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# Entanglement between Collective Operators in the Linear Harmonic Chain
## I Introduction
Quantum entanglement is a physical phenomenon in which the quantum states of two or more systems can only be described with reference to each other, even though the individual systems may be spatially separated. This leads to correlations between observables of the systems that cannot be understood on the basis of classical (local realistic) theories Bell1964 . Its importance today exceeds the realm of the foundations of quantum physics and entanglement has become an important physical resource, like energy, that allows performing communication and computation tasks with efficiency which is not achievable classically Niel1998 .
In the near future we will certainly see more and more experiments on entanglement of increasing complexity. Moving to higher entangled systems or entangling more systems with each other, will eventually push the realm of quantum physics well into the macroscopic world. It will be therefore important to investigate under which conditions entanglement within or between ”macroscopic” objects, each consisting of a sample containing a large number of the constituents, can arise.
Recently, it was shown that macroscopic entanglement can arise ”naturally” between constituents of various complex physical systems. Examples of such systems are chains of interacting spin systems Arne2001 ; Niel1998 , harmonic oscillators Aude2002 ; Sera2005 and quantum fields Rezn2003 . Entanglement can have an effect on the macroscopic properties of these systems Gosh2003 ; Bruk2005 ; Wies2005 and can be in principle extractable from them for quantum information processing Rezn2003 ; Pate2004 ; deCh2005 ; Retz2005 .
With the aim of better understanding macroscopical entanglement we will investigate entanglement between collective operators in this paper. A simple and natural system is the ground state of a linear chain of harmonic oscillators furnished with harmonic nearest-neighbor interaction. The mathematical entanglement properties of this system were extensively investigated in Aude2002 ; Sera2005 ; Pate2005 ; Bote2004 . Entanglement was computed in the form of logarithmic negativity for general bisections of the chain and for contiguous blocks of oscillators that do not comprise the whole chain. It was shown that the log-negativity typically decreases exponentially with the separation of the groups and that the larger the groups, the larger the maximal separation for which the log-negativity is non-zero Aude2002 . It also was proven that an area law holds for harmonic lattice systems, stating that the amount of entanglement between two complementary regions scales with their boundary Cram2005 .
In a real experimental situation, however, we are typically not able to determine the complete mathematical amount of entanglement (as measured, e.g., by log-negativity) which is non-zero even if two blocks share only one arbitrarily weak entangled pair of oscillators. Our measurement apparatuses normally cannot resolve single oscillators, but rather interact with a whole bunch of them in one way, potentially even in non-contiguous regions, thus measuring certain global properties. Here we will study entanglement between ”physical blocks” of harmonic oscillators — existing only if there is entanglement between the collective operators defined on the entire blocks — as a function of their size, relative distance and the coupling strength. Our aim is to quantify (experimentally accessible) entanglement between global properties of two groups of harmonic oscillators. Surprisingly, we will see that such collective entanglement can be demonstrated even in the case where none of the oscillators from one block is entangled with an oscillator from the other block (i.e., it cannot be understood as a cumulative effect of entanglement between pairs of oscillators), which is in agreement with Aude2002 . This shows the existence of bipartite entanglement between collective operators.
Because of the area law Cram2005 the amount of entanglement is relatively small in the first instance. We suggest a way to overcome this problem by allowing the collective blocks to consist of a periodic sequence of subblocks. Then the total boundary region between them is increased and we verify that indeed a larger amount of entanglement is found for periodic blocks, where the entanglement scales at most with the total boundary region. We give an analytical approximation of this amount of entanglement and motivate how it can in principle be extracted from the chain Rezn2003 ; Pate2004 ; deCh2005 ; Retz2005 .
Methodologically, we will quantify the entanglement between collective operators of two blocks of harmonic oscillators by using a measure for continuous variable systems based on the Peres-Horodecki criterion Pere1996 ; Horo1997 ; Simo2000 ; Kim2002 . The collective operators will be defined as sums over local operators for all single oscillators belonging to the block. The infinite harmonic chain is assumed to be in the ground state and since the blocks do not comprise the whole chain, they are in a mixed state.
## II Linear Harmonic Chain
We investigate a linear harmonic chain, where each of the $`N`$ oscillators is situated in a harmonic potential with frequency $`\omega `$ and each oscillator is coupled with its neighbors by a harmonic potential with the coupling frequency $`\mathrm{\Omega }`$. The oscillators have mass $`m`$ and their positions and momenta are denoted as $`\overline{q}_i`$ and $`\overline{p}_i`$, respectively. Assuming periodic boundary conditions ($`\overline{q}_{N+1}\overline{q}_1`$), the Hamiltonian thus reads Schw2003
$$H=\underset{j=1}{\overset{N}{}}(\frac{\overline{p}_j^2}{2m}+\frac{m\omega ^2\overline{q}_j^2}{2}+\frac{m\mathrm{\Omega }^2(\overline{q}_j\overline{q}_{j1})^2}{2}).$$
(1)
We canonically go to dimensionless variables: $`q_jC\overline{q}_j`$ and $`p_j\overline{p}_j/C`$, where $`C\sqrt{m\omega (1+2\mathrm{\Omega }^2/\omega ^2)^{1/2}}`$ Bote2004 . By this means the Hamiltonian becomes
$$H=\frac{E_0}{2}\underset{j=1}{\overset{N}{}}(p_j^2+q_j^2\alpha q_jq_{j+1}),$$
(2)
with the abbreviations $`\alpha 2\mathrm{\Omega }^2/(2\mathrm{\Omega }^2+\omega ^2)`$ and $`E_0\sqrt{2\mathrm{\Omega }^2+\omega ^2}`$. The coupling constant is restricted to values $`0<\alpha <1`$, where $`\alpha 0`$ in the weak coupling limit ($`\mathrm{\Omega }/\omega 0`$) and $`\alpha 1`$ in the strong coupling limit ($`\mathrm{\Omega }/\omega \mathrm{}`$).
In the language of second quantization the positions and momenta are converted into operators ($`q_j\widehat{q}_j`$, $`p_j\widehat{p}_j`$) and are expanded into modes of their annihilation and creation operators, $`\widehat{a}`$ and $`\widehat{a}^{}`$, respectively:
$`\widehat{q}_j`$ $`={\displaystyle \frac{1}{\sqrt{N}}}{\displaystyle \underset{k=0}{\overset{N1}{}}}{\displaystyle \frac{1}{\sqrt{2\nu (\theta _k)}}}[\widehat{a}(\theta _k)\text{e}^{\text{i}\theta _kj}+\text{H.c.}],`$ (3)
$`\widehat{p}_j`$ $`={\displaystyle \frac{\text{i}}{\sqrt{N}}}{\displaystyle \underset{k=0}{\overset{N1}{}}}\sqrt{{\displaystyle \frac{\nu (\theta _k)}{2}}}[\widehat{a}(\theta _k)\text{e}^{\text{i}\theta _kj}\text{H.c.}].`$ (4)
Here $`\theta _k2\pi k/N`$ (with $`k=0,1,\mathrm{},N1`$) is the dimensionless pseudo-momentum and $`\nu (\theta _k)\sqrt{1\alpha \mathrm{cos}\theta _k}`$ is the dispersion relation. The annihilation and creation operators fulfil the well known commutation relation $`[\widehat{a}(\theta _k),\widehat{a}^{}(\theta _k^{})]=\delta _{kk^{}}`$, since $`[\widehat{q}_i,\widehat{p}_j]=`$i$`\delta _{ij}`$ has to be guaranteed. The ground state (vacuum), denoted as $`|0`$, is defined by $`\widehat{a}(\theta _k)|0=0`$ holding for all $`\theta _k`$. The two-point vacuum correlation functions
$`g_{|ij|}`$ $`0\left|\widehat{q}_i\widehat{q}_j\right|0\widehat{q}_i\widehat{q}_j,`$ (5)
$`h_{|ij|}`$ $`0\left|\widehat{p}_i\widehat{p}_j\right|0\widehat{p}_i\widehat{p}_j,`$ (6)
are given by $`g_l=(2N)^1_{k=0}^{N1}\nu ^1(\theta _k)\mathrm{cos}(l\theta _k)`$ and $`h_l=(2N)^1_{k=0}^{N1}\nu (\theta _k)\mathrm{cos}(l\theta _k)`$, where $`l|ij|`$. In the limit of an infinite chain ($`N\mathrm{}`$) — which we will study below — and for $`l<N/2`$ they can be expressed in terms of the hypergeometric function $`{}_{2}{}^{}F_{1}^{}`$ Bote2004 $`g_l=(z^l/2\mu )\left(\genfrac{}{}{0pt}{}{l1/2}{l}\right)_2F_1(1/2,l+1/2,l+1,z^2)`$, $`h_l=(\mu z^l/2)\left(\genfrac{}{}{0pt}{}{l3/2}{l}\right)_2F_1(1/2,l1/2,l+1,z^2)`$, where $`z(1\sqrt{1\alpha ^2})/\alpha `$ and $`\mu 1/\sqrt{1+z^2}`$.
## III Defining Collective Operators
In the following, we are interested in entanglement between two ”physical blocks” of oscillators, where the blocks are represented by a specific form of collective operators which are normalized sums of individual operators. By means of such a formalism we seek to fulfil experimental conditions and constraints, since finite experimental resolution implies naturally the measurement of, e.g., the average momentum of a bunch of oscillators rather than the momentum of only one. On the other hand, this formalism can easily take account of blocks that consist of non-contiguous regions, leading to interesting results which will be shown below. We want to point out that this convention of the term block is not the same as it is normally used in the previous literature. In contrast to the latter, for which one allows any possible measurement, our simulation of realizable experiments already lacks some information due to the averaging.
Let us now consider two non-overlapping blocks of oscillators, $`A`$ and $`B`$, within the closed harmonic chain in its ground state, where each block contains $`n`$ oscillators. The blocks are separated by $`d0`$ oscillators (Fig. 1). We assume $`n,dN`$ and $`N\mathrm{}`$ for the numerical calculations of the two-point correlation functions.
By a Fourier transform we map the $`n`$ oscillators of each block onto $`n`$ (”orthogonal”) frequency-dependent collective operators
$`\widehat{Q}_A^{(k)}`$ $`{\displaystyle \frac{1}{\sqrt{n}}}{\displaystyle \underset{jA}{}}\widehat{q}_j\text{e}^{{\scriptscriptstyle \frac{2\pi \text{i}jk}{n}}},`$ (7)
$`\widehat{P}_A^{(k)}`$ $`{\displaystyle \frac{1}{\sqrt{n}}}{\displaystyle \underset{jA}{}}\widehat{p}_j\text{e}^{{\scriptscriptstyle \frac{2\pi \text{i}jk}{n}}},`$ (8)
with the frequencies $`k=0,\mathrm{},n1`$, and analogously for block $`B`$. The commutator of the collective position and momentum operators is
$$[\widehat{Q}_A^{(k)},\widehat{P}_A^{(k^{})}]=\text{i}\delta _{kk^{}}.$$
(9)
This means that collective operators for different frequencies $`kk^{}`$ commute. For different blocks the commutator vanishes: $`[\widehat{Q}_A^{(k)},\widehat{P}_B^{(k^{})}]=0`$.
If the individual positions and momenta of all oscillators are written into a vector
$$\widehat{𝐱}(\widehat{q}_1,\widehat{p}_1,\widehat{q}_2,\widehat{p}_2,\mathrm{},\widehat{q}_N,\widehat{p}_N)^\text{T},$$
(10)
then there holds the commutation relation
$$[\widehat{x}_i,\widehat{x}_j]=\text{i}\mathrm{\Omega }_{ij}$$
(11)
with $`𝛀`$ the $`n`$-fold direct sum of $`2\times 2`$ symplectic matrices:
$$𝛀\underset{j=1}{\overset{n}{}}(\begin{array}{cc}0& 1\\ 1& 0\end{array}).$$
(12)
A matrix $`𝐒`$ transforms $`\widehat{𝐱}`$ into a vector of collective (and uninvolved individual) oscillators:
$$\widehat{𝐗}𝐒\widehat{𝐱}=(\{\widehat{Q}_A^{(k)},\widehat{P}_A^{(k)}\}_k,\{\widehat{Q}_B^{(k)},\widehat{P}_B^{(k)}\}_k,\{\widehat{q}_j,\widehat{p}_j\}_j)^\text{T}.$$
(13)
Here $`\{\widehat{Q}_A^{(k)},\widehat{P}_A^{(k)}\}_k=(\widehat{Q}_A^{(0)},\widehat{P}_A^{(0)},\mathrm{},\widehat{Q}_A^{(n1)},\widehat{P}_A^{(n1)})`$ denotes all collective oscillators of block $`A`$ and analogously for block $`B`$, whereas $`\{\widehat{q}_j,\widehat{p}_j\}_j`$ denotes the $`2(N2n)`$ position and momentum entries of those $`N2n`$ oscillators which are not part of one of the two blocks. The matrix $`𝐒`$ corresponds to a Gaussian operation Eise2003 . It has determinant det$`𝐒=1`$ and preserves the symplectic structure
$$𝛀=𝐒^\text{T}𝛀𝐒,$$
(14)
and hence
$$[\widehat{X}_i,\widehat{X}_j]=\text{i}\mathrm{\Omega }_{ij}$$
(15)
for all $`i,j`$, in particular verifying (9). This means that the Gaussianness of the ground state of the harmonic chain (i.e., the fact that the state is completely characterized by its first and second moments, see below) was preserved by the (Fourier) transformation to the frequency-dependent collective operators.
## IV Quantifying Entanglement between Collective Operators
In reality, we are typically not capable of single particle resolution measurements and only of measuring the collective operators with one frequency, namely $`k=0`$. Note that in general the correlations of higher-frequency collective operators, e.g., $`(\widehat{Q}_A^{(k)})^2`$ or $`\widehat{Q}_A^{(k)}\widehat{Q}_B^{(k)}`$ with $`k0`$, are not real numbers. Therefore, as a natural choice, we denote as the collective operators
$`\widehat{Q}_A`$ $`\widehat{Q}_A^{(0)}={\displaystyle \frac{1}{\sqrt{n}}}{\displaystyle \underset{jA}{}}\widehat{q}_j,`$ (16)
$`\widehat{P}_A`$ $`\widehat{P}_A^{(0)}={\displaystyle \frac{1}{\sqrt{n}}}{\displaystyle \underset{jA}{}}\widehat{p}_j,`$ (17)
and analogously for block $`B`$. It seems to be a very natural situation that the experimenter only has access to these collective properties and we are interested in the amount of (physical) entanglement one can extract from the system if only the collective observables $`\widehat{Q}_{A,B}`$ and $`\widehat{P}_{A,B}`$ are measured.
Reference Simo2000 derives a separability criterion which is based on the Peres-Horodecki criterion Pere1996 ; Horo1997 and the fact that — in the continuous variables case — the partial transposition allows a geometric interpretation as mirror reflection in phase space. Following largely the notation in the original paper, we introduce the vector
$$\widehat{\xi }(\widehat{Q}_A,\widehat{P}_A,\widehat{Q}_B,\widehat{P}_B)$$
(18)
of collective operators. The commutation relations have the compact form $`[\widehat{\xi }_\alpha ,\widehat{\xi }_\beta ]=`$i$`K_{\alpha \beta }`$ with $`𝐊_{j=1}^2\left(\genfrac{}{}{0.0pt}{}{0}{1}\genfrac{}{}{0.0pt}{}{1}{0}\right)`$. The separability criterion bases on the covariance matrix (of first and second moments)
$$V_{\alpha \beta }\frac{1}{2}\mathrm{\Delta }\widehat{\xi }_\alpha \mathrm{\Delta }\widehat{\xi }_\beta +\mathrm{\Delta }\widehat{\xi }_\beta \mathrm{\Delta }\widehat{\xi }_\alpha ,$$
(19)
where $`\mathrm{\Delta }\widehat{\xi }_\alpha \widehat{\xi }_\alpha \widehat{\xi }_\alpha `$ with $`\widehat{\xi }_\alpha =0`$ in our case (state around the origin of phase space).
The covariance matrix $`𝐕`$ is real (which would not be the case for higher-frequency collective operators) and symmetric: $`\widehat{Q}_A\widehat{Q}_B=\widehat{Q}_B\widehat{Q}_A`$ and $`\widehat{P}_A\widehat{P}_B=\widehat{P}_B\widehat{P}_A`$, coming from the fact that the two-point correlation functions (5) and (6) only depend on the absolute value of the position index difference. On the other hand, using (3) and (4), we verify that $`\widehat{q}_i\widehat{p}_j=`$i$`(2N)^1_{k=0}^{N1}\mathrm{exp}[`$i$`\theta _k(ij)]`$ and $`\widehat{p}_j\widehat{q}_i=`$i$`(2N)^1_{k=0}^{N1}\mathrm{exp}[`$i$`\theta _k(ji)]`$. For $`ij`$ both summations vanish ($`\theta _k2\pi k/N`$ and $`i,j`$ integer) and for $`i=j`$ they are the same but with opposite sign. Thus, in all cases $`\widehat{q}_i\widehat{p}_j=\widehat{p}_j\widehat{q}_i`$. These symmetries also hold for the collective operators and hence we obtain
$$𝐕=(\begin{array}{cccc}G& 0& G_{AB}& 0\\ 0& H& 0& H_{AB}\\ G_{AB}& 0& G& 0\\ 0& H_{AB}& 0& H\end{array}).$$
(20)
The matrix elements are
$`G`$ $`\widehat{Q}_A^2=\widehat{Q}_B^2={\displaystyle \frac{1}{n}}{\displaystyle \underset{jA}{}}{\displaystyle \underset{iA}{}}g_{|ji|},`$ (21)
$`H`$ $`\widehat{P}_A^2=\widehat{P}_B^2={\displaystyle \frac{1}{n}}{\displaystyle \underset{jA}{}}{\displaystyle \underset{iA}{}}h_{|ji|},`$ (22)
$`G_{AB}`$ $`\widehat{Q}_A\widehat{Q}_B={\displaystyle \frac{1}{n}}{\displaystyle \underset{jA}{}}{\displaystyle \underset{iB}{}}g_{|ji|},`$ (23)
$`H_{AB}`$ $`\widehat{P}_A\widehat{P}_B={\displaystyle \frac{1}{n}}{\displaystyle \underset{jA}{}}{\displaystyle \underset{iB}{}}h_{|ji|}.`$ (24)
To quantify entanglement between two collective blocks we use the degree of entanglement $`\epsilon `$, given by the absolute sum of the negative eigenvalues of the partially transposed density operator: $`\epsilon `$Tr$`|\rho ^{\text{T}_B}|1`$, i.e., by measuring how much the mirror reflected state fails to be positive definite. This measure (negativity) is based on the Peres-Horodecki criterion Pere1996 ; Horo1997 and was shown to be an entanglement monotone Lee2000 ; Vida2002 . For covariance matrices of the form (20) it reads Kim2002
$$\epsilon =\mathrm{max}(0,\frac{(\delta _1\delta _2)_0}{\delta _1\delta _2}1),$$
(25)
where $`\delta _1G|G_{AB}|`$ and $`\delta _2H|H_{AB}|`$. In general, the numerator is defined by the square of the Heisenberg uncertainty relation
$$(\delta _1\delta _2)_0\left(\frac{1}{2}|[\widehat{Q}_{A,B},\widehat{P}_{A,B}]|\right)^2,$$
(26)
with $`(\delta _1\delta _2)_0=1/4`$ due to (9). We note that $`\epsilon `$ is a degree of entanglement (in the sense of necessity and sufficiency) only for Gaussian states which are completely characterized by their first and second moments, as for example the ground state of the harmonic chain we are studying. However, we left out the higher-frequency collective operators (and all the oscillators which are not part of the blocks) and therefore, the entanglement $`\epsilon `$ has to be understood as the Gaussian part of the amount of entanglement which exists between (and can be extracted from) the two blocks when only the collective properties $`\widehat{Q}_{A,B}`$ and $`\widehat{P}_{A,B}`$, as defined in (16) and (17), are accessible.
There also exists an entanglement witness in form of a separability criterion based on variances, where $`\mathrm{\Delta }(\widehat{Q}_A\widehat{Q}_B)^2+(\widehat{P}_A+\widehat{P}_B)^2=2(GG_{AB}+H+H_{AB})<2`$ is a sufficient condition for the state to be entangled Duan2000 . We note that the above negativity measure (25) is ”stronger” than this witness in the whole parameter space ($`\alpha ,n`$). In particular, there are cases where $`\epsilon >0`$ although $`\mathrm{\Delta }2`$. This is in agreement with the finding that the variance criterion is weaker than a generalized negativity criterion Shch2005 .
We further note that the amount of entanglement (25) is invariant under a change of potential redefinitions of the collective operators, e.g., $`\widehat{Q}_A_{jA}\widehat{q}_j`$ or $`\widehat{Q}_A(1/n)_{jA}\widehat{q}_j`$, as then the modified scaling in the correlations ($`G`$, $`G_{AB}`$, $`H`$, and $`H_{AB}`$) is exactly compensated by the modified scaling of the Heisenberg uncertainty in the numerator.
Figure 2 shows the results for $`d=0`$ and $`d=1`$. In the first case — if the blocks are neighboring — there exists entanglement for all possible coupling strengths $`\alpha `$ and block sizes $`n`$. In the latter case — if there is one oscillator between the blocks — due to the strongly decaying correlation functions $`g`$ an $`h`$ there is no entanglement between two single oscillators ($`n=1`$), but entanglement for larger blocks (up to $`n=4`$, depending on $`\alpha `$). The statement that entanglement can emerge by going to larger blocks was also found in Aude2002 . But there the blocks were abstract objects, containing all the information of their constituents. In the case of collective operators, however, increasing the block size (averaging over more oscillators) is also connected with a loss of information. In spite of this loss and the mixedness of the state, two blocks can be entangled, although non of the individual pairs between the blocks is entangled — indicating true bipartite entanglement between collective operators. For $`d2`$, however, no entanglement can be found anymore.
These results are in agreement with the general statement that entanglement between a region and its complement scales with the size of the boundary Cram2005 . In the present case of two blocks in a one-dimensional chain (Fig. 1) the boundary is constant and as the blocks are made larger, the entanglement decreases since it is distributed over more and more oscillators. We therefore propose to increase the number of boundaries by considering two non-overlapping blocks, where we allow a periodic continuation of the situation above, i.e. a sequence of $`m1`$ subblocks, separated by $`d`$ oscillators and each consisting of $`s1`$ oscillators, where $`ms=n`$ (Fig. 3).
The degree of entanglement between two periodic blocks of non-separated ($`d=0`$) one-particle subblocks ($`s=1`$) is larger for stronger coupling constant $`\alpha `$ and grows with the overall number of oscillators $`n`$ (Fig. 4a). For given $`\alpha `$ and $`n`$ and no separation between the subblocks ($`d=0`$) the entanglement is larger for the case of small subblocks, as then there are many of them, causing a large total boundary (Fig. 4b). Entanglement can be even found for larger separation ($`d=1,2`$) with a more complicated dependence on the size $`s`$ of the subblocks. There is a trade-off between having a large number of boundaries and the fact that one should have large subblocks as individual separated oscillators are not entangled (Fig. 4c,d). For $`d3`$ no entanglement can be found anymore. (In a realistic experimental situation, where the separation $`d`$ is not sharply defined, e.g., where there are weighted contributions for $`d=0,1,\mathrm{},d_{\text{max}}`$, entanglement can persist even for $`d_{\text{max}}3`$, depending on the weighting factors.)
For the sake of completeness we give a rough approximation of the entanglement between two periodic blocks. Let us assume that the subblocks are directly neighbored, $`d=0`$. Furthermore, we consider couplings $`\alpha `$ such that we may neglect higher than next neighbor correlations ($`\alpha 0.5`$), i.e., we only take into account $`g_0`$, $`g_1`$, $`h_0`$ and $`h_1`$. The correlations read
$`G`$ $`={\displaystyle \frac{1}{n}}{\displaystyle \underset{jA}{}}{\displaystyle \underset{iA}{}}g_{|ji|}g_0+{\displaystyle \frac{2m(s1)}{n}}g_1,`$ (27)
$`G_{AB}`$ $`={\displaystyle \frac{1}{n}}{\displaystyle \underset{jA}{}}{\displaystyle \underset{iB}{}}g_{|ji|}{\displaystyle \frac{1}{n}}(2m1)g_1,`$ (28)
and analogously for $`H`$ and $`H_{AB}`$. The first equation reflects that there are $`n`$ self-correlations and $`m(s1)`$ nearest neighbor pairs (which are counted twice) within one block, i.e., $`s1`$ pairs per subblock. The second equation represents the fact that there are $`2m1`$ boundaries where blocks $`A`$ and $`B`$ meet. Using $`s=n/m`$, the entanglement (25) becomes (note that $`g_1>0`$ and $`h_1<0`$)
$$\epsilon \frac{1}{4[g_0+(2\frac{4m1}{n})g_1][h_0+(2\frac{1}{n})h_1]}1.$$
(29)
For given $`n`$ this approximation obviously increases with the total number of boundaries, $`m`$. It can be considered as an estimate for a situation like in Fig. 4b, if a smaller coupling is used such that the neglect of higher correlations becomes justified.
We close this section by annotating that the entanglement (25) between collective blocks of oscillators — being the Gaussian part — can in principle (for sufficient control of the block separation $`d`$) be transferred to two remote qubits via a Jaynes–Cummings type interaction Rezn2003 ; Pate2004 ; Retz2005 . For the interaction with periodic blocks ”gratings” have to be employed in the experimental setup. The interaction Hamiltonian is of the form
$`\widehat{H}_{\text{int}}`$ $`(\text{e}^{\text{ i }\omega _1t}\widehat{\sigma }_1^++\text{e}^{+\text{ i }\omega _1t}\widehat{\sigma }_1^+)\widehat{Q}_A`$
$`+(\text{e}^{\text{ i }\omega _2t}\widehat{\sigma }_2^++\text{e}^{+\text{ i }\omega _2t}\widehat{\sigma }_2^+)\widehat{Q}_B,`$ (30)
where $`\omega _i`$ is the Rabi frequency and $`\widehat{\sigma }_i^+=(\widehat{\sigma }_i^{})^{}=|e_i_ig|`$ is the bosonic operator (with $`|g_i`$ and $`|e_i`$ the ground and the excited state) of the $`i`$-th qubit ($`i=1,2`$).
## V Collective Operators for Scalar Quantum Fields
The continuum limit of the linear harmonic chain is the (1+1)-dimensional Klein–Gordon field $`\varphi (x,t)`$ with the canonical momentum field $`\pi (x,t)=\dot{\varphi }(x,t)`$. It satisfies the Klein–Gordon equation (in natural units $`\mathrm{}=c=1`$) with mass $`m`$
$$\ddot{\varphi }^2\varphi +m^2\varphi =0.$$
(31)
With the canonical quantization procedure $`\varphi `$ and $`\pi `$ become operators satisfying the non-trivial commutation relation $`[\widehat{\varphi }(x,t),\widehat{\pi }(x^{},t)]=`$i$`\delta (xx^{})`$. The field operator can be expanded into a Fourier integral over elementary plane wave solutions Bjo2003
$`\widehat{\varphi }(x,t)`$ $`={\displaystyle \frac{\text{d}k}{\sqrt{4\pi \omega _k}}\left[\widehat{a}(k)\text{e}^{\text{i}kx\text{i}\omega _kt}+\text{H.c.}\right]},`$ (32)
$`\widehat{\pi }(x,t)`$ $`=\text{i}{\displaystyle \frac{\text{d}k\omega _k}{\sqrt{4\pi }}\left[\widehat{a}(k)\text{e}^{\text{i}kx\text{i}\omega _kt}\text{H.c.}\right]},`$ (33)
where $`k`$ is the wave number and $`\omega _k=+\sqrt{k^2+m^2}`$ is the dispersion relation. The annihilation and creation operators fulfil $`[\widehat{a}(k),\widehat{a}^{}(k^{})]=\delta (kk^{})`$. We write the field operator as a sum of two contributions $`\widehat{\varphi }=\widehat{\varphi }^{(+)}+\widehat{\varphi }^{()}`$, where $`\widehat{\varphi }^{(+)}`$ ($`\widehat{\varphi }^{()}`$) is the contribution with positive (negative) frequency. Thus, $`\widehat{\varphi }^{(+)}`$ corresponds to the term with the annihilation operator in (32). The vacuum correlation function is given by the (equal-time) commutator of the positive and the negative frequency part:
$$0\left|\widehat{\varphi }(x,t)\widehat{\varphi }(y,t)\right|0=[\widehat{\varphi }^{(+)}(x,t),\widehat{\varphi }^{()}(y,t)].$$
(34)
It is a peculiarity of the idealization of quantum field theory that for $`x=y`$ this propagator diverges in the ground state:
$$0\left|\widehat{\varphi }^2(x,t)\right|0\mathrm{}.$$
(35)
The same is true for $`0\left|\widehat{\pi }^2(x,t)\right|0`$ and hence we cannot easily build an entanglement measure like for the harmonic chain, since the analogs of the two-point correlation functions $`g_0`$ and $`h_0`$, (5) and (6), are divergent now. Automatically, we are motivated to study the more physical situation and consider extended space-time regions, which means that we should integrate the field (and conjugate momentum) over some spatial area. We define the collective operators
$`\widehat{\mathrm{\Phi }}_L(x_0,t)`$ $`{\displaystyle \frac{1}{\sqrt{L}}}{\displaystyle _{L/2}^{L/2}}\widehat{\varphi }(x+x_0,t)\text{d}x,`$ (36)
$`\widehat{\mathrm{\Pi }}_L(x_0,t)`$ $`{\displaystyle \frac{1}{\sqrt{L}}}{\displaystyle _{L/2}^{L/2}}\widehat{\pi }(x+x_0,t)\text{d}x,`$ (37)
Therefore, $`\widehat{\mathrm{\Phi }}_L(x_0,t)`$ and $`\widehat{\mathrm{\Pi }}_L(x_0,t)`$ are equal-time operators which are spatially averaged over a length $`L`$, centered at position $`x_0`$. The commutator is
$$[\widehat{\mathrm{\Phi }}_L(x_0,t),\widehat{\mathrm{\Pi }}_L(x_0,t)]=\frac{1}{L}_{L/2}^{L/2}_{L/2}^{L/2}\text{i}\delta ^{(3)}(xy)=\text{i},$$
(38)
which is in complete analogy to (9). If $`\widehat{\mathrm{\Phi }}_L`$ and $`\widehat{\mathrm{\Pi }}_L`$ correspond to separated regions without overlap, i.e., $`|x_0y_0|>L`$, then of course $`[\widehat{\mathrm{\Phi }}_L(x_0,t),\widehat{\mathrm{\Pi }}_L(y_0,t)]=0`$. The spatial integration in (36) and (37) can be carried out analytically:
$`\widehat{\mathrm{\Phi }}_L(x_0,t)`$ $`={\displaystyle \frac{1}{\sqrt{\pi L}}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}{\displaystyle \frac{\text{d}k}{k\sqrt{\omega _k}}}\mathrm{sin}\left({\displaystyle \frac{kL}{2}}\right)`$
$`\times [\widehat{a}(k)\text{e}^{\text{i}kx_0\text{i}\omega _kt}+\text{H.c.}],`$ (39)
$`\widehat{\mathrm{\Pi }}_L(x_0,t)`$ $`={\displaystyle \frac{\text{i}}{\sqrt{\pi L}}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}{\displaystyle \frac{\text{d}k\sqrt{\omega _k}}{k}}\mathrm{sin}\left({\displaystyle \frac{kL}{2}}\right)`$
$`\times [\widehat{a}(k)\text{e}^{\text{i}kx_0\text{i}\omega _kt}\text{H.c.}].`$ (40)
The final step is to calculate the propagators of the field and the conjugate momentum. We find
$`D_{\widehat{\mathrm{\Phi }},L}(r)`$ $`0\left|\widehat{\mathrm{\Phi }}_L(x_0,t)\widehat{\mathrm{\Phi }}_L(y_0,t)\right|0`$ (41)
$`={\displaystyle \frac{1}{\pi L}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}{\displaystyle \frac{\text{d}k}{k^2\sqrt{k^2+m^2}}}\mathrm{sin}^2\left({\displaystyle \frac{kL}{2}}\right)\mathrm{cos}(kr),`$
$`D_{\widehat{\mathrm{\Pi }},L}(r)`$ $`0\left|\widehat{\mathrm{\Pi }}_L(x_0,t)\widehat{\mathrm{\Pi }}_L(y_0,t)\right|0`$ (42)
$`={\displaystyle \frac{1}{\pi L}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}{\displaystyle \frac{\text{d}k\sqrt{k^2+m^2}}{k^2}}\mathrm{sin}^2\left({\displaystyle \frac{kL}{2}}\right)\mathrm{cos}(kr),`$
with $`r|x_0y_0|`$ the distance between the centers of the two regions, reflecting the spatial symmetry. Thus $`D_{\widehat{\mathrm{\Phi }},L}(0)`$ and $`D_{\widehat{\mathrm{\Pi }},L}(0)`$ are the analogs of $`\widehat{Q}_{A,B}^2`$ and $`\widehat{P}_{A,B}^2`$ (intra-block correlations within the same block), respectively, whereas $`D_{\widehat{\mathrm{\Phi }},L}(r>L)`$ and $`D_{\widehat{\mathrm{\Pi }},L}(r>L)`$ correspond to $`\widehat{Q}_A\widehat{Q}_B`$ and $`\widehat{P}_A\widehat{P}_B`$ (inter-block correlations between separated blocks).
The expressions (41) and (42) are finite, especially for $`r=0`$. Mathematically, the integration over a finite spatial region $`L`$ corresponds to a cutoff, which removes the divergence we faced in (35). However, the expressions are ill defined for $`L0`$.
Applying the entanglement measure (25) with $`G=D_{\widehat{\mathrm{\Phi }},L}(0)`$, $`H=D_{\widehat{\mathrm{\Pi }},L}(0)`$, $`G_{AB}=D_{\widehat{\mathrm{\Phi }},L}(r)`$, and $`H_{AB}=D_{\widehat{\mathrm{\Pi }},L}(r)`$ does not indicate entanglement for any choice of $`L`$ and $`r>L`$. The same is true for the generalized case of blocks consisting of periodic subregions of space, showing an inherent difference between the harmonic chain and its continuum limit. We believe this is due to the fact, that any spatial integration immediately corresponds to an infinitely large block in the discrete harmonic chain and that the information loss (compared to the mathematical indeed existing exponentially small entanglement Rezn2003 ) due to the collective operators already is too large. Nonetheless, defining collective operators like in (36) and (37) and use of the measure (25) may reveal entanglement between spatially separated regions for other quantum field states, which is the subject of future research.
## VI Conclusion
Our results have importance for investigating the conditions under which entanglement can be detected by measuring collective observables of blocks consisting of a large number of harmonic oscillators. This has relevance for schemes of extracting entanglement where the probe particles normally interact with whole (periodic) groups of oscillators rather than single oscillators. The results are also relevant for the transition from the quantum to the classical domain as they suggest that entanglement between collective operators (global properties) may persist even in the limit of a large number of particles. It is obvious that our approach of collective observables can be extended to more dimensions. Furthermore, we demonstrated its potential application to scalar quantum field theory.
## Acknowledgement
We thank A. Ferreira and M. Wieśniak for stimulating discussions. J. K. and Č. B. acknowledge financial support by the FWF (Austrian Science Fund), SFB project P06, as well as the European Commission under the Integrated Project Qubit Applications QAP funded by the IST Directorate as contract number 015846. V. V. acknowledges support from the Engineering and Physical Sciences Council, the British Council in Austria and the European Union.
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# Spectral representation of the effective dielectric constant of graded composites
## I Introduction
In graded materials Milton1 the physical properties may vary continuously in space making them distinctly different from homogeneous materials. Hence, composite media consisting of graded inclusions have attracted much interest in various engineering applications Yamanouchi , such as reduced residual and thermal barrier coatings of high temperature components in gas turbines, surface hardening for tribological protection, and graded interlayers used in multilayered microelectronic and optoelectronic components Holt ; Cherradi .
Like graded materials, thin films are of great interest in many practical applications and often possess different optical properties Kammler in comparison to bulk materials. Recently, it was found experimentally that graded thin films may have high relative dielectric permittivity as well as a flatter temperature characteristic of permittivity Kawasaki than single-layer films LuAPL03 .
The traditional theories used to deal with the homogeneous materials Jackson ; AIP , however fail to deal with composites of graded inclusions. To treat these composites, we have recently developed a first-principles approach Dong ; Gu-JAP and a differential effective dipole approximation Huang .
This work has been motivated by a recent study of the optical absorption spectrum of a graded metallic film Huang1 . In that work, a broad surface plasmon absorption band was observed in addition to a strong Drude absorption peak at zero frequency. Such a broad absorption band has been shown to be responsible for the enhanced nonlinear optical response as well as an attractive figure of merit (the degree of optical absorption). Yuen et al. Yuen1 pointed out that such an absorption spectrum, being related to the imaginary part of the effective dielectric constant, should equally well be reflected in the Bergman-Milton spectral representation of the effective dielectric constant Bergman ; Milton .
Bergman-Milton spectral representation was originally developed for calculating the effective dielectric constant and other response functions of two-component composites Bergman ; Milton . However, the two concerned components are all homogeneous. Therefore, it is worth extending the spectral representation to graded composite materials. The work on graded films is just a simple example of a more general graded composite in three dimensions. One of the main purpose of this work help to identify the physical origin of the broad absorption band. It turns out that, unlike in the case of homogeneous materials, the characteristic function of a graded composite is a continuous function because of the continuous variation of the dielectric function within the constituent component.
Moreover, we apply our theory to a special case of graded composites, i. e., multilayer material, which is more convenient to fabricate in practice than graded material Hobson , and many algorithms are now available for designing of multilayer coatings Martin ; Verly . Thus, the present work is necessary in the sense that we shall discuss the multilayer effect as the number of layers inside the material increases. In this regard, this work should be expected to have practical relevance. As the number of layers $`N`$ increases, we shall show a gradual transition from sharp peaks to a broad continuous band until the graded composite results are recovered by the limit of $`N\mathrm{}`$.
The paper is organized as follows. In Sec. II, the general derivation of the spectral representation for graded composites is presented. In Sec. III, we describe the model and present analytical results for the spectral representation of the effective dielectric constant of a graded film with an interface, as well as a graded sphere. In those cases, the Bergman-Milton formalism has been modified for graded composites. We further obtain an analytic form for the spectral density function of a multilayer film and a multilayer sphere in Sec. IV. Numerical results are presented in Sec. V, and discussion and conclusion are given in Sec. VI.
## II Formalism
We consider a two-component composite in which graded inclusions of dielectric constant $`ϵ_1(𝐫)`$ are embedded in a homogeneous host medium of dielectric constant $`ϵ_2`$. It is noted that the dielectric constant $`ϵ_1(𝐫)`$ is a gradation profile as a function of the position $`𝐫`$. And we will restrict our discussion and calculation to the quasi-static approximation, i. e., $`dc/\omega 1`$, where $`d`$ is the characteristic size of the inclusion, $`c`$ is the speed of light in vacuum and $`\omega `$ is the frequency of the applied field. In the quasi-static approximation, the whole graded structure can be regarded as an effective homogeneous one with effective (overall) linear dielectric constant defined as Stroud
$$ϵ_e=\frac{1}{V}\frac{𝐄𝐃}{E_0^2}𝑑V,$$
(1)
where $`E_0`$ is the applied electric field along $`z`$ direction, $`𝐄`$ and $`𝐃`$ are the local electric field and local displacement, respectively.
The object of the present section is to solve the Laplace’s equation
$$(ϵ(𝐫)\varphi (𝐫))=0$$
(2)
subject to the boundary condition $`\varphi _0=E_0z`$. The dielectric function $`ϵ(𝐫)`$ varies from component to component but has a fixed mathematical expression for a given component. It can be expressed as Bergman
$$ϵ(𝐫)=ϵ_2\left[1\frac{1}{s}\eta (𝐫)\right],$$
(3)
where $`s=[1ϵ_{\mathrm{ref}}/ϵ_2]^1`$ is the material parameter and $`ϵ_{\mathrm{ref}}`$ is some reference dielectric constant in the graded component. The characteristic function $`\eta (𝐫)`$ is may be written in terms of a real function $`f(𝐫)`$ as
$$\eta (𝐫)=\{\begin{array}{cc}1+f(𝐫)& \text{ in inclusion},\\ 0& \text{in host},\end{array}$$
which accords for the microstructure of graded composites. The function $`f(𝐫)`$ depends on the specific variation of the dielectric constant in the inclusion component. For homogeneous constituent component, i. e., $`f(𝐫)=0`$, $`\eta (𝐫)=1`$ in the inclusion component, while $`\eta (𝐫)=0`$ in the host medium. For graded systems, $`\eta (𝐫)`$ can be a continuous function in the inclusion component because of the continuous variation of the dielectric function within the inclusion component. Thus, Eq. (1) can be solved
$$\varphi (𝐫)=E_0z+\frac{1}{s}𝑑V^{}\eta (𝐫^{})^{}G(𝐫𝐫^{})^{}\varphi (𝐫^{}),$$
(4)
where $`G(𝐫𝐫^{})`$ is a Green’s function satisfying:
$$\{\begin{array}{cc}^2G(𝐫𝐫^{})=\delta ^3(𝐫𝐫^{})\text{for }\text{r}\text{ in V},& \\ G=0\text{ for }\text{r}\text{ on the boundary.}& \end{array}$$
In order to obtain a solution for Eq. (2), we introduce an integral-differential Hermitian operator $`\widehat{\mathrm{\Gamma }}`$, which satisfies
$$\widehat{\mathrm{\Gamma }}𝑑V^{}\eta (𝐫^{})^{}G(𝐫𝐫^{})^{},$$
and define an inner product as
$$\varphi |\psi =𝑑V\eta (𝐫)\varphi ^{}\psi .$$
(5)
With the above definitions, Eq. (4) can be simplified to
$$\varphi (𝐫)=E_0z+\frac{1}{s}\widehat{\mathrm{\Gamma }}\varphi (𝐫).$$
Let $`s_n`$ and $`|\varphi _n`$ be the $`n`$th eigenvalue and eigenfunction of operator $`\widehat{\mathrm{\Gamma }}`$. Then, the generalized eigenvalue problem becomes
$$(\eta (𝐫)\varphi _n)=s_n^2\varphi _n.$$
The potential $`|\varphi `$ can be expanded in series of eigenfunctions,
$$|\varphi \underset{n}{}\left(\frac{s}{s_ns}\right)\frac{|\varphi _n\varphi _n|z}{\varphi _n|\varphi _n},$$
(6)
where we choose $`E_0=1`$ for convenience. Since $`\eta (𝐫)`$ is a real function, the eigenvalues $`s_n`$ will be real. Also, for graded component, $`\eta (𝐫)`$ is a continuous function, which will cover the full region, i. e., $`\mathrm{}\eta (𝐫)\mathrm{}`$. Therefore, the eigenvalues $`s_n`$, which depend on the continuously graded microstructure $`\eta (𝐫)`$, do not lie within the interval $`[0,1]`$ but extend to $`\mathrm{}s_n\mathrm{}`$ as first pointed by Gu and Gong GuYing for three-component composites case. However, eigenvalues $`s_n`$ still lie in $`[0,1]`$ for $`0\eta (𝐫)1`$.
We are now in the position to find an analytical representation for the effective dielectric constant $`ϵ_e`$ according to Eq. (1). We take advantage of Green’s theorem, the boundary condition $`\varphi _0=z`$, and the Maxwell equation $`𝐃=0`$ to obtain the effective dielectric constant
$`{\displaystyle \frac{ϵ_e}{ϵ_2}}`$ $`=`$ $`{\displaystyle \frac{1}{ϵ_2V}}{\displaystyle \left(\varphi \right)𝐃𝑑V}`$
$`=`$ $`{\displaystyle \frac{1}{V}}{\displaystyle \widehat{𝐳}\left[\left(1\frac{1}{s}\eta (𝐫)\right)\varphi \right]𝑑V}`$
$`=`$ $`1+{\displaystyle \frac{1}{sV}}z|\varphi .`$
If we now introduce the reduced response Bergman
$$F(s)=1\frac{ϵ_e}{ϵ_2},$$
(8)
and substitute Eq. (6) into Eq. (II) we find
$$F(s)=\frac{1}{V}\underset{n}{}\frac{\left|z|\varphi _n\right|^2}{\varphi _n|\varphi _n}\left(\frac{1}{ss_n}\right).$$
We can now express the effective dielectric constant as
$$ϵ_e=ϵ_2\left(1\underset{n}{}\frac{f_n}{ss_n}\right),$$
(9)
where $`f_n`$ is given by
$$f_n=\frac{1}{V}\frac{\left|z|\varphi _n\right|^2}{\varphi _n|\varphi _n}.$$
Using the above equations, we obtain the following sum rule
$`{\displaystyle \underset{n}{}}f_n`$ $`=`$ $`{\displaystyle \frac{1}{V}}z|z`$ (10)
$`=`$ $`{\displaystyle \frac{1}{V}}{\displaystyle 𝑑V\eta (𝐫)zz}`$
$`=`$ $`{\displaystyle \frac{1}{V}}{\displaystyle 𝑑V\eta (𝐫)}.`$
It is worth noting that the sum rule will not equal to the volume fraction of inclusion. This is different from the Bergman-Milton spectral representation for two homogeneous systems, in which the sum rule equals to the volume fraction of the inclusion.
When the operator $`\widehat{\mathrm{\Gamma }}`$ has a continuous spectrum, Eq. (9) should be replaced with the integral form
$$ϵ_e=ϵ_2\left(1𝑑s^{}\frac{m(s^{})}{ss^{}}\right),$$
(11)
where $`m(s^{})`$ is the spectral density function. Then, the reduced response becomes
$$F(s)=𝑑s^{}\frac{m(s^{})}{ss^{}}.$$
(12)
If we write $`s`$ as $`s+i0^+`$, the right side of Eq. (12) becomes
$$P𝑑s^{}\frac{m(s^{})}{ss^{}}i\pi m(s),$$
and thus, $`m(s^{})`$ is given through the limiting process
$$m(s^{})=\frac{1}{\pi }\mathrm{Im}[F(s^{}+i0^+)].$$
(13)
This final result is identical in form to Bergman’s expression for the analogous function in scalar composite materials. However, there are differences in the derivation, namely, the definition of the inner product Eq. (5), the continuous graded microstructure $`\eta (𝐫)`$, the sum rule, as well as the range of eigenvalues $`s_n`$.
From Eq. (11) it is evident that if the spectral density function $`m(s^{})`$ is known, the effective dielectric constant can be obtained accurately, and vice versa. The spectral representation has been used to analyze the effective dielectric properties of composites. Recently, Levy and Bergman Levy also used it in their study of nonlinear optical susceptibility. In this regard, Sheng and coworkers Ma developed a practical algorithm for calculating the effective dielectric constants based on the spectral representation. In what follows, we restrict ourselves to a graded composite both in one dimension and three dimensions, as well as corresponding multilayer composites.
## III Spectral density function of graded composites
### III.1 Spectral density function of a graded film
We consider a graded dielectric film of width $`L`$ , in which two media meet at a planar interface as shown in Fig. 1 (a). The first medium $`ϵ_1(z)`$ varies along $`z`$axis, while the second medium $`ϵ_2`$ is homogeneous. We define the graded microstructure as
$$\eta (z)=\{\begin{array}{cc}1+az& 0<zh,\\ 0& h<z<L,\end{array}$$
(14)
where $`a`$ and $`h`$ are real constants. They can be varied to describe different graded films. Thus, according to Eq. (3), the dielectric function of graded film can be expressed as
$$ϵ(z)=ϵ_2\left(1\frac{\eta (z)}{s}\right).$$
(15)
Owing to the simple geometry of a graded film, we can use the equivalent capacitance of a series combination to calculate the effective dielectric constant as
$$\frac{1}{ϵ_e}=\frac{1}{L}_0^L\frac{1}{ϵ(z)}𝑑z.$$
(16)
Substituting Eqs. (14) and (15) into Eq. (16), we obtain
$$\frac{1}{ϵ_e}=\frac{1h}{ϵ_2}+\frac{s\left[\mathrm{ln}\left(1{\displaystyle \frac{\eta (0)}{s}}\right)\mathrm{ln}\left(1{\displaystyle \frac{\eta (h)}{s}}\right)\right]}{aϵ_2},$$
with the assumption $`L=1`$.
We are now in a position to extend the Bergman-Milton spectral representation of the effective dielectric constant Bergman ; Milton to a graded film. For a graded system, $`\eta (z)`$ can be a continuous function in the inclusion medium. Using Eqs. (11)$``$(13), we obtain the spectral density function for a graded film as
$$m(s^{})=\frac{as^{}\mathrm{arg}\left({\displaystyle \frac{s1}{sah1}}\right)}{\pi \left[\left(s^{}\mathrm{arg}\left({\displaystyle \frac{s1}{sah1}}\right)\right)^2+\left(a(h1)s^{}\mathrm{ln}\left({\displaystyle \frac{s^{}1}{sah1}}\right)\right)^2\right]},$$
where $`s^{}=\mathrm{Re}[s]`$ and $`\mathrm{arg}[\mathrm{}]`$ denote the arguments of complex functions.
### III.2 Spectral density function of a graded sphere
The above theory can be generalized to graded composites in three dimensions. We consider a graded sphere with dielectric constant $`ϵ_1(r)`$ embedded into a homogeneous host medium with dielectric constant $`ϵ_2`$. The dielectric constant of the graded sphere $`ϵ_1(r)`$ varies along the radius $`r`$. We can obtain the effective dielectric constant of a graded sphere using the spectral representation. We consider the graded microstructure as
$$\eta (r)=\{\begin{array}{cc}1+ar& 0<rR,\\ 0& r>R,\end{array}$$
where $`R`$ is the radius of the graded sphere. Thus, from Eq. (3) the dielectric constant in the graded sphere is given by
$$ϵ(r)=ϵ_2\left(1\frac{\eta (r)}{s}\right).$$
(17)
In the dilute limit the effective dielectric constant of a small volume fraction $`p`$ of graded spheres embedded in a host medium is given by YuHui ; Zhang
$$ϵ_e=ϵ_2+3ϵ_2\mathrm{𝑝𝑏},$$
(18)
where $`b`$ is the dipole factor of graded spheres embedded in a host as given in Ref. Dong . Using Eq. (8) and Eq. (18), the reduced response can be obtained as
$$F(s)=3ϵ_2\mathrm{𝑝𝑏}.$$
(19)
Thus, the spectral density function of a graded sphere can be given through a numerical evaluation of Eq. (13).
## IV spectral density function of multilayer composites
A multilayer composite is a special case of graded composites. The gradation becomes continuous as the number of layers approaches infinity. To investigate the multilayer effect, we shall use a finite difference approximation for the graded profile (Eqs. (15) and (17)) for a finite number of layers. To mimic a multilayer system, we divide the interval $`[0,1]`$ into $`N`$ equally spaced sub-intervals, $`[0,z_1]`$, $`(z_1,z_2]`$, $`\mathrm{}`$, $`(z_N1,1]`$. Then we adopt the midpoint value of $`ϵ(z)`$ for each sub-interval as the dielectric constant of that sublayer. In this way, we calculate the effective dielectric constant, eigenvalues, as well spectral density function for each $`N`$. It is worth noting that the results of $`N\mathrm{}`$ (e. g., $`N=1024`$ recovers the results of graded composites.
In addition to multilayer films, we can use the above approach to study the much simpler problem of a two-layer film. In this system, we have two layers of dielectric constants $`ϵ_1`$, $`ϵ_2`$, and host $`ϵ_0`$. Thickness are $`hy`$, $`h(1y)`$, and $`1h`$, respectively, where $`y`$ is the length ratio between component $`ϵ_1`$ and component $`ϵ_2`$. We also define two microstructure parameters, $`\eta _1`$ and $`\eta _2`$. If we let $`s=1/(1ϵ_1/ϵ_0)`$, then $`\eta _1=1`$, and $`\eta _2=(ϵ_0ϵ_2)/(ϵ_0ϵ_1)`$. According to Eq. (15), the effective dielectric constant of the two-layer film is now given by
$$\frac{1}{ϵ_e}=\frac{hy}{ϵ_1}+\frac{h(1y)}{ϵ_2}+\frac{1h}{ϵ_0}.$$
According to Eq. (8), the reduced response can be given by
$$F(s)=\frac{F_1}{ss_1}+\frac{F_2}{ss_2},$$
(20)
where
$`F_1`$ $`=`$ $`{\displaystyle \frac{h(s_1(yy\eta +\eta )\eta )}{s_1s_2}},`$
$`F_2`$ $`=`$ $`{\displaystyle \frac{h(s_2(yy\eta +\eta )\eta )}{s_1s_2}},`$
$`s_1`$ $`=`$ $`{\displaystyle \frac{1}{2}}[1h(yy\eta +\eta )+\eta `$
$`\sqrt{4\eta (1+h)+(1h(yy\eta +\eta )+\eta )^2}],`$
$`s_2`$ $`=`$ $`{\displaystyle \frac{1}{2}}[1h(yy\eta +\eta )+\eta `$
$`+\sqrt{4\eta (1+h)+(1h(yy\eta +\eta )+\eta )^2}].`$
From the sums of $`F_1`$ and $`F_2`$ and the integral of graded microstructure $`\eta (z)`$ given by Eq. (14), we can check that the sum rule expressed by Eq. (10) is obeyed. It should also be noted that there are two poles in the expression for the reduced response corresponding to two peaks in the spectral density function. If $`h=1`$, then $`s_1=0`$, that is, one peak is located at zero, which is explicitly shown in Fig. 2(a).
Similarly, we can also apply our graded spectral representation to a single-shell sphere of core dielectric constant $`ϵ_1`$, covered by a shell of $`ϵ_2`$, and suspended in a host of $`ϵ_0`$. In this example, we can also define two microstructure parameters $`\eta _1`$ and $`\eta _2`$. If we let $`s=1/(1ϵ_1/ϵ_0)`$, then $`\eta _1=1`$, and $`\eta _2=(ϵ_0ϵ_2)/(ϵ_0ϵ_1)`$. The dipole factor of single-shell sphere is given YuHui ; Zhang
$$b=\frac{ϵ_2ϵ_0+(ϵ_0+2ϵ_2)xf^3}{ϵ_2+2ϵ_0+2(ϵ_2ϵ_0)xf^3},$$
where $`f`$ is the ratio between radius core and radius shell, and $`x`$ is given by
$$x=\frac{ϵ_1ϵ_2}{ϵ_1+2ϵ_2}.$$
Then, we can also write Eq. (19) similarly to Eq. (20), where the residues and eigenvalues are given by,
$`F_1`$ $`=`$ $`{\displaystyle \frac{3ps_1[(1+\eta )y^3\eta ]\eta p[12(1+\eta )y^3+2\eta ]}{3(s_1s_2)}},`$
$`F_2`$ $`=`$ $`{\displaystyle \frac{3ps_2[(1+\eta )y^3\eta ]+\eta p[12(1+\eta )y^3+2\eta ]}{3(s_1s_2)}},`$
$`s_1`$ $`=`$ $`{\displaystyle \frac{1}{6}}\left[1+3\eta \sqrt{1+(28y^3)\eta +(1+8y^3)\eta ^2}\right],`$
$`s_2`$ $`=`$ $`{\displaystyle \frac{1}{6}}\left[1+3\eta +\sqrt{1+(28y^3)\eta +(1+8y^3)\eta ^2}\right].`$
Analysis shows that the spectral representation for $`N=2`$ contains two simple poles corresponding to two peaks in the spectral density function. Therefore, we draw the conclusion that, $`N`$ peaks are a result of $`N`$ layers. Moreover, $`N1`$ peaks will accumulate into a continuous broad absorption spectrum when $`N`$ tends toward infinity, which can be seen from Fig. 4(f) and 5(f).
## V Numerical results
We are now in a position to do some numerical calculations of the spectral density function from Eqs. (11) and (13). A small but finite imaginary part in the complex parameter has been used in the calculations. Without any loss of generality, we choose $`L=1`$ and $`R=1`$ for convenience. We show the effect of different graded profiles, as well as the effect of the thickness of the inclusion. It should be noted, that in all figures the range of $`s`$ is limited to $`[0,1]`$, because we chose $`1<a<0`$ which limits the value of $`\eta `$ into $`[0,1]`$.
Fig. 1 displays the dielectric profile of a graded film (Fig. 1(a)) and a graded sphere (Fig. 1(b)). This figure obviously shows that the dielectric constant varies with the position in inclusion while a constant in host medium. Also, different values of $`a`$ accord with different graded materials.
In Fig. 2(a), we plot the spectral density function $`m(s)`$ of a graded film without an interface against the spectral parameter for various graded microstructures $`\eta (z)`$. It is evident that there is always a broad continuous band in the spectral density function. Both the strength as well as the width of the continuous part of $`m(s)`$ increase with the gradient of the dielectric profile. Thus, the previous results of the broad surface-plasmon band can be expected. Note that there is a sharp peak at $`s=0`$, which is also present in a homogeneous film. In Fig. 2(b), we plot the spectral density function of a graded film meeting a homogeneous medium at an interface for various graded microstructure $`\eta (z)`$. Again, there is always a broad continuous band in the spectral density function. However, the sharp peak has now shifted to a finite value of $`s`$, which is also present in a homogeneous film.
In Fig. 3, the spectral density function of graded sphere is displayed for a volume fraction $`p=0.1`$. In this case, the interface always exists. It is clear that a broad continuous function in the spectral density function is always observed, as well as the shift of the sharp peak. However, the decrease of the broad continuous function is more abrupt for graded sphere than for graded film with increasing $`s`$.
Figures 4 and 5 display the spectral density function for a multilayer film and a sphere, respectively. It is clear that there are always $`N`$ sharp peaks for $`N`$ layers. Moreover, it is worth noting that there occurs a transition from sharp peaks to a broad continuous band with increasing $`N`$ (see Fig. 4(f) and Fig. 5(f)), that is, the graded results are recovered by the limit results of $`N\mathrm{}`$. In particular, we had obtained the analytical expression of spectral density function for $`N=2`$. There are two resonances corresponding to the two peaks in Fig. 4(a) and Fig. 5(a).
## VI Discussion and Conclusion
We have investigated a graded composite film and a sphere by means of the Bergman-Milton spectral representation. It has been shown that the spectral density function can be obtained analytically for a graded system. However, unlike in the case of homogeneous constituent components, the characteristic function is a continuous function due to the presence of gradation. Moreover, the derivation as well as some salient properties, namely, the sum rule, the definition of inner product, the definition of the integral-differential operators, and the range of spectral parameters, do change because of the continuous variation of the dielectric profile within the constituent components. It should be noted that in graded composite, the eigenvalues are not limited to $`[0,1]`$, and they can be extended to $`\mathrm{}s_n\mathrm{}`$ for the full region $`\eta `$, i. e., $`\mathrm{}\eta \mathrm{}`$. In this work, however for simplicity, we investigated the spectral density function in $`0s1`$ by choosing $`1<a<0`$ to limit the value of $`\eta `$ into $`[0,1]`$.
We also study multilayer composites and calculated the spectral density function versus the number of layers, to explicitly demonstrate that the broad continuous spectrum arises from the accumulation of poles when the number of layers tends to infinity. This finding coincides with the broad surface-plasmon absorption band associated with the optical properties of graded composites.
To sum up, we have investigated the spectral density function of graded film and graded sphere, as well as multilayer cases. There is always a broad continuous function in the spectral density function in graded composite, but simple poles in multilayer composite, the number of pole depends on the number of layers. Moreover, there is a gradual transition from sharp peaks to a broad continuous band until the graded composite results recover in the limit of $`N\mathrm{}`$.
###### Acknowledgements.
This work was supported by the Research Grant Council of Hong Kong SAR Government Earmarked Grant (K. W. Y.) and by the Academy of Finland (L. D., M. K.).
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# Three-dimensional stability of the solar tachocline
## 1 Motivation
Helioseismology has provided us with information on the internal rotation of the Sun. The angular velocity depends on latitude as well as radius. The dependence is mainly a latitudinal in the bulk of the convection zone, whereas the solar radiative core rotates nearly uniform with an angular velocity of the convection zone at about $`30^{}`$ latitude. The transition from the differential rotation in the convection zone to the uniform rotation in the core is thin – probably thinner than 5% of the solar radius – and is called the tachocline.
The convection zone is thermally overcritical and the stability of shear flows is not an issue there. The shear in the tachocline however is latitudinal and radial and may be subject to shear-flow instabilities. If the tachocline is turned into a turbulent layer, a problem arises with the mixing of elements, most notably with Lithium which will be destroyed in nuclear fusion up to 0.68 solar radii ($`R_{}`$), just below the convection zone. A turbulent tachocline mixing the Lithium into its fusion zone would contradict the observed relatively high Lithium abundance at the surface of the Sun.
The situation in the solar tachocline was described by Spiegel & Zahn (1992) as exhibiting horizontal turbulence, excited by hydrodynamic shear-flow instability. It was argued that the two-dimensional turbulence provides angular momentum transport but inhibits too strong mixing of Lithium below $`0.68R_{}`$.
The hydrodynamic stability of the latitudinal shear in the tachocline was studied by Watson (1981). The dependence on colatitude $`\theta `$ was $`\mathrm{\Omega }=\mathrm{\Omega }_{\mathrm{eq}}(1\alpha _2\mathrm{cos}^2\theta )`$, where $`\mathrm{\Omega }_{\mathrm{eq}}`$ is the equatorial angular velocity at the bottom of the convection zone and $`\alpha _2`$ is a parameter for the relative equator-to-pole difference of the angular velocity. Watson found instability for a differential rotation with $`\alpha _2>0.286`$. At that time, this was supposed to be in the vicinity of the value of solar differential rotation, and the result was not definitely deciding between a turbulent and a stable tachocline. The investigation was two-dimensional arguing that the stable stratification will not allow significant radial flows, but thereby the radial shear is neglected, too.
Charbonneau et al. (1999a) extended the linear stability analysis to various rotation profiles of the form
$$\mathrm{\Omega }=\mathrm{\Omega }_{\mathrm{eq}}(1\alpha _2\mathrm{cos}^2\theta \alpha _4\mathrm{cos}^4\theta ).$$
(1)
Different layers of the tachocline were associated with different results from helioseismology in terms of $`\alpha _2`$ and $`\alpha _4`$. The upper layer which is thought to be penetrated by convective overshooting was found to be unstable to the shear, whereas the lower layer with smaller latitudinal shear turned out to be stable. In a step towards the three-dimensional stability of the tachocline, Dikpati & Gilman (2001) studied the two-dimensional system allowing for deformations in the third – the radial – dimension. As a result, the critical differential rotation for instability was reduced to 11% in the overshoot part of the tachocline.
In order to address the entire tachocline and since the tachocline is a place where latitudinal and radial shear meet, we investigate the stability of the three-dimensional rotation profile. Although it is very reasonable to assume that radial flows will be weak, the variation of the latitudinal differential rotation with radius across weakly coupled spherical layers could provide different results for the stability of the tachocline. We will briefly summarize the numerical background in Section 2, provide details of the computational results in Section 3 with different rotatation profiles, and summarize our findings in Section 4. We find that the tachocline is hydrodynamically stable in all the configurations studied. The paper is thus intended to be a preparatory step towards fully three-dimensional studies of the MHD stability of the tachocline.
## 2 Computational setup
The rotational profile depends on both latitude and radius in this study. Between the inner and outer radius of the tachocline, $`r_\mathrm{i}=0.65`$ and $`r_\mathrm{o}=0.7`$ respectively, the angular velocity is defined by
$$\mathrm{\Omega }(r,\theta )=\mathrm{\Omega }_{\mathrm{eq}}[1\alpha _2\mathrm{cos}^2\theta \alpha _2(\frac{1}{4}\mathrm{cos}^2\theta )\frac{r_\mathrm{o}r}{r_\mathrm{o}r_\mathrm{i}}],$$
(2)
where $`\theta `$ denotes the colatitude in the spherical shell, $`r`$ is the radial coordinate, and $`\mathrm{\Omega }_{\mathrm{eq}}`$ is the equatorial angular velocity. The profile implies that the rotation velocity at the inner boundary of the computational domain, $`r=r_\mathrm{i}`$, is the one at $`\theta =60^{}`$ (or $`30^{}`$ heliographic latitude). This appears to be a valid assumption in agreement with various helioseismological inversions (recently e.g. Schou et al. 2002), whereas core rotation may be adopted at larger latitudes higher up in the convection zone. The formation of the tachocline rotation profile is supposed to be caused by rotating convection on top of it and reduction of differential rotation by magnetic fields in the solar core at the bottom. Such a profile causes a meridional circulation reaching a steady-state in competition with the aforementioned (or other tachocline-forming) effects (see e.g. Sule et al. 2005). In our linear analysis, this flow has no effect on the stability of the non-axisymmetric modes investigated in this Paper.
We employ the incompressible, viscous Navier-Stokes equation and linearize the problem. We can then separate the axisymmetric background rotation $`U`$ from the nonaxisymmetric flow $`u`$. The latter is evolved by numerical computations. The normalised equation of motion reads
$`{\displaystyle \frac{u}{t}}`$ $`=`$ $`u\times \times U+U\times \times u`$ (3)
$`p(uU)+\mathrm{}u,`$
and the continuity equation $`u=0`$ holds. The equation is evolved with the spectral spherical code by Hollerbach (2000). We look for exponentially decaying or growing solutions in order to find the critical differential rotation $`\alpha _2`$ of the marginal case. The actual integration employs the radial components of the curl and the curl-curl of the equation, thereby eliminating the gradient terms. The linear (viscous) part of the full equation of motion is evolved implicitly, while the nonlinear parts are integrated explicitly with the advection term and forces being computed in real space. We kept this splitting even in our linearized problem, since a fully implicit scheme would have required a substantial modification of the code. The $`\mathrm{}u`$ is thus treated implicitly, the other rhs terms (two remaining after curling) are computed in real space and are used for a second-order Runge-Kutta integration.
The normalisation of the equation with the viscous time $`\tau _\nu =r_\mathrm{o}^2/\nu `$ and the length scale $`r_\mathrm{o}`$ leads to the Reynolds number
$$\mathrm{Re}=\frac{r_\mathrm{o}^2\mathrm{\Omega }_{\mathrm{eq}}}{\nu }$$
(4)
as a free parameter which is essentially a variation of the viscosity $`\nu `$ since radius and $`\mathrm{\Omega }_{\mathrm{eq}}`$ are sufficiently well known. The solar Reynolds number in the tachocline is – in terms of the definition of (4) – about $`10^{14}`$. We try to achieve time series for numerically demanding $`\mathrm{Re}>10^4`$. By comparison with known results from inviscid two-dimensional analyses, we find that the critical viscous differential rotation at $`\mathrm{Re}>10^3`$ or $`10^4`$ is already sufficiently close to the inviscid value.
The velocity is decomposed into toroidal and poloidal potentials, $`u=\times (e\widehat{r})+\times \times (f\widehat{r})`$, where $`\widehat{r}`$ is the unit vector in radial direction. The potentials are again decomposed into radial Chebyshev polynomials and spherical harmonics. Since it is the potentials being evolved, the continuity equation is fulfilled automatically. The density is constant throughout the computational domain. In a thin shell of 5% of the solar radius, this is a reasonably good approximation of the true situation. We also do not take any deformation of the tachocline into account. Top and bottom radius of the tachocline are constant over latitude. Charbonneau et al. (1999b) found a prolate tachocline whose equatorial part is located at slightly smaller radius than the polar end. We assume that the difference of 3.5% in location of the tachocline has negligible effect on the results, but may be an issue of future investigations.
The radial boundary conditions for the velocity perturbations are stress-free at both $`r_\mathrm{i}`$ and $`r_\mathrm{o}`$. At Reynolds numbers of $`\mathrm{Re}>10^4`$, high spectral resolution was necessary to obtain reliable results. Up to 80 Chebyshev and 80 Legendre polynomials were used to resolve the flow properly.
The azimuthal modes of the problem described by (3) are decoupled, and we can study the stability of individual $`m`$-modes separately. Moreover, even and odd latitudinal modes (symmetric and antisymmetric modes with respect to the equator) decouple, and we will have a look into the critical differential rotation for the excitation of instability of the two kinds separately. The radial modes all couple and do not provide results on the stability of individual radial wavelengths.
## 3 Results
### 3.1 Stability of various solutions
In a set of fiducial computations, we applied a purely latitudinal profile of the angular velocity. An $`m=1`$ mode is evolved with the profile of (1) where $`\alpha _4=0`$ and $`\alpha _2`$ is varied. Since we solve a viscous problem, the critical differential rotation, $`\alpha _2^{\mathrm{crit}}`$, depends on the Reynolds number. The result is already very close to Watson’s inviscid result for $`\mathrm{Re}1000`$. This is good reason to assume that the numerical solutions reaching $`\mathrm{Re}\mathrm{30\hspace{0.17em}000}`$ are suitable approximations for the solar plasma.
Despite allowing for radial motions, the evolution provides solutions which are nearly toroidal and do not show significant radial flows. They are surface flows forming two cells on each hemisphere with stream lines through the poles. Figure 2 shows a representation of the flow in the spherical surface. The graph has to assume that the poloidal component of the velocity is zero though, which is not entirely true.
The second step involved the rotation profile of (2) for which the critical steepness of the differential rotation, $`\alpha _2^{\mathrm{crit}}`$, is again sought for various Reynolds numbers. Figure 3 shows the lines of marginal stability, i.e. the critical differential rotation, versus Reynolds number for the symmetric $`m=1`$ mode. The solid line refers to profile (2), the dashed line is the latitudinal profile and converges to the result by Watson (1981) for $`\mathrm{Re}\mathrm{}`$.
The most easily excited patterns of $`m=1`$ are always symmetric with respect to the equator. We can also look for the stability of antisymmetric patterns and find the results shown in Fig. 4. They are more stable than the symmetric configurations with an $`\mathrm{\Omega }(r,\theta )`$ profile. The antisymmetric solutions from the $`\mathrm{\Omega }(\theta )`$ profile have also higher critical differential rotation values than their symmetric counterparts.
The patterns drift with a certain velocity in azimuthal direction around the solar axis. Since the equations hold for the nonrotating system, we can directly convert the pattern rotation into physical times. The pattern rotation period for the two-dimensional and the three-dimensional profile of $`\mathrm{\Omega }`$ is shown in Fig. 5 by a dashed and a solid line, respectively. The actual solar rotation periods are also plotted. The reference equatorial period of 25.44 d ($`\mathrm{\Omega }_{\mathrm{eq}}=455`$ nHz) is given as dash-dot line. The pattern rotation periods are determined at marginal stability. Since the marginal case gives us a value for the differential rotation, we also plot the polar rotation period $`P_{\mathrm{pole}}`$ versus Re. The short-dash line is for the two-dimensional $`\mathrm{\Omega }(\theta )`$ profile, the dotted line is for the three-dimensional case described by (2).
The pattern rotation periods are always between the equatorial and polar rotation periods, in agreement with the 2D results by Charbonneau et al. (1999a). While the patterns from the 2D-$`\mathrm{\Omega }`$ profile are close to the polar rotation period, the patterns for the 3D profile rotate with nearly the average rotation period between the polar and equatorial ones.
We can compute the time after which the pattern is passed by a given point on the equator. This time is often called lap time. Assuming an equatorial rotation period of 25.44 d ($`\mathrm{\Omega }_{\mathrm{eq}}=455`$ nHz), we find a lap time of 91 d for the 2D case, and a lap time of 78 d for the 3D case.
Modes with higher azimuthal mode numbers require significantly higher differential rotation for instability. The antisymmetric $`m=2`$ mode which is symmetric with respect to the equator was found to be stable even in the entire parameter range covered by Fig. 3. This is in agreement with the inviscid, two-dimensional stability analysis by Charbonneau et al. (1999a). The stability lines for the antisymmetric $`m=2`$ mode are shown in Fig. 6. We could not find instability for any $`m=3`$ mode in the range covered by Fig. 3.
### 3.2 Effects of buoyancy
Little influence is expected from the stable temperature gradient in the tachocline. Since we do have the chance to prove this in our three-dimensional simulations, we demonstrate the effect of a negative buoyancy force on the stability of the differential rotation. The Navier-Stokes equation is extended by the buoyancy force and reads in the non-dimensional form:
$`{\displaystyle \frac{u}{t}}`$ $`=`$ $`u\times \times U+U\times \times u`$ (5)
$`+\mathrm{Ra}\mathrm{\Theta }rp(uU)+\mathrm{}u,`$
$`{\displaystyle \frac{\mathrm{\Theta }}{t}}`$ $`=`$ $`U\mathrm{\Theta }uT+{\displaystyle \frac{1}{\mathrm{Pr}}}\mathrm{}\mathrm{\Theta },`$ (6)
with a background temperature profile of
$$T=\frac{r_\mathrm{i}}{r_\mathrm{o}r_\mathrm{i}}\left(\frac{r_\mathrm{o}}{r}1\right)$$
(7)
and the Prandtl number being the ratio between viscosity and thermal diffusivity, $`\mathrm{Pr}=\nu /\chi `$. The Rayleigh number in (6) is
$$\mathrm{Ra}=\frac{g\alpha (T_\mathrm{i}T_\mathrm{o})r_\mathrm{o}^3}{\nu ^2},$$
(8)
where $`g`$ is the gravitational acceleration, $`\alpha `$ is the coefficient of volume expansion, and $`T_\mathrm{i}`$ and $`T_\mathrm{o}`$ are the temperatures at the two boundaries. In the Boussinesq formulation used here, the presence of a sub-adiabatic temperature gradient actually translates into a negative value of Ra. We set our version of the Rayleigh number to a value as small (“as negative”) as $`\mathrm{Ra}=10^8`$ in order to see any notable effect on the flow. The Prandtl number is set to unity.
The critical differential rotation for a growing symmetric $`m=1`$ mode at $`\mathrm{Re}=10^4`$ increases slightly to 53.3% as compared with the non-buoyant value of 52%. This is in line with the fact that the solutions contain nearly horizontal motions.
### 3.3 Effects of higher-degree terms
The differential rotation has been expressed by a $`\mathrm{cos}^2\theta `$ dependence in latitude and a linear $`r`$ dependence over radius. We are studying the stability of the three-dimensional tachocline setup with higher-degree dependences in this Section, such as the $`\mathrm{cos}^4\theta `$ term and a nonlinear radial dependence of $`\mathrm{\Omega }`$.
Including the $`\mathrm{cos}^4\theta `$ term in (2) yields an angular velocity profile of the form
$`\mathrm{\Omega }(r,\theta )`$ $`=`$ $`\mathrm{\Omega }_{\mathrm{eq}}\{1\alpha _2\mathrm{cos}^2\theta \alpha _4\mathrm{cos}^4\theta {\displaystyle \frac{r_\mathrm{o}r}{r_\mathrm{o}r_\mathrm{i}}}`$ (9)
$`[\alpha _2({\displaystyle \frac{1}{4}}\mathrm{cos}^2\theta )+\alpha _4({\displaystyle \frac{1}{16}}\mathrm{cos}^4\theta )]\}`$
The computations for the full profile could not easily be extended beyond $`\mathrm{Re}=5000`$. But the results for the possible Re and for the easier two-dimensional $`\mathrm{\Omega }(\theta )`$ show no decreased critical differential rotation when the $`\mathrm{cos}^4\theta `$ term is included. This also holds for the extreme case of $`\alpha _4`$ carrying the shear alone ($`\alpha _2=0`$).
Starting from (2), a modified $`\mathrm{\Omega }`$-profile was constructed in order to find any influence of the particular radial dependence of the latitudinal shear on the results. We used
$`\mathrm{\Omega }(r,\theta )`$ $`=`$ $`\mathrm{\Omega }_{\mathrm{eq}}[1\alpha _2\mathrm{cos}^2\theta `$ (10)
$`\alpha _2({\displaystyle \frac{1}{4}}\mathrm{cos}^2\theta ){\displaystyle \frac{1}{2}}(1\mathrm{cos}({\displaystyle \frac{r_\mathrm{o}r}{r_\mathrm{o}r_\mathrm{i}}}\pi ))]`$
The radial profile introduces two inflection points. It does not apply the usual error function for the radial $`\mathrm{\Omega }`$-step in order to obtain an exact $`\mathrm{\Omega }/r=0`$ at both inner and outer boundary. A graphic representation is shown in Fig. 7.
The critical differential rotation for instability is shown versus Reynolds number in Fig. 8. We were able to reach extraordinary Reynolds numbers of $`10^5`$, probably because of the vanishing $`\mathrm{\Omega }/r`$ at the boundaries. Instability does not occur at reduced differential rotation, and the line of marginal stability actually converges to the simpler profile (2) for high Re. The additional effect of buoyancy on the stability of the profile (10) with inflection points is also shown as a dashed line.
## 4 Summary
A fully three-dimensional, linear analysis of the stability of the solar tachocline was carried out. If radial variation of the angular velocity is included in the model, the maximum pole-equator difference of the angular velocity can be as large as 52% for a symmetric $`m=1`$ mode, before instability sets in. Two-dimensional analyses have delivered much lower critical values. The difference of 3D versus 2D is not radial flows emerging from the extension in the third dimension, but it is changed stability conditions emerging from the radial shear and radial dependence of the differential rotation.
Other modes such as higher $`m`$ or different flow symmetries do not get unstable at lower critical differential rotation values under the influence of a three-dimensional rotation profile.
The stabilizing effect of the temperature gradient has been added, but since all the unstable modes are very nearly horizontal, the influence is small. The assumption that horizontal motions dominate is valid even without a stabilizing temperature gradient. However the assumption that spherical shells of infinitesimal thickness do not interact with each other is not applicable, according to our results. One may argue that the viscosity in the computer simulations is much too high, but the variation of the results is small at $`\mathrm{Re}>1000`$. This is an indication that computations with $`\mathrm{Re}=10^4`$ or higher are a good approximation of the near-inviscid solar case. We conclude that all parts of the tachocline not being affected by convective overshooting are stable. We believe that the answer on the question of the dynamical state of the tachocline will be given by MHD computations of a three-dimensional domain.
###### Acknowledgements.
We are grateful to Peter Gilman for his valuable comments on the topic. AS thanks the Deutsche Forschungsgemeinschaft for the support by grant No. Ru 488/15-1.
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# Quark model description of quasi-elastic pion knockout from the proton at JLAB
## 1 Introduction
In Ref. the JLab F$`\pi `$1 Collaboration extracted from the data the charged pion form factor $`F_\pi (Q^2)`$ using a Regge model for high energy meson electroproduction . Along with this the recent JLab data on the Rosenbluth separation of longitudinal and transverse cross sections for the reaction $`p(e,e^{}\pi ^+)n`$ were presented for $`Q^2=`$ 0.6, 0.75, 1 and 1.6 GeV<sup>2</sup>/c<sup>2</sup> and for the invariant mass $`W=`$ 1.95 GeV.
The procedure for reggeization according to is quite natural for the meson ($`\pi `$, $`\rho `$, $`b_1`$, …) $`t`$-pole diagrams, however, the authors of Ref. have been forced to add a nucleon $`s`$-pole term to the Regge amplitude to ensure gauge invariance. In general, the full sum of $`s`$-channel resonances taken with their proper form factors should be dual to the full sum of reggeon exchanges (see, e.g. the recent review ), but actually the form factors for the transitions $`\gamma ^{}+NN^{}`$ are only known with large uncertainties, and thus one can stay on a phenomenological level and use only some separate terms from both sums.
The aim of this letter is to show that an alternative (microscopic) description of quasi-elastic pion knockout from the proton can be obtained on the basis of a constituent quark model (CQM). Then with the fixed parameters of model (which are a few in number) one can predict the absolute value and the $`t`$-dependence of longitudinal $`d\sigma _L/dt`$ and transverse $`d\sigma _T/dt`$ cross sections for the reaction of pion quasi-elastic knockout from nucleon by electrons in a wide interval of intermediate energies. We start from two assumptions, which are natural in the framework of the CQM:
(a) At energies $`W`$ 2.1 - 2.3 GeV and $`Q^2`$ 1 - 2 GeV<sup>2</sup>/c<sup>2</sup> characteristic of the JLab new measurements there are no resonance peaks in the cross section, and thus one need not to account for the details of the intermediate nucleon excitations. Only a combined effect of all the $`s`$\- and $`u`$-channel contributions should be evaluated. In our opinion, a quark approach for such evaluation would be more convenient than the traditional baryon-resonance approach which characterized by many coupling constants and form factors.
The $`t`$-pole contributions to the quasi-elastic pion knockout correspond to trivial quark diagrams shown in Fig. 1a. Our hypothesis is that in the kinematical region discussed the full sum of nucleon and excited-nucleon $`s`$\- and $`u`$-pole terms can be represented by the quark diagrams in Figs. 1b,c. Moreover, in the region of the pion $`t`$-pole, $`|t|m_\pi ^2`$, the sum of all the $`s`$\- and $`u`$-channel contributions should be a small correction to the t-pole contribution $`(tm_\pi ^2)^1.`$ Therefore, one can neglect the highest order corrections $`(m_q^2/Q^2)^n`$, e.g. generated by diagrams with one- or two-gluon exchange, in comparison to the dominant $`s`$\- and $`u`$-pole terms in Figs. 1b,c.
(b) The $`s`$\- and $`u`$-pole quark diagrams in Figs. 1b,c at large momentum transfer $`Q^2`$ imply using the strong and e.-m. quark form factors. In the CQM one can take the Fourier transform of the pion wave function for the strong $`\pi qq`$ vertex form factor (e.g. this can be shown in terms of the $`{}_{}{}^{3}P_{0}^{}`$ model ). Then, the quark amplitude for the $`s`$\- or $`u`$-channel transition $`\gamma ^{}+qq+\pi `$ becomes proportional to the product of two form factors, the strong one $`F_{\pi qq}`$ and the electromagnetic form factor of the constituent quark $`F_q(Q^2)`$<sup>1</sup><sup>1</sup>1It implies that the constituent quark is an extended object and has its own electromagnetic form factor, e.g. $`F_q(Q^2)=1/(1+Q^2/\mathrm{\Lambda }_q^2)`$ . The parameter $`\mathrm{\Lambda }_q`$ is believed to be set by the chiral symmetry scale $`\mathrm{\Lambda }_\chi 4\pi f_\pi `$ 1 GeV.. In the considered region of large $`Q^2`$ and $`W`$ this product is equivalent to the pion electromagnetic form factor $`F_\pi (Q^2)F_q(Q^2)F_{\pi qq}(Q^2)`$.
As this relationship has not been considered earlier on, here in Appendix we obtain it with a simple calculation in terms of two diagrams pictured in Fig. 2. The quark diagram for pion form factor $`F_\pi `$ is shown in Fig. 2b with the same notations of momentum variables as for the discussed $`s`$-pole diagram represented in Fig. 2a (with the scalar $`\overline{q}q`$ vertex $`v`$ of $`{}_{}{}^{3}P_{0}^{}`$ model). The solid line in both pictures shows the propagation of a large momentum $`𝐪𝐤`$ transfered from photon to pion by a highly excited quark. It is intuitively clear that at fixed $`q^2`$ the contribution of this deeply off-shell quark<sup>2</sup><sup>2</sup>2In the CQM a quark has a fixed mass $`m_qM_N/3`$ and all the non-trivial $`q^2`$-dependence of the quark propagator takes effectively into account through the $`q^2`$-dependent vertex form factors $`F_q`$ and $`F_{\pi qq}`$. into the both processes should be similar in value independently on details in diagrams in Fig.2a and Fig.2b. This assumption is formally confirmed by a simple calculation in terms of Gaussian wave functions (see Appendix). As a result, in the considered kinematical region of quasi-elastic pion knockout both the $`t`$\- and $`s(u)`$-contributions to the amplitude become proportional to a common form factor $`F_\pi (Q^2)`$ which can be factorized from the sum of $`t`$\- and $`s(u)`$-pole diagram contributions. Note that in a phenomenological approach such common form factor was introduced by hand to preserve gauge invariance of the calculated amplitude.
In this work we obtain an evaluation of a common effect of $`s`$\- and $`u`$-pole contributions to the cross section. The overall contribution of such terms is $`Q^2`$\- and $`W`$-dependent and vanishes with increasing $`W`$ and $`Q^2`$ in direct proportion to the factor $`(QM_N/W^2)F_\pi (Q^2)`$ (the final expression looks like a generalization of the Kroll-Ruderman contact term by introduction of a $`Q^2`$\- and $`W`$-dependent form factor into it). Our evaluation shows that for JLab data at lower $`Q^2`$ (0.6 and $`0.75`$ GeV<sup>2</sup>/c<sup>2</sup>) this contribution is important and cannot be neglected<sup>3</sup><sup>3</sup>3Unfortunately at small $`Q^2`$, comparable with the cutoff parameter $`\mathrm{\Lambda }`$ 0.7 GeV/c in $`\gamma ^{}NN^{}`$ vertices, one cannot exclude the direct contribution of intermediate baryon resonances to the cross section. Formally such small $`Q^2`$ are out of the region of the quasi-elastic kinematics., while at higher values of $`Q^2`$ and $`W`$ it becomes considerably smaller and practically invisible in the cross section at $`W`$ 3 GeV. Here we predict the $`t`$-dependence of the longitudinal and transverse cross section for a new JLab measurement at $`W=`$2.1 - 2.3 GeV.
## 2 The $`t`$\- and $`s(u)`$-channel contributions in terms of the CQM
The possibility of describing the high-energy pion electroproduction in terms of the pion $`t`$-channel mechanism has been discussed in the literature over many years . An important argument is that the contribution of the $`s`$-channel diagrams in Fig.1b is suppressed by a factor $`1/Q^4`$ in comparison to the contribution of the $`t`$-channel diagrams. In Refs. , the pion-exchange ($`t`$-channel diagram) was calculated using the light-front dynamics. In principle, this made it possible to extract the strong $`\pi NN`$ form factor $`F_{\pi NN}(t)`$ from comparison of a calculated longitudinal cross section $`d\sigma _L/dt`$ with the data at high momentum transfer $`Q^2`$ 2 -3 GeV<sup>2</sup>/c<sup>2</sup>. But in practice, the accessible data on high-quality Rosenbluth separation are limited by too small $`Q^2`$ and $`W`$, where the $`s(u)`$-contributions cannot be neglected.
In contrast to those studies, in Ref. the process $`p(e,e^{}\pi ^+)n`$ was considered in the laboratory frame. In this case one is able to single out, in a natural way, the kinematical region where the recoil momentum of the final nucleon $`P^\mu =\{P_0^{},𝐤\}`$ with $`P_0^{}M_N+𝐤^2/(2M_N)`$ is small and where only the momentum $`𝐤^{}=𝐪+𝐤`$ of the knocked out meson is large. In that (quasi-elastic) region at $`|t|0.10.2`$ GeV<sup>2</sup> 4-momentum $`k=PP^{}`$ transfered to the nucleon can be related to the 3-momentum $`𝐤`$, $`k^\mu \{𝐤^2/(2M_N),𝐤\}`$. Therefore, both the initial $`|N(P)>`$ and the final-state $`|N(P^{})>`$ nucleons are nonrelativistic ones. In the CQM each of them can be described by a nonrelativistic wave function
$`|N_{3q}(P)`$ $`=`$ $`\mathrm{\Phi }_N(\rho \rho \text{ }\rho \rho \text{ }_1,\rho \rho \text{ }\rho \rho \text{ }_2)|[2^3]_C[3]_{SI}S_z,I_z.`$ (1)
Here, $`|[2^3]_C[3]_{SI}S_z,I_z>`$ is the color$`(C)`$-spin$`(S)`$-isospin$`(I)`$ part and
$`\mathrm{\Phi }_N(\rho \rho \text{ }\rho \rho \text{ }_1,\rho \rho \text{ }\rho \rho \text{ }_2))=𝒩^1_Ne^{\rho \rho \text{ }\rho \rho \text{ }_1^2/4b^2\rho \rho \text{ }\rho \rho \text{ }_2^2/3b^2}`$ (2)
is the internal wave function normalized as $`𝑑\rho \rho \text{ }\rho \rho \text{ }_1𝑑\rho \rho \text{ }\rho \rho \text{ }_2|\mathrm{\Phi }_N(\rho \rho \text{ }\rho \rho \text{ }_1,\rho \rho \text{ }\rho \rho \text{ }_2)|^2=1`$, where $`\rho \rho \text{ }\rho \rho \text{ }_1`$ and $`\rho \rho \text{ }\rho \rho \text{ }_2`$ are the Jacobi coordinates, $`𝒩_N`$ is a normalization constant and $`b`$ is a “quark radius” of the nucleon.
The $`t`$-pole amplitude for the $`\gamma ^{}qq^{}\pi ^+`$ process on the $`i`$-th quark (Fig.1a) is
$`_{q(t)}^{(i)\mu }=ie\widehat{G}_{\pi qq}^{(i)}{\displaystyle \frac{F_\pi (Q^2)}{tm_\pi ^2}}(k+k^{})^\mu ,\widehat{G}_{\pi qq}^{(i)}=g_{\pi qq}\tau _{}^{(i)}\sigma \sigma \text{ }\sigma \sigma \text{ }^{(i)}𝐤,`$ (3)
where $`F_\pi (Q^2)`$ is the pion electromagnetic form factor and $`\widehat{G}_{\pi qq}^{(i)}`$ is the $`\pi qq`$ vertex operator for the $`i`$-th quark, which is related to the pion-nucleon form factor $`G_{\pi NN}(𝐤^2)`$ as
$`G_{\pi NN}(𝐤^2)\tau _{}\sigma \sigma \text{ }\sigma \sigma \text{ }𝐤N_{3q}(P^{})|{\displaystyle \underset{i=1}{\overset{3}{}}}\widehat{G}_{\pi qq}^{(i)}|N_{3q}(P)={\displaystyle \frac{5\tau _{}}{3}}g_{\pi qq}\sigma \sigma \text{ }\sigma \sigma \text{ }𝐤e^{𝐤^2b^2/6}.`$ (4)
In the following, for convenience, we proceed with the normalized pion-nucleon form factor $`F_{\pi NN}(𝐤^2)G_{\pi NN}(𝐤^2)/g_{\pi NN}`$ with $`g_{\pi NN}G_{\pi NN}(0)`$. We use a simple form of the operator $`\widehat{G}_{\pi qq}^{(i)}`$ neglecting the momentum dependence since the exchanged pion and the constituent quarks are approximately on their mass shells.
The quark-level amplitude $`_{q(t)}^\mu `$ gives rise to the $`t`$-pole matrix element for $`\pi ^+`$ electroproduction from the nucleon
$`_{N(t)}^\mu =N_{3q}(P^{})|{\displaystyle \underset{i=1}{\overset{3}{}}}_{q(t)}^{(i)\mu }|N_{3q}(P),`$ (5)
which is proportional to the strong form factor
$`F_{\pi NN}(t)={\displaystyle 𝑑\rho \rho \text{ }\rho \rho \text{ }_1𝑑\rho \rho \text{ }\rho \rho \text{ }_2e^{i\frac{2}{3}𝐤\rho \rho \text{ }\rho \rho \text{ }_2}|\mathrm{\Phi }_N(\rho \rho \text{ }\rho \rho \text{ }_1,\rho \rho \text{ }\rho \rho \text{ }_2)|^2}=e^{𝐤^2b^2/6},t𝐤^2,`$ (6)
defined above. At $`b=0.6`$ fm, which is a typical scale in the CQM (at small $`|t|0.2`$ GeV<sup>2</sup> this Gaussian is very close to the monopole form factor with $`\mathrm{\Lambda }_{\pi NN}0.7`$ GeV), the matrix element (5) gives a reasonable description of the $`t`$-dependence of the differential cross sections in the JLab experiment (see, e.g., Ref. ). However, since the $`s`$\- and $`u`$\- channel quark contributions (Fig. 1b,c), which are required for gauge invariance, have not been taken into account, the results of such simple quark evaluation (dashed lines in the upper panels of Fig. 3) are not so good in comparison with the Regge parametrization used in Refs. (dash-dotted lines in Fig. 3).
Here we shall evaluate these contributions, starting from a nonlocal (nl) variant of the pseudoscalar $`\pi qq`$ coupling $`g_{\pi qq}^{\mathrm{nl}}`$, which differs from the local coupling $`g_{\pi qq}`$ by the presence of the quark form factor: $`g_{\pi qq}^{\mathrm{nl}}(p^2)g_{\pi qq}F_{\pi qq}(p^2)`$. Hence, the vertex operator $`\widehat{G}_{\pi qq}^{(i)}`$ is modified accordingly. In this case the kinematics is the following. Two particles, pion and ingoing (outgoing) quark, are on mass shell: $`k_{}^{}{}_{}{}^{2}=m_\pi ^2`$, $`p_i^2=m_q^2`$ ($`p_{i}^{}{}_{}{}^{2}=m_q^2`$), where $`m_q=M_N/3`$ is the mass of the constituent quark. But the intermediate quark with a large momentum $`p_{i(s)}^{\prime \prime }=p_i+q`$ ($`p_{i(u)}^{\prime \prime }=p_i^{}q`$) obtained by absorption (emission) of a momentum transfer from the $`\gamma ^{}`$ should be considerably off its mass shell: $`p_{i}^{\prime \prime }{}_{}{}^{2}Q^2`$, (see solid lines in Figs.1b,c).
The contribution of the quark diagrams in Figs.1b,c to the matrix element for absorption of the longitudinal virtual photon ($`ϵ^{(0)\mu }=\frac{1}{Q}\{|𝐪|,q_0\widehat{𝐪}\}`$, $`Q=\sqrt{q^2}`$)
$`_{q(s+u)}^{(i)\mu }`$ $`=`$ $`i\sqrt{2}g_{\pi qq}F_{\pi qq}(Q^2)F_q(Q^2)`$ (7)
$`\times `$ $`\left\{e_u\gamma ^5{\displaystyle \frac{1}{\overline{)}p_i^{}+\mathit{}^{}m_q}}\gamma ^\mu +e_d\gamma ^\mu {\displaystyle \frac{1}{\overline{)}p_i\mathit{}^{}m_q}}\gamma ^5\right\}^{(i)}`$
was evaluated in the approximation $`W^2(Q^2M_N^2),M_N^2`$. Here $`e_u(e_d)`$ is the quark charge and $`F_q(Q^2)`$ is the quark electromagnetic form factor. In the considered region of a large momentum $`k^\mu =(k^0,𝐤^{})`$ transfered to the pion (with $`k^0|𝐤^{}||𝐪|`$) and a small momentum $`k^\mu =\{t/(2M_N),𝐤\}\{0,𝐤\}`$ transfered to the nucleon we can omit all contributions to the amplitude proportional to small parameters
$`{\displaystyle \frac{Q^2M_N^2}{W^2}},{\displaystyle \frac{|𝐤|}{2M_N}},{\displaystyle \frac{m_q^2}{W^2}},{\displaystyle \frac{m_\pi ^2}{M_N^2}},\mathrm{}1.`$ (8)
Then, for the 4-momenta of intermediate (off-shell) quarks in the diagrams in Fig. 1b,c one can approximately write
$`p_{i(s)}^{\prime \prime \mathrm{\hspace{0.17em}2}}2m_qk_0^{}\left(1+{\displaystyle \frac{|\varkappa \varkappa \text{ }\varkappa \varkappa \text{ }_2|}{m_q}}cos\theta \right)Q^2,`$ $`p_{i(u)}^{\prime \prime \mathrm{\hspace{0.17em}2}}p_{i(s)}^{\prime \prime \mathrm{\hspace{0.17em}2}}Q^2,`$ (9)
where $`\varkappa \varkappa \text{ }\varkappa \varkappa \text{ }_2`$ is a relative momentum conjugated to the Jacobi coordinate $`\rho \rho \text{ }\rho \rho \text{ }_2`$, and $`cos\theta =\widehat{𝐤}^{}\widehat{\varkappa \varkappa \text{ }\varkappa \varkappa \text{ }}_2\widehat{𝐪}\widehat{\varkappa \varkappa \text{ }\varkappa \varkappa \text{ }}_2`$. By making use of the approximate equality $`|𝐪|/q_01+Q^2/(2q_0^2)`$, which holds within the conditions (8), one obtains the approximation $`(\mathit{}^{}\gamma ^5\gamma ^\mu +\gamma ^5\gamma ^\mu \mathit{}^{})\epsilon _\mu ^{(0)}k_0^{}Q/q_0\mathrm{diag}\{\sigma \sigma \text{ }\sigma \sigma \text{ }\widehat{𝐪},\sigma \sigma \text{ }\sigma \sigma \text{ }\widehat{𝐪}\}`$ useful for reducing the r.h.s. of Eq. (7) (the approximation implies that $`\widehat{𝐤}^{}\widehat{𝐪}`$, which is true for the forward pion knockout).
Finally, in the lowest order of expansion in the small parameters $`M_N^2/W^2`$ and $`Q^2/W^2`$ the matrix element (7) reduces to an effective $`\gamma \pi qq`$ interaction<sup>4</sup><sup>4</sup>4At the photon point $`Q^2=`$0, $`ϵ_\mu ^{(\lambda )}=\{0,\widehat{𝐧}_{}\}`$ and in the low-energy limit $`k_0^{}m_\pi `$ the same quark calculation gives exactly rise to the Kroll-Ruderman term for the threshold pion photoproduction.
$`_{q(s+u)}^{(i)\mu }ϵ_\mu ^{(\lambda =0)}`$ $`=`$ $`i{\displaystyle \frac{eg_{\pi qq}}{2m_q}}{\displaystyle \frac{M_NQ}{W^2}}{\displaystyle \frac{\tau _{}^{(i)}(\sigma \sigma \text{ }\sigma \sigma \text{ }^{(i)}\widehat{𝐪})}{(1+\frac{|\varkappa \varkappa \text{ }\varkappa \varkappa \text{ }_2|}{m_q}cos\theta )}}F_q(Q^2)F_{\pi qq}(Q^2),`$ (10)
in which one can use the relation
$`F_q(Q^2)F_{\pi qq}(Q^2)F_\pi (Q^2),`$ (11)
discussed in the Section 1 (see also Appendix). Note that the $`\varkappa \varkappa \text{ }\varkappa \varkappa \text{ }_2`$-dependent denominator in Eq. (10), being integrated with Gaussian functions (1), introduces only a small correction to the integral, and thus we shall omit it in the final expression (see next section) for the nucleon matrix element (but it has not been omitted in the numerical calculation).
Eq. (11) is very useful to check gauge invariance of the sum of $`t`$\- and $`s(u)`$-pole amplitudes in the final expression, where these amplitudes destructively interfere (see Sect. 3). It should be noted that gauge invariance can also be restored in the general case, if one uses different vertex form factors (see detailed discussion in Refs. ).
## 3 Longitudinal and transeverse cross section
We calculate the longitudinal part of differential cross section
$`{\displaystyle \frac{d\sigma _L}{dt}}={\displaystyle \frac{M_N^2\overline{|J_\pi ^\mu \epsilon _\mu ^{(\lambda =0)}|^2}}{4\pi (W^2M_N^2)\sqrt{(W^2Q^2M_N^2)^2+4W^2Q^2}}}`$ (12)
for quasi-elastic pion knockout taking into account both the $`t`$-pole diagram in Fig.1a and the $`s`$\- and $`u`$-pole contribution (10),
$`\overline{|J_\pi ^\mu \epsilon _\mu ^{(\lambda =0)}|^2}={\displaystyle \frac{1}{2}}{\displaystyle \underset{spin}{}}|N_{3q}(P^{})|{\displaystyle \underset{j=1}{\overset{3}{}}}[_{q(t)}^{(j)\mu }+_{q(s+u)}^{(j)\mu }]\epsilon _\mu ^{(\lambda =0)}|N_{3q}(P)|^2,`$ (13)
and compare the results to the JLab data and to the Regge model predictions (see Fig. 3).
The CQM calculation of the right-hand side of Eq. (13) leads to a simple matrix element for a nonrelativistic nucleon wave function $`\mathrm{\Phi }_N(\rho \rho \text{ }\rho \rho \text{ }_1,\rho \rho \text{ }\rho \rho \text{ }_2)`$. Using Eqs. (5) and (10) we obtain in the laboratory frame the following matrix element of the longitudinal hadron current for the transition $`\gamma ^{}+pn+\pi ^+`$:
$`J_\pi ^\mu \epsilon _\mu ^{(0)}=i\tau _{}{\displaystyle \frac{eg_{\pi NN}}{2M_N}}F_\pi (Q^2)F_{\pi NN}(t)\left[{\displaystyle \frac{2\epsilon ^{(0)}k^{}\sigma \sigma \text{ }\sigma \sigma \text{ }𝐤}{tm_\pi ^2}}{\displaystyle \frac{M_NQ}{W^2}}\sigma \sigma \text{ }\sigma \sigma \text{ }\widehat{𝐪}\right].`$ (14)
Here the contribution of the $`b_1`$-meson pole (the P-wave excitation of the pion in the CQM) may be also taken into account. Then, the third term, $`+(g_{b_1NN}/g_{\pi NN})(g_{b_1\pi \gamma }k_0Q/(2M_Nm_{b_1}))(𝐤\widehat{𝐪}\sigma \sigma \text{ }\sigma \sigma \text{ }𝐤/(tm_{b_1}^2))`$, should be inserted into the square brackets in the r.h.s. of Eq. (14). However, the $`b_1`$-pole contribution is too small when compared to the $`\pi `$-pole contribution. The $`\rho `$-meson pole is also not included as in the CQM it does not contribute to the longitudinal part of the cross section (in the CQM the quark spin-flip amplitude $`\rho ^++\gamma ^{}(M1)\pi ^+`$ only contributes to the transverse part of the cross section where it is of prime importance ).
The quark model calculation leading to the r.h.s. of Eq. (14) has the advantage that the parameters of the vertices $`g_{\pi NN}`$, $`g_{b_1NN}`$, $`g_{b_1\pi \gamma }`$ and their form factors $`F_{\pi NN}(t)`$, $`F_{b_1NN}(t)`$, $`F_{b_1}(Q^2)`$ are not free, but related to each other by the following relationships: (i) The integral part of Eq. (13) defines a nonrelativistic vertex form factor (6) $`F_{mNN}(t)=F_{\pi NN}(t)=F_{b_1NN}(t)`$ common to all the terms in Eq. (14). (ii) The relative phases of all the amplitudes in the r.h.s of Eq. (14) are fixed by the results of quark model calculations. Therefore, the negative sign of the interference term between $`t`$\- and $`s`$-pole contributions (the destructive interference) is unambiguously determined: the spin average of the product of the first term in the squared bracket of Eq. (14) with the second one has a negative value: $`1/2_{S_zS_z^{}}<S_z^{}|\sigma \sigma \text{ }\sigma \sigma \text{ }\widehat{𝐤}|S_z><S_z|\sigma \sigma \text{ }\sigma \sigma \text{ }\widehat{𝐪}|S_z^{}>=\mathrm{cos}(\widehat{\mathrm{𝐤𝐪}}).`$ Recall that $`\mathrm{cos}(\widehat{\mathrm{𝐤𝐪}})`$-1 in the quasielastic pion knockout kinematics, where $`𝐤`$ is the momentum of the nucleon-spectator in the lab frame.
The $`\rho `$-pole contribution to the transverse part of differential cross section $`d\sigma _T/dt`$ is proportional to the value $`\overline{|J_\rho ^\mu \epsilon _\mu ^{(\lambda =1)}|^2}`$, where
$`J_\rho ^\mu \epsilon _\mu ^{(\lambda )}=i\tau _{}eg_{\rho \pi \gamma }{\displaystyle \frac{M_\rho }{M_\pi }}F_{\rho \pi \gamma }(Q^2){\displaystyle \frac{1+\varkappa _\rho }{2M_N}}g_{\rho NN}F_{\rho NN}(t)|𝐪|{\displaystyle \frac{i[\sigma \sigma \text{ }\sigma \sigma \text{ }\times 𝐤]ϵϵ\text{ }ϵϵ\text{ }^{(\lambda )}}{tM_\rho ^2}}`$ (15)
and $`\epsilon ^{(\lambda )\mu }=\{0,ϵϵ\text{ }ϵϵ\text{ }^{(\lambda )}\}`$ for $`\lambda =\pm `$ 1. For coupling constants and form factors in Eq. (15) we use the results of calculations in terms of $`{}_{}{}^{3}P_{0}^{}`$ model from Ref. : $`g_{\rho NN}=\frac{1}{5}g_{\pi NN}\sqrt{\frac{M_\rho }{M_\pi }}`$, $`1+\varkappa _\rho =`$ 5, $`F_{\rho NN}(t)=F_{\pi NN}(t)`$. For $`g_{\rho \pi \gamma }`$ coupling constant we take the value $`g_{\rho \pi \gamma }=0.103`$ which is fixed by the the experimental value of the decay width $`\mathrm{\Gamma }_{\rho \pi \gamma }=`$ 67 KeV. The momentum dependence of $`F_{\rho \pi \gamma }(Q^2)`$ will be discussed in the next section.
## 4 Results and outlook
In our calculation we use a standard (monopole-like) representation for the transition form factors $`F_\pi (Q^2)`$ and $`F_{\rho \pi \gamma }(Q^2)`$:
$`F_\pi (Q^2)={\displaystyle \frac{1}{1+Q^2/\mathrm{\Lambda }_\pi ^2}},F_{\rho \pi \gamma }(Q^2)={\displaystyle \frac{1}{1+Q^2/\mathrm{\Lambda }_{\rho \pi }^2}}.`$ (16)
For the pion charge form factor we use $`\mathrm{\Lambda }_\pi ^2=`$0.54 GeV<sup>2</sup>/c<sup>2</sup> which is close to the recent theoretical evaluation and correlates well with the recent JLab data . In the case of the $`F_{\rho \pi \gamma }(Q^2)`$ form factor we vary $`\mathrm{\Lambda }_{\rho \pi }^2`$ from $`\mathrm{\Lambda }_\pi ^2=0.54`$ GeV<sup>2</sup>/c<sup>2</sup> to 0.7 GeV<sup>2</sup>/c<sup>2</sup>. We think that an accurate analysis of the transverse part of the differential cross section can shed light on the value of $`\mathrm{\Lambda }_{\rho \pi }`$. From Fig. 3 one can see, that the value $`\mathrm{\Lambda }_{\rho \pi }^2=0.7`$ GeV<sup>2</sup>/c<sup>2</sup> is more appropriate in description of the transverse cross section. Apart from this value we vary only one free parameter in the standard representation of the strong $`\pi NN`$ form factor $`F_{\pi NN}(Q^2)=\mathrm{\Lambda }_{\pi NN}^2/(\mathrm{\Lambda }_{\pi NN}^2+Q^2)`$ where the range parameter $`\mathrm{\Lambda }_{\pi NN}`$ 0.6 - 0.7 GeV/c corresponds to the reasonable value $`b`$ 0.5 - 0.6 fm for the radius of the three-quark configuration $`s^3`$ used in the CQM nucleon wave function (1). At realistic values of $`\mathrm{\Lambda }_{\pi NN}=0.7`$ GeV/c and $`g_{\pi NN}=13.5`$ our results (the solid lines in Fig.3) are in a better agreement with the data than the simplified model (dashed lines in two upper panels) taking into account only the pion $`t`$-pole not satisfying gauge invariance.
In our calculation we model the contribution of baryon resonances in the $`s`$\- and $`u`$-channel by a nonlocal extension of the Kroll-Rudermann contact term. This contribution is negative, which is essential to improve the description of the measured cross section. At intermediate values of $`Q^2`$ 1 GeV<sup>2</sup>/c<sup>2</sup>, which correspond to a quasi-elastic mechanism of pion knockout, the calculated cross section is close to the experimental data. At smaller $`Q^2`$ 0.7 GeV<sup>2</sup>/c<sup>2</sup> our results are close to the cross section calculated in Ref (dash-dotted lines) on the basis of a model which takes into account both the Reggeon-pole exchange and the nucleon-pole contribution.
For a new JLab measurement of $`d\sigma _L/dt`$ and $`d\sigma _T/dt`$ at higher invariant mass (the data analysis is presently underway) in Fig. 4 we give our prediction for the $`t`$-dependence of the cross sections at $`W=`$ 2.1 - 2.3 GeV and $`Q^2`$ centered at 1.6 and 2.45 GeV<sup>2</sup>/c<sup>2</sup>. Here we use fixed (above defined) parameters of the model. Because of the large value of $`W`$ the contribution of the effective contact term (10), which is proportional to the factor $`QM_N/W^2`$, becomes too small and not shown in Fig. 4.
Our calculations show that the longitudinal cross section practically does not depend on the contribution of the $`b_1`$-meson pole, and thus the data cannot constrain the $`b_1NN`$ and $`b_1\pi \gamma `$ vertices. In contrast, data on the transverse cross section should be critically dependent on the $`\rho `$-pole contribution and on the $`\rho \pi \gamma ^{}`$ spin-flip amplitude. This was first shown in Ref. and then supported in Ref. . Our results (Figs. 3 and 4) confirm this statement and show that the data on the $`d\sigma _T/dt`$ at high $`Q^2`$ and $`W`$ can be used for the direct measurement of the $`F_{\rho \pi \gamma }(Q^2)`$ form factor. On the other hand, a Regge description of the transverse cross section is not as good as for the longitudinal case.
Acknowledgements
The authors thank the Fpi2 Collaboration \[Experiment E01-004 at JLab\] for the interest to our work and informative discussion. We thank Vladimir Neudatchin, Nikolai Yudin, Garth Huber and Tanja Horn for fruitful discussions and suggestions. This work was supported by the DFG under contracts FA67/25-3 and GRK683 and the DFG grant 436 RUS 113/790/. This research is also part of the EU Integrated Infrastructure Initiative Hadronphysics project under contract number RII3-CT-2004-506078, President grant of Russia ”Scientific Schools” No. 1743.2003 and grants of Russia RFBR No. 05-02-04000, 05-02-17394 and 03-02-17394.
## Appendix A Appendix
The contribution of the diagram in Fig. 2a is proportional to the overlap integral
$`F_{diag}(𝐪,𝐤)`$ $`=`$ $`F_q(𝐪^2){\displaystyle \frac{1}{𝒩}}{\displaystyle \frac{d\varkappa \varkappa \text{ }\varkappa \varkappa \text{ }_2}{(2\pi )^3}\mathrm{\Phi }_N(\varkappa \varkappa \text{ }\varkappa \varkappa \text{ }_2)\mathrm{\Phi }_N\left(\varkappa \varkappa \text{ }\varkappa \varkappa \text{ }_2+\frac{2}{3}𝐤\right)\mathrm{\Phi }_\pi \left(\varkappa \varkappa \text{ }\varkappa \varkappa \text{ }_2+\frac{𝐪𝐤}{2}\right)}`$ (17)
$`=`$ $`e^{𝐤^2b^2/6}F_q(𝐪^2)\mathrm{exp}\left[{\displaystyle \frac{(𝐪\frac{𝐤}{3})^2b_\pi ^2}{4(1+2x_\pi ^2/3)}}\right],`$
where $`\mathrm{\Phi }_N`$ is a wave function of the 3-rd quark in the nucleon and $`\mathrm{\Phi }_\pi `$ is a pion wave function in the momentum representation (here it plays the role of a strong $`\pi qq`$ form factor). In this simple calculation we use $`\mathrm{\Phi }_N(\varkappa \varkappa \text{ }\varkappa \varkappa \text{ }_2)e^{3\varkappa \varkappa \text{ }\varkappa \varkappa \text{ }_2^2b^2/4}`$ and $`\mathrm{\Phi }_\pi (\varkappa \varkappa \text{ }\varkappa \varkappa \text{ })e^{\varkappa \varkappa \text{ }\varkappa \varkappa \text{ }^2b_\pi ^2}`$, where $`\varkappa \varkappa \text{ }\varkappa \varkappa \text{ }_2=(𝐩_1+𝐩_22𝐩_3)/3`$. Here $`b`$ and $`b_\pi `$ are the quark radii in nucleon and pion correspondingly with $`x_\pi =b_\pi /b`$. In the CQM one can obtain a very similar expression for the diagram in Fig. 2b:
$`F_\pi (𝐪^2)`$ $`=`$ $`F_q(𝐪^2){\displaystyle \frac{d𝐩}{(2\pi )^3}\mathrm{\Phi }_\pi \left(𝐩\frac{𝐤}{2}\right)\mathrm{\Phi }_\pi \left(𝐩+\frac{𝐪𝐤}{2}\right)}=F_q(𝐪^2)e^{𝐪^2b_\pi ^2/8}.`$ (18)
One can see that the strong $`\pi NN`$ form factor (in the CQM it is $`F_{\pi NN}(𝐤^2)=e^{𝐤^2b^2/6}`$) appears as a multiplier in the second line of Eq. (17) and the reminder of this expression coincides with the pion form factor (18) in the limit $`𝐪𝐤`$ at a specific value $`x_\pi =\sqrt{3/2}`$, which is not very different from the realistic value $`x_\pi `$ 1 characteristic of the CQM. Really the factor $`\sqrt{3/2}`$ appears because of a specific procedure of center-of-mass motion eliminating, which is model dependent. We shell consider a small difference between $`F_\pi `$ and $`F_qF_{\pi qq}`$ as a small model-dependent correction which can be neglected in line of the assumption (a) given in Section 1.
It could also happen that the $`\gamma ^{}`$ is absorbed by one quark while the pion is emitted from a different quark in the nucleon. Then a large momentum $`𝐪`$ should be firstly transfered to the $`N`$-$`3q`$ vertex. In the CQM this amplitude can be evaluated as well. One can insert an intermediate nucleon (baryon) state between two points in the diagram where a large momentum $`𝐪`$ or $`|𝐤^{}||𝐪|`$ is absorbed or emitted. It is essential that the nucleon propagator and form factors (both electromagnetic and strong) depend on the large momentum. This leads to the following contribution to the amplitude of pion quasi-elasic knockout, which scales as:
$`F_{nondiag}(𝐪,𝐤^{})F_q(𝐪^2){\displaystyle \frac{d\varkappa \varkappa \text{ }\varkappa \varkappa \text{ }_2}{(2\pi )^3}\mathrm{\Phi }_N(\varkappa \varkappa \text{ }\varkappa \varkappa \text{ }_2)\mathrm{\Phi }_N\left(\varkappa \varkappa \text{ }\varkappa \varkappa \text{ }_2+\frac{2}{3}𝐪\right)\frac{M_N\sqrt{Q^2}}{Q^2M_N^2}}`$
$`\times {\displaystyle }{\displaystyle \frac{d\varkappa \varkappa \text{ }\varkappa \varkappa \text{ }_2^{}}{(2\pi )^3}}\mathrm{\Phi }_N(\varkappa \varkappa \text{ }\varkappa \varkappa \text{ }_2^{})\mathrm{\Phi }_N(\varkappa \varkappa \text{ }\varkappa \varkappa \text{ }_2^{}{\displaystyle \frac{2}{3}}𝐤^{})G_E(Q^2){\displaystyle \frac{M_N}{\sqrt{Q^2}}}F_{\pi NN}(Q^2),`$ (19)
where $`G_E`$ is a nucleon electric form factor. At large $`Q^2`$ 1 - 2 GeV<sup>2</sup>/c<sup>2</sup> characteristic of the new F$`\pi `$2 measurement the ratio $`|F_{nondiag}(𝐪,𝐤^{})/F_{diag}(𝐪,𝐤)|^2`$ is too small. Hence one can consider the amplitude (19) as a second order correction to the $`\pi `$-pole amplitude, while the amplitude (17) is the first order correction.
LIST OF FIGURES
Fig.1: Quark diagrams for the $`t`$-, $`s`$\- and $`u`$-channel mechanism of pion knockout. Thick lines indicate large $`Q^2`$ transition.
Fig.2: Quark diagrams: (a) for the $`s`$-pole mechanism of pion emission (in terms of $`{}_{}{}^{3}P_{0}^{}`$ model), (b) for pion form factor.
Fig.3: Longitudinal and transverse cross sections for the $`p(e,e^{}\pi ^+)n`$ process. The $`F\pi 1`$ data for $`W`$ 0.95 GeV. 1) The upper row of panels: $`Q^2=`$ 0.6 GeV<sup>2</sup>/c<sup>2</sup> (left panel) and 0.75 GeV<sup>2</sup>/c<sup>2</sup> (right panel). Dashed lines: the pion $`t`$-pole only. Solid lines: the total sum of the pion $`t`$-pole and the contact ($`s+u`$)-term with common e.-m. and strong form factors. Dotted lines: the same sum, but without form factors in the contact term. Dash-dotted lines: the Regge model prediction. 2) The lower two rows of panels: $`Q^2=`$ 1.0 GeV<sup>2</sup>/c<sup>2</sup> (upper panels) and 1.6 GeV<sup>2</sup>/c<sup>2</sup> (lower panels). Left panels ($`d\sigma _L/dt`$): the same notations as in the first two panels. Right panels ($`d\sigma _T/dt`$): the $`\pi `$-pole contribution only (dotted); the sum of $`\pi `$\- and $`\rho `$-pole + ($`s+u`$) contributions with a monopole $`\rho \pi \gamma `$ form factor, $`\mathrm{\Lambda }_{\rho \pi }^2=`$ 0.54 GeV<sup>2</sup>/c<sup>2</sup> (solid) and 0.7 GeV<sup>2</sup>/c<sup>2</sup> (dashed).
Fig.4: Predicted cross sections for new JLab measurements at a higher value of $`W=`$ 2.1 - 2.3 GeV and $`Q^2`$ centered at 1.6 GeV<sup>2</sup>/c<sup>2</sup> (upper curves) and 2.45 GeV<sup>2</sup>/c<sup>2</sup> (lower curves). Left panel ($`d\sigma _L/dt`$): the total sum of the pion $`t`$-pole and the contact ($`s+u`$)-term (solid); the Regge model prediction (dash-dotted). Right panel ($`d\sigma _T/dt`$): the $`\pi `$-pole contribution only (dotted); the sum of $`\pi `$\- and $`\rho `$-pole + ($`s+u`$) contributions with a monopole $`\rho \pi \gamma `$ form factor, $`\mathrm{\Lambda }_{\rho \pi }^2=`$ 0.54 GeV<sup>2</sup>/c<sup>2</sup> (solid) and 0.7 GeV<sup>2</sup>/c<sup>2</sup> (dashed).
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# Extraction of scattering lengths from final-state interactions
## I Introduction
The scattering length provides not only an important measure for the strength of the interaction in a specific hadronic two-body system Joach but often allows to draw further more general and thus even more interesting conclusions. For example in the case of the proton-proton and neutron-neutron systems the corresponding scattering lengths in the $`{}_{}{}^{1}S_{0}^{}`$ partial wave provide a very sensitive test of charge symmetry in the strong interaction Miller . SU(3) symmetry can be tested by comparing the $`{}_{}{}^{1}S_{0}^{}`$ scattering lengths in the neutron-proton and $`\mathrm{\Sigma }^+p`$ systems, which should fulfil the relation $`a_{np}=a_{\mathrm{\Sigma }^+p}`$ in case SU(3) symmetry holds rigorously Dover . In the chiral limit the $`\pi N`$ S-wave scattering lengths vanish and therefore any deviation from that value is a direct measure for how strongly this symmetry is broken. Here especially the isoscalar component is of high interest due to its close link to the sigma term of the nucleon xx . On a more phenomenological level the $`\eta N`$ scattering length is interesting because the magnitude of its real part is directly linked with speculations about the existence of $`\eta `$-mesonic hadronic bound states such as $`\eta ^3He`$ Wycech ; Belyaev1 ; Belyaev2 ; Rakityansky ; Fix1 ; Niskanen ; Niskanen1 or $`\eta ^4He`$ Haider2 .
Unfortunately, a direct determination of the scattering length is only feasible in a few cases. It can be done with scattering experiments sufficiently close to the reaction threshold so that the effective range expansion can be utilized for extracting the scattering length. But in practice such experiments are only possible for charged and also (quasi) stable particles, as it is the case with, e.g., $`pp`$, $`\pi ^+p`$ or $`K^+p`$ scattering. One can also extract the scattering length from a study of the hadronic level shifts of atoms rus like $`\pi ^{}p`$, $`\pi ^{}d`$ or also $`\overline{p}p`$ Gotta . For the majority of the hadronic two-body systems information about the scattering length is only accessible via an investigation of the final-state interaction of systems which have at least three particles in the final state.
While in the former type of experiments the accuracy of the scattering length is directly connected with the precision of the data more detailed and often sophisticated considerations are necessary in order to estimate additional uncertainies that arise when the scattering length is extracted from final-state effects Gibbs . However, in some cases such an estimation is facilitated by the fact that the reaction mechanism is known. For example, the reaction $`nd(nn)p`$, one of the prime sources of the $`nn`$ scattering length, can be analysed by means of rigorous Faddeev calculations Trotter ; Huhn ; Deng . Reactions involving the pion such as $`\pi d`$, $`\pi ^{}d\gamma nn`$ or $`\gamma d\pi ^+nn`$, which can be used to extract the $`\pi N`$ and $`nn`$ scattering lengths, respectively, can be tackled by chiral perturbation theory in a well controlled way in the relevant near-threshold regime xxx ; Lensky .
In a recent publication Gasparyan2003 we argued that also large-momentum transfer reactions such as $`ppK^+p\mathrm{\Lambda }`$ C11\_1 ; Bilger ; jan or $`\gamma dK^+n\mathrm{\Lambda }`$ Renard ; Kerbikov ; Mecking ; Adel ; Li ; Yama are excellent candidates for extracting information about the scattering lengths. In reactions with large momentum transfer the production process is necessarily of short-ranged nature. As a consequence the results are basically insensitive to details of the production mechanism and therefore a reliable and general error estimation can be given. Indeed, in Gasparyan2003 a formalism based on dispersion theory was presented that relates spectra from large-momentum transfer reactions, such as $`ppK^+p\mathrm{\Lambda }`$ or $`\gamma dK^+n\mathrm{\Lambda }`$, directly to the scattering length of the interaction of the final state particles. The theoretical error of the method was estimated to be 0.3 fm or even less, which is comparable to the error quoted in the context of the determination of $`a_{nn}`$ Howell , say. This estimate was confirmed by comparing results obtained with the proposed formalism to those of microscopic model calculations for the specific reaction $`ppK^+\mathrm{\Lambda }p`$. The arguments but also the formalism of Ref. Gasparyan2003 are, of course, valid for any production or decay process which is of short-ranged nature, i.e. also for investigation of hadronic two-particle subsystems resulting from the decay of the $`J/\mathrm{\Psi }`$ or $`B`$ mesons Psi ; Belle .
In the present paper we want to investigate further aspects of extracting scattering lengths from final-state interactions which were not addressed in our earlier work. One of those topics is the presence of the Coulomb interaction. In many interesting hadronic two-particle systems both particles carry charges, like in the already mentioned $`\mathrm{\Sigma }^+p`$ channel whose scattering length could be extracted from the reaction $`ppK^0\mathrm{\Sigma }^+p`$. Then the production amplitude acquires additional singularities, due to the long-range nature of the Coulomb forces, and the formalism developed in Ref. Gasparyan2003 is no longer directly applicable. We will derive the modifications that are necessary in order to adapt the dispersion-relation method to the situation when the Coulomb force is present in the final-state interaction. We also demonstrate in a toy model calculation how one has to proceed in a concrete application to data.
In addition we present a more detailed examination of the accuracy of the method proposed in Ref. Gasparyan2003 . A test based on one specific model calculation, namely for the reaction $`ppK^+\mathrm{\Lambda }p`$, has been already performed in that paper. However, here we want to put this investigation on a broader basis by considering final-state interactions of varying strengths, corresponding to a much larger range of values of the scattering length. In addition we take a look at the effective range $`r_e`$ as well which can be also extracted by the proposed dispersion-integral method. For the effective range a sensible error estimation is not possible, as was already pointed out in Ref. Gasparyan2003 , but it is still interesting to examine in concrete applications whether meaningful results could be achieved. Finally, and equally important, we want to compare the present method with the performance of other, approximative treatments of the final-state interaction that are commonly used in the literature to extract information on the scattering length and also the effective range. This concerns in particular the Jost-function approach book based on the effective-range approximation and an even simpler approach that relies simply on utilizing the effective range approximation itself migdal . Thereby, we will show that the latter methods lead to (partly drastic) systematic deviations from the true values and therefore one has to be rather cautious in the interpretation of results achieved with those methods.
The paper is structured in the following way: In the subsequent section we give a short review of the dispersion integral method for extracting the scattering length from final-state interactions. In section 3 results of an examination of the accuracy of this method are presented. Thereby, we consider various (singlet and triplet) $`S`$ wave $`YN`$ and $`NN`$ interactions and compare the scattering lengths and effective range extracted with the dispersion-integral method from appropriately generated final-state effects with the ones predicted by the models. We also apply two approximative methods for treating final-state effects, namely the Jost-function approach based on the effective range approximation (Jost-ERA) as well as the effective range approximation itself, and compare their performance with the one of our method. In section 4 we generalize the dispersion-integral method to the case where a repulsive Coulomb interaction is present in the final state. Test calculations for a final-state interaction with Coulomb are then presented in sect. 5 and it is discussed in detail how one has to procede in a practical application. The paper ends with a short summary.
## II Formalism
Our method, which goes back to an idea of Geshkenbein Geshkenbein1969 ; Geshkenbein1998 , is based on using the dispersion relation technique. Consider the production amplitude $`A_S`$ of a $`23`$ reaction. To be concrete we discuss $`ppK^+p\mathrm{\Lambda }`$, or $`\gamma dK^+n\mathrm{\Lambda }`$, with the $`\mathrm{\Lambda }N`$ system being in an $`L=0`$ partial wave and a specific spin state $`S`$ ($`{}_{}{}^{1}S_{0}^{}`$ or $`{}_{}{}^{3}S_{1}^{}`$). This amplitude depends on the total energy squared $`s=(p_1+p_2)^2`$, the invariant mass squared of the outgoing $`\mathrm{\Lambda }N`$ system $`m^2=(p_N+p_\mathrm{\Lambda })^2`$ and the momentum transfer $`t=(p_1p_{K^+})^2`$, where $`p_1`$, $`p_2`$, $`p_N`$, $`p_\mathrm{\Lambda }`$, and $`p_{K^+}`$ are the 4-momenta of the two initial particles, final nucleon, lambda, and kaon, respectively. Then one can write down a dispersion relation for this amplitude with respect to $`m^2`$ at fixed $`s`$ and $`t`$
$`A_S(s,t,m^2)={\displaystyle \frac{1}{\pi }}{\displaystyle _{\mathrm{}}^{\stackrel{~}{m}^2}}{\displaystyle \frac{D_S(s,t,m^{}{}_{}{}^{2})}{m^{}{}_{}{}^{2}m^2}}dm^{}{}_{}{}^{2}+{\displaystyle \frac{1}{\pi }}{\displaystyle _{m_0^2}^{\mathrm{}}}{\displaystyle \frac{D_S(s,t,m^{}{}_{}{}^{2})}{m^{}{}_{}{}^{2}m^2}}dm^{}{}_{}{}^{2},`$ (1)
where $`\stackrel{~}{m}^2`$ is the upper boundary of the lefthand cut, $`m_0^2=(m_N+m_\mathrm{\Lambda })^2`$, and
$`D_S(s,t,m^2)={\displaystyle \frac{1}{2i}}(A_S(s,t,m^2+i0)A_S(s,t,m^2i0))`$ (2)
is the discontinuity of the amplitude along the cuts. We neglect here the contributions from possible kaon-baryon interactions. In case they are not small, they still can be considered as constant (weakly mass dependent) if one chooses the kinematics such that the excess energy of the reaction is significantly larger than the typical range of the $`\mathrm{\Lambda }N`$ interaction, cf. the discussion in Ref. Gasparyan2003 . The index $`S`$ denoting the spin state will be suppressed in the following to simplify the notation.
For a purely elastic $`\mathrm{\Lambda }N`$ system, the discontinuity along the righthand cut would be given by
$`D(s,t,m^2)=A(s,t,m^2)e^{i\delta }\mathrm{sin}\delta ,`$ (3)
where $`\delta `$ is the $`\mathrm{\Lambda }N`$ ($`{}_{}{}^{1}S_{0}^{}`$ or $`{}_{}{}^{3}S_{1}^{}`$) scattering phase shift. Then the solution of Eq. (1) in the physical region reads (see Refs. Muskhelishvili1953 ; Omnes1958 ; Frazer1959 )
$`A(s,t,m^2)=\mathrm{exp}[{\displaystyle \frac{1}{\pi }}{\displaystyle _{m_0^2}^{\mathrm{}}}{\displaystyle \frac{\delta (m^{}{}_{}{}^{2})}{m^{}{}_{}{}^{2}m^2i0}}dm^{}{}_{}{}^{2}]\mathrm{\Phi }(s,t,m^2),`$ (4)
where $`\mathrm{\Phi }(s,t,m^2)`$ contains only lefthand singularities and therefore is a slowly varying function of $`m^2`$. In order to ensure this requirement to be fulfilled it is important that the momentum transfer $`t`$ is large. We assume also that there is no bound state in the $`\mathrm{\Lambda }N`$ system.
Consider now a realistic situation where inelastic channels are present – as it is the case with $`\mathrm{\Lambda }N`$ due to the coupling to the $`\mathrm{\Sigma }N`$ channel, say. Then one can write down a formula similar to Eq. (1), but with the integration performed over a finite range of masses Gasparyan2003 :
$`A(m^2)=\mathrm{exp}[{\displaystyle \frac{1}{\pi }}{\displaystyle _{m_0^2}^{m_{max}^2}}{\displaystyle \frac{\delta (m^{}{}_{}{}^{2})}{m^{}{}_{}{}^{2}m^2i0}}dm^{}{}_{}{}^{2}]\stackrel{~}{\mathrm{\Phi }}(m^2),`$ (5)
where $`\stackrel{~}{\mathrm{\Phi }}(m^2)`$ is again a slowly varying function of $`m^2`$ given the phase shift $`\delta `$ is sufficiently small in the vicinity of $`m_{max}`$ Gasparyan2003 . The upper limit $`m_{max}`$ has to be chosen in such a way that the corresponding relative momentum of the $`\mathrm{\Lambda }N`$ system, $`p_{max}`$, is of the order of the typical scale of the $`\mathrm{\Lambda }N`$ interaction, i.e. in the order of $`1/a`$ or $`1/r`$. The equation (5) can be solved with respect to the $`\delta `$ Gasparyan2003 :
$`{\displaystyle \frac{\delta (m^2)}{\sqrt{m^2m_0^2}}}=`$
$`{\displaystyle \frac{1}{2\pi }}𝐏{\displaystyle _{m_0^2}^{m_{max}^2}}{\displaystyle \frac{\mathrm{log}|A(m^{}{}_{}{}^{2})/\stackrel{~}{\mathrm{\Phi }}(m_{max}^2,m^{}{}_{}{}^{2})|^2}{\sqrt{m^{}{}_{}{}^{2}m_0^2}(m^{}{}_{}{}^{2}m^2)}}\sqrt{{\displaystyle \frac{m_{max}^2m^2}{m_{max}^2m^{}^2}}}dm^{}{}_{}{}^{2}.`$ (6)
If one neglects the mass dependence of $`\stackrel{~}{\mathrm{\Phi }}(m^2)`$ and uses the relation between the partial cross section $`\sigma _S`$ and the amplitude
$$\frac{d^2\sigma _S}{dm^{}{}_{}{}^{2}dt}p^{}|A_S(s,t,m^{}{}_{}{}^{2})|^2,$$
then one obtains the expression for the scattering length in terms of observables
$`a_S`$ $`=`$ $`\underset{m^2m_0^2}{lim}{\displaystyle \frac{1}{2\pi }}\left({\displaystyle \frac{m_\mathrm{\Lambda }+m_N}{\sqrt{m_\mathrm{\Lambda }m_N}}}\right)𝐏{\displaystyle _{m_0^2}^{m_{max}^2}}𝑑m^{}{}_{}{}^{2}\sqrt{{\displaystyle \frac{m_{max}^2m^2}{m_{max}^2m^{}^2}}}`$ (7)
$`\times {\displaystyle \frac{1}{\sqrt{m^{}{}_{}{}^{2}m_0^2}(m^{}{}_{}{}^{2}m^2)}}\mathrm{log}\left\{{\displaystyle \frac{1}{p^{}}}\left({\displaystyle \frac{d^2\sigma _S}{dm^{}{}_{}{}^{2}dt}}\right)\right\},`$
and analogously for the effective range $`r_e`$.
## III Accuracy of the method and comparison with other approaches
The most important advantage of the method proposed by us Gasparyan2003 is that a reliable estimate for the uncertainty of the extracted scattering length can be given. The sources for the uncertainty are trifold: (i) A possible influence of the final-state interaction in the other outgoing channels. For the reaction $`ppK^+\mathrm{\Lambda }p`$ considered in Ref. Gasparyan2003 this concerns the $`K\mathrm{\Lambda }`$ and $`KN`$ systems. (ii) The adopted value for $`m_{max}^2`$, the upper limit chosen for the dispersion integral in Eq. (7). (iii) A sensitivity to left-hand cuts of the production operation. A detailed analysis of the issues (ii) and (iii) , based on general arguments, presented in Ref. Gasparyan2003 suggests that the error in the scattering length should be typically in the order of 0.3 fm or less. The role of issue (i) cannot be quanitified theoretically but has to be investigated by performing experiments and corresponding analyses at different beam momenta Gasparyan2003 .
In this section we want to present a thorough examination of the accuracy of the proposed method and, in particular, to corroborate the error estimate, by means of concrete model calculations. A test based on one specific model calculation, namely for the reaction $`ppK^+\mathrm{\Lambda }p`$, has been already performed in Ref. Gasparyan2003 . However, here we want to put this investigation on a broader basis by considering final-state interactions of varying strengths, corresponding to a much larger range of values of the scattering length. Furthermore, and equally important, we want to compare the present method with the performance of other, approximative treatments of the final-state interaction that are commonly used in the literature to extract information on the scattering length and the effective range and have been applied to $`ppK^+\mathrm{\Lambda }p`$ Bale ; Hinter .
One of those approximative treatments follows from the assumption that the phase shifts are given by the first two terms in the effective range expansion,
$`p\mathrm{cot}(\delta (m^2))={\displaystyle \frac{1}{a}}+{\displaystyle \frac{r_e}{2}}p^2,`$ (8)
usually called the effective range approximation (ERA), over the whole energy range. Here $`p`$ is the relative momentum of the final-state particles under consideration in their center of mass system, corresponding to the invariant mass $`m^2`$. In this case the relevant integrals (4) can be evaluated in closed form as book
$`A(m^2){\displaystyle \frac{(p^2+\alpha ^2)r_e/2}{1/a+(r_e/2)p^2ip}},`$ (9)
where $`\alpha =1/r_e(1+\sqrt{12r_e/a})`$. Because of its simplicity Eq. (9) is often used for the treatment of the final-state interaction (FSI).
A further simplification can be made if one assumes that $`ar_e`$. This situation is practically realized in the $`{}_{}{}^{1}S_{0}^{}`$ partial wave of the $`NN`$ system. Then the energy dependence of the quantity in Eq. (9) is given by the energy dependence of the elastic amplitude
$`A(m^2){\displaystyle \frac{1}{1/a+(r_e/2)p^2ip}},`$ (10)
as long as $`p1/r_e`$. Therefore one expects that, at least for small kinetic energies, $`NN`$ elastic scattering and meson production in $`NN`$ collisions with a $`NN`$ final state exhibit the same energy dependence book ; migdal , which indeed was experimentally confirmed. This treatment of FSI effects is often referred to as Migdal-Watson (MW) approach migdal .
In order to examine the reliability of the three methods described above we took different $`YN`$ models from the literature Holz ; Melni2 ; NijV ; Haiden and calculated the production amplitude $`A(m^2)`$ utilizing the meson exchange model from Ref. model . Then this amplitude was used for extracting the scattering length by means of the dispersion integral Eq. (7) or from the approximative prescriptions given by Eqs. (9) and (10). For comparison we considered also the $`{}_{}{}^{1}S_{0}^{}`$ partial wave of the $`np`$ system of the Argonne potential Argonne . In this case $`A(m^2)`$ was set equal to the scattering wave function $`\mathrm{\Psi }(p,r)`$ at the origin, more precisely to $`\mathrm{\Psi }^{}(p,0)^{}`$, which corresponds to the assumption that the production operator is point-like.
Some selective results (for the Nijmegen NSC97 NijV and Jülich 01 Melni2 $`YN`$ models and the Argonne v14 Argonne $`NN`$ potential) are summarized in Table 1. The second column contains the correct scattering length evaluated directly from the potential model. One can see that the extraction of the scattering length via the dispersion integral (7) yields results pretty close to the original values for all considered potentials. In fact, in most cases the deviation is significantly smaller than the uncertainty of the method, estimated in Ref. Gasparyan2003 to be 0.3 fm. The results of the Jost-ERA approach Eq. (9) exhibit a systematic offset in the order of 0.3 fm. The situation is much worse for the Migdal-Watson approach Eq. (10) where a similar offset is found though now in the order of 0.6 fm. As a consequence, the extracted values differ by 50 % or more from the correct scattering lengths. Only for the $`{}_{}{}^{1}S_{0}^{}`$ $`np`$ partial wave the disagreement is still in the order of 5 %. Here the reliability of the Jost-ERA and Migdal-Watson approaches are comparable. This is in agreement with the expectations mentioned above.
The systematic offset inherent in the Jost-ERA approach as well as in the Migdal-Watson prescription can be best seen in Fig. 1, where we shown the difference between the scattering lengths predicted by various models and the values extracted via the dispersion integral (circles), the Jost-ERA method (squares) and the Migdal-Watson prescription (triangles).
While the Jost-ERA approach might still be a reasonable tool for getting a first rough estimate of the scattering length for a particular two-body interaction one should be rather cautious when using it for more quantitative analyses. In particular, its application in a combined fit to elastic scattering data and invariant mass spectra, e.g. to $`\mathrm{\Lambda }p`$ and $`ppK^+\mathrm{\Lambda }p`$, is rather problematic and can easily cause misleading results. Because of the offset in the scattering length in applications to final-state effects it is clear that a combined fit cannot converge to a unique (the “true”) $`\mathrm{\Lambda }p`$ scattering length. Only the elastic data will favour values close to the “true” scattering length whereas the production data tend to support larger (negative) values. This is obvious from the corresponding Jost-ERA results presented in Table 1 and also from Fig. 1. We believe that the analysis of Hinterberger and Sibirtsev presented in Ref. Hinter is an instructive exemplification of this dilemma. Employing the Jost-ERA approach to low energy total $`\mathrm{\Lambda }p`$ cross sections Alex ; Sechi and to experimental results for the missing mass spectrum of the reaction $`ppK^+X`$ Siebert separately, they derived (spin averaged) scattering lengths of $`a=1.81_{0.21}^{+0.18}`$ fm and $`a=2.57_{0.23}^{+0.20}`$ fm, respectively. Taking into account the error bars this is roughly the difference we would expect from the offset (of around $`0.3`$ fm) seen in our test calculations and, therefore, one must consider the results as being practically consistent with each other. But the authors of Ref. Hinter attempted to “reconcile” the results even more by introducing a spin-dependence in the fitting procedure. Indeed, with the relative magnitude of singlet to triplet contribution in the production reaction as free parameter (their relative strength in the elastic channel is fixed at 1:3 by the spin weight!) a “fully consistent” description of the combined data could be achieved Hinter and apparently the spin-singlet as well as spin-triplet $`\mathrm{\Lambda }p`$ $`S`$-wave scattering lengths could be determined from spin-averaged observables. Our experience with the Jost-ERA approach reported above, however, strongly suggests that the sensitivity to the spin seen in this analysis is most likely just an artifact of the method applied.
Let us now come to the effective range $`r_e`$. Since the dispersion relations yield only an integral representation for the product $`a^2((2/3)ar_e)`$ but not for the effective range $`r_e`$ alone Gasparyan2003 it follows that the attainable accuracy of $`r_e`$ is always limited roughly by twice the relative error on $`a`$. Still, it is interesting to see what values one gets for $`r_e`$ from the dispersion integrals. Corresponding results are presented in Table 2 and in Fig. 2. Evidently, the values extracted via the dispersion integral agree much better with the original results as one might have expected. In fact, in practically all cases the deviation is in the order of only 5 % or even less. This suggests that one could use the dispersion integrals also to extract the effective range $`r_e`$ from data. But one should keep in mind that, unlike the case of the scattering length, now one cannot rely on a solid and general estimate of the uncertainty. As far as the Jost-ERA approach is concerned it is clear from Table 2 that it yields rather poor results. In case of the Migdal-Watson prescription (10) it turned out that the fit always prefers an effective range $`r_e`$ equal to zero. This is due to the term proportional to $`r_e^2p^4`$ in the denominator of the $`A(m^2)`$ that make the production cross section decrease too fast as compared to the data (or to our calculations with realistic models). Therefore we don’t show any results of the Migdal-Watson fit for $`r_e`$.
One should note here that the upper limit in the dispersion integrals was always taken such that $`p_{max}=205`$ MeV/c (as in Gasparyan2003 ) that corresponds to $`ϵ_{max}m_{max}m_040`$ MeV for the $`\mathrm{\Lambda }N`$ and $`\mathrm{\Sigma }N`$ systems and $`ϵ_{max}45`$ MeV for the $`NN`$ system. In the latter case it is interesting to see what happens if one varies the range of integration, since the energy structure in the $`NN`$ interaction is much narrower due to the large $`NN`$ scattering length. For $`ϵ_{max}=10`$ MeV and $`ϵ_{max}=20`$ MeV as upper limits of the integration one gets the scattering lengths $`a=22.62`$ fm and $`a=23.17`$ fm, respectively – which are in principle still close to the original value. For the effective range, however, the calculation yields $`r_e=4.78`$ fm and $`3.72`$ fm, respectively. The reason why the agreement for $`ϵ_{max}>40`$ MeV is so good is that the $`NN`$ $`{}_{}{}^{1}S_{0}^{}`$ phase shift becomes sufficiently small at such energies, which implies a small uncertainty according to the error estimation in Ref. Gasparyan2003 .
## IV Dispersion relation in the presence of Coulomb repulsion
In the case when both baryons in the final state carry charges (for example in the reaction $`ppK^0p\mathrm{\Sigma }^+`$) there is a Coulomb interaction between them. Then the production amplitude $`A(m^2)`$ acquires additional singularities at $`p=0`$, due to the long-range nature of the Coulomb forces, and the formalism developed in Sect. 2 is no longer applicable directly. In this section we want to describe the modifications that are necessary in order to adapt the dispersion-relation method to the situation when the Coulomb force is present in the final-state interaction. We restrict ourselves to the case of a repulsive Coulomb interaction so that no bound states are present.
In order to elucidate the principle idea we start out from the case of elastic (two-body) scattering. Here the problem can be most conveniently dealt with by applying the Gell-Mann–Goldberger two-potential formalism Gell . Let us assume that the total potential $`V=V_c+V_s`$ is given by the sum of a short-ranged hadronic potential $`V_s`$ and the Coulomb interaction $`V_c`$. Then the total reaction amplitude $`T`$ can be written as $`T=T_c+T_{cs}`$, where $`T_c`$ is the Coulomb amplitude and $`T_{cs}`$ is defined by
$`T_{cs}=(1+T_cG_0)t_{cs}(1+G_0T_c),`$ (11)
where $`t_{cs}`$ fulfils a Lippmann-Schwinger equation,
$`t_{cs}=V_s+V_sG_ct_{cs},`$ (12)
with the short-range potential $`V_s`$ as driving term. To obtain the physical on-shell amplitudes one needs to project the corresponding $`T`$-operators on the so-called Coulombian asymptotic states $`|p_{\mathrm{}}\pm `$ which are related to the Coulomb scattering states (with fixed angular momentum – in our case $`l=0`$) $`|p\pm _c`$ via $`|p\pm _c=|p_{\mathrm{}}\pm +G_0^\pm T_c^\pm |p_{\mathrm{}}\pm `$ vanHaeringen1976 . Here $`p`$ denotes the center of mass momentum in the baryon-baryon system. In this way one obtains in particular
$`{}_{c}{}^{}p\left|t_{cs}\right|p+_{c}^{}=\frac{1}{\pi \mu }f_{cs}(p),`$ (13)
where $`f_{cs}`$ is the so-called Coulomb-modified nuclear scattering amplitude and $`\mu `$ is the reduced mass. Its relation to the phase shift $`\delta _{cs}`$ is the following
$`f_{cs}={\displaystyle \frac{e^{2i\delta _c}(e^{2i\delta _{cs}}1)}{2ip}},`$ (14)
with $`\delta _c`$ denoting the pure Coulomb $`S`$-wave phase shift given by $`\delta _c=\mathrm{arg}(\mathrm{\Gamma }(1+i\eta ))`$ with $`\eta =\frac{\mu e^2}{p}`$.
It has been shown in Ref. Hamilton1973 under rather general assumptions that the modified amplitude $`\stackrel{~}{f}(p)=e^{2i\delta _c}f_{cs}(p)/C^2(p)`$ is free of the Coulomb singularities on the physical sheet and possesses only the singularities caused by dynamical cuts (see also Refs. Cornille1962 ; vanHaeringen1977 ; Heller1966 ; Scotti1965 ). In addition below the inelastic cuts and above the two-baryon threshold the modified unitarity relation reads
$`\stackrel{~}{f}(s+i0)\stackrel{~}{f}(si0)=2ip\stackrel{~}{f}(p)\stackrel{~}{f}^{}(p)C^2(p)`$ (15)
with $`C^2(p)=\frac{2\pi \eta }{e^{2\pi \eta }1}`$ being the Coulomb penetration factor.
Furthermore an effective range function, modified for the presence of the Coulomb interaction, can be defined as well. Is is given by
$`S(p)pC^2(p)\mathrm{cot}\delta _{cs}(p)+Q(p)=1/a_{cs}+r_ep^2/2+..,`$ (16)
where $`Q(p)\mu e^2[\psi (i\eta )+\psi (i\eta )2\mathrm{ln}\eta ],\psi (z)=\mathrm{\Gamma }^{}(z)/\mathrm{\Gamma }(z)`$.
Coming back now to the production reaction it can be analogously shown that also the modified production amplitude
$`\stackrel{~}{A}(m^2)=e^{i\delta _c}A(m^2)/C(p)`$ (17)
is free of the aforementioned singularities Hamilton1973 . Therefore a dispersion relation similar to Eq. (1) can be written down
$`\stackrel{~}{A}(m^2)={\displaystyle \frac{1}{\pi }}{\displaystyle _{\mathrm{}}^{\stackrel{~}{m}^2}}{\displaystyle \frac{\stackrel{~}{D}(m^{}{}_{}{}^{2})}{m^{}{}_{}{}^{2}m^2}}dm^{}{}_{}{}^{2}+{\displaystyle \frac{1}{\pi }}{\displaystyle _{m_0^2}^{\mathrm{}}}{\displaystyle \frac{\stackrel{~}{D}(m^{}{}_{}{}^{2})}{m^{}{}_{}{}^{2}m^2}}dm^{}{}_{}{}^{2}.`$ (18)
Unitarity implies that the discontinuity for the elastic cut is
$`\stackrel{~}{D}(m^2)=\stackrel{~}{A}(m^2)e^{i\delta _{cs}}\mathrm{sin}\delta _{cs}.`$ (19)
The solution to Eq. (18) is found in complete analogy to the case without the presence of the Coulomb interaction,
$`\stackrel{~}{A}(m^2)=\mathrm{exp}[{\displaystyle \frac{1}{\pi }}{\displaystyle _{m_0^2}^{m_{max}^2}}{\displaystyle \frac{\delta _{cs}(m^{}{}_{}{}^{2})}{m^{}{}_{}{}^{2}m^2i0}}dm^{}{}_{}{}^{2}]\stackrel{~}{\mathrm{\Psi }}(m^2),`$ (20)
where $`\stackrel{~}{\mathrm{\Psi }}(m^2)`$ is some function slowly varying with $`m^2`$. If one neglects the the weak $`m^2`$ dependence present in $`\stackrel{~}{\mathrm{\Psi }}(m^2)`$ the expression for the phase shift $`\delta _{cs}`$ in terms of the differential cross section becomes
$`{\displaystyle \frac{\delta _{cs}(m^2)}{\sqrt{m^2m_0^2}}}={\displaystyle \frac{1}{2\pi }}𝐏{\displaystyle _{m_0^2}^{m_{max}^2}}{\displaystyle \frac{\mathrm{log}\left[{\displaystyle \frac{1}{p^{}C^2(p^{})}}{\displaystyle \frac{d^2\sigma }{dm^{}{}_{}{}^{2}dt}}\right]}{\sqrt{m^{}{}_{}{}^{2}m_0^2}(m^{}{}_{}{}^{2}m^2)}}\sqrt{{\displaystyle \frac{m_{max}^2m^2}{m_{max}^2m^{}^2}}}dm^{}{}_{}{}^{2}.`$ (21)
Using the effective range expansion (16) one can then extract the scattering length $`a_{cs}`$ from this dispersion integral.
## V Test of the method for the Coulomb case
One of the obvious reactions for applying the formalism with Coulomb is $`ppK^0\mathrm{\Sigma }^+p`$ where one could determine the $`\mathrm{\Sigma }N`$ scattering length for the isospin 3/2 state. Note that the $`\mathrm{\Sigma }^+p`$ channel does not couple to the $`\mathrm{\Lambda }N`$ system and is therefore free of inelastic cuts (that start already on the left-hand side) as required for the applicability of the disperson integral method. For this reaction one could perform a model calculation analogous to the one for $`ppK^+\mathrm{\Lambda }p`$ model which we used for testing the dispersion-integral method in the absence of the Coulomb interaction Gasparyan2003 . However, the implementation of Coulomb effects into our momentum-space code is technically complicated and requires also some approximations Hanhart . Thus, for the present test calculation we adopt a different strategy. First, instead of the momentum-space $`YN`$ models of Refs. Holz ; Melni2 ; Haiden we take the r-space Argonne ($`NN`$) potential, however, with parameters modified in such a way that the effective range parameters are similar to those predicted by realistic $`YN`$ potentials Holz ; NijV ; Melni2 ; Haiden for the $`\mathrm{\Sigma }N`$ $`I=3/2`$ $`{}_{}{}^{1}S_{0}^{}`$ partial wave. In particular we prepared two models with Coulomb modified scattering length of $`a_{cs}=3.24`$ fm (model 1) and $`a_{cs}=1.86`$ fm (model 2), respectively. The corresponding scattering lengths without Coulomb interaction are $`4.11`$ fm and $`2.01`$ fm, respectively. For the transition amplitude we use the scattering wave function calculated from those potential models and evaluated at the origin. This corresponds to the assumption that the production operator is point-like, which is reasonable as long as we are interested only in the mass dependence of the production amplitude. The corrections stemming from a possible mass dependence of the production operator were discussed in Ref. Gasparyan2003 .
The results of applying Eq. (7) with $`m_{max}m_0=40`$ MeV ($`p_{max}=205`$ MeV/c) are shown in Fig. 3, where we plot the function $`1/S(p)`$ which should coincide with the scattering length $`a_{cs}`$ at $`p=0`$. Obviously, there is a strongly nonanalytic behavior of the extracted inverse effective range function when approaching the threshold – which, however, can be easily understood. It is clear from Eq. (16) that the threshold behavior of the Coulomb modified phase shift is $`\delta _{cs}a_{cs}pC^2(p)`$, i.e. $`\delta _{cs}`$ goes to zero very rapidly. To obtain such a behavior on the left-hand side of Eq. (21) one needs to have a very precise cancellation in the integral on the right-hand side of Eq. (21) that is, of course, impossible if one truncates the integral at a finite momentum. But still one can expect for Eq. (21) to work for momenta not too close to the threshold, namely above the typical Coulomb scale of $`2\pi \alpha /\mu 25`$ MeV/c where the factors $`C^2(p)`$ and $`Q(p)`$ that appear in the effective range function (Eq. (16)) become smoother. This is indeed the case, as can be seen from Fig. 3. Thus, a natural and practical step here would be to extrapolate the extracted $`S(p)`$ to the threshold from above. Using a 4-th order polynomial of the type $`1/a_{cs}+r_ep^2/2Pr_e^3p^4`$ and fixing the coefficients in the region $`50100`$ MeV/c, i.e. well above the Coulomb structure, one can then extrapolate $`S(p)`$ to the threshold, cf. the dash-dotted lines in Fig. 3. In this case one gets a satisfactory agreement between the true and extracted scattering lengths. In fact, the deviations are not worse then in the case when we consider the same potentials without Coulomb interaction and they are also within the theoretical error of 0.3 fm estimated in Ref. Gasparyan2003 . The extracted values are $`a_{cs}=3.10`$ fm and $`a_{cs}=1.86`$ fm for models 1 and 2, respectively, with Coulomb, and $`a_s=4.05`$ fm and $`a_s=2.06`$ fm when the Coulomb interaction is switched off. We also checked that the sensitivity of the result to the region of interpolation of the effecitve range function $`S(p)`$ is rather low. For instance if one shifts the lower bound of this region to $`70`$ MeV/c the corresponding change in the scattering length will be less then $`0.05`$ fm.
If one wishes to consider a more realistic situation, one needs to deal with mass distributions with finite statistical errors, finite mass resolution, and finite binning. To examine also this situation we have generated two data sets, corresponding to the models 1 and 2 as shown on Fig. 4. We have chosen the binning as well as the mass resolution to be equal to $`2`$ MeV (the same as in the experiment Siebert that was analyzed in Gasparyan2003 ), and the statistics to be rather high in order to minimize the influence of the statistical error bars on the results. The excess energy was set to $`40`$ MeV to simplify the simulation. In a realistic situation larger values are preferable to minimize the influence of the meson-baryon interactions, cf. the corresponding remarks in Ref. Gasparyan2003 . In our test calculation such meson-baryon interactions are neglected anyway.
We start here with the procedure suggested in appendix A of Gasparyan2003 , namely by fitting the cross section with an exponential parameterization of the type
$`{\displaystyle \frac{d^2\sigma }{dm^2dt}}=C^2(p)\times \mathrm{exp}\left[C_0+{\displaystyle \frac{C_1^2}{(m^2C_2^2)}}\right]\times \text{phase space}.`$ (22)
This formula fits the generated cross section with the $`\chi ^2`$ per degree of freedom of $`\chi _{dof}^21`$, cf. Fig. 4. A new problem that arises here is that the production amplitude contains a very narrow structure (of the size $`2\pi \alpha /\mu `$) close to threshold as can be seen in Fig. 5 (solid lines). Clearly, this structure can not be reproduced after the fitting procedure as it gets smeared out by the mass resolution and binning (their size is much larger than the scale of the structure). The amplitudes coming from the fit are depicted in Fig. 5 by dash-dotted lines. However the fit can be improved if one notes that the structure comes mostly from the part of the dispersion integral (20) containing the leading term in the $`\delta _{cs}`$ expansion near the threshold, namely (in the nonrelativistic case)
$`\mathrm{exp}[{\displaystyle \frac{1}{\pi }}{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{a_{cs}C^2(p^{})}{p^{}{}_{}{}^{2}p^2i0}}\left({\displaystyle \frac{p^2}{p^2}}\right)dp^{}{}_{}{}^{2}]=\mathrm{exp}[a_{cs}Q(p)],`$ (23)
where we made a subtraction at $`p=0`$ to render the integral convergent. This does not change the energy dependence of the resulting exponent. Indeed the structure disappears after dividing the production amplitude by the factor $`\mathrm{exp}[a_{cs}Q(p)]`$.
Obviously the scattering length is unknown before its extraction! But what can be done is to resort to an iterative procedure including first the extraction of the unimproved scattering length, then putting it into the fit function,
$`{\displaystyle \frac{d^2\sigma }{dm^2dt}}=C^2(p)\times \mathrm{exp}\left[C_0+{\displaystyle \frac{C_1^2}{(m^2C_2^2)}}2a_{cs}Q(p)\right]\times \text{phase space},`$ (24)
and repeating this step until the procedure converges. Fortunately the iterations converge already after three or four iterations, and the resulting amplitudes are shown in Fig. 5 by the dashed lines. The improvement of the fit is quite obvious. Finally we applied the combination of the extrapolation and iteration procedures to obtain the scattering lengths from the pseudodata. Since the data have a statistical uncertainty we generated a sample of 1000 mass distributions and looked at the average value of the scattering lengths. They turned out to be $`2.89\pm 0.06`$ fm for model 1 and $`1.82\pm 0.05`$ for model 2. Note that the deviation from the correct values is now a bit larger, but still reasonable ($`0.35`$ fm in the worst case). Fig. 6 shows how the extracted inverse effective range function approaches to the correct one given by the model by the example of model 1.
## VI Summary
In a recent publication Gasparyan2003 we have presented a formalism based on dispersion relations that allows one to relate spectra from large-momentum transfer reactions, such as $`ppK^+p\mathrm{\Lambda }`$ or $`\gamma dK^+n\mathrm{\Lambda }`$, directly to the scattering length of the interaction of the final-state particles. An estimation of the systematic uncertainties of that method, relying on general arguments, led to the conclusion that the theoretical error in the extracted scattering length should be less than 0.3 fm. This finding was corroborated in an application of the method to results of a microscopic model calculation for $`ppK^+p\mathrm{\Lambda }`$.
In the present paper this dispersion theoretical method was generalized to the case where a repulsive Coulomb force is present in the final-state interaction. As an example let us mention the reaction $`ppK^0p\mathrm{\Sigma }^+`$ which could be used to extract the $`p\mathrm{\Sigma }^+`$ scattering length. Though the generalization of the formalism itself is straight forward it turned out that there are some additional features due to the Coulomb interaction which need to be taken into account in concrete applications of the method to data. These practical aspects were thoroughly discussed and it was shown how to circumvent the difficulties. In a test calculation utilizing potential models with effective range parameters similar to those of realistic $`YN`$ interactions the extracted values for the scattering lengths were found to agree within 0.3 fm with those predicted by the models. Thus, the accuracy of the dispersion theoretical method for extracting the scattering lengths from final-state interactions with Coulomb force is comparable to the case where no Coulomb interaction is present.
We presented also a more detailed examination of the accuracy of the dispersion-integral method than in Ref. Gasparyan2003 . In particular we considered final-state interactions of varying strengths, corresponding to a much larger range of values of the scattering length. These investigations confirmed the reliability of the general error estimate provided in Ref. Gasparyan2003 . Indeed in most of the considered cases the deviation of the extracted scattering length from the true value was significantly smaller than the uncertainty of 0.3 fm derived in that paper. In addition we studied the effective range $`r_e`$ which can be also extracted by the proposed dispersion-integral method. For most of the interaction models considered the extracted values of $`r_e`$ agreed remarkably good with the true results. Thus, it might be sensible to use the proposed method to extract the effective range from data - though one should always keep in mind that for this quantity a generally valid error estimation is not possible Gasparyan2003 .
Finally, we compared the present method with the performance of other, approximative treatments of the final-state interaction that are commonly used in the literature to extract information on the scattering length and also the effective range. In particular we tested the Jost-function approach based on the effective-range approximation (Jost-ERA) and an even simpler approach that relies simply on utilizing the effective range approximation itself. Thereby, we showed that the latter methods lead to systematic deviations from the true values of the scattering lengths in the order of 0.3 fm (Jost-ERA) and even 0.7 fm (direct effective-range approximation). This suggests that one should be rather cautious in the interpretation of results achieved with those methods.
###### Acknowledgements.
A.G. would like to acknowledge finanical support by the grant No. 436 RUS 17/75/04 of the Deutsche Forschungsgemeinschaft. Furthermore he thanks the Institut für Kernphysik at the Forschungszentrum Jülich for its hospitality during the period when the present work was carried out.
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# Completely Quantized Collapse and Consequences
## I Introduction
Quantum theory as a theory of measurement was promulgated by BohrBohr , and this point of view has its adherents todayPeres-Fuchs . But, the success of the theory suggests that it could be more, a theory of reality. Einstein wrote SchrödingerEinsteinletter that, if there were no more complete description of nature than that envisaged by Bohr, “…then physics could only claim the interest of shopkeepers and engineers; the whole thing would be a wretched bungle.” Bell made much the same point, writingBell “ To restrict quantum mechanics to be exclusively about piddling laboratory experiments is to betray the great enterprise.”
The lack in standard quantum theory has been called the “measurement problem,” but I prefer to call it the “reality problem.” I believe it is most precisely characterized as a failure to give a satisfactory response to the following operational requirement for a well-defined quantum theory of reality. You are given, at time $`t=0`$, a state vector for an arbitrary physical system, together with its Hamiltonian, but no other information. For $`t>0`$, specify (i.e., give a procedure for how to identify) the realizable states and their probabilities of realization.
Not only Bohr’s version of quantum theory, but more recent sophisticated attempts, such as many worldslev , environmental decoherencezehzurek , consistent historiesgriffithshartlegellmannomnes , have not succeeded so far at providing such a specificationmekent . They are successfully applied to many important and interesting examples, but these examples require additional information (this is the nature of the system, that is the apparatus, there is the environment…) which lie outside the theory.
The reality problem has been solved by enlarging quantum theory so that it describes state vector collapse as a physical process. The Continuous Spontaneous Localization (CSL) modelme ; GPR ; myreview provides a modified Schrodinger equation. It describes how, in addition to their usual interactions, the particles of the world interact with a posited classical field $`w(𝐱,t)`$ resulting in a nonunitary state vector evolution (which changes the norm of the state vector). What are the realizable states? Each such state vector, corresponding to a different $`w(𝐱,t)`$ (drawn from the class of white noise “functions), is realizable. What are the associated probabilities? CSL provides a second equation, the prescription that the probability of occurrence of a state vector (the probability of occurrence of its associated $`w(𝐱,t)`$) is proportional to its squared norm. In this straightforward manner, the operational requirement cited above is satisfied. Each of these CSL states accurately describes (except for occasional brief moments) the macroscopic world we see around us. The predictions agree with all presently performed experiments (but there are testable differences from standard quantum theory). The CSL model is reviewed in section II.
I have previously presented a modelWays ; mekent ; myreview called the “Index” model, but here called the “Completely Quantized Collapse” (CQC) model. It is reviewed in section III. It utilizes a formalism invented to model continuous measurementscontmeas . CQC is completely equivalent to CSL in its specification of the realizable states and their probabilities, in the following manner.
In CQC, the classical field is replaced by a quantized field $`W(𝐱,t)`$ which, like the classical field $`w(𝐱,t)`$ of CSL, commutes with itself, $`[W(𝐱,t),W(𝐱^{},t^{})]=0`$. There is a conjugate field, $`\mathrm{\Pi }(𝐱,t)`$, which also commutes with itself, but which satisfies $`[W(𝐱,t),\mathrm{\Pi }(𝐱^{},t^{})]=i\delta (tt^{})\delta (𝐱𝐱^{})`$. The basis of eigenstates of $`W(𝐱,t)`$ for every $`𝐱`$, $`t`$, is denoted $`|w`$. It satisfies $`W(𝐱,t)|w=w(𝐱,t)|w`$, i.e., there is a different eigenstate for every possible $`w(𝐱,t)`$ function. There is a single state vector, called here the “ensemble vector,” which evolves unitarily (the Hamiltonian is given). When the ensemble vector is expanded as a Schmidt decomposition in the “preferred basis” $`|w`$, the particle state in the direct product with $`|w`$ is the CSL state associated with $`w(𝐱,t)`$: schematically,
$$|\text{ensemble vector},t=\underset{\text{all}w}{}|w|\psi _{\text{csl}},t_w.$$
The usual quantum rule for probabilities (the squared norm of a state in the superposition) gives the CSL probability of realization of that state.
The measurement (reality) problem in standard quantum theory has often been phrased in terms of difficulties associated with the ill-defined collapse postulate, nowhere more succinctly than Bell’s Boolean phrase“And is not Or,” i.e., Schrödinger’s equation gives a sum of terms but we observe one term or another. CSL provides a resolution of the problem by individually, dynamically, providing each term we can observe. However, CQC gives a sum of terms: how can that be considered as a resolution of the problem?
A colleague once observed to me that construction of CQC goes against everything I have worked for, the resolution of the reality problem by dynamical collapse. However, what I have worked for is to make a well-defined quantum theory of realityPearle'67 , i.e., one which satisfies the aforementioned operational requirement. CSL satisfies this requirement, but so does CQC. The terms in the superposition which make up the ensemble vector are precisely defined (by the orthogonal preferred basis $`|w`$), do not interfere (they obey a superselection rule), are macroscopically sensible (because the associated CSL states are), and occur with the correct amplitudes (which give the correct probabilities). In his famous “cat paradox” paperSchr , Schrödinger referred to the state vector as the “prediction catalog” or the “expectation catalog” because, in Bohr’s quantum theory, it lists what can be expected to occur if a measurement is performed. The CQC state vector may, likewise, be called a “reality catalog.”
However, the justification here is greater than Schrödinger’s. His usage requires judgment outside the theory as to when the state vector has evolved sufficiently to be considered a catalog, and as to which are the appropriate states in the superposition (i.e., which basis) to consider as making up the catalog. In CQC, these criteria lie inside the theory, since the ensemble-vector may always be considered a catalog, and the appropriate states are labeled by the preferred basis $`|w`$. Roughly speaking, CSL provides a vertical catalog (Or), and CQC provides a horizontal catalog (And), but it is the same catalog. If one has a precise rule for identifying a superposition of states at time $`t`$ which never interfere thereafter, it is a matter of indifference whether one removes them from the superposition and follows their separate evolutions that way, or leaves them in the superposition and follows their joint evolution that way.
As previously pointed outmekent , one may regard CQC as a model for a satisfactory resolution of other attempted interpretations of quantum theoryGS . The superposed states making up the ensemble vector could be thought of as many worlds (although it would seem that there would then be no reason to insist upon such an interpretation). If the $`W`$-field is regarded as the environment, then CQC could be regarded as describing environmental decoherence. The evolving superposed states making up the ensemble vector could be taken to represent a single family of consistent historieskent . The Schmidt decomposition of the ensemble vector is precisely what the modal interpretation requires. What is currently missing from these approaches, and what CQC provides, is a subsystem with a well-defined “universal basis” (here the $`|w`$ basis vectors) whose Schmidt decomposition partners make macroscopic sense. One might hope that such a basis may be found arising in a natural way from a physical entity, whence perhaps CSL/CQC and all these interpretations may merge (see the last two paragraphs of this introduction for further remarks on this).
Now I turn to the topics which will occupy the rest of this paper. Because collapse narrows wave functions, the energy of particles in CSL increases with time. This growth of energy of particles provides for experimental tests of the modelexptlpapers but, from a theoretical point of view, it is desirable to have a conserved energy. Recently, I showed how to define energy associated with the $`w`$-field in CSL so that the first moment of the total energy (particles plus $`w`$-field) is conservedenergy . However, energy conservation requires all moments of the energy to be conserved. The formalism of CQC facilitates that demonstration (which can be converted to CSL language with less transparency). Because the ensemble vector evolution in CQC is a Hamiltonian evolution, the Hamiltonian is a conserved energy associated with the ensemble.
The Hamiltonian does not commute with $`W(𝐱,t)`$ since it is composed of the $`W`$ and $`\mathrm{\Pi }`$-field operators. Thus one cannot associate a conserved energy with an individual realizable state labeled by eigenvalues of $`W`$ . Nonetheless, it is of interest to see how energy fares for the ensemble of possible states, for example, how the ensemble energy of the $`W`$-field can go negative to compensate for the energy increase of the particles whose wavefunctions are narrowed by collapse. The discussion of energy conservation takes place in section IV.
Another example of the utility of the CQC model is considered in section V. Because the $`W`$-field energy has the whole real line for its spectrum, it is possible to construct its self-adjoint conjugate, a time-operator $`T`$, built solely out of the $`W`$ and $`\mathrm{\Pi }`$-field operators. One may consider that time should be intimately related to the realization of physical states, since without such a realization there are no events and without events there is no time. This, and the fact that, like the Matterhorn, it is “there” (i.e., very few viable physical theories possess an unbounded Hamiltonian) motivate looking at $`T`$.
$`T`$’s ensemble probability distribution is examined. $`T`$’s variance diminishes with time. $`T`$’s mean, roughly speaking, is a “center of time,” i.e., the average time over all space, and over time from 0 to $`t`$, weighted by the square of the particle number density. Thus, if the particle Hamiltonian vanishes so the only particle evolution is collapse, which does not change the ensemble spatial particle density, the mean of $`T`$ is $`t/2`$. If, however, the particles have a high clumping rate (for example, due to gravitational interaction), the spatial average of the square of the particle number density is largest most recently, and the mean of $`T`$ is closer to $`t`$.
One of the motivations of this work is to provide a bridge from dynamical collapse models to other interpretations which seek a well-defined quantum theory based upon the standard quantum formalism alone (although they may be satisfied with just accounting for observation, not reality), rather than the altered formalism of collapse models. Now, what I consider to be a pressing problem associated with CSL is that it is phenomenological. (Some have said “ad hoc,” but that is not appropriate terminology since ad hoc means “for this case only,” whereas CSL applies to all cases, all physical situations: in fact, the label “ad hoc” more appropriately belongs to these other interpretations to the extent that they require additional information appropriate to each physical situation). It could be that CSL dynamics may arise from a larger structure which generates quantum theory itself, as has been suggested by AdlerAdler . In this paper, we consider another possibility.
CSL depends upon interaction of a posited field’s interaction with particles, but no identification of that field with a physical entity has been made. A purpose in discussing CQC is that, since it is formulated completely quantum mechanically, that may make it easier to identify the universal $`W`$-field with a quantized physical entity which arises naturally in another context. In that case, CQC would appear as based upon the standard quantum formalism alone.
However, it is not so easy to find a candidate for a universal fluctuating field with the desired properties. There is experimental evidenceexptlpapers that the coupling of $`W`$ is to the mass-density of particles, which strongly suggests a gravitational connection, a proposal which has a long historygravity ; pearlesquires . This sugggests looking for a mechanism whereby gravitational excitations could naturally act like the $`W`$-field.
By way of illustration, Section VI contains two models for $`W`$ with gravitational overtones, neither to be taken too seriously. In one, $`W`$ is constructed out of many scalar quantum fields (of Planck mass), each interacting only for a brief interval (the Planck time) with the particle’s mass-density. The other utilizes aspects of a previous gravitational classical modelpearlesquires of $`W`$ and crudely pays homage to spin-network modelssmolinseth . It assumes that space consists of Planck-volume cells, each containing a bound spin 1/2 entity with two possible energy states equal to $`\pm `$Planck mass. A spin occasionally breaks free to briefly interact with a thermal bath, as well as gravitationally with the particle’s mass-density, before it is bound once more. In this case, $`W(𝐱,t)`$ is identified with the free mean spin, in a small volume surrounding $`𝐱`$, during the small time interval surrounding $`t`$,. These examples suggest that what is required for $`W`$ is a quantity which has the freedom to vary independently over all space-time and for which the state vector describes its prior space-time values . This is unlike a conventional field, which can vary independently (in magnitude and time derivative) only on a single spacelike hypersurface, and whose present values are all that is described by the state vector.
## II CSL
### II.1 The Simplest Model
The simplest CSL model, which will supply most of the illustrative calculations in this paper, describes collapse toward an eigenstate of the operator $`A`$. The state vector evolution, governed by a function $`w(t)`$ of white noise class, is given by
$$|\psi ,t_S=𝒯e^{_0^t𝑑t^{}\{iH_A+(4\lambda )^1[w(t^{})2\lambda A]^2\}}|\varphi .$$
(1)
In Eq.(1), $`𝒯`$ is the time-ordering operator, $`H_A`$ is the (usual) Hamiltonian of the system, $`|\varphi =|\psi ,0`$ is the initial state vector, the subscript $`S`$ denotes the Schrödinger-picture state vector, and the parameter $`\lambda `$ governs the collapse rate. Eq. (1) is the solution of the modified Schrödinger equation:
$$d|\psi ,t_S/dt=\{iH_A(4\lambda )^1[w(t)2\lambda A]^2\}|\psi ,t_S.$$
However, only the solution form (1) shall be utilized in this paper.
CSL is completed by specifying the probability that a particular $`w`$, lying in the interval $`(w(t),w(t)+dw),`$ is realized in nature:
$$P_t(w)Dw=_S\psi ,t|\psi ,t_SDw.$$
(2)
In Eq.(2), and for all integrations, one may discretize time, i.e., with $`t_jjdt`$, each $`w(t_j)`$ is considered to be an independent variable ($`\mathrm{}<w(t_j)<\mathrm{}`$), and $`Dw_{j=1}^{t/dt}\{dw(t_j)/[dt/2\pi \lambda ]^{1/2}\}`$, with the eventual limit $`dt0`$.
In the “collapse interaction picture,” which is what is mostly employed in this paper, the state vector is
$`|\psi ,t`$ $``$ $`\mathrm{exp}iH_At|\psi ,t_S`$ (3a)
$`=`$ $`𝒯e^{(4\lambda )^1_0^t𝑑t^{}[w(t^{})2\lambda A(t^{})]^2}|\varphi ,`$ (3b)
where (3b) follows from Eq.(1), and $`A(t^{})\mathrm{exp}(iH_At)A\mathrm{exp}(iH_At)`$. From Eq.(3a), the probability rule, Eq.(2), is the same, except that the subscript $`S`$ is dropped.
The density matrix describing the ensemble in the interaction picture follows from Eqs.(2), (3b):
$`\rho (t)`$ $`{\displaystyle P_t(w)Dw\frac{|\psi ,t\psi ,t|}{\psi ,t|\psi ,t}}={\displaystyle Dw|\psi ,t\psi ,t|}`$ (4b)
$`=𝒯e^{(\lambda /2)_0^t𝑑t^{}[A_L(t^{})A_R(t^{})]^2}|\varphi \varphi |`$
where, in an expansion of the exponent in Eq.(4b), $`A_L(t^{})`$ $`(A_R(t^{})`$) appears to the left (right) of $`|\varphi \varphi |`$, and is time-ordered (reverse-time-ordered) by $`𝒯`$.
### II.2 Proof of Collapse When $`H_A=0`$
If $`H_A=0`$ (or if the collapse rate is so fast that, over the time interval $`(0,t)`$, we can well approximate $`A(t)A`$), then the state vector dynamically collapses, as $`t\mathrm{}`$, to one of the eigenstates of A. Let $`|\varphi =_n\alpha _n|a_n`$, where $`A|a_n=a_n|a_n`$. An intimation of the collapse behavior in this case may be seen from the density matrix (4b):
$$\rho (t)=\underset{n,m}{}\alpha _n\alpha _m^{}e^{(\lambda t/2)(a_na_m)^2}|a_na_m|,$$
since the off-diagonal elements vanish as $`t\mathrm{}`$, and the diagonal element magnitudes are the ones obtained from $`|\varphi `$ using the “collapse postulate” of standard quantum theory. However, this ensemble density matrix behavior does not prove collapse in the individual case, since many different ensembles can have the same density matrix.
Here is a proof of collapse not previously given. Eqs.(3),(2) become respectively:
$`|\psi ,t`$ $`=`$ $`{\displaystyle \underset{n}{}}\alpha _ne^{(4\lambda )^1_0^t𝑑t^{}[w(t^{})2\lambda a_n]^2}|a_n`$ (5a)
$`P_t(w)`$ $`=`$ $`{\displaystyle \underset{n}{}}|\alpha _n|^2e^{(2\lambda )^1_0^t𝑑t^{}[w(t^{})2\lambda a_n]^2}`$ (5b)
It shall now be shown that the only $`w(t)`$’s of non-zero measure are those for which the asymptotic ($`t\mathrm{}`$) time-average of $`w(t)`$ is $`2\lambda a_n`$. To see this, consider the joint probability of $`w`$ and $`a(2\lambda t)^1_0^t𝑑t^{}w(t^{})`$:
$$P_t(w,a)Dwda=2\lambda tda\delta \left[_0^t𝑑t^{}w(t^{})2\lambda ta\right]P_t(w)Dw,$$
(6)
where $`P_t(w)Dw`$ is the probability of $`w`$ and $`2\lambda tda\delta `$ is the conditional probability, given $`w`$, that $`w`$’s time-average is $`2\lambda a`$. To find the probability distribution $`P_t(a)`$ of $`a`$, one must integrate (6) over $`Dw`$. To do so, make the replacement $`\delta (z)=(2\pi )^1𝑑\omega \mathrm{exp}i\omega z`$, and then the resulting gaussian integrals over $`Dw`$ may be performed. Next, the (gaussian) integral over $`\omega `$ can be performed, resulting in:
$$P_t(a)da=\underset{n}{}|\alpha _n|^2\frac{e^{2\lambda t(aa_n)^2}}{\sqrt{\pi /2\lambda t}}da_\stackrel{}{t\mathrm{}}\underset{n}{}|\alpha _n|^2\delta (aa_n)da.$$
(7)
Thus, the asymptotic ensemble probability is completely accounted for by $`a`$ taking on only the values $`a_n`$, and the probability of an $`a_n`$ is $`|\alpha _n|^2`$.
It remains to show that (5a) asymptotically describes collapse. Let $`w(t)`$ satisfy the constraint $`C_a`$ defined as $`lim_T\mathrm{}(2\lambda T)^1_0^T𝑑t^{}w(t^{})=a`$. Taking the asymptotic limit of (6) and then integrating over an infinitesimal range about $`a`$ followed by integrating over $`w`$, it follows from (6), (7) that
$$|\alpha _n|^2\delta _{a,a_n}=Dw_{aϵ}^{a+ϵ}𝑑aP_{\mathrm{}}(w,a)=_{C_a}DwP_{\mathrm{}}(w).$$
Applying this result to (5b) yields
$$_{C_a}Dwe^{(2\lambda )^1_0^{\mathrm{}}𝑑t^{}[w(t^{})2\lambda b]^2}=\delta _{a,b}.$$
Since each term in the integral here is non-negative, if $`ab`$, for each $`w(t^{})`$ obeying the constraint $`C_a`$ (i.e., for each $`w(t^{})`$ which asymptotically equals $`2\lambda a`$ plus a fluctuating term which averages out to 0), then $`exp(2\lambda )^1_0^{\mathrm{}}𝑑t^{}[w(t^{})2\lambda b]^2=0`$.
This result can be applied to Eq.(5a), where the gaussians are of the same form (except $`4\lambda `$ replaces $`2\lambda `$). Thus, for each such $`w(t)`$ which obeys the constraint $`C_a`$, with $`aa_n`$, the state vector asymptotically vanishes: that doesn’t matter since, by Eq.(5b), the probability of its occurrence vanishes also. When $`w(t)`$ obeys the constraint $`C_{a_m}`$, all terms in (5a) but the $`m`$th asymptotically vanish, $`|\psi ,\mathrm{}|a_m`$: collapse occurs.
### II.3 Full-Blown CSL
The CSL proposal for a serious collapse model, applicable to all nonrelativisitic systems, involves multiple $`A`$’s, each with its own $`w`$. That is, collapse is toward the joint eigenstates of an operator $`A(𝐱)`$ for all $`𝐱`$, and $`w(t)`$ is replaced by the random field $`w(𝐱,t)`$. The choice of $`A`$ isme ; pearlesquires :
$$A(𝐱)\underset{i}{}\frac{m_i/m_0}{(\pi a^2)^{3/4}}𝑑𝐳e^{(2a^2)^1(𝐱𝐳)^2}\xi _i^{}(𝐳)\xi _i(𝐳).$$
(8)
In Eq.(8), $`\xi _i(𝐳)`$ is the annihilation operator at location $`𝐳`$ for a particle of mass $`m_i`$ ($`m_0`$ is the proton mass), and $`a`$ (and $`\lambda `$) are parameters from the seminal GRW modelGRW , with chosen values $`a10^5`$cm and $`\lambda 10^{16}`$sec<sup>-1</sup>. Thus, since $`\xi _i^{}(𝐳)\xi _i(𝐳)`$ is the particle number density of the $`i`$th particle type at $`𝐳`$, $`A(𝐱)`$ is essentially proportional to the mass density inside a sphere of radius $`a`$ centered at x. The CSL dynamical equation which replaces Eq.(3b) is:
$$|\psi ,t=𝒯e^{(4\lambda )^1_0^t𝑑t^{}𝑑𝐱[w(𝐱,t^{})2\lambda A(𝐱,t^{})]^2}|\varphi $$
(9)
(where $`A(𝐱,t^{})\mathrm{exp}(iH_At)A(𝐱)\mathrm{exp}(iH_At)`$), i.e., time is replaced by space-time. The probability rule is unchanged, $`P_t(w)Dw=\psi ,t|\psi ,tDw`$, but space-time is discretized for integration, so $`Dw_{j,k}\{dw(𝐱_k,t_j)/[dtd𝐱/2\pi \lambda ]^{1/2}\}`$,
## III CQC
### III.1 Simplest Model
Start by introducing creation and annihilation operators $`b^{}(\omega )`$, $`b(\omega )`$ for $`\mathrm{}<\omega <\mathrm{}`$, which obey the usual commutation rules,
$$[b(\omega ),b(\omega ^{})]=0,[b^{}(\omega ),b^{}(\omega )^{}]=0,[b(\omega ),b^{}(\omega ^{})]=\delta (\omega \omega ^{}).$$
Then, the operator $`W(t)`$ (the quantized equivalent of the CSL $`w(t)`$) and its conjugate $`\mathrm{\Pi }(t)`$ may be defined as
$`W(t)`$ $``$ $`(\lambda /2\pi )^{1/2}{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\omega [e^{i\omega t}b(\omega )+e^{i\omega t}b^{}(\omega )],`$ (10a)
$`\mathrm{\Pi }(t)`$ $``$ $`i/(8\pi \lambda )^{1/2}{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\omega [e^{i\omega t}b(\omega )+e^{i\omega t}b^{}(\omega )]`$ (10b)
which are readily shown to satisfy
$$[W(t),W(t^{})]=0,[\mathrm{\Pi }(t),\mathrm{\Pi }(t^{})]=0,[W(t),\mathrm{\Pi }(t^{})]=i\delta (tt^{})$$
(11)
The vanishing of the first two commutators in Eqs.(11) result from cancellation of the positive frequency contribution in Eq.(10) by the negative frequency contribution. (Incidentally, a free tachyonic field also has equal positive and negative energy spectra, so it too can be quantized to have vanishing self-commutatortachyon ).
The joint eigenstate of $`W(t)`$ ($`\mathrm{\Pi }(t)`$) is denoted $`|w`$ ($`|\pi `$), where the label represents a particular function $`w(t)`$ ($`\pi (t)`$) of white noise class. It satisfies $`W(t)|w=w(t)|w`$ ($`\mathrm{\Pi }(t)|\pi =\pi (t)|\pi `$). The range of $`t`$ is $`(\mathrm{},\mathrm{})`$, and the states are normalized so that $`_{\mathrm{}}^{\mathrm{}}D^{}w|ww|=1`$ ($`_{\mathrm{}}^{\mathrm{}}D^{}\pi |\pi \pi |=1`$). $`D^{}`$ differs from the $`D`$ employed for CSL in that it involves the product of $`dw(t)`$’s for the infinite range of $`t`$. To discretize, write $`|w=_{j=\mathrm{}}^{\mathrm{}}|w(t_j`$, with $`w(t_j)|w^{}(t_j)=\delta [w(t_j)w^{}(t_j)]\sqrt{2\pi \lambda /dt}`$.
It follows from Eqs.(11) that, for an arbitrary vector $`|\mathrm{\Phi }`$,
$$w|\mathrm{\Pi }(t)|\mathrm{\Phi }=\frac{1}{i}\frac{\delta }{\delta w(t)}w|\mathrm{\Phi }.$$
(12)
Here, $`\delta /\delta w(t)`$ is the functional derivative. In time-discretized calculations, $`\delta /\delta w(t)(dt)^1/w(t_j)`$.
The vacuum state $`|0`$ satisfies $`b(\omega )|0=0`$ as usual. Since, from Eqs.(10), $`W(t)+2i\lambda \mathrm{\Pi }(t)`$ is an integral over $`b(\omega )`$, Eqs.(10), (12) imply
$$w|W(t)+2i\lambda \mathrm{\Pi }(t)|0=[w(t)+2\lambda \delta /\delta w(t)]w|0=0$$
and a similar equation for $`\pi |0`$, whose solutions are
$$w|0=e^{(4\lambda )^1_{\mathrm{}}^{\mathrm{}}𝑑t^{}w^2(t^{})},\pi |0=e^{\lambda _{\mathrm{}}^{\mathrm{}}𝑑t^{}\pi ^2(t^{})}.$$
(13)
The normalization in Eqs.(13) ensure that $`D^{}w|w|0|^2=D^{}\pi |\pi |0|^2=0|0=1`$.
The $`W`$-field Hamiltonian, the energy operator which generates time-translations of $`W(t)`$, $`\mathrm{\Pi }(t)`$, is
$$H_w_{\mathrm{}}^{\mathrm{}}𝑑\omega \omega b^{}(\omega )b(\omega )=_{\mathrm{}}^{\mathrm{}}𝑑t^{}\dot{W}(t^{})\mathrm{\Pi }(t^{}),$$
(14)
which follows from Eqs.(10). (By adding an infinitesimal imaginary part to $`\omega `$ in Eqs.(10), one can arrange that $`W(t)`$, $`\mathrm{\Pi }(t)`$ vanish at $`t\pm \mathrm{}`$, so one can freely integrate Eq.(14) by parts, to put the time derivative on $`\mathrm{\Pi }`$. One can also change the order of $`\dot{W}(t^{})`$ and $`\mathrm{\Pi }(t^{})`$ in Eq.(14), since $`𝑑t^{}[\dot{W}(t^{}),\mathrm{\Pi }(t^{})]=0`$.)
The “ensemble vector” of CQC in the collapse interaction picture is defined as:
$`|\psi ,t`$ $``$ $`𝒯e^{2i\lambda _0^t𝑑t^{\prime \prime }A(t^{\prime \prime })\mathrm{\Pi }(t^{\prime \prime })}|0|\varphi `$ (15a)
$`=`$ $`{\displaystyle D^{}w|ww|𝒯e^{2i\lambda _0^t𝑑t^{\prime \prime }A(t^{\prime \prime })\mathrm{\Pi }(t^{\prime \prime })}|0|\varphi }`$ (15b)
$`=`$ $`{\displaystyle D^{}w|w𝒯e^{2\lambda _0^t𝑑t^{\prime \prime }A(t^{\prime \prime })(\delta /\delta w(t^{\prime \prime }))}e^{(4\lambda )^1_{\mathrm{}}^{\mathrm{}}𝑑t^{}w^2(t^{})}|\varphi }`$ (15c)
$`=`$ $`{\displaystyle D^{}w|w𝒯e^{(4\lambda )^1_0^t𝑑t^{}[w(t^{})2\lambda A(t^{})]^2}|\varphi e^{(4\lambda )^1\{_{\mathrm{}}^0𝑑t^{}w^2(t^{})+_t^{\mathrm{}}𝑑t^{}w^2(t^{})\}}}`$ (15d)
A way to look at this is as a continuous measurement sequence. The change of the wave function at time $`t`$ over the interval $`dt`$, according to (15a), resides solely in
$$e^{2i\lambda A(t)i^1/w(t)}e^{(4\lambda )^1dtw^2(t)}=e^{(4\lambda )^1dt[w(t)2\lambda A(t)]^2},$$
which has the form of the von Neumann measurement of $`A(t)`$ by the pointer $`w(t)`$. However, here the pointer uncertainty $`(dt)^1`$ is quite large. It is only the accumulation of a large number of measurements which can distinguish the eigenvalues of $`A(t)`$. One can visualize the collection of pointers for all t (and for all $`(𝐱,t)`$ in the case of full-blown CSL), each waiting, patiently, for its brief opportunity to perform its measurement.
### III.2 Equivalence of CSL and CQC
The time-ordered term in Eq.(15d) is precisely the CSL expression (3b) for the state vector which evolves under the classical noise $`w(t^{})`$ for $`0t^{}t`$. The remaining exponentials are unaltered parts of the vacuum state, one for $`\mathrm{}<t^{}<0`$ which never changes, and one for $`t<t^{}<\mathrm{}`$ which evolves.
In standard quantum theory, a superselection rulesuper obtains at time $`t`$ if thereafter the state vector may be written as a superposition of orthogonal states, each evolving independently of all the others (i.e., the matrix elements of the Hamiltonian between these states vanishes). In such a case, the state vector may be regarded as a ”catalog,” as discussed in the Introduction, merely a convenient device for listing the different states and their probabilities, since any state’s evolution after time t can be separately described.
In CQC, the ensemble vector admits a superselection rule at any time. Consider the superselected states at time $`t`$. The labels which distinguish them are the values of $`w(t^{})`$ for $`0t^{}t`$. The differently labeled states evolve independently thereafter. Each state corresponds to a different CSL state. The values of $`w(t^{})`$ for $`\mathrm{}<t^{}<0`$ and $`t<t^{}\mathrm{}`$ are irrelevant as far as the particle physics over the interval $`(0,t)`$ is concerned. The former never becomes relevant. For the latter, as time increases, say from $`t`$ to $`t_1`$, the values of $`w(t^{})`$ for $`t<t^{}<t_1`$ become relevant: a superselected state at time $`t`$ splits into many superselected states over the interval $`(t,t_1)`$ (in a way that a many-worlds advocate might admire). Thus, as far as the particle physics is concerned, the ensemble vector may be regarded as merely a convenient device for listing the different CSL states, and this is the sense in which CSL and CQC are equivalent.
The evolution (15) is unitary and guarantees conservation of probability if the usual quantum rule is employed, that the probability of realization of the state $`|ww|\psi ,t`$ is $`D^{}w|w|\psi ,t|^2`$. As far as the particle physics is concerned at time $`t`$, one may integrate the probability of a superselected set over $`w(t^{})`$ for $`t^{}`$ outside the range $`0t^{}t`$: the result is the CSL probability $`P_t(w)Dw`$.
### III.3 Full-Blown CQC
Full-blown CQC corresponding to full-blown CSL described in Eqs.(8),(9), simply replaces time integrals by space-time integrals, e.g.,
$$W(x)(\lambda )^{1/2}/(2\pi )^2d^4k[e^{ikx}b(k)+e^{ikx}b^{}(k)]$$
$$|\psi ,t𝒯e^{2i\lambda _0^td^4x^{}A(x^{})\mathrm{\Pi }(x^{})}|0|\varphi $$
replace Eqs.(10a), (15a), etc.
### III.4 Nonwhite Noise CQC
A CSL collapse evolution which generalizes Eq.(3b)nonwhite is
$$|\psi ,t=𝒯e^{(4\lambda )^1_0^t𝑑t^{}𝑑t^{\prime \prime }[w(t^{})2\lambda A(t^{})]g(t^{}t^{\prime \prime })[w(t\mathrm{"})2\lambda A(t\mathrm{"})]}|\varphi ,$$
where $`g(t)`$ is real, symmetric in $`t`$ and has a positive spectrum, i.e., $`g(t)=𝑑\omega \mathrm{exp}(i\omega t)h^2(\omega )`$ with $`h(\omega )=h(\omega )=h^{}(\omega )`$.
A CQC counterpart requires defining $`W(t)`$ and $`\mathrm{\Pi }(t)`$ as in Eqs.(10), except that the square bracket in the integral of Eq.(10a) (Eq.(10b)) is multiplied by $`h^1(\omega )`$ ($`h(\omega )`$). This replaces the equation before (13) and Eq.(13a) respectively by
$$\left[_{\mathrm{}}^{\mathrm{}}𝑑t^{}g(tt^{})w(t^{})+2\lambda \delta /\delta w(t)\right]w|0=0,w|0=e^{(4\lambda )^1_{\mathrm{}}^{\mathrm{}}𝑑t^{}𝑑t^{\prime \prime }w(t^{})g(t^{}t^{\prime \prime })w(t^{\prime \prime })}$$
With the evolution equation the same as Eq.(15a), Eq.(15d) is replaced by
$$|\psi ,t=𝒯e^{(4\lambda )^1_{\mathrm{}}^{\mathrm{}}𝑑t^{}𝑑t^{\prime \prime }[w(t^{})2\lambda A(t^{})\mathrm{\Theta }_{0,t}(t^{})]g(t^{}t^{\prime \prime })[w(t^{\prime \prime })2\lambda A(t^{\prime \prime })\mathrm{\Theta }_{0,t}(t^{\prime \prime })]}|\varphi $$
where $`\mathrm{\Theta }_{0,t}(t^{})=1`$ for $`0t^{}t`$ and $`=0`$ otherwise. For general $`h(\omega )`$, this equation is not the same as its CSL counterpart written above, except in the S-matrix limit where the interval $`(0,t)`$ is replaced by $`(\mathrm{},\mathrm{})`$. This generalization shall not be considered any further in this paper.
## IV Energy
### IV.1 Pictures
The CSL Schrödinger picture state vector (1) has a corresponding CQC ensemble vector:
$`|\psi ,t_S`$ $``$ $`e^{iH_At}|\psi ,t`$ (16a)
$`=`$ $`𝒯e^{i_0^t𝑑t^{}[H_A+2\lambda A\mathrm{\Pi }(t^{})]}|0|\varphi ,`$ (16b)
where (16b) follows from putting (15a) into (16a). Although $`A`$’s time dependence in (15a) is removed in (16b), this is still not the Schrödinger picture for CQC’s ensemble vector (nonetheless, we’ll keep the $`S`$ subscript on it), since the $`\mathrm{\Pi }`$ operator still has time dependence. It can be removed in what we’ll call the sub-Schrödinger picture:
$`|\psi ,t_{SS}`$ $``$ $`e^{iH_wt}|\psi ,t_S`$ (17a)
$`=`$ $`e^{it[H_w+H_A+2\lambda A\mathrm{\Pi }(0)]}|0|\varphi ,`$ (17b)
Thus, since this is a unitary evolution, there is a conserved energy, the Hamiltonian in Eq.(17b). In the collapse-interaction picture it is
$$H(t)=e^{i(H_A+H_W)t}[H_A+H_W+2\lambda A\mathrm{\Pi }(0)]e^{i(H_A+H_W)t}=H_A+H_W+2\lambda A(t)\mathrm{\Pi }(t).$$
It should be emphasized that, since $`H_w`$ and the interaction term depend upon $`\mathrm{\Pi }`$ which has nonvanishing matrix elements between different eigenstates of $`W`$, the conserved energy belongs to the ensemble of realizable states, not to an individual state. A single collapse to the realized state generally will not preserve the energy distribution. (Of course, this is also the case in standard quantum theory with the collapse postulate.) Nonetheless, it is interesting to consider the probability distributions of the ensemble total energy and its parts. A conserved quantity gives insight to the behavior of the ensemble of possibilities.
For example, one may see how the the ensemble average increase of particle energy, due to the localization mechanism of CSL, is compensated by the ensemble average decrease of the random field energy. And, laboratory experiments often do involve an ensemble of (approximately) identical initial states, for which therefore energy is (approximately) conserved for the ensemble of outcome states. Also interesting are circumstances where one can speak of conservation of energy associated with the ensemble of states that collapses to a distinctly identifiable outcome. For example, when considering a situation when the initial state vector is a superposition of particles in widely separated locations, the ongoing orthogonality of the differently located particle states ensures energy conservation for the ensemble of states which describe collapse to particles in one of the locations.
### IV.2 Simplest Model, $`H_A=0`$
We display here generating functions and probability distributions for the total energy and its components, for the model of section IIa, in the simplest case, $`H_A=0`$.
Using the expression for $`|\psi ,t`$ given in Eq.(15a), the generating function for the moments of the total energy is
$`\psi ,t|e^{iH(t)\beta }|\psi ,t`$ $`=`$ $`{\displaystyle \underset{n}{}}|\alpha _n|^20|e^{2i\lambda a_n_0^tdt^{}\mathrm{\Pi }(t^{})]}e^{i[H_w+2\lambda a_n\mathrm{\Pi }(t)]\beta }e^{2i\lambda a_n_0^tdt^{}\mathrm{\Pi }(t^{})]}|0`$ (18a)
$`=`$ $`{\displaystyle \underset{n}{}}|\alpha _n|^20|e^{i[H_w+2i\lambda a_n_0^t𝑑t^{}i\dot{\mathrm{\Pi }}(t^{})+2\lambda a_n\mathrm{\Pi }(t)]\beta }|0`$ (18b)
$`=`$ $`{\displaystyle \underset{n}{}}|\alpha _n|^20|e^{i[H_w+2\lambda a_n\mathrm{\Pi }(0)]\beta }|0`$ (18c)
$`=`$ $`{\displaystyle \underset{n}{}}|\alpha _n|^20|e^{iH_w\beta }e^{2i\lambda a_n_0^\beta 𝑑t^{}\mathrm{\Pi }(t^{})}|0`$ (18d)
$`=`$ $`{\displaystyle \underset{n}{}}|\alpha _n|^2{\displaystyle D^{}\pi e^{2\lambda _{\mathrm{}}^{\mathrm{}}\pi ^2(t^{})}e^{2i\lambda a_n_0^\beta 𝑑t^{}\pi (t^{})}}`$ (18e)
$`=`$ $`{\displaystyle \underset{n}{}}|\alpha _n|^2e^{(\lambda /2)a_n^2|\beta |}.`$ (18f)
(18b) is the result of performing the unitary transformation on $`e^{iH(t)\beta }`$ in (18a), the integration in the exponent of (18b) results in (18c) which confirms that the generating function is time-independent, (18d) expresses the unitary operator in (18c) in a different form, (18e) follows from $`H_w|0=0`$ and insertion of the identity $`1=D^{}\pi |\pi \pi |`$, (18f) follows by discretization and performance of the gaussian integrals.
The moments of the energy follow from Eq.(18f) using the formal Taylor series expansion, $`|\beta |=0+\beta ϵ(0)+(\beta ^2/2!)2\delta (0)+(\beta ^3/3!)2\delta ^{}(0)+\mathrm{}`$ (where $`ϵ(x)x/|x|`$ for $`x0`$ and $`ϵ(0)0`$). Thus the mean energy is 0 and the higher moments of the energy are infinite. The associated probability distribution more transparently entails this result:
$`𝒫(E)`$ $``$ $`{\displaystyle \frac{1}{2\pi }}{\displaystyle _{\mathrm{}}^{\mathrm{}}}e^{iE\beta }\psi ,t|e^{iH(t)\beta }|\psi ,t`$ (19a)
$`=`$ $`{\displaystyle \frac{1}{\pi }}{\displaystyle \underset{n}{}}|\alpha _n|^2{\displaystyle \frac{\lambda a_n^2/2}{E^2+(\lambda a_n^2/2)^2}}.`$ (19b)
Although the total energy probability distribution is time-independent, the energy $`H_w`$ is not. The generating function here is
$`\psi ,t|e^{iH_w\beta }|\psi ,t`$ $`=`$ $`{\displaystyle \underset{n}{}}|\alpha _n|^20|e^{2i\lambda a_n_0^t𝑑t^{}\mathrm{\Pi }(t^{})}e^{2i\lambda a_n_0^t𝑑t^{}\mathrm{\Pi }(t^{}\beta )}|0`$ (20a)
$`=`$ $`{\displaystyle \underset{n}{}}|\alpha _n|^2{\displaystyle D^{}\pi e^{2\lambda _{\mathrm{}}^{\mathrm{}}\pi ^2(t^{})}e^{2i\lambda a_n[_0^t𝑑t_\beta ^{t\beta }]dt^{}\pi (t^{})}}`$ (20c)
$`=`$ $`{\displaystyle \underset{n}{}}|\alpha _n|^2\{[\mathrm{\Theta }(\beta t)+\mathrm{\Theta }(t\beta )]e^{\lambda a_n^2t}`$
$`+\mathrm{\Theta }(t\beta )\mathrm{\Theta }(t+\beta )e^{\lambda a_n^2|\beta |}\}.`$
where (20a) utilizes (15a) and the time-translation capability of $`H_w`$, (20b) involves going to the $`|\pi `$ representation, and (20c) follows from performing the gaussian integrals. From the Taylor series expansion of (20c), one sees that the mean value of $`H_w`$ is zero and the higher moments are infinite.
The probability distribution following from Eq.(20c) confirms this result:
$`𝒫_w(E)`$ $``$ $`{\displaystyle \frac{1}{2\pi }}{\displaystyle _{\mathrm{}}^{\mathrm{}}}e^{iE\beta }\psi ,t|e^{iH_w\beta }|\psi ,t`$ (21b)
$`=`$ $`{\displaystyle \frac{1}{\pi }}{\displaystyle \underset{n}{}}|\alpha _n|^2\{e^{\lambda a_n^2t}[\pi \delta (E){\displaystyle \frac{\mathrm{sin}Et}{E}}+{\displaystyle \frac{E\mathrm{sin}Et\lambda a_n^2\mathrm{cos}Et}{E^2+(\lambda a_n^2)^2}}]`$
$`+{\displaystyle \frac{\lambda a_n^2}{E^2+(\lambda a_n^2)^2}}\}.`$
Eq. (21) shows how the $`W`$-field energy, starting at $`t=0`$ with the vacuum energy $`E=0`$, evolves toward a stationary energy distribution associated with the completely collapsed state. (Note that the asymptotic value of (21b) is not the same as (19b), although it is similar in form.)
The interaction energy’s probability distribution is
$`𝒫_I(E)`$ $``$ $`{\displaystyle \frac{1}{2\pi }}{\displaystyle _{\mathrm{}}^{\mathrm{}}}e^{iE\beta }\psi ,t|e^{iH_I\beta }|\psi ,t`$ (22a)
$`=`$ $`{\displaystyle \frac{1}{2\pi }}{\displaystyle _{\mathrm{}}^{\mathrm{}}}e^{iE\beta }{\displaystyle \underset{n}{}}|\alpha _n|^2{\displaystyle D^{}\pi e^{2\lambda _{\mathrm{}}^{\mathrm{}}𝑑t^{}\pi ^2(t^{})}e^{2i\lambda a_n\pi (t)\beta }}`$ (22b)
$`=`$ $`{\displaystyle \underset{n}{}}{\displaystyle \frac{|\alpha _n|^2}{\sqrt{2\pi \lambda a_n^2/dt}}}e^{E^2/(2\lambda a_n^2/dt)}`$ (22c)
Since the interaction energy is proportional to the white noise operator $`\mathrm{\Pi }(t)`$, it has a white noise energy spectrum, with all energies equally likely.
### IV.3 Simplest Model, $`H_A0`$
#### IV.3.1 Moment Generating Functions
As in the previous section, also in this more general case, the random field variables can be integrated out of the expressions for the moment generating functions. For the total energy, corresponding to Eq.(18), one obtains:
$`\psi ,t|e^{iH(t)\beta }|\psi ,t`$ $`=`$ $`\varphi |0|e^{i[H_A+H_w+2\lambda \mathrm{\Pi }(0)A]\beta }|0|\varphi `$ (23a)
$`=`$ $`\varphi |0|e^{i[H_A+H_w]\beta }𝒯e^{2i\lambda _0^\beta 𝑑t^{}\mathrm{\Pi }(t^{})A(t^{})}|0|\varphi `$ (23b)
$`=`$ $`\varphi |e^{iH_A\beta }{\displaystyle D^{}\pi e^{2\lambda _{\mathrm{}}^{\mathrm{}}𝑑t^{}\pi ^2(t^{})}𝒯e^{2i\lambda _0^\beta 𝑑t^{}\pi (t^{})A(t^{})}}|\varphi `$ (23c)
$`=`$ $`\varphi |e^{iH_A\beta }𝒯e^{(\lambda /2)_0^\beta 𝑑t^{}A^2(t^{})}|\varphi `$ (23d)
$`=`$ $`\varphi |e^{[iH_A+(\lambda /2)A^2(0)]\beta }|\varphi .`$ (23e)
((23a) is expressed in the manifestly time-independent sub-Schrödinger picture, and the remaining steps are ones taken before.) For $`H_A`$, the moment generating function is
$`\psi ,t|e^{iH_A\beta }|\psi ,t`$ $`=`$ $`\varphi |0|𝒯e^{2i\lambda _0^t𝑑t^{}\mathrm{\Pi }(t^{})A_L(t^{})}e^{iH_A\beta }e^{2i\lambda _0^t𝑑t^{}\mathrm{\Pi }(t^{})A_R(t^{})}|0|\varphi `$ (24a)
$`=`$ $`\varphi |𝒯\left\{e^{(\lambda /2)_0^t𝑑t^{}[A_L(t^{})A_R(t^{})]^2}e^{iH_A\beta }\right\}|\varphi ,`$ (24b)
where $`A_R`$ ($`A_L`$) acts to the right (left) of $`\mathrm{exp}iH_A\beta `$, and is time-ordered (time–reverse-ordered) by $`𝒯`$.
For $`H_w`$, the calculation is
$`\psi ,t|e^{iH_w\beta }|\psi ,t`$ $`=\varphi |0|𝒯e^{2i\lambda _0^t𝑑t^{}\mathrm{\Pi }(t^{})A_L(t^{})}e^{iH_w\beta }e^{2i\lambda _0^t𝑑t^{}\mathrm{\Pi }(t^{})A_R(t^{})}|0|\varphi `$ (25d)
$`=\varphi |0|𝒯e^{2i\lambda _0^t𝑑t^{}\mathrm{\Pi }(t^{}+\beta /2)A_L(t^{})}e^{2i\lambda _0^t𝑑t^{}\mathrm{\Pi }(t^{}\beta /2)A_R(t^{})}|0|\varphi `$
$`=\varphi |{\displaystyle D^{}\pi e^{2\lambda _{\mathrm{}}^{\mathrm{}}𝑑t^{}\pi ^2(t^{})}}`$
$`𝒯e^{2i\lambda _{\beta /2}^{t+\beta /2}𝑑t^{}\pi (t^{})A_L(t^{}\beta /2)}e^{2i\lambda _{\beta /2}^{t\beta /2}𝑑t^{}\pi (t^{})A_R(t^{}+\beta /2)}|\varphi `$
$`=\varphi |𝒯e^{(\lambda /2)_0^t𝑑t^{}[A_L^2(t^{})+A_R^2(t^{})]}[\mathrm{\Theta }(\beta t)+\mathrm{\Theta }(\beta t)`$
$`+\mathrm{\Theta }(t\beta )\mathrm{\Theta }(t+\beta )e^{\lambda _{|\beta |/2}^{t|\beta |/2}𝑑t^{}A_L(t^{}\beta /2)A_R(t^{}+\beta /2)}]|\varphi `$
#### IV.3.2 First Moments
Here is shown conservation of energy in the mean, reproducing more simply a previously obtained result (all that that had been proved about conservation of energy for CSLenergy ). Equating the terms proportional to the first power of $`\beta `$ in Eq,(24b) produces
$$\psi ,t|H_A|\psi ,t=\varphi |𝒯e^{(\lambda /2)_0^t𝑑t^{}[A_L(t^{})A_R(t^{})]^2}H_A|\varphi .$$
(26)
To find the first moment of $`H_w`$, the exponent in (25d) must be written for small $`\beta `$:
$`{\displaystyle _{|\beta |/2}^{t|\beta |/2}}𝑑t^{}A_L(t^{}\beta /2)A_R(t^{}+\beta /2)={\displaystyle _0^t}𝑑t^{}A_L(t^{})A_R(t^{})`$
$`(|\beta |/2)[A_L(t)A_R(t)+A_L(0)A_R(0)]+(\beta /2){\displaystyle _0^t}𝑑t^{}[A_L(t^{})\dot{A}_R(t^{})\dot{A}_L(t^{})A_R(t^{})]+\mathrm{}`$ (27)
The zeroth and first order terms in the Taylor series expansion of (27) are the same as (27), except that the second term vanishes, since $`d|\beta |/dt|_{\beta =0}=ϵ(0)=0`$. The step-functions in (25d) likewise do not contribute to the Taylor series expansion’s first two terms, so we get from (25d):
$`\psi ,t|H_w|\psi ,t`$ $`=`$ $`i(\lambda /2)\varphi |𝒯e^{(\lambda /2)_0^t𝑑t^{}[A_L(t^{})A_R(t^{})]^2}`$ (28a)
$`{\displaystyle _0^t}dt^{}[A_L(t^{})\dot{A}_R(t^{})\dot{A}_L(t^{})A_R(t^{})]|\varphi `$
$`=`$ $`(\lambda /2)\varphi |𝒯{\displaystyle _0^t}𝑑t^{}e^{(\lambda /2)_0^t^{}𝑑t^{\prime \prime }[A_L(t^{\prime \prime })A_R(t^{\prime \prime })]^2}[A(t^{}),i\dot{A}(t^{})]|\varphi `$ (28b)
$`=`$ $`(\lambda /2)\varphi |𝒯{\displaystyle _0^t}𝑑t^{}e^{(\lambda /2)_0^t^{}𝑑t^{\prime \prime }[A_L(t^{\prime \prime })A_R(t^{\prime \prime })]^2}[A(t^{}),[A(t^{}),H_A]]|\varphi `$ (28c)
In going from (28a) to (28b) we have utilized
$$𝒯e^{(\lambda /2)_0^t𝑑t^{}[A_L(t^{})A_R(t^{})]^2}B_L(t_1)C_R(t_1)=𝒯e^{(\lambda /2)_0^{t_1}𝑑t^{}[A_L(t^{})A_R(t^{})]^2}B(t_1)C(t_1)$$
for $`t_1<t`$ and arbitrary operators $`B_L(t_1)`$, $`C_R(t_1)`$ (the subscripts $`R`$, $`L`$ may be dispensed with when operators having the same time argument are adjacent and in the correct order).
Since the first moment of the interaction energy vanishes, from Eqs.(26), (28c) follows conservation of energy in the mean:
$`(d/dt)\psi ,t|H_w|\psi ,t`$ $`=`$ $`(d/dt)\psi ,t|H_A|\psi ,t`$ (29)
$`=`$ $`(\lambda /2)\varphi |𝒯e^{(\lambda /2)_0^t𝑑t^{}[A_L(t^{})A_R(t^{})]^2}[A(t),[A(t),H_A]]|\varphi `$
#### IV.3.3 Large t
The large t-behavior of the probability density distribution of $`H_w`$ is due to the third term in the bracket of Eq.(25d). Because this term’s exponent is largest for small $`\beta `$, and also because $`\mathrm{exp}iE\beta `$ oscillates rapidly for large $`\beta `$, the small $`\beta `$ approximation of (25d) should be a good approximation to the large t limit. Putting (27) into (25d) results in
$`𝒫_w(E)={\displaystyle \frac{1}{2\pi }}{\displaystyle _{\mathrm{}}^{\mathrm{}}}e^{iE\beta }\psi ,t|e^{iH_w\beta }|\psi ,t`$ (30a)
$`{\displaystyle \frac{1}{2\pi }}{\displaystyle _t^t}𝑑\beta e^{iE\beta }\varphi |𝒯e^{(\lambda /2)_0^t𝑑t^{}[A_L(t^{})A_R(t^{})]^2}e^{(\lambda /2)|\beta |[A^2(t)+A^2(0)]}e^{i(\lambda /2)\beta _0^t𝑑t^{}C_{LR}(t^{})}|\varphi `$ (30b)
$`=\varphi |𝒯e^{(\lambda /2)_0^t𝑑t^{}[A_L(t^{})A_R(t^{})]^2}{\displaystyle \frac{1}{\pi }}{\displaystyle \frac{\lambda [A^2(t)+A^2(0)]/2}{[E(\lambda /2)_0^t𝑑t^{}C_{LR}(t^{})]^2+[\lambda (A^2(t)+A^2(0))/2]^2}}|\varphi `$
where $`C_{LR}(t^{})A_L(t^{})[A_R(t^{}),H_A][A_L(t^{}),H_A]A_R(t^{})`$. Note that, when $`H_A=0`$, the large t-limit of (30c) is the same as (21b). Also, note that, when $`H_A0`$, the mean value of $`E`$ given by (30c) is the same as (28c).
For example, in the case of a free particle of mass $`m`$ which collapses toward a position eigenstate, where $`H=p^2/2m`$ and $`A=x`$, the term in (30c) $`[E(\lambda /2)_0^t𝑑t^{}C_{LR}(t^{})]^2=[E+\lambda \mathrm{}t/2m]^2`$. Because the ensemble energy is conserved, the mean $`W`$-field energy decreases with time to compensate for the increase of the particle’s mean energy $`=\lambda \mathrm{}t/2m`$ due to the narrowing of wavepackets by the collapse mechanism.
## V Time Operator
Unlike the usual Hamiltonian which is semi-boundedtimeref , $`H_w`$ has an unbounded spectrum, providing an opportunity to define a self-adjoint time operator:
$`T`$ $``$ $`{\displaystyle \frac{_{\mathrm{}}^{\mathrm{}}𝑑\omega b^{}(\omega )\frac{1}{i}\frac{}{\omega }b(\omega )}{_{\mathrm{}}^{\mathrm{}}𝑑\omega b^{}(\omega )b(\omega )}}{\displaystyle \frac{B}{N}}`$ (31a)
$`=`$ $`{\displaystyle \frac{_{\mathrm{}}^{\mathrm{}}𝑑t^{}t^{}[(4\lambda )^1W^2(t^{})+\lambda \mathrm{\Pi }^2(t^{})]}{_{\mathrm{}}^{\mathrm{}}𝑑t^{}[(4\lambda )^1W^2(t^{})+\lambda \mathrm{\Pi }^2(t^{})]}},`$ (31b)
where $`[B,N]=0`$. It is simple to show that $`[H_w,T]=i`$, where $`H_w`$ is given by Eq.(14). The action of $`T`$ on $`|0`$, is undefined, but one can take $`T=B/(N+ϵ)`$ so that $`T|0=0/ϵ=0`$. Of course, $`T`$ plus any function of $`H_w`$ and $`N`$ is also conjugate to $`H_w`$, but (31) is the simplest choice.
### V.1 Eigenfunctions
$`T`$’s eigenvalues are highly degenerate. A complete set of (unnormalized) eigenfunctions in the $`|w`$ basis is:
$`w|{\displaystyle \underset{t}{}}[{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\omega e^{i\omega t}b^{}(\omega )]^{n(t)}|0`$ $``$ $`{\displaystyle \underset{t}{}}[w(t)2\lambda \delta /\delta \omega ]^{n(t)}e^{(4\lambda )^1_{\mathrm{}}^{\mathrm{}}𝑑t^{}w^2(t^{})}`$
$``$ $`{\displaystyle \underset{t}{}}H_{n(t)}[w(t)(dt/2\lambda )^{1/2}]e^{(4\lambda )^1_{\mathrm{}}^{\mathrm{}}𝑑t^{}w^2(t^{})}`$
where $`n(t)=0,1,2\mathrm{}`$ and the $`H_n`$ are Hermite polynomials. The associated eigenvalues $`_{t=\mathrm{}}^{\mathrm{}}tn(t)/_{t=\mathrm{}}^{\mathrm{}}n(t)`$ have a “center of time” form.
### V.2 Moment Generating Function
To find the moment generating function for $`T`$, first observe, using (15a) and (10b), that
$`e^{i\beta T}|\psi ,t=e^{i\beta T}𝒯e^{(\lambda /2\pi )^{1/2}_0^t𝑑t^{}A_R(t^{})_{\mathrm{}}^{\mathrm{}}𝑑\omega e^{i\omega t^{}}b^{}(\omega )}|0e^{(\lambda /2)_0^t𝑑t^{}A_R^2(t^{})}|\varphi `$
$`=𝒯\left[1+{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}e^{i\beta B/n}{\displaystyle \frac{(\lambda /2\pi )^{n/2}}{n!}}\left[{\displaystyle _0^t}𝑑t^{}A_R(t^{}){\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\omega e^{i\omega t^{}}b^{}(\omega )\right]^n\right]|0e^{(\lambda /2)_0^t𝑑t^{}A_R^2(t^{})}|\varphi .`$
Since $`\mathrm{exp}(i\beta B/n)b^{}(\omega )\mathrm{exp}(i\beta B/n)=b^{}(\omega +\beta /n)`$, one gets:
$`\psi ,t|e^{i\beta T}|\psi ,t`$ $`=`$ $`\varphi |𝒯e^{(\lambda /2)_0^t𝑑t^{}[A_R^2(t^{})+A_L^2(t^{})]}0|\left[1+{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(\lambda /2\pi )^{n/2}}{n!}}\left[{\displaystyle _0^t}𝑑t^{}A_L(t^{}){\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\omega e^{i\omega t^{}}b(\omega )\right]^n\right]`$ (32a)
$`\left[1+{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(\lambda /2\pi )^{n/2}}{n!}}\left[{\displaystyle _0^t}𝑑t^{}A_R(t^{}){\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\omega e^{i(\omega n^1\beta )t^{}}b^{}(\omega )\right]^n\right]|0|\varphi `$
$`=`$ $`\varphi |𝒯e^{(\lambda /2)_0^t𝑑t^{}[A_R^2(t^{})+A_L^2(t^{})]}\left[1+{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(\lambda )^n}{n!}}\left[{\displaystyle _0^t}𝑑t^{}A_L(t^{})A_R(t^{})e^{i\beta t^{}/n}\right]^n\right]|\varphi `$ (32b)
As a check, when $`\beta =0`$, (32b) becomes:
$$\varphi |𝒯e^{(\lambda /2)_0^t𝑑t^{}[A_R^2(t^{})+A_L^2(t^{})]}e^{\lambda _0^t𝑑t^{}[A_R(t^{})A_L(t^{})]}|\varphi =\varphi |𝒯e^{(\lambda /2)_0^t𝑑t^{}[A_R(t^{})A_L(t^{})]^2}|\varphi =1.$$
### V.3 First Moments
Now, extract the first two moments of $`T`$ from (32b). From the coefficient of $`\beta `$ in (32b) we obtain:
$`\psi ,t|T|\psi ,t`$ $`=`$ $`\varphi |𝒯{\displaystyle \frac{_0^t𝑑t^{}t^{}A_L(t^{})A_R(t^{})}{_0^t𝑑t^{}A_L(t^{})A_R(t^{})}}`$ (33)
$`\left[e^{(\lambda /2)_0^t𝑑t^{}[A_L(t^{})A_R(t^{})]^2}e^{(\lambda /2)_0^t𝑑t^{}[A_L^2(t^{})+A_R^2(t^{})]}\right]|\varphi `$
To get a feel for the meaning of (33), consider the case $`H_A=0`$:
$`\psi ,t|T|\psi ,t={\displaystyle \underset{k}{}}|\alpha _k|^2{\displaystyle \frac{_0^t𝑑t^{}t^{}}{_0^t𝑑t}}[1e^{\lambda ta_k^2}]={\displaystyle \frac{t}{2}}[1{\displaystyle \underset{k}{}}|\alpha _k|^2e^{\lambda ta_k^2}].`$ (34)
Thus, asymptotically, $`\psi ,t|T|\psi ,t`$ approaches the “center of time,” a weighted average (here uniformly weighted) of the time over the interval $`(0,t)`$. For $`H_A0`$ much the same behavior prevails, according to Eq.(33). For example, set $`A_R(t^{})=A_L(t^{})=At^s`$. Then, from (33), asymptotically, $`\psi ,t|T|\psi ,tt(2s+1)/(2s+2)`$: for $`s`$ large, the mean value of $`T`$ is almost $`t`$.
For full-blown CQC, the expression that replaces (33) is the same expression with $`A(t^{})A(𝐱,t^{})`$ and $`dt^{}d𝐱dt^{}`$. $`A`$ (given by Eq. (8)) is $``$ the mass density in a sphere of radius $`a`$. Then the appropriate “center of time ” expression appearing in the equivalent of (33) is
$$_0^t𝑑t^{}𝑑𝐱t^{}A_L(𝐱,t^{})A_R(𝐱,t^{})/_0^t𝑑t^{}𝑑𝐱A_L(𝐱,t^{})A_R(𝐱,t^{}),$$
i.e., $`t^{}`$, roughly speaking, is weighted by the squared mass density integrated over all space and over time from 0 to $`t`$. Should there be a high rate of agglomeration of matter, e.g., due to gravitational clumping, then $`\psi ,t|T|\psi ,t`$ will approach $`t`$. It should be emphasized that gravitational collapse has this effect, not wave function collapse, as Eq. (34) attests.
The second moment of $`T`$ is likewise found from (32b):
$`\psi ,t|T^2|\psi ,t`$ $`=`$ $`\varphi |𝒯\{\left[{\displaystyle \frac{_0^t𝑑t^{}t^{}A_L(t^{})A_R(t^{})}{_0^t𝑑t^{}A_L(t^{})A_R(t^{})}}\right]^2`$ (35)
$`\left[e^{(\lambda /2)_0^t𝑑t^{}[A_L(t^{})A_R(t^{})]^2}e^{(\lambda /2)_0^t𝑑t^{}[A_L^2(t^{})+A_R^2(t^{})]}[1+f(Z_{LR})]\right]`$
$`+{\displaystyle \frac{_0^t𝑑t^{}t^2A_L(t^{})A_R(t^{})}{_0^t𝑑t^{}A_L(t^{})A_R(t^{})}}e^{(\lambda /2)_0^t𝑑t^{}[A_R^2(t^{})+A_L^2(t^{})]}f(Z_{LR})\}|\varphi ,`$
where $`Z_{LR}\lambda _0^t𝑑t^{}A_L(t^{})A_R(t^{})`$ and $`f(z)_0^z𝑑z^{}[\mathrm{exp}z^{}1]/z^{}`$. For the example $`A_R(t^{})=A_L(t^{})=At^s`$ ($`H_A=0`$ corresponds to $`s=0`$), Eq.(35) asymptotically becomes (using $`f(z)z^1\mathrm{exp}z`$):
$$\psi ,t|T^2|\psi ,t\left[\frac{t(2s+1)}{2s+2}\right]^2\left[1+\frac{1}{\lambda t^{2s+1}}\frac{(2s+2)^2}{2s+3}\underset{k}{}\frac{|\alpha _k|^2}{a_k^2}\right].$$
(36)
so the fractional deviation decreases over time, $`\overline{\mathrm{\Delta }T^2}/\overline{T}^2t^{(2s+1)}`$.
### V.4 Large t
More generally, an approximation to the large $`t`$ probability distribution of $`T`$ may be had by employing the small $`\beta `$ approximation in (32b),
$`[{\displaystyle _0^t}𝑑t^{}A_L(t^{})A_R(t^{})e^{i\beta t^{}/n}]^n`$ $``$ $`[{\displaystyle _0^t}𝑑t^{}A_L(t^{})A_R(t^{})(1i\beta t^{}/n)]^n`$
$``$ $`[{\displaystyle _0^t}𝑑t^{}A_L(t^{})A_R(t^{})]^ne^{i\beta _0^t𝑑t^{}t^{}A_L(t^{})A_R(t^{})/_0^t𝑑t^{}A_L(t^{})A_R(t^{})},`$
with the result
$`𝒫(\tau )`$ $``$ $`(2\pi )^1{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\beta e^{i\beta \tau }\psi ,t|e^{i\beta \tau }|\psi ,t`$ (37)
$``$ $`\varphi |𝒯e^{(\lambda /2)_0^t𝑑t^{}[A_L(t^{})A_R(t^{})]^2}\delta \left[\tau {\displaystyle \frac{_0^t𝑑t^{}t^{}A_L(t^{})A_R(t^{})}{_0^t𝑑t^{}A_L(t^{})A_R(t^{})}}\right]|\varphi .`$
To summarize, there is a natural clock associated with the random field, whose time reading depends upon global dynamics and whose deviance decreases with time. In future work, rather than collapse depending upon time, one may look to express the attractive idea that the notion of time may depend upon collapse, since no collapse implies no events implies no time.
## VI Two examples
Collapse models are phenomenological, a guess (based upon a perceived flaw in standard quantum theory) that nature utilizes this description of dynamical behavior. To be generally accepted, such models need experimental verificationexp and grounding within a larger physical framework. One of the main aims of this paper has been to emphasize that CSL can be formulated as CQC, within an essentially quantum framework, in order to encourage thoughts about how collapse might arise from a physical structure which has a natural quantum theory expression. In that spirit, this section contains two examples of CQC dynamics arising from other quantum conceptual structures.
### VI.1 Quantum Fields
Consider a collection of scalar quantum fields, associated with particles of large mass $`M`$ which, for definiteness, take to be the Planck mass, so large that $`\sqrt{𝐤^2+M^2}M`$ for any so-far experimentally achievable $`𝐤`$. Each field (more precisely, its time derivative, which is its conjugate field) interacts with the mass density $`A(𝐱,t)`$ (given in Eq.(8) et. seq.) over a brief interval which, for definiteness, we shall imagine is the Planck time $`\tau =M^1`$. Call the field $`\mathrm{\Phi }_t(𝐱)`$ (and call its conjugate $`\mathrm{\Pi }_t(𝐱)`$) which is responsible for the interaction over the interval centered at time $`t`$;
$`\mathrm{\Phi }_t(𝐱)`$ $``$ $`C{\displaystyle \frac{1}{(2\pi )^{3/2}\sqrt{2M}}}{\displaystyle \mathrm{𝐝𝐤}[e^{i(𝐤𝐱Mt)}b_t(𝐤)+e^{i(𝐤𝐱Mt)}b_t^{}(𝐤)]}`$ (38a)
$`\mathrm{\Pi }_t(𝐱)`$ $``$ $`C^1{\displaystyle \frac{i}{(2\pi )^{3/2}}}\sqrt{{\displaystyle \frac{M}{2}}}{\displaystyle \mathrm{𝐝𝐤}[e^{i(𝐤𝐱Mt)}b_t(𝐤)+e^{i(𝐤𝐱Mt)}b_t^{}(𝐤)]}`$ (38b)
where the constant $`C`$ is yet to be chosen.
Define the eigenstates by $`\mathrm{\Phi }_t(𝐱)|\phi _t=\phi _t(𝐱)|\phi _t`$ and the associated vacuum state by $`b_t(𝐤)|0_t=0`$. Then, Eqs.(38) imply
$$\left[\frac{M}{C^2}\phi _t(𝐱)+\frac{\delta }{\delta \phi _t(𝐱)}\right]\phi _t(𝐱)|0_t=0,\phi _t(𝐱)|0_t=e^{\frac{M}{2C^2}{\scriptscriptstyle 𝑑𝐱\phi _t(𝐱)^2}}.$$
(39)
The joint vacuum state for all the fields is $`|0_t|0_t`$ and a joint eigenstate of all the fields is $`|\phi _t|\phi _t(𝐱)_t`$. To make their scalar product close to the scalar product in Eq.(13) (suitably generalized to space-time), choose $`C=M\sqrt{2\lambda }`$ so, from (39),
$$\phi |0=e^{(4\lambda )^1_t\tau {\scriptscriptstyle 𝑑𝐱\phi _t(𝐱)^2}}e^{(4\lambda )^1_{\mathrm{}}^{\mathrm{}}𝑑t{\scriptscriptstyle 𝑑𝐱\phi _t(𝐱)^2}}$$
(40)
The evolution equation is chosen as:
$`\phi |\psi ,t`$ $``$ $`𝒯e^{2i\lambda _{t^{}=0}^t{\scriptscriptstyle 𝑑𝐱^{}\mathrm{\Pi }_t^{}(𝐱)A(𝐱^{},t^{})}}|0|\varphi e^{2i\lambda _{t^{}=0}^t𝑑t^{}{\scriptscriptstyle 𝑑𝐱^{}M\mathrm{\Pi }_t^{}(𝐱^{})A(𝐱^{},t^{})}}|0|\varphi `$ (41a)
$`=`$ $`𝒯e^{(4\lambda )^1_{t^{}=0}^t\tau {\scriptscriptstyle 𝑑𝐱^{}[\phi _t^{}(𝐱^{})2\lambda A(𝐱^{},t^{})]^2}}|\varphi e^{(4\lambda )^1[_{t^{}=\mathrm{}}^0+_{t^{}=t}^{\mathrm{}}]\tau {\scriptscriptstyle d(𝐱^{})\phi _t^{}^2(𝐱^{})}}`$ (41b)
To complete the demonstration that this model is equivalent to CQC, identify the fields
$$\stackrel{~}{W}(𝐱,t)\mathrm{\Phi }_t(𝐱),\stackrel{~}{\mathrm{\Pi }}(𝐱,t)M\mathrm{\Pi }_t(𝐱)$$
(42)
from which follows
$$[\stackrel{~}{W}(𝐱,t),\stackrel{~}{\mathrm{\Pi }}(𝐱^{},t^{})]=\delta (𝐱𝐱^{})\delta _{t,t^{}}/\tau \delta (𝐱𝐱^{})\delta (tt^{}).$$
It is immediately seen, using Eq.(42) for $`\stackrel{~}{\mathrm{\Pi }}(𝐱,t)`$, that the statevector evolution equation (41a) is identical to Eq.(15a) (generalized to space-time).
To see that the expressions for $`\stackrel{~}{W}`$ and $`\stackrel{~}{\mathrm{\Pi }}`$ in terms of particle creation and annihilation operators closely approximate the similar expressions for $`W`$ and $`\mathrm{\Pi }`$, note that
$$[e^{iMt}b_t(𝐤)/\sqrt{\tau },e^{iMt}b_t^{}(𝐤)/\sqrt{\tau }]=\delta (𝐤𝐤^{})\delta _{t,t^{}}/\tau \delta (𝐤𝐤^{})\delta (tt^{}).$$
(43)
Therefore, define
$$\stackrel{~}{b}(𝐤,\omega )(2\pi )^{1/2}𝑑t^{}e^{i\omega t^{}}e^{iMt^{}}b_t^{}(𝐤^{})/\sqrt{\tau },$$
(44)
so $`[\stackrel{~}{b}(𝐤,\omega ),\stackrel{~}{b}^{}(𝐤^{},\omega ^{}]=\delta (𝐤𝐤^{})\delta (\omega \omega ^{})`$ follows from Eqs.(43),(44). Inserting the inverse of Eq.(44), for $`b_t(𝐤)`$ in terms of $`\stackrel{~}{b}(𝐤,\omega )`$, into Eqs.(38),(42) gives for $`\stackrel{~}{W}`$ and $`\stackrel{~}{\mathrm{\Pi }}`$ expressions whose form is identical to (10)’s expressions (generalized to space-time) for $`W`$ and $`\mathrm{\Pi }`$.
### VI.2 Many Spins
The model described here borrows from a gravitationally based semi-classical collapse modelpearlesquires and spin-network models of space-timesmolinseth , without being faithful to either. Suppose that space consists of “Planck cells” (each of volume $`\mathrm{}^3`$, where $`\mathrm{}`$ is the Planck length), each containing a spin which is normally rigidly enmeshed with the other spins, but which occasionally breaks loose. Then, for a duration equal to the Planck time $`\tau `$, it interacts both with a thermal bath, (characterized by the energy $`\beta ^1`$) and gravitationally with nearby particles (more on this shortly). Following this brief dynamics it is frozen, no longer interacting with either bath or particles, once more enmeshed with other spins.
Characterize the spin in cell i located at $`𝐱_i`$, which interacts over the interval $`(t\tau ,t)`$ and then freezes at time $`t`$, by $`\sigma _{i,t}`$ (a Pauli z-spin matrix) and mass $`\mu \sigma _{i,t},`$ where $`\mu `$ is the Planck mass. Suppose that a spin-mass at $`𝐱_i`$ interacts gravitationally only with the particles in a sphere of radius $`a`$ around it, and that the interaction is as if the particles are uniformly smeared over that sphere, with resulting mass density represented by the operator $`\rho (𝐱_i)`$. For simplicity, take the particle Hamiltonian to vanish, so that $`\rho `$ is independent of $`t`$. Then, since any initial statevector for the particles can be written as the sum of joint eigenstates of $`\rho (𝐱_i)`$ for all $`𝐱_i`$, calculations for the statevector need consider only the special case of one such term, for which $`\rho (𝐱_i)`$ is a c-number, with the more general case being reconstructed afterward as the sum. Then, for one such term, the Hamiltonian for the gravitational interaction of all the spins which have interacted with the particles over the interval $`(0,t)`$ (up to a numerical factor which, for simplicity, is taken equal to 1) is
$$H=G\mu a^2\underset{i,t^{}=\tau }{\overset{t}{}}\sigma _{i,t^{}}\rho (𝐱_i).$$
(45)
Because all the spins act independently, focus upon a group of spins excited over the time interval $`(t^{},t^{}+\mathrm{\Delta }t)`$ which lie within a volume $`\mathrm{\Delta }V`$. Take $`\mathrm{\Delta }t`$ and $`\mathrm{\Delta }V`$ to be large ($`\mathrm{\Delta }t/\tau >>1`$, $`\mathrm{\Delta }V/\mathrm{}^3>>1`$), large enough to contain $`N`$ interacting spins where $`N>>1`$, but still small on the scale of particle physics (e.g., $`\rho (𝐱_i)`$ is essentially constant over $`\mathrm{\Delta }V`$). Denote by $`p`$ the probability that a spin in a cell is freed up to interact over a time interval $`\tau `$ so, on average,
$$N=p\frac{\mathrm{\Delta }V}{\mathrm{}^3}\frac{\mathrm{\Delta }t}{\tau }.$$
(46)
For simplicity, we shall take (46) to be the exact expression for $`N`$, not just for its average.
Now consider the (interaction picture) statevector $`|\psi ,t`$. It may be written as the direct product of states associated with all volumes $`\mathrm{\Delta }V`$ and all time intervals $`\mathrm{\Delta }t`$ within $`(0,t)`$. The state which is the contribution of one such volume and interval has the form
$$|\chi =\underset{s_1=1}{\overset{1}{}}\mathrm{}\underset{s_N=1}{\overset{1}{}}|s_1\mathrm{}s_N|\text{bath},s_1\mathrm{}s_NC_{s_1\mathrm{}s_N}(\rho ).$$
(47)
The sum in Eq.(47) contains all the spin states entangled with the associated bath states which interacted with the spins and brought them to thermal equilibrium: assume for simplicity that the bath states are mutually orthogonal.
Then, $`|\chi `$’s contribution to the density matrix describing the spins alone is
$$\text{Tr}_{\text{bath}}|\chi \chi |=\underset{s_1=1}{\overset{1}{}}\mathrm{}\underset{s_N=1}{\overset{1}{}}|s_1\mathrm{}s_Ns_1\mathrm{}s_N||C_{s_1\mathrm{}s_N}(\rho )|^2.$$
On the other hand, because these states are all in thermal equilibrium, using Eq.(45)’s Hamiltonian for these spins, the thermal density matrix is:
$$\underset{s_1=1}{\overset{1}{}}\mathrm{}\underset{s_N=1}{\overset{1}{}}|s_1\mathrm{}s_Ns_1\mathrm{}s_N|\frac{e^{\beta G\rho \mu a^2_{i=1}^Ns_i}}{_{s_1=1}^1\mathrm{}_{s_N=1}^1e^{\beta G\rho \mu a^2_{i=1}^Ns_i}}.$$
Thus, by comparison of these two equations,
$$|C_{s_1\mathrm{}s_N}(\rho )|^2=\frac{e^{\beta G\rho \mu a^2_{i=1}^Ns_i}}{_{s_1=1}^1\mathrm{}_{s_N=1}^1e^{\beta G\rho \mu a^2_{i=1}^Ns_i}}.$$
(48)
Eventually, $`W`$ associated with this space-time region will be taken $`S_{i=1}^N\sigma _i`$. The statevector which is the sum of all states in (47) which have the same eigenvalue of $`S`$ and which is normalized to 1 is:
$$|s\underset{\{s_i\},{\scriptscriptstyle s_i}=s}{}|s_1\mathrm{}s_N|\text{bath},s_1\mathrm{}s_N\frac{C_{s_1\mathrm{}s_N}(\rho )}{\sqrt{_{\{s_i\},{\scriptscriptstyle s_i}=s}|C_{s_1\mathrm{}s_N}|^2}},$$
so, using (48),
$$s|\chi =\left[\underset{\{s_i\},{\scriptscriptstyle s_i}=s}{}|C_{s_1\mathrm{}s_N}(\rho )|^2\right]^{1/2}=\left[\frac{_{\{s_i\},{\scriptscriptstyle s_i}=s}e^{\beta G\rho \mu a^2_{i=1}^Ns_i}}{_{\{s_i\}}e^{\beta G\rho \mu a^2_{i=1}^Ns_i}}\right]^{1/2}.$$
(49)
Evaluation of (49) is well known but, for completeness, it is sketched here. The partition function in the denominator of (49) is (writing $`C=G\rho \mu a^2`$ and taking $`N`$ even)
$$\text{Tr}e^{\beta CS}=\underset{k=0}{\overset{N}{}}e^{\beta C(N2k)}\frac{N!}{k}!(Nk)!=(2\mathrm{cosh}\beta C)^N,$$
(50)
so the expression in (49) under the square root is
$`{\displaystyle \underset{\{s_i\},{\scriptscriptstyle s_i}=s}{}}|C_{s_1\mathrm{}s_N}(\rho )|^2`$ $`=`$ $`{\displaystyle \frac{e^{\beta Cs}}{(2\mathrm{cosh}\beta C)^N}}{\displaystyle \frac{N!}{[(Ns)/2]![(N+s)/2]!}}`$ (51a)
$``$ $`{\displaystyle \frac{1}{\sqrt{2\pi N\mathrm{cosh}^2\beta C}}}e^{[sN\mathrm{tanh}\beta C]^2/2N\mathrm{cosh}^2\beta C}`$ (51b)
$``$ $`{\displaystyle \frac{1}{\sqrt{2\pi N}}}e^{[sN\beta C]^2/2N}.`$ (51c)
In (51b) the well-known gaussian approximation to the binomial distribution has been employedFeller which is valid for large $`N`$ and for s within a few standard deviations of its mean. The approximation in (51c) depends upon validity of the inequality
$$\beta G\rho \mu a^2<<1$$
(52)
which shall need to be checked.
Now it is possible to put together the contributions of all volumes $`\mathrm{\Delta }V`$ and all time intervals $`\mathrm{\Delta }t`$ within $`(0,t)`$ to construct the statevector $`|\psi ,t`$. Define $`|\stackrel{~}{s}_{\mathrm{\Delta }V,\mathrm{\Delta }t}|s_{i,j}`$, the normalized basis vector which is the joint eigenvector of $`S(𝐱_i,t_j)`$ for all $`𝐱_i`$ and for $`0t_jt`$. The scalar product $`\stackrel{~}{s}|\psi ,t`$ is the direct product of the expressions (49) so, from (51c) (without the normalization factors, which are tucked into the element of integration):
$`\stackrel{~}{s}|\psi ,t`$ $`=`$ $`e^{(4N)^1_{𝐱_i,t_j=0}^t[s_{i,j}N\beta G\rho \mu a^2]^2}`$ (53a)
$`=`$ $`e^{4^1N(\beta G\mu a^2)^2_{𝐱_i,t_j=0}^t[s_{i,j}/N\beta G\mu a^2\rho _i]^2}`$ (53b)
$``$ $`e^{4^1(p/\mathrm{}^3\tau )(\beta G\mu a^2)^2{\scriptscriptstyle 𝑑𝐱_0^t𝑑t^{}[s^{}(𝐱,𝐭^{})\rho (𝐱)]^\mathrm{𝟐}}},`$ (53c)
where (53c) follows from (46), and we have defined:
$$s^{}(𝐱,t)s_{i,j}/N\beta G\mu a^2.$$
(54)
Now, compare (53c) with CQC’s comparable expression. Note from (8) that $`A(𝐱)=m_0^1a^{3/2}\rho (𝐱)`$ (up to a numerical factor which shall be ignored) so the CQC expression comparable to Eq.(9) is
$$w|\psi ,t=e^{(4\lambda )^1{\scriptscriptstyle 𝑑𝐱_0^t𝑑t^{}[w(𝐱,t^{})2\lambda m_0^1a^{3/2}\rho (𝐱)]^\mathrm{𝟐}}}.$$
(55)
The model’s Eq.(53c) and the CQC Eq.(55) are identical if
$`W(𝐱,t)`$ $`=`$ $`{\displaystyle \frac{2\lambda a^{3/2}}{m_0}}S^{}(𝐱,t)={\displaystyle \frac{2\lambda }{(\beta m_0c^2)\mathrm{}a^{1/2}}}{\displaystyle \frac{S(𝐱,t)}{N}},`$ (56a)
$`4(\lambda \tau )(\mathrm{}/a)`$ $`=`$ $`p(\beta m_0c^2)^2,`$ (56b)
where $`G\mu =\mathrm{}c^2`$ has ben used, and $`S^{}𝐱,t)`$ is the operator whose eigenvalues are $`s^{}(𝐱,t)`$ given in Eq.(54).
Thus, this model produces CSL dynamics.
It remains to check whether the numerical values in Eqs.(52), (56) are reasonable. Consider two possible temperatures for the bath, the cosmic radiation temperature $`\beta ^1210^4`$ eV. and the Planck temperature $`\beta ^110^{28}`$ eV. For normal mass densities, $`\rho 1\text{gm/cc}510^{33}\text{eV/cc}`$, one gets respectively $`\beta G\rho \mu a^210^6`$ and $`10^{38}`$, so the inequality of Eq.(52) is satisfied for both.
For these two temperatures, (56b) gives $`p10^{112}`$ and $`10^{49}`$ respectively. The second number is too small for the conceptual picture we have outlined since, after $`10^{49}\tau 10^5`$sec, the spin in every cell is likely to have interacted and frozen, so the process would cease. However, one may change the conceptual picture. It is possible for a spin in a cell to interact repeatedly, without changing the dynamical equation because the bath states preserve the past orthogonality of the realizable states. Although they are represented by bath states, they are labeled by the past values of W(x,t). Each of these orthogonal realizable states at time $`t`$ evolves independently of the other states over the next $`\mathrm{\Delta }t`$, via the CSL evolution.
Lastly, look at Eq.(56a). Note that $`WS/N`$ is intensive, proportional to the mean value of the spin in a cell, so this result is independent of the size of $`\mathrm{\Delta }𝐱`$, $`\mathrm{\Delta }t`$ as it should be. For the numerical coefficient in (56a), it is more informative to discuss $`S^{}(𝐱,t)`$’s relation to $`S(𝐱,t)`$ rather than $`W(𝐱,t)`$’s relation to either, because $`S^{}(𝐱,t)`$’s mean value is simply $`\rho `$ (see (53c)). One obtains from Eq.(56a), for a cosmic bath and for a Planck bath respectively, $`S^{}(S/N)10^{39}`$eV/cc and $`S^{}(S/N)10^{71}`$eV/cc. Both factors are large compared to normal mass density $`10^{34}`$eV/cc, so that the mean spin value $`S/N`$ for such a density is respectively $`10^5`$ and $`10^{37}`$.
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# Resonances, Unstable Systems and Irreversibility: Matter Meets Mind
## 1 Introduction
Usually dynamical systems are considered to be time-reversible as their equations of motion are time-reversal symmetric under the time inversion operator $`R:(\stackrel{}{x},t)(\stackrel{}{x},t)`$. This means that if $`\varphi (t)`$ is a solution of the equations of motion, then so is $`R\varphi (t)`$. Such systems should then be reversible in the sense that if they exhibit a temporal succession of state transitions $`\varphi _1`$, $`\varphi _2`$, $`\varphi _3`$,…, $`\varphi _n`$, they can also exhibit the reverse temporal sequence $`R\varphi _n`$, $`R\varphi _{n1}`$, $`R\varphi _{n2}`$,…, $`R\varphi _1`$. In quantum mechanics these evolutions typically are described by one-parameter unitary groups of operators.
Resonances appear in a number of dynamical systems, both classical and quantum (e.g., Antoniou and Prigogine 1993; Antoniou and Tasaki 1993; Bohm et al. 1997) and are prototypical irreversible processes (e.g. scattering resonances). When the number of resonances in dynamical systems is sufficiently large, the dynamics is extremely unstable (e.g., exhibiting sensitive dependence on initial conditions), and becomes irreversible. For such unstable systems time arrows for the dynamics can be clearly defined (e.g., Bishop 2004a; Bishop 2004b; Bishop 2005a). In the rigged Hilbert space framework for quantum mechanics (e.g., Bohm and Gadella 1989; Bohm et al. 1997), such time arrows are represented by semigroups. It is typically the case that there are two possible semigroups for such dynamics, one defined in the forward direction in time and one defined in the backward direction. However, if the evolutions of such resonance phenomena as scattering resonances and quasistable particles are irreversible, then there appears to be no physical relevance to the mathematical descriptions of the time-reversed states and evolutions.
After presenting the background of the rigged Hilbert space (RHS) framework for quantum mechanics (QM) in section 2, I will review its application to resonance states for scattering (section 3). This will be followed by a brief review of the extended Galilean group of Wigner and its application to resonance states in the RHS framework (section 4). I will then give an interpretation of the time-reversed resonance states and evolutions as abstract representations of mental systems as opposed to material systems (section 5).
## 2 Rigged Hilbert Space Quantum Mechanics
An RHS may be briefly characterized as follows. Let $`\mathrm{\Psi }`$ be an abstract linear scalar product space and complete $`\mathrm{\Psi }`$ with respect to two topologies. The first topology is the standard Hilbert space (HS) topology $`\tau _{}`$ defined by the norm
$$h=\sqrt{(h,h)}$$
(1)
where $`h`$ is an element of $`\mathrm{\Psi }`$. The second topology $`\tau _\mathrm{\Phi }`$ is defined by a countable set of norms
$$\varphi _n=\sqrt{(\varphi ,\varphi )_n},n=0,1,2,\mathrm{}$$
(2)
where $`\varphi `$ is also an element of $`\mathrm{\Psi }`$ and the scalar product in (2) is given by
$$(\varphi ,\varphi ^{})_n=(\varphi ,(\mathrm{\Delta }+1)^n\varphi ^{}),n=0,1,2,\mathrm{}$$
(3)
where $`\mathrm{\Delta }`$ is the Nelson operator $`\mathrm{\Delta }=`$ $`_i\chi _i^2`$. The $`\chi _i`$ are the generators of an enveloping algebra of observables for the system in question and they form a basis for a Lie algebra (Nelson 1959; Bohm et al. 1999). In the case of the harmonic oscillator, for example, the $`\chi _i`$ would be the position and momentum operators or, alternatively, the raising and lowering operators. Furthermore if the operator $`\mathrm{\Delta }+1`$ is nuclear then the space $`\mathrm{\Phi }`$ defined by (2) is a nuclear space (Treves 1967).
A Gel’fand triplet is obtained by completing $`\mathrm{\Psi }`$ with respect to $`\tau _\mathrm{\Phi }`$ to obtain $`\mathrm{\Phi }`$ and with respect to $`\tau _{}`$ to obtain $``$. In addition there are the dual spaces of continuous linear functionals $`\mathrm{\Phi }^\times `$ and $`^\times `$ respectively. Since $``$ is self dual, we obtain
$$\mathrm{\Phi }\mathrm{\Phi }^\times \text{ .}$$
(4)
The Nelson operator fully determines the space $`\mathrm{\Phi }`$. However, there are many inequivalent irreducible representations of an enveloping algebra of a group characterizing a physical system (e.g. Bohm et al. 1999). Therefore further restrictions may be required to obtain a realization for $`\mathrm{\Phi }`$, e.g., due to the convergence properties desired for test functions in $`\mathrm{\Phi }`$. In general one chooses the weakest topology such that the algebra of operators for the physical problem is continuous and $`\mathrm{\Phi }`$ is nuclear. The physical symmetries of the system play an important role in such choices (Bohm et al. 1999).
In RHS QM, the observables form an algebra on the entire space of physical states (including $`\mathrm{\Phi }^\times `$, where Dirac kets reside), so a RHS contains observables with continuous or even complex eigenvalues, whereas a HS does not. This means that the basis vector expansion of eigenvectors (Dirac’s spectral decomposition) can be given a rigorous foundation resulting in the nuclear spectral theorem:
$$|\varphi =\underset{n}{}|E_n)(E_n|\phi )+|EE|\phi d\mu (E)\text{.}$$
(5)
Here the rounded bras and kets denote elements elements of $``$ and the summation in (5) represents the discrete part of the spectrum. The angular bras, $`\phi |`$, denote elements defined in $`\mathrm{\Phi }`$, while the angular kets, $`|E`$, denote elements defined in $`\mathrm{\Phi }^\times `$; hence, the integral in (5) represents the continuous part of the spectrum.
## 3 States, Observables and Resonances in Scattering
A typical scattering experiment consists of an accelerator, which prepares a projectile in a particular state, a target and detectors. The total Hamiltonian modeling the interaction of the particle with the target is, therefore, $`H`$ = $`H_o`$ \+ $`V`$, where $`H_o`$ represents the free particle Hamiltonian and $`V`$ the potential in the interaction region. The vectors representing growing and decaying states are associated with the resonance poles of the analytically continued S-matrix (Lax and Phillips 1967).
Following the Bohm group, a time arrow emerges in scattering resonances through imposing the preparation/registration arrow of time (Bohm et al. 1994; Bishop 2004b). The key intuition behind this arrow is that no observable properties of a state can be measured unless the state has first been prepared. Following Ludwig (1983; 1985), an in-state of a particular quantum system (considered as an ensemble of individual systems such as elementary particles) is prepared by a preparation apparatus (considered macrophysical). The detector (considered macrophysical) registers so-called out-states of post-interaction particles. In-states are taken to be elements $`\varphi \mathrm{\Phi }_{}`$ and observables are taken to be elements $`\psi \mathrm{\Phi }_+`$. (Resonance states, such as the Dirac, Lippman, Schwinger kets and Gamow vectors, are elements of $`\mathrm{\Phi }_\pm ^\times `$). This leads to a distinction between prepared states, on the one hand, and observables, each described by a separate RHS (Bohm and Gadella 1989; Bohm et al. 1997):
$`\mathrm{\Phi }_{}`$ $`\mathrm{\Phi }_{}^\times `$ (6a)
$`\mathrm{\Phi }_+`$ $`\mathrm{\Phi }_+^\times \text{ ,}`$ (6b)
where $`\mathrm{\Phi }_{}`$ is the Hardy space of the lower complex energy half-plane intersected with the Schwartz class functions and $`\mathrm{\Phi }_+`$ is the Hardy space of the upper complex energy half-plane intersected with the Schwartz class functions. As Bohm and Gadella (1989) demonstrate, some elements of the generalized eigenstates in $`\mathrm{\Phi }_{}^\times `$ and $`\mathrm{\Phi }_+^\times `$ correspond to exponentially growing and decaying states respectively. The semigroups governing these states are<sup>1</sup><sup>1</sup>1If $`U(t)`$ is a unitary operator on $``$ and $`\mathrm{\Phi }\mathrm{\Phi }^\times `$, then $`U^{}`$ can be extended to $`\mathrm{\Phi }^\times `$ provided that (1) $`U`$ leaves $`\mathrm{\Phi }`$ invariant and (2) $`U`$ is continuous on $`\mathrm{\Phi }`$ with respect to the topology $`\tau _\mathrm{\Phi }`$. The operator $`U^\times `$ denotes the extension of the HS operator $`U^{}`$ to $`\mathrm{\Phi }^\times `$ and is defined by $`U\varphi |F=\varphi |U^\times F`$ for all $`\varphi \mathrm{\Phi }`$ and $`F\mathrm{\Phi }^\times `$. When the group operator $`U^{}`$ is extended to $`\mathrm{\Phi }^\times `$, continuity requirements force the operators $`U^\times `$ to be semigroups defined only on the temporal half-domains (Bohm and Gadella 1989).
$`\varphi |U^\times |Z_R^{}`$ $`=e^{iE_Rt}e^{\frac{\mathrm{\Gamma }}{2}t}\varphi |Z_R^{}\text{ }t0\text{}t:\mathrm{}0`$ (7a)
$`\psi |U^\times |Z_R`$ $`=e^{iE_Rt}e^{\frac{\mathrm{\Gamma }}{2}t}\psi |Z_R\text{ }t0\text{}t:0\mathrm{}\text{,}`$ (7b)
where $`E_R`$ represents the total resonance energy, $`\mathrm{\Gamma }`$ represents the resonance width, $`Z_R`$ represents the pole at $`E_Ri\frac{\mathrm{\Gamma }}{2}`$, $`Z_R^{}`$ represents the pole at $`E_R+i\frac{\mathrm{\Gamma }}{2}`$, $`|Z_R^{}\mathrm{\Phi }_{}^\times `$ represents a growing Gamow vector and $`|Z_R\mathrm{\Phi }_+^\times `$ represents a decaying Gamow vector. The $`t<0`$ semigroup is identified as future-directed along with $`|Z_R^{}`$ as a forming/growing state. The $`t>0`$ semigroup is identified as future-directed along with $`|Z_R`$ as a decaying state.
## 4 Time-reversed States and Observables
Following Wigner (1964), the time-reversal operator, $`R(t)`$, is the HS representation of the physical spacetime transformation
$$R:(\stackrel{}{x},t)(\stackrel{}{x},t)\text{.}$$
(8)
Therefore, $`R`$ is an element of a co-representation of the extended Galilei symmetry group (Cariñena and Santander 1981) for a nonrelativistic spacetime (extended Poincaré group for a relativistic spacetime). These representations must be unitary and linear except for $`R`$, which is antilinear.
Wigner originally derived the properties of $`R`$ for the spacetime symmetry group extended by time inversions and studied the parity inversion operator $`\mathrm{\Sigma }`$ and the total inversion operator $`T`$ in combination with $`R`$ (Wigner 1964). The parity inversion operator is unitary so its phase can be chosen such that $`\mathrm{\Sigma }^2`$ $`=I`$ (the identity operator), while $`T`$ and $`R`$ are both anti-unitary, so that the associative law for group multiplication then dictates that $`R^2=\epsilon _RI`$ and $`T^2=\epsilon _TI`$, where $`\epsilon _R=\pm 1`$ and $`\epsilon _T=\pm 1`$. The phase of $`T`$ can be chosen so that $`T=\mathrm{\Sigma }R`$ (where the order of application of $`\mathrm{\Sigma }`$ and $`R`$ is physically immaterial). The extension of the Galilei spacetime symmetry group is summarized in Table I.
$`\underset{\text{Table I. Properties of the Galilei spacetime symmetry group.}}{\begin{array}{ccccc}\epsilon _R& \epsilon _T& \mathrm{\Sigma }& R& T\\ (1)^{2j}& (1)^{2j}& 1& C& C\\ (1)^{2j}& (1)^{2j}& \left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right)& \left(\begin{array}{cc}0& C\\ C& 0\end{array}\right)& \left(\begin{array}{cc}0& C\\ C& 0\end{array}\right)\\ (1)^{2j}& (1)^{2j}& \left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right)& \left(\begin{array}{cc}0& C\\ C& 0\end{array}\right)& \left(\begin{array}{cc}0& C\\ C& 0\end{array}\right)\\ (1)^{2j}& (1)^{2j}& \left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right)& \left(\begin{array}{cc}0& C\\ C& 0\end{array}\right)& \left(\begin{array}{cc}0& C\\ C& 0\end{array}\right)\end{array}}`$
The index $`j`$ refers to the spin of the particle being considered while $`C`$ is an operator whose $`(2j+1)`$-dimensional matrix has the elements $`c_{\mu ,\nu }=(1)^{j+\mu }\delta _{\mu ,\nu }`$, where $`j\mu `$ and $`\nu j`$. In the first representation, where $`\epsilon _R=\epsilon _T=(1)^{2j}`$, there are no changes to the underlying vector space. This is the typical case discussed in QM (and relativistic quantum field theory). The other three representations, however, exhibit a doubling of the vector spaces (note the block matrices in the last three columns of Table I). In order to track this space doubling, let the index $`r=0,1`$ label the rows and columns of the matrices in Table I.
Although no quantum fields have been constructed for representations two and three of Table I (indeed they are highly problematic), Bohm and co-workers have constructed models for the fourth representation by applying $`R`$ to the states and observables in (7) (Bohm 1995; Bohm and Wickramasekara 1997). First, consider the growing Gamow vectors for, $`\varphi ^{r=0,\times }\mathrm{\Phi }_{}^{r=0,\times }`$. Applying $`R`$ yields
$$R\varphi ^{r=0,\times }=\psi ^{r=1,\times }\mathrm{\Phi }_+^{r=1,\times }\text{.}$$
(9)
Similarly for the decaying Gamow vectors, $`\psi ^{r=0,\times }\mathrm{\Phi }_+^{r=0,\times }`$, applying $`R`$ yields
$$R\psi ^{r=0,\times }=\varphi ^{r=1,\times }\mathrm{\Phi }_{}^{r=1,\times }\text{.}$$
(10)
The transformation properties of $`R`$ may be summarized as $`R:\mathrm{\Phi }_\pm ^{r=0,\times }\mathrm{\Phi }_{}^{r=1,\times }`$. The temporal evolution of these time-reversed vectors is also given by semigroups. Identify $`r=0`$ with the scattering experiment as normally carried out in the laboratory and $`r=1`$ with the extended spacetime transformed situation (”time-reversed counterparts”). Then $`U^\times (t)\varphi ,r=0|Z_R^{},r=0\mathrm{\Phi }_{}^{r=0,\times }`$, a growing Gamow vector representing a preparable state for $`t0`$, is transformed under $`R`$ into $`U^\times (t)\psi ,r=1|Z_R,r=1\mathrm{\Phi }_+^{r=1,\times }`$, where
$$e^{iE_Rt}e^{\frac{\mathrm{\Gamma }}{2}t}\psi ,r=1|Z_R,r=1$$
(11)
is restricted to the time domain $`t0`$ by continuity requirements. In the case of $`|Z_R^{},r=0`$, time counts up from $`\mathrm{}`$ to $`0`$; in contrast, for $`|Z_R,r=1`$, time counts down from $`\mathrm{}`$ to $`0`$, meaning that it represents a Gamow vector that increases as $`t`$ decreases. Similarly, $`U^\times (t)\psi ,r=0|Z_R,r=0\mathrm{\Phi }_+^{r=0,\times }`$, a decaying Gamow vector representing observables for $`t0`$, is transformed under $`R`$ into $`U^\times (t)\varphi ,r=1|Z_R^{},r=1\mathrm{\Phi }_{}^{r=1,\times }`$, where
$$e^{iE_Rt}e^{\frac{\mathrm{\Gamma }}{2}t}\varphi ,r=1|Z_R^{},r=1$$
(12)
is restricted to the time domain $`t0`$ by continuity requirements. In the case of $`|Z_R,r=0`$, time counts up from $`0`$ to $`\mathrm{}`$; in contrast, for $`|Z_R^{},r=1`$, time counts down from $`0`$ to $`\mathrm{}`$, meaning that it represents a Gamow vector that decays as $`t`$ increases.
## 5 Matter Meets Mind
Comparing eqs. (7) with (11) and (12), we can see that in the $`r=0`$ regime the association of prepared states with growing eigenvectors and of detected observables with decaying eigenvectors is quite natural. On the other hand, the $`r=1`$ regime has no natural association with physical phenomena (to apply the eigenstates in this regime “straightforwardly” within our framework would lead to identifying the growing eigenvectors with “prepared observables” and the decaying eigenvectors with “detected states,” counterintuitive to say the least).<sup>2</sup><sup>2</sup>2The question of interpreting the time-reversed states and observables was first suggested to me as an interesting problem by Arno Bohm.
Suppose we consider an alternative interpretation of the states and observables of the $`r=1`$ regime as an abstract representation of mental rather than material systems. The semigroups in this regime carry vectors from the future to the past. This could be taken as an abstract representation of final causation, appropriate to teleological or goal-directed behavior. For example, suppose I have a particular vision of the kind of person I want to become, say a more humble person; or suppose I have a particular goal I want to achieve, say landing a top-flight permanent academic position. These would be examples of final causation at work in everyday decisions and actions.<sup>3</sup><sup>3</sup>3It would be interesting, though difficult, to connect the framework proposed here with analyses of goal-directed behavior in cognitive psychology, and cognitive science more broadly, as the latter perspectives tend to transmute the apparent final-causal nature of everyday goal-directed behavior into mechanisms of efficient causation (e.g. Bishop 2005b). Hence, establishing the desired connection is not straightforward without at least extending the current cognitive paradigms. Drawing on the analogy with final causation as a backwards-directed influence, an eigenvector growing in the backwards time direction might represent the formation (“preparation” or “excitation”) of such a goal or vision of the future. This could be taken as representing the building influence of the goal or vision of the future on the present decision. Similarly, an eigenvector decaying in the backwards time direction might represent the decision state (“registration” or “de-excitation”) resulting in concrete action toward the goal. It is plausible that decision states decay back to some kind of “ready state” after action is initiated so that a new decision state can be created for the next set of goals and actions. The rate of decay could be slower or faster depending on whether the intended action required more effort of will to “stay on track” as it were to completion or not. The resonance state might be taken as a representation of the decision itself.
The $`r=1`$ regime could, then, serve as an abstract model of goal-directed decision and action. Moreover, both regimes together would play a role in the abstract description of mental and material systems and their relations. The $`r=0`$ regime would correspond to material systems while the $`r=1`$ regime would correspond to mental systems. We would, then, have a unified abstract description of mental-material systems.
Such an abstract description could be deployed to represent a “dualistic” distinction between material and mental domains, emerging from a “monistic” domain without such a distinction. It has been proposed that this emergence is related to some temporal symmetry breaking (Atmanspacher 2003; Primas 2004) in the spirit of ideas of Pauli and Jung (Pauli and Enz 2001), where physical and psychical aspects originate in a psychophysically neutral domain. The symmetry breaking envisaged need not be a unique, one-time event, but is perhaps best understood as an ongoing process due to a number of contingent conditions giving rise to the mental-material distinction.<sup>4</sup><sup>4</sup>4The role of contingent conditions in the emergence of properties is discussed in (Bishop and Atmanspacher 2005). Furthermore, this symmetry breaking can lead to the Cartesian distinction of the dualistic approach while still allowing for correlations or forms of interaction emerging from the neutral domain, perhaps leading to resolution of a number of problems plaguing the dualistic approach.
To be a bit more precise, suppose the neutral domain is characterized by states $`\omega `$ and a unitary symmetry (continuum order, automorphic dynamics, etc.). The dynamics of this domain would then exhibit the time-reversal symmetry described in section 1 and might be characterized by a one-parameter unitary group of bounded operators on $``$. Some, as yet unspecified, symmetry breaking leads to the generation of time-asymmetric dynamics characterized by semigroups governing the two regimes, $`r=0,1`$. As originally characterized, the states $`\omega `$ are neutral with respect to mental-material aspects, whereas after the symmetry breaking, the states are differentiated into material ($`r=0`$) and mental ($`r=1`$) states and processes. It is at the level of symmetry breaking that the states and observables discussed above emerge. The characterization of observables in the unitary domain is left unaddressed here.<sup>5</sup><sup>5</sup>5As a referee helpfully pointed out, one could also consider a more general kind of interpretation of the $`r=1`$ regime, namely as representing observational systems, where mental systems are a very important special case. Space does not permit consideration of this interesting possibility.
The fact that the $`r=0`$ and $`r=1`$ regimes are related to each other via a time-reversal operator suggests the possibility that there is some form of intertwining relation among the states and observables of the two regimes. If so, then the relationship among the elements of the mental and material domains would not be so starkly disjoint as in Descartés’ view, where the two domains are conceived as distinctly different kinds of substances. Therefore, the distinction between mental and material domains need not imply Descartés’ metaphysical distinction nor the kinds of interaction problems encountered in that view.
The RHS framework for QM allows for the description of both time-symmetric and time-asymmetric phenomena. In particular, it is well-suited for the description of resonances and other kinds of unstable states. If the interpretation sketched here makes the time-reversed states and observables of the $`r=1`$ regime plausible, then the RHS framework is also well-suited for such an abstract representation of mental and material states. The unitary neutral domain might be related to $``$ while the $`r=0`$ (material) and $`r=1`$ (mental) regimes are related to $`\mathrm{\Phi }`$ and $`\mathrm{\Phi }^\times `$.
Although one might wonder about the propriety of using concrete models to motivate the framework and then subsequently throwing those models to the side to apply the framework to more abstract questions, this way of proceeding represents a well-established use of models in mathematical physics (Redhead 1980). There are also technical questions about the application of Wigner’s ideas to observables as well as to semigroup representations, but these question are fairly straightforward. What is not so straightforward are questions such as the emergence of time, or the kinds of contingent conditions leading to the symmetry breaking generating the two regimes. The abstract RHS framework proposed here appears to be promising as one avenue for exploring such topics.
## 6 Concluding Summary
One way the fundamental time-reversal invariance of dynamical systems might be broken is through the presence of resonances and their interactions giving rise to unstable dynamical systems. When time-reversal invariance is broken, this results in two well-defined time arrows associated with semigroups bearing time orientations, one directed toward the future and one directed toward the past. Scattering resonances provide an example where time-reversal invariance is broken. The resulting forward-directed semigroups and states correspond to the processes of resonance formation, decay and detection, but the backward-directed semigroups and states are thought to have perhaps only mathematical significance. Here, I have sketched a possible interpretation of these latter semigroups and states as corresponding abstractly to the domain of mental systems, while the forward-directed semigroups and states would correspond to the domain of material systems. The crucial idea is that these two domains might emerge from a more fundamental domain that is neutral with respect to any mental-material distinction and that, hence, various possibilities exist for relations between the emergent domains that are typically precluded by traditional Cartesian dualisms. The RHS framework seems well-suited for describing and exploring these possibilities.
###### Acknowledgement 1
I would like to thank H. Atmanspacher and A. Bohm for helpful discussions. Financial support from the Alexander von Humboldt Foundation as well as the Federal Ministry of Education and Research and the German government’s Program for the Investment in the Future of the are gratefully acknowledged.
## 7 References
I. Antoniou and S. Tasaki, Int. J. Quant. Chem., 46 (1993) 425
I. Antoniou and I. Prigogine, Physica A, 192 (1993) 443.
H. Atmanspacher, BioSystems, 68 (2003) 19.
R. C. Bishop, Studies in History and Philosophy of Modern Physics 35 (2004a) 1.
R. C. Bishop, Int. J. Theor. Phys., 43 (2004b) 1675.
R. C. Bishop, Int. J. Theor. Phys., (2005a) in press.
R. C. Bishop, “Cognitive Psychology: Hidden Assumptions,” in Critical Thinking About Psychology: Hidden Assumptions and Plausible Alternatives, B. Slife, J. Reber and F. Richardson (eds.), Washington, D. C.: American Psychological Association Books, (2005b) 151.
R. C. Bishop and H. Atmanspacher, (2005), “Contextual Emergence in the Description of Properties,” submitted.
A. Bohm, Phys. Rev. A, 51 (1995) 1758.
A. Bohm, I. Antoniou and P. Kielanowski, Phys. Lett. A, 189 (1994) 442.
A. Bohm and M. Gadella, Dirac Kets, Gamow Vectors, and Gel’fand Triplets, Lecture Notes in Physics, vol. 348, Springer, Berlin, (1989).
A. Bohm, M. Gadella and S. Wickramasekara, “Some Little Things about Rigged Hilbert Spaces and Quantum Mechanics and All That,” in I. Antoniou and G. Lumer (eds.) Generalized Functions, Operator Theory, and Dynamical Systems. Boca Raton, FL: Chapman & Hall/CRC, (1999) 202.
A. Bohm, S. Maxson, M. Loewe and M. Gadella, Physica A, 236 (1997) 485.
A. Bohm and S. Wickramasekara, Found. of Phys., 27 (1997) 969.
J. Cariñena and M. Santander, J. Math. Phys., 22 (1981) 1548.
P. Lax and R. Phillips, Scattering Theory, Academic Press, New York, (1967).
G. Ludwig, Foundations of Quantum Mechanics, Vol. I, Springer, Berlin, (1983).
G. Ludwig, Foundations of Quantum Mechanics, Vol. II, Springer, Berlin, (1985).
E. Nelson, Ann. Math., 70 (1959) 572.
W. Pauli and C. P. Enz, Atom and Archetype: The Pauli/Jung Letters, 1932-1958. Princeton: Princeton University Press, (2001).
H. Primas, Mind and Matter, 1, (2004) 81.
M. Redhead, Brit. J. Phil. Sci., 31, (1980) 145.
F. Treves, Topological Vector Spaces, Distributions and Kernels. New York: Academic Press, (1967).
E. Wigner, “Unitary Representations of the Inhomogeneous Lorentz Group Including Reflections,” in Group Theoretical Concepts and Methods in Elementary Particle Physics, F. Gürsey (ed.), Gordon and Breach, Science Publishers, New York, (1964).
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# Cluster Baryon Fraction and Structure from the Convergence/SZ Effect Diagram
## 1. Introduction
Clusters of galaxies, the largest virialized systems known, are key tracers of the matter distribution in the Universe. Clusters are composed of the hot, diffuse intracluster medium (ICM), luminous galaxies (stars), and a dominant mass component in the form of the unseen dark matter. The mass ratio of hot gas to dark matter in clusters (the gas fraction) is believed to be a fair estimate of the universal baryon fraction (White et al. 1993). However, non-gravitational processes associated with cluster formation, such as radiative gas cooling or pre-heating, would break the self-similarities in cluster properties and hence cause the gas fraction to acquire some mass dependence (Kravtsov et al. 2005; Bialek et al. 2001). Weak gravitational lensing on background galaxies probes directly the projected cluster mass distribution (Bartelmann & Schneider 2001). The thermal Sunyaev-Zel’dovich (SZ) effect on the cosmic microwave background spectrum, on the other hand, measures the projected thermal electron pressure (Birkinshaw 1999). Therefore, a combination of weak-lensing and SZ measurements will allow us to measure directly the cluster gas fractions when the ICM temperature is available (Myers et al. 1997; Holder, Carlstrom, & Evrard 2000).
Both blind SZ surveys, such as that to be done with AMiBA (Ho et al. 2004), and gravitational weak lensing surveys, such as the CFHT Legacy Survey<sup>1</sup><sup>1</sup>1http://www.cfht.hawaii.edu/Science/CFHTLS, aim to detect massive clusters of $`\stackrel{>}{}10^{14}M_{}`$ to study the formation of structure through the abundance and properties of the detected clusters. A comparison of the two types of survey should reveal much about the different structures of hot gas and dark matter in clusters of different mass. Here we investigate this possibility using $`N`$-body simulations.
## 2. Cluster gas fractions from weak-lensing and SZ observations
The lensing convergence $`\kappa `$ is essentially the surface mass density projected on the sky,
$$\kappa =\mathrm{\Sigma }_m/\mathrm{\Sigma }_{\mathrm{crit}}(z_d,z_s),$$
(1)
where $`\mathrm{\Sigma }_m`$ is the surface mass density of a halo acting as gravitational lens of redshift $`z_d`$, and $`\mathrm{\Sigma }_{\mathrm{crit}}`$ is the lensing critical surface mass density as a function of lens redshift $`z_d`$ and source redshift $`z_s`$ specified for a given background cosmology (e.g., Bartelmann & Schneider 2001).
The SZ effect is described by the Compton $`y`$ parameter defined as a line-of-sight integral of the temperature-weighted thermal electron density (Birkinshaw 1999):
$$y=\frac{\sigma _T}{m_ec^2}𝑑ln_ek_\mathrm{B}T_g\frac{1+X}{2}\frac{\sigma _T}{m_p}\frac{k_\mathrm{B}T_g}{m_ec^2}\mathrm{\Sigma }_g.$$
(2)
where $`n_e`$ is the electron number density, $`T_g`$ the gas temperature, $`m_e`$ the electron mass, $`\sigma _\mathrm{T}`$ the Thomson cross section, $`\mathrm{\Sigma }_g`$ the surface gas mass density, and $`X=0.76`$ the hydrogen abundance. Note that we have assumed isothermality in the second equality in Eq. (2). The presence of temperature gradients in the ICM still remains controversial (see Piffaretti et al. 2005 and references therein). In the paper, we assume the ICM is isothermal as a zeroth-order approximation; we found a temperature variation within the virial radius by a factor two at most from our sample of massive halos (see §4).
In general, the lensing $`\kappa `$ and the Compton $`y`$ reconstructed from noisy observations are smoothed by a certain window function in the map making process. In interferometric observations, the synthesized beam (or PSF) is well approximated by a Gaussian window. In weak lensing surveys, Gaussian smoothing is often adopted to find clusters as local maxima in the reconstructed convergence map (e.g., Miyazaki et al. 2002). We hence express the convergence field smoothed with a Gaussian window $`W_\mathrm{G}(\theta )=\mathrm{exp}(\theta ^2/\theta _\mathrm{G}^2)/\pi \theta _\mathrm{G}^2`$ as $`\kappa _\mathrm{G}(𝜽)`$; $`\theta _\mathrm{G}`$ is related to FWHM of the Gaussian window $`\theta _{\mathrm{fwhm}}`$ by $`\theta _\mathrm{G}=\theta _{\mathrm{fwhm}}/\sqrt{4\mathrm{ln}(2)}`$. Similarly, we define the Gaussian smoothed Compton $`y`$ parameter, $`y_\mathrm{G}(𝜽)`$. We then compare the reconstructed $`\kappa _\mathrm{G}`$ and $`y_\mathrm{G}`$ values at a certain point $`𝜽`$.Taking the ratio of smoothed $`\kappa _\mathrm{G}`$ and $`y_\mathrm{G}`$, we have
$`\eta _g{\displaystyle \frac{y_\mathrm{G}}{\kappa _\mathrm{G}}}`$ $`=`$ $`8.9h\times 10^4\left({\displaystyle \frac{k_\mathrm{B}T_g}{10\mathrm{k}\mathrm{e}\mathrm{V}}}\right)\left({\displaystyle \frac{\mathrm{\Sigma }_{\mathrm{crit}}}{1h\mathrm{g}/\mathrm{cm}^2}}\right)\left({\displaystyle \frac{f_g}{0.1}}\right)`$ (3)
where $`f_g(𝜽)=\mathrm{\Sigma }_{g,\mathrm{G}}(𝜽)/\mathrm{\Sigma }_{m,\mathrm{G}}(𝜽)`$ is the local gas fraction defined with the Gaussian smoothed surface mass densities. In order to extract information on the cluster gas fraction from weak lensing and SZ observations, we thus need additional information on the ICM temperature $`T_g`$, the cluster redshift $`z_d`$, and the redshift distribution of background galaxies. Such information can be in principle available from X-ray observations and photometric redshift measurements of background galaxies.
We can define a global estimator for the cluster gas fraction based on weak-lensing and SZ observation as
$$\eta _g\frac{_iw_i\eta _{g,i}}{_iw_i}\mathrm{\Sigma }_{\mathrm{crit}}(z_d,z_s)f_gT_g$$
(4)
where $`\eta _{g,i}=\eta _g(𝜽_i)`$ and $`w_i=w(𝜽_i)`$ is the weight for the $`i`$th pixel $`𝜽_i`$. Accordingly, the global gas fraction of a halo can be estimated from $`\kappa `$ and $`y`$ maps as $`f_g\eta _gT_g^1\mathrm{\Sigma }_{\mathrm{crit}}^1`$ if one assumes the isothermality of the ICM. The case with $`w_i=\kappa _\mathrm{G}(𝜽_i)`$ corresponds to a commonly used definition of the gas fraction, $`\eta _g^{(1)}y/\kappa \mathrm{\Sigma }_g/\mathrm{\Sigma }_m`$, as the ratio of the mean surface mass densities. However, the mean convergence $`\kappa `$ is ill constrained from weak-lensing observations: Weak lensing $`\kappa `$ reconstructions based solely on image distortions suffer from the mass sheet degeneracy (Schneider & Seitz 1995) unless additional information such as the magnification bias is taken into account (Broadhurst, Taylor & Peacock 1995; Broadhurst et al. 2005). On the other hand, interferometric observations are insensitive to the DC signal, and hence do not constrain the total flux in the field-of-view (FoV). Bolometers also suffer from a high level of contamination by atmospheric and other environmental signals, and hence rely on the differencing scheme to extract the sky signals. We therefore consider an alternative weight $`w_i=\kappa _\mathrm{G}^2(𝜽_i)`$ that weights high density regions more strongly than $`w_i=\kappa _\mathrm{G}(𝜽_i)`$. Then we have the gas fraction estimator of the form:
$$\eta _g^{(2)}\frac{_i\kappa _\mathrm{G}(𝜽_i)y_\mathrm{G}(𝜽_i)}{_i\kappa _\mathrm{G}(𝜽_i)\kappa _\mathrm{G}(𝜽_i)}=\frac{\kappa y}{\kappa ^2}.$$
(5)
Suppose that $`\kappa `$ and $`y`$ satisfy a linear bias relation $`y(𝜽)=a\kappa (𝜽)`$. Then errors in the Compton parameter of the form $`yy+ϵ`$ (where $`y|ϵ|`$) lead to errors in the two estimators of $`\eta _g`$ of $`\mathrm{\Delta }\eta _g^{(1)}=ϵ/\kappa `$ and $`\mathrm{\Delta }\eta _g^{(2)}=ϵ\kappa /\kappa ^2`$. The difference between these errors $`|\mathrm{\Delta }\eta _g^{(1)}||\mathrm{\Delta }\eta _g^{(2)}|=\left(|ϵ|/\kappa \right)\left(\mathrm{Var}[\kappa ]/\kappa ^2\right)0`$, with equality only if variance in $`\kappa `$ vanishes. Thus the $`\eta _g^{(2)}`$ estimator is always better than the $`\eta _g^{(1)}`$ estimator. Similarly, an error on the convergence of the form $`\kappa \kappa +ϵ`$ (where $`\kappa |ϵ|)`$ leads to $`|\mathrm{\Delta }\eta ^{(1)}||\mathrm{\Delta }\eta ^{(2)}|\left(a|ϵ|/\kappa \right)\left(\mathrm{Var}[\kappa ]/\kappa ^2\right)0`$, again showing that the cross-correlation based estimator $`\eta _g^{(2)}`$ is more robust than $`\eta _g^{(1)}`$.
In practical observations, noise contributions to the estimator (4) must be taken into account. Since the noise properties between the $`\kappa `$ and $`y`$ maps will not be correlated, the required correction is in the denominator of Eq. (5), where we replace $`\kappa ^2`$ by $`\kappa ^2\sigma _\kappa ^2`$ with $`\sigma _\kappa ^2`$ being the noise variance in the smoothed convergence field, $`\kappa _\mathrm{G}`$ (§3.3).
## 3. Mock observations
We use cosmological simulation data by Lin et al. (2004) to demonstrate what useful information we can extract from the cross correlation between weak-lensing and SZ surveys. Specifically, we simulate the forthcoming AMiBA experiment as an illustration of such combined surveys.
### 3.1. $`N`$-body/hydrodynamic simulations
To make sky maps with realistic SZ and weak lensing signals, we use results from preheating cosmological simulations of a $`\mathrm{\Lambda }`$CDM model $`(\mathrm{\Omega }_m=0.34,\mathrm{\Omega }_\mathrm{\Lambda }=0.66,\mathrm{\Omega }_b=0.044,h=0.66,\sigma _8=0.94)`$ in a $`100h^1`$Mpc co-moving box generated with an N-body/hydrodynamics code GADGET (Springel et al. 2001), which reproduce the observed cluster $`M_X`$-$`T_X`$ and $`L_X`$-$`T_X`$ relations at $`z=0`$ (Lin et al. 2004). We construct a total of $`29`$ SZ sky maps, each $`1\mathrm{d}\mathrm{e}\mathrm{g}^2`$ in size and containing $`1024^2`$ pixels, by projecting the electron pressure through the randomly displaced and oriented simulation boxes, separated by $`100h^1`$ Mpc, along a viewing cone out to the redshift of $`z=2`$. Similarly, 29 $`\kappa `$-maps are constructed by projecting the distance-weighted mass overdensity $`\delta \rho `$ out to a source plane at a fixed redshift, $`z=z_s`$.
### 3.2. AMiBA SZ cluster survey
As an AMiBA specification, we adopt a close-packed hexagonal configuration of $`19\times 1.2`$m dishes on a single platform. AMiBA operates at $`95`$ GHz, and this array configuration yields a synthesized beam of $`\theta _{\mathrm{fwhm}}2^{}`$ ($`0.6h^1`$Mpc at $`z=0.8`$), which is optimized to detect high-$`z`$ clusters (Zhang et al. 2002). The primary beam has an FoV of $`\mathrm{FWHM}11^{}`$. We assume the following telescope system: bandwidth $`\mathrm{\Delta }\nu =20`$ GHz, system temperature $`T_{\mathrm{sys}}=70`$ K, system efficiency of $`0.7`$, and dual polarizations.
AMiBA will operate as a drift scan interferometer to remove ground contamination to high order (Pen et al. 2002<sup>2</sup><sup>2</sup>2http://www.cita.utoronto.ca/~pen/download/Amiba; Zhang et al. 2002; Park et al. 2003). We approximate the resulting synthesized beam in the AMiBA map by a Gaussian with $`\theta _{\mathrm{fwhm}}=2^{}`$, and we convolve the simulated Compton $`y`$ maps with this Gaussian beam to generate a set of noise-free AMiBA $`y_\mathrm{G}`$-maps. The expected rms noise over a sky area of $`\mathrm{\Omega }_s`$ with a total integration time of $`t_{\mathrm{int}}`$ is (Pearson et al. 2003),
$$\sigma _{y_\mathrm{G}}=2.0\times 10^6(\frac{T_{\mathrm{sys}}}{70\mathrm{K}})\left(\frac{\mathrm{\Delta }\nu }{20\mathrm{G}\mathrm{H}\mathrm{z}}\right)^{1/2}\left(\frac{\mathrm{\Omega }_\mathrm{s}}{1\mathrm{d}\mathrm{e}\mathrm{g}^2}\right)^{1/2}\left(\frac{t_{\mathrm{int}}}{240\mathrm{h}\mathrm{r}\mathrm{s}}\right)^{1/2}.$$
(6)
### 3.3. Weak lensing cluster survey
We assume a weak lensing survey with a mean galaxy number density of $`n_g=40`$ arcmin<sup>-2</sup>, an rms amplitude of the intrinsic ellipticity distribution of $`\sigma _ϵ=0.3`$, and a mean galaxy redshift of $`z_s=1`$, which are close to typical values for a ground based optical imaging survey with a magnitude limit of $`R_{\mathrm{lim}}\stackrel{>}{}25.5`$ mag in a sub-arcsecond seeing condition (e.g., Hamana, Takada, & Yoshida 2004). The rms noise in a Gaussian-smoothed convergence field, $`\kappa _\mathrm{G}`$, is given as
$$\sigma _{\kappa _\mathrm{G}}=1.1\times 10^2\left(\frac{\sigma _ϵ}{0.3}\right)\left(\frac{n_g}{40\mathrm{a}\mathrm{r}\mathrm{c}\mathrm{m}\mathrm{i}\mathrm{n}^2}\right)^{1/2}\left(\frac{\theta _{\mathrm{fwhm}}}{2^{}}\right)^1$$
(7)
(e.g., van Waerbeke 2000). We choose the same smoothing scale of $`\theta _{\mathrm{fwhm}}=2^{}`$ as for the AMiBA SZ survey and convolve the $`\kappa `$ maps from $`N`$-body simulations (see §3.1) with the Gaussian beam to produce a set of noise-free $`\kappa _\mathrm{G}`$ maps. Note that the adopted value of $`\theta _{\mathrm{fwhm}}`$ is close to the optimal smoothing scale of $`\theta _{\mathrm{fwhm}}1\stackrel{}{\mathrm{.}}7`$ (or $`\theta _\mathrm{G}1^{}`$) for an efficient survey of massive halos with $`M\stackrel{>}{}10^{14}h^1M_{}`$ found by Hamana et al. (2004) based on cosmological $`N`$-body simulations.
## 4. Cluster gas fractions from cross correlation between weak-lensing and SZ signals
We carried out a statistical analysis on the simulated weak lensing $`\kappa _\mathrm{G}`$ and SZ Compton $`y_\mathrm{G}`$ maps of $`29\times 1\mathrm{d}\mathrm{e}\mathrm{g}^2`$ for the AMiBA experiment. Figure 1 shows the cross correlation between the noise-free $`\kappa _\mathrm{G}`$ and $`y_\mathrm{G}`$ sky maps of the entire $`29\mathrm{d}\mathrm{e}\mathrm{g}^2`$ as a 2D scatter plot histogram. The scatter plot is presented in units of rms dispersions $`\kappa _{}`$ and $`y_{}`$ of the smoothed noise-free fields $`\kappa _\mathrm{G}(𝜽)`$ and $`y_\mathrm{G}(𝜽)`$, respectively; $`\kappa _{}=1.4\times 10^2`$ and $`y_{}=2.2\times 10^6`$. Halo peaks with $`\kappa _\mathrm{G}\stackrel{>}{}4.5\kappa _{}`$ and $`y_\mathrm{G}\stackrel{>}{}6y_{}`$ ($`3.3\mathrm{mJy}/\mathrm{beam}`$ at $`95`$ GHz) are indicated in Fig. 1 according to their redshifts. These thresholds effectively select massive halos with $`M_{200}\stackrel{>}{}2\times 10^{14}h^1M_{}`$ (see Fig. 2 of Umetsu et al. 2004); $`M_{\mathrm{\Delta }_c}(4\pi /3)\mathrm{\Delta }_c\rho _cr_{\mathrm{\Delta }_c}^3`$ is the mass enclosed by $`r_{\mathrm{\Delta }_c}`$ with the critical density of the universe $`\rho _c`$. Figure 1 clearly exhibits different capabilities of SZ and weak-lensing in cluster detections: weak lensing depends strongly on $`z_d`$ through $`\mathrm{\Sigma }_{\mathrm{crit}}`$, while the surface brightness of SZ effect is independent of $`z_d`$. In particular, gravitational lensing is less efficient for high-$`z`$ lensing halos due to pure geometrical effects. As a consequence, high-$`z`$ halos tend to appear in low-$`\kappa `$ and high-$`y`$ regions.
Another important feature in Fig. 1 is a number of butterfly-wing-shaped structures extending from the origin to points corresponding to the halo peaks. Each blade of a butterfly wing corresponds to an individual cluster. Hereafter, we call the $`\kappa /y`$ cross-correlation diagram the butterfly diagram. For illustrative purposes, we show in Fig. 2 a butterfly diagram from the noise-free result for the most massive halo, which has $`M_{200}=7.0\times 10^{14}h^1M_{}`$ at $`z=0.142`$. The SZ signal for this halo is more extended than the weak lensing signal in this cluster region. Furthermore, we can see a clear disparity between the shapes of the contours in the $`\kappa `$ and $`y`$ maps.
The butterfly diagram of a single cluster shows a wing that contains rich information on the relative distributions of dark matter and hot gas. The tip of a butterfly wing shows the surface mass densities of both dark matter and hot gas in the cluster core, and the bulk of the wing reveals the degree to which the hot gas traces dark matter. The wing width serves as an indicator of the difference between projected distributions of dark matter and hot gas. Most halos shown in Fig. 1 have significant wing widths. These widths arise because the baryon distribution is closer to axisymmetric than the dark matter distribution due to collisional relaxation of the baryons. The local slope $`\eta _g(𝜽)`$ of a wing in conjunction with the cluster redshift and temperature can be an indicator of the local gas fraction $`f_g(𝜽)`$ (see Eq. ). Note that the estimator $`\eta _g^{(2)}`$ defined in Eq. (5) can be interpreted as a global slope of a wing in the least-square sense.
We use the halo sample with $`\kappa _\mathrm{G}\stackrel{>}{}4.5\kappa _{}`$ and $`y_\mathrm{G}\stackrel{>}{}6y_{}`$ derived from our noise-free sky maps of $`29\mathrm{d}\mathrm{e}\mathrm{g}^2`$ to compare the halo-based cumulative gas fractions $`f_g(<r)=M_g(r)/M(r)`$ with their projected estimates $`f_g^{(2)}(<\theta )\eta _g^{(2)}(<\theta )T_g^1\mathrm{\Sigma }_{\mathrm{crit}}^1(z_d,z_s)`$ from cross correlation between $`\kappa `$ and $`y`$; $`\theta =r/D_A(z_d)`$ with $`D_A(z)`$ being the angular diameter distance to redshift $`z`$. We make this comparison inside $`r_{500}`$, and calculate for each halo a gas fraction using Eqs. (4) and (5) within an aperture of projected radius $`\theta _{500}=r_{500}/D_A(z_d)`$. We assume cluster isothermality and take the emission-weighted mean temperature $`T_{500}`$ inside $`r_{500}`$ as the observable X-ray temperature. The weight $`w_i`$ in Eq. (4) is set to zero if $`\kappa _i0`$ or $`y_i0`$ to exclude low-signal regions where cross correlation measurements are of low significance due to projection effects and/or observational noise. Further, we remove those halos whose centers are located outside the inner $`50^{}\times 50^{}`$ region of a $`1\mathrm{d}\mathrm{e}\mathrm{g}^2`$ sky map, in order to avoid systematics in the gas fraction estimates due to the boundaries. The effective survey area is thus about $`20\mathrm{deg}^2`$
Figure 3 shows the distribution of the ratio $`f_g^{(2)}(<\theta _{500})/f_g(<r_{500})`$ obtained from the noise-free maps as a function of $`M_{200}`$ for three redshift subsamples: (a) $`z0.3`$ (where we find 8 halos); (b) $`0.3z<0.5`$ (16 halos); and (c) $`0.5z<1`$ (13 halos). The sample means and their standard errors of $`f_g^{(2)}(<\theta _{500})/f_g(<r_{500})`$ are (a) $`0.96\pm 0.03`$, (b) $`1.10\pm 0.04`$, and (c) $`0.95\pm 0.05`$. Note that the quoted errors represent the intrinsic scatter of the estimator $`f_g^{(2)}`$ from the noise-free signal maps. On the whole, the estimator $`f_g^{(2)}(<\theta _{500})`$ for the gas fraction is close to unbiased, but there is a slight ($`2\sigma `$) over-estimation bias of $`f_g`$ in $`0.3z<0.5`$. We found that most halos with overestimated values of $`f_g`$ are associated with mergers. The shock heating of gas in mergers boosts the temperature in the region between the merging halos. However, we ascribe the lower temperature $`T_{500}`$ to this gas, based on the emission-weighted temperature of the entire halo, and so overestimate the value of $`f_g`$ for the system as a whole. This overestimate of $`f_g`$ appears to be smaller at higher $`z`$ because such halos have smaller angular sizes and so suffer larger beam smearing effects. If the 5 interacting/merger systems in subsample (b) are rejected, then the sample mean $`f_g^{(2)}(<\theta _{500})/f_g(<r_{500})=1.01\pm 0.03`$, showing that the bias is eliminated. Such a debias should be possible based on X-ray observations of the survey field.
On the other hand, the projection effect in $`\kappa `$ due to intervening matter (Metzler et al. 2001) also affects cluster $`f_g`$ estimates made using weak lensing data. This is more severe for higher-$`z`$ clusters due to their lower geometrical lensing efficiency and smaller angular sizes. If an intervening mass concentration is projected onto a cluster, then the gas fraction of that cluster could be under-estimated. Furthermore, the gradients in 3D cluster gas-fraction and temperature profiles (Ettori et al. 2004) can also induce a systematic bias in projected $`f_g^{(2)}`$ estimates, especially for spatially-resolved low-$`z`$ clusters. Intrinsic gradients in $`T_g`$ and $`f_g`$ are manifest as a curvature of the wing in a butterfly diagram, and they cannot be distinguished unless detailed information is available on the temperature distribution in the ICM.
The error bars in Fig. (3) are based on 50 realizations of Monte-Carlo simulated noise maps which are linearly added to the Gaussian smoothed signal maps of $`\kappa _\mathrm{G}`$ and $`y_\mathrm{G}`$. Following van Waerbeke (2000), a correlated noise map of $`\kappa _\mathrm{G}`$ is generated for each realization using the noise auto-correlation function specified by the smoothing scale $`\theta _{\mathrm{fwhm}}`$ and the observational noise rms $`\sigma _\kappa `$ (Eq. ). A similar correlated noise map for $`y_\mathrm{G}`$ uses the noise rms $`\sigma _y`$ (Eq. ). For a fixed halo mass, the random error in $`f_g^{(2)}`$ is smaller for lower-$`z`$ halos with larger angular scales, $`\theta _{500}`$, over which pixel noise is averaged, but systematic errors will be dominant for these clusters.
## 5. Conclusions
By simulating the AMiBA experiment, we have shown that the butterfly diagram that emerges from a comparison of the SZ and weak-lensing signals contains rich information on the distribution of the cluster gas fraction within an individual cluster and between different clusters (§4). The cross-correlation based estimator $`f_g^{(2)}(<\theta )`$ should be used for a robust estimate of the cluster gas fraction. Overall, this estimator is close to unbiased, but some halos show intrinsic scatter that exceeds the observational errors (§3) due to various sources of systematics (§4). The bias for the the intermediate-$`z`$ subsample of halos can be eliminated by removing interacting/merger systems. Careful selection of clusters should minimize systematic uncertainties in $`f_g`$. Detailed information regarding the temperature distribution in clusters, derived from X-ray observations, improves the calibration of the gas fraction estimator, and yields valuable information on the thermodynamic history of the ICM.
We acknowledge Kai-Yang Lin and Masahiro Takada for useful discussions.
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# Correlations from p-p collisions at √𝑠=200 GeV
## 1 Introduction
QCD theory predicts that an abundance of low-$`Q^2`$ scattered gluons (minijets) should be produced in relativistic nuclear collisions at RHIC energies. Those gluons are believed to drive formation of the colored medium in heavy ion collisions . If so, we may discover remnants of low-$`Q^2`$ ($`Q/2`$ 1 - 5 GeV) partons in the correlation structure of final-state hadrons. Initial studies of two-particle correlations in p-p collisions emphasized momentum subspace $`(\eta ,\varphi )`$ (pseudorapidity and azimuth) described at smaller $`p_t`$ in terms of string fragmentation . Angular correlations of fragments from hard-scattered partons (jets) were observed on $`(\eta ,\varphi )`$ at larger $`p_t`$ and with increasing $`\sqrt{s}`$ . The two-component model of p-p collisions, with string and parton fragments forming distinguishable soft and hard components, describes particle multiplicity distributions at large $`\sqrt{s}`$ and has also been introduced to the heavy-ion context .
Two issues arise in a conventional study of parton fragment distributions: 1) the distribution of fragment momenta along the jet thrust axis (fragmentation function) and 2) the angular distribution of fragments relative to the thrust axis. The thrust axis (estimating the parton direction) may be inferred from a collection of hadrons (jet) identified with the scattered parton, or (especially in heavy ion collisions) a ‘high-$`p_t`$ leading particle’ may be used to estimate the parton momentum direction and magnitude. In the present analysis we adopt no a priori jet or factorization hypothesis. We study minimum-bias two-particle distributions on transverse rapidity space $`(y_{t1},y_{t2})`$ to obtain fragment distributions (not fragmentation functions) and on angle space $`(\eta _1,\eta _2,\varphi _1,\varphi _2)`$ to obtain fragment angular correlations. Particles in pairs are treated symmetrically, as opposed to asymmetric ‘trigger’ and ‘associated’ particle combinations.
We observe that hard-component correlations obtained with this minimum-bias analysis, in contrast to jet correlations obtained with the conventional leading- or trigger-particle approach, represent the majority of parton fragment pairs, those with $`p_{t1}p_{t2}1`$ GeV/c. Analysis of nonperturbative low-$`Q^2`$ parton fragment correlations has motivated some novel techniques, including use of transverse rapidity $`y_t`$ rather than momentum $`p_t`$ and formation of joint (2D) angular autocorrelations. Symmetric analysis of fragment pairs requires a generalized treatment of fragmentation functions and jet angular correlations. An immediate consequence has been observation of substantially asymmetric angular correlations about the thrust axis for low-$`Q^2`$ parton collisions. While we attempt to understand QCD in A-A collisions we should also revisit its manifestations in elementary collisions, where novel phenomena are still emerging.
## 2 String and low-$`Q^2`$ parton fragment correlations on $`(y_{t1},y_{t2})`$ and $`(\eta _\mathrm{\Delta },\varphi _\mathrm{\Delta })`$
In a study of $`p_t`$ spectra from p-p collisions at $`\sqrt{s}=200`$ GeV we observed that the spectra can be separated into a soft component (string fragments - longitudinal fragmentation) described by a Lévy distribution on transverse mass $`m_t`$ and a ‘semi-hard’ component (parton fragments - transverse fragmentation) described by a gaussian distribution on transverse rapidity $`y_t`$, based on systematic variation with event multiplicity $`n_{ch}`$. The hard component is interpreted as fragments of minimum-bias (mainly low-$`Q^2`$) partons. When a soft-component Lévy distribution (defined as the limiting case of the $`n_{ch}`$ dependence) is subtracted from $`p_t`$ spectra for ten multiplicity classes we obtain distributions in Fig. 1 (first panel), described by hard-component model distributions (solid curves) with gaussian shape essentially independent of $`n_{ch}`$. We have transformed to transverse rapidity $`y_t\mathrm{ln}\{(m_t+p_t)/m_0\}`$ with $`m_0`$ assigned as the pion mass to provide a ‘native’ description of jets as hadron fragments from a moving source.
That novel single-particle result motivated a follow-up study of two-particle correlations on transverse rapidity . The minimum-bias distribution in Fig. 1 (second panel) represents all event multiplicities and charge combinations within the STAR $`(\eta ,\varphi )`$ detector acceptance. Separate soft and hard components are evident, as is the correspondence with hard-component gaussians in the first panel. The $`(y_{t1},y_{t2})`$ space can thus be used as a cut space to study trends of corresponding angular correlations on pseudorapidity $`\eta `$ and azimuth $`\varphi `$ for soft (string fragment) and hard (parton fragment) components, as shown in the third and fourth panels of Fig. 1.
A sequence of several analysis steps has led from study of single-particle $`p_t`$ spectra in a two-component context to isolation of soft and hard spectrum components to corresponding two-particle correlations on transverse rapidity to jet-like angular correlations obtained without imposition of a jet hypothesis. We thus achieve a model-independent analysis of parton scattering and fragmentation. Angular autocorrelations are not dependent on a leading or trigger particle and provide unprecedented access to low-$`Q^2`$ partons. The lower limit on fragment $`p_t`$ is determined by the fragmentation process itself, not the analysis method. We observe jet correlations in p-p collisions down to $`p_t`$ = 0.35 GeV/c for both particles. The main subject of this paper is the properties of fragment distributions on rapidity and angle from low-$`Q^2`$ partons in p-p collisions.
## 3 Analysis method and data
Correlation measure $`\mathrm{\Delta }\rho /\sqrt{\rho _{ref}}`$ is closely related to Pearson’s normalized covariance or correlation coefficient . For event-wise particle counts $`n_a`$ and $`n_b`$ in histogram bins $`a`$ and $`b`$ on single-particle space $`x`$ and pair counts in corresponding bin $`(a,b)`$ on two-particle space $`(x_1,x_2)`$ we obtain Pearson’s coefficient $`r_{ab}\overline{(n\overline{n})_a(n\overline{n})_b}/\sqrt{\overline{(n\overline{n})_a^2}\overline{(n\overline{n})_b^2}}`$ averaged over an event ensemble. Pearson’s coefficient is approximated by density ratio $`\mathrm{\Delta }\rho /\sqrt{\rho _{ref}}1/ϵ_x\overline{(n\overline{n})_a(n\overline{n})_b}/\sqrt{\overline{n}_a\overline{n}_b}`$, where $`ϵ_x`$ is the histogram bin size on $`x`$ and Poisson values of the number variances in the denominator have been substituted. $`\mathrm{\Delta }\rho /\sqrt{\rho _{ref}}`$, estimating the density of correlated pairs per particle, is our basic correlation measure for two-particle distributions on $`(y_{t1},y_{t2})`$, $`(\eta _1,\eta _2)`$ and $`(\varphi _1,\varphi _2)`$. To view angular correlations compactly we combine spaces $`(\eta _1,\eta _2)`$ and $`(\varphi _1,\varphi _2)`$ by means of a joint autocorrelation, providing projection (by averaging) to a lower-dimensional space with essentially no information loss. An autocorrelation for 1D primary space $`x`$ is constructed by averaging coefficients $`r_{a,a+k}`$ over index $`a`$ along the $`k^{th}`$ diagonal on $`(x_1,x_2)`$. For angle space $`(\eta ,\varphi )`$ we average simultaneously along diagonals on $`(\eta _1,\eta _2)`$ and $`(\varphi _1,\varphi _2)`$ to obtain a joint autocorrelation, a projection by averaging of the full two-particle angle space onto its difference variables $`\eta _\mathrm{\Delta }\eta _1\eta _2`$ and $`\varphi _\mathrm{\Delta }\varphi _1\varphi _2`$.
Data for this analysis were selected from 12 million p-p minimum-bias collision events at $`\sqrt{s_{NN}}=200`$ GeV obtained with the STAR detector at RHIC . Non-single-diffractive collisions were defined by a trigger coincidence between beam-beam counters positioned symmetrically in $`3.5|\eta |5.0`$ . Particle tracks were reconstructed with the STAR Time Projection Chamber (TPC) operating within a uniform 0.5 T magnetic field parallel to the beam axis. Tracks were accepted within full azimuth, $`|\eta |1.0`$ and $`0.15p_t6`$ GeV/c. TPC tracks and transverse beam position were used to reconstruct the primary or collision vertex.
## 4 p-p number correlations on $`(y_{t1},y_{t2})`$
Use of transverse rapidity $`y_t`$ to study parton fragmentation insures more compact distributions and relates directly to fundamental QCD aspects such as the double log approximation. It also provides a common basis, with longitudinal rapidity, for comparing string and parton fragmentation. Particle pairs from nuclear collisions can be separated on azimuth difference variable $`\varphi _\mathrm{\Delta }\varphi _1\varphi _2`$ into same-side (SS) ($`|\varphi _\mathrm{\Delta }|<\pi /2`$) and away-side (AS) ($`|\varphi _\mathrm{\Delta }|>\pi /2`$) pairs. Fig. 2 shows same-side (left two panels) and away-side (right two panels) correlations of the form $`\mathrm{\Delta }\rho /\sqrt{\rho _{ref}}`$ (normalized covariance) distributed on transverse rapidity subspace $`(y_{t1},y_{t2})`$. Each pair of panels represents like-sign (LS) and unlike-sign (US) charge combinations (left and right respectively). Structure can be separated into a soft component ($`y_t<2`$ for both pair partners) and a hard component ($`y_t>2`$ for both partners). Both components are strongly dependent on charge-sign combination and azimuth opening angle (same-side or away-side pairs). Description of the correlation structure is simplest in terms of sum and difference variables $`y_{t\mathrm{\Sigma }}y_{t1}+y_{t2}`$ and $`y_{t\mathrm{\Delta }}y_{t1}y_{t2}`$.
In the left (SS) panels, the LS soft component is interpreted as quantum correlations (HBT). The LS hard component along the diagonal is small and may itself be dominated by quantum correlations (from parton fragmentation). The US hard component is a peak at $`y_t2.8`$ ($`p_t`$ 1 GeV/c) elongated along $`y_{t\mathrm{\Sigma }}`$. The US hard component runs continuously into the US soft component at lower $`y_t`$, which is suppressed due to transverse momentum conservation. In the right (AS) panels, the hard-component peaks for LS and US pairs are nearly symmetric about their centers and much broader on $`y_{t\mathrm{\Delta }}`$ than their SS US counterpart, with rapid falloff below $`y_{t\mathrm{\Sigma }}4`$ (hadron $`p_t0.5`$ GeV/c) and nearly equal amplitudes. The large US soft component represents longitudinal string fragmentation, subject to local momentum and charge conservation (charge ordering on $`z`$ and therefore contributing negligibly to AS LS correlations.
## 5 p-p number correlations on $`(\eta _\mathrm{\Delta },\varphi _\mathrm{\Delta })`$
The $`(y_{t1},y_{t2})`$ correlations in the previous section are directly related to angular correlations on $`(\eta ,\varphi )`$. To isolate soft and hard components of p-p angular correlations we define soft pairs by $`y_t<2`$ ($`p_t<0.5`$ GeV/c) and hard pairs by $`y_t>2`$ for each particle. Fig. 3 shows minimum-bias correlations (all pairs for all multiplicity classes) of the form $`\mathrm{\Delta }\rho /\sqrt{\rho _{ref}}`$ on $`(\eta _\mathrm{\Delta },\varphi _\mathrm{\Delta })`$ for the soft component (left two panels) and hard component (right two panels), with charge combinations like-sign (LS) and unlike-sign (US) (left and right respectively). The first panel is dominated by a 2D gaussian peak at the origin representing quantum correlations. The US combination in the next panel is dominated by a 1D gaussian peak on $`\eta _\mathrm{\Delta }`$ arising from local charge conservation during string fragmentation (producing charge ordering along the collision axis). That trend is suppressed near the origin (the depression on azimuth of the $`\eta _\mathrm{\Delta }`$ gaussian) due to local transverse-momentum conservation. The narrow peak at the origin represents electron-positron pairs from photon conversions. Except for the HBT contribution all structure in the soft component is apparently a consequence of local measure conservation during string fragmentation. The resulting ‘canonical suppression’ is dependent on local particle density (event multiplicity).
The hard-component correlations in the right panels consist of a same-side peak at the origin and an away-side ridge. The LS same-side peak may be due to quantum correlations rather than jet fragmentation per se. The US same-side peak represents angular correlations of parton fragments (jet cone). The away-side hard-component correlations for LS and US pairs are essentially identical in shape and amplitude and reflect momentum conservation between scattered partons (dijets), including lack of structure on $`\eta _\mathrm{\Delta }`$ due to the broad distribution of parton-collision centers of momentum. These hard-component $`(\eta ,\varphi )`$ systematics, fully consistent with conventional expectations for high-$`p_t`$ jet angular correlations, are observed in this study for pairs of particles with both $`p_t`$s as low as 0.35 GeV/c ($`y_t1.6`$), much lower than previously observed with leading-particle methods. In what follows we emphasize low-$`Q^2`$ parton fragmentation.
## 6 Parton fragmentation: linear vs logarithmic presentation
We first consider parton fragments distributed on transverse momentum. A fragmentation function describes the conditional distribution of fragment momenta along a jet thrust axis given the parton momentum (as estimated directly from collision kinematics or from the observed fragment distribution or jet). Hadron momentum is conventionally normalized by the parton momentum/energy (e.g., $`x_p=p_t/E_{jet}`$), represented by $`E_{jet}`$ (reconstructed jets), $`\sqrt{s}/2`$ (e<sup>+</sup>-e<sup>-</sup>) or $`Q/2`$ (e-p deep-inelastic scattering or DIS). Alternatively, relative fragment momentum is measured logarithmically by $`\xi \mathrm{ln}(1/x)`$.
One semilog plotting convention is a logarithmic frequency distribution on linear transverse momentum, or normalized ratio $`x_p=p_t/E_{jet}`$, as shown in Fig. 4 (first panel) from e-p deep-inelastic scattering . Those distributions are similar to the hard component $`H_0`$ of the p-p single-particle $`p_t`$ spectrum when plotted on $`p_t`$ shown as the dashed curve in Fig. 4 (second panel). The fragmentation function plotted on linear $`p_t`$ is approximately exponential. An alternative convention is a linear frequency distribution on logarithmic transverse-momentum variable $`\xi \mathrm{ln}(1/x_p)\mathrm{ln}(E_{jet}/p_t)`$ as shown in Fig. 4 (third panel) . Those fragmentation functions (for two $`Q^2`$ intervals) can be compared with the hard components obtained from RHIC p-p spectra plotted on transverse rapidity in the fourth panel (the sense of increasing $`p_t`$ is reversed in the two cases). The conditional (on parton $`Q/2`$) fragmentation functions on $`\xi `$ in the third panel are approximately gaussian (further discussed below), as are the minimum-bias fragment distributions (no parton condition) in the fourth panel.
## 7 Fragmentation functions on transverse rapidity
The comparisons above motivated us to replot measured fragmentation functions on transverse rapidity $`y_t`$ as an infrared safe logarithmic momentum variable. In Fig. 5 (first two panels) we have replotted fragmentation functions from several collisions systems on transverse rapidity $`y_t(p_t;m_0)\mathrm{ln}\{(\sqrt{p_t^2+m_0^2}+p_t)/m_0\}`$, with $`m_0`$ assigned as the pion mass for all particles. We find that fragmentation functions so plotted are well described by beta distribution $`\beta (x;p,q)=x^{p1}(1x)^{q1}/B(p,q)`$, with beta function $`B(p,q)\mathrm{\Gamma }(p+q)/\mathrm{\Gamma }(p)\mathrm{\Gamma }(q)`$. The standard beta distribution defined on $`x[0,1]`$ is rescaled to $`\beta (y_t;p,q)`$ on $`[y_{t,min},y_{t,max}]`$. We observe that $`y_{t,min}`$ is 0.5 ($`p_t0.075`$ GeV/c) for e-e collisions and 1.5 ($`p_t0.25`$ GeV/c) for p-p collisions, the latter endpoint presumably reflecting the presence of the spectator parton system in the p-p case. The parton momentum is represented by $`y_{t,max}(X;m_0)\mathrm{ln}\{(\sqrt{X^2+m_0^2}+X)/m_0\}`$, with $`X=E_{jet}`$, $`\sqrt{s}/2`$ or $`Q/2`$ depending on the collision system. $`y_{t,max}`$ values for each fragmentation function in Fig. 5 are marked with short vertical lines. For e-e collisions the beta parameters are $`p=2.75`$ and $`q=3.5`$ for the cases examined; for p-p collisions the values are $`p2.8`$ and $`q4.2`$. The hard component from the STAR two-component analysis of p-p collisions is included in the first panel as the small gaussian curve (describing minimum-bias parton fragments with no $`y_{t,max}`$ condition). The energies in that panel are dijet energies $`E_{jj}=2E_{jet}`$. The data in the second panel are also plotted on $`\mathrm{ln}(p)`$ (comparable to $`\xi `$ but with the sense of increasing $`p_t`$ reversed) in the third panel. The solid curves are MLLA predictions , which reflect momentum conservation and QCD branching above the peaks but deviate strongly from data at lower momenta. In contrast, the beta distribution provides a precise description of the entire fragment distributions on transverse rapidity in the left panels.
The distribution systematics for e-e collisions in the second panel, interpreted in the context of the double-log approximation (DLA), imply that the e-e jet multiplicity trend can be described by $`n_{ch}=1+A(y_{t,max}y_{t,min})^2`$, whereas the modified leading-log approximation (MLLA) prediction is $`\mathrm{exp}\{a\mathrm{ln}(Q/2\mathrm{\Lambda })\}`$ . The fourth panel shows measured jet multiplicities vs jet energy ($`y_{t,max}`$) for e-e collisions (points) , our parameterization with $`A=0.8`$ (solid curve) and the MLLA prediction with $`a`$ chosen to provide the best fit (dashed curve). The hatched region corresponds to the low-$`Q^2`$ partons which dominate a minimum-bias p-p analysis. By a similar argument we can represent the most probable point (mode) as $`y_{t,peak}=B(y_{t,max}y_{t,min})+y_{t,min}`$, where for the beta distribution and e-e collisions $`B(p1)/(p+q2)=0.45`$, whereas the MLLA gives $`\xi _{peak}\mathrm{ln}(Q/2\mathrm{\Lambda })`$. Fragment distributions on transverse rapidity are consistent with conventional treatments (e.g., $`\xi `$) but are ‘infrared safe’ as $`p_t0`$.
## 8 Symmetrized two-particle fragment distributions
The conventional asymmetric treatment of parton fragments in terms of trigger and associated particles cannot access the low-$`Q^2`$ partons of greatest interest to us. We therefore symmetrize the analysis of fragment correlations, first for two-particle fragment distributions on $`y_t`$ and then for angular correlations. The same-side unlike-sign (US) $`(y_t,y_t)`$ correlations in Sec. 4 can be interpreted as a two-particle jet fragment distribution which we now model: we combine what is known about single-particle fragmentation functions with expectations for two-particle correlations to sketch the parameterization of a two-particle fragment distribution.
In the previous section we presented single-particle fragmentation functions plotted on transverse rapidity $`y_t`$. An example of fragmentation functions plotted on $`\xi `$ from the H1 experiment at HERA is shown in Fig. 4 (third panel) for two intervals on $`Q/2`$. Based on those trends we sketch a 2D distribution relating parton momentum $`Q/2`$ and fragment momentum $`p_t`$, both in rapidity format, as shown in the first two panels of Fig. 6. The first panel shows the locus of modes (most probable points, dashed line with smaller slope) for a continuum of fragmentation functions varying with parton rapidity on the horizontal axis, with fragment rapidity on the vertical axis. To construct the two-particle distribution we require a curve smoothly joining that line of modes and the diagonal with unit slope (solid line) given by the solid curve. The plot is then symmetrized. In the second panel we sketch the corresponding 2D distribution of fragmentation functions vs parton rapidity $`y_{t,max}`$. The diagonal dotted line represents parton momentum $`Q/2`$ plotted as a rapidity on the fragment axis. The dashed line is the line of modes. The fragmentation function for a given parton momentum would be a vertical slice from this distribution.
In the third panel we symmetrize the distribution from the second panel because in our $`(y_t,y_t)`$ analysis we correlate fragments with fragments symmetrically, rather than fragments with partons. That sketch can be compared with data in the fourth plot for unlike-sign same-side pairs – the minimum-bias two-particle fragment distribution. For low-$`Q^2`$ partons, which dominate the minimum-bias parton distribution, the two-particle fragment distribution becomes symmetric about the sum diagonal (and the fragment number $`2`$). This parton-fragment distribution is biased by the requirement of fragment pairs: ‘fragmentation’ to single hadrons cannot contribute to this plot. Because of the symmetry, the more appropriate variables to describe the distribution are sum and difference rapidities $`y_{t\mathrm{\Sigma }}y_{t1}+y_{t2}`$ and $`y_{t\mathrm{\Delta }}y_{t1}y_{t2}`$ as illustrated in the first panel. The comparison between sketch and data is informative. The small excess yield at low $`y_t`$ in the fourth panel is string fragments, partially suppressed by transverse momentum conservation.
## 9 Jet morphology for low-$`Q^2`$ partons
In the previous section we described parton fragment correlations on transverse rapidity. Jet structure is also characterized by two-particle angular correlations of hadron fragments, both intra-jet (within one jet) and inter-jet (between opposing dijets). The conventional method to describe angular correlations of fragments in the absence of full jet reconstruction is as a conditional distribution relative to a leading or trigger particle (e.g., the highest-$`p_t`$ particle in an event) used as an estimator for the parton momentum. The distribution of ‘associated’ particles relative to the trigger particle is an approximation to full jet reconstruction. Angular correlations relative to the trigger reveal 2D same-side peaks with (possibly different) widths on $`\eta `$ and $`\varphi `$. From those widths the mean momenta of fragments transverse to the jet thrust axis in the two angular directions $`(\eta ,\varphi )`$ are inferred.
In this minimum-bias study we encounter fragment distributions with a most-probable $`p_t`$ of about 1 GeV/c. The most probable fragment multiplicity is 2, with approximately equal fragment momenta, and the parton $`Q/2`$ is thus somewhat more than 2 GeV—comparable to the intrinsic parton $`k_t`$. For jets from low-$`Q^2`$ partons we cannot differentiate trigger and associated particles. We want a symmetrized formulation of angular correlations relative to which the leading-particle conditional treatment is a special case. For intra-jet correlations we define a symmetrized expression relating angular correlation width and transverse momentum $`j_t`$ relative to thrust as
$`\overline{j_{t\varphi }^2}_{12}{\displaystyle \frac{\sqrt{p_{t,1}^2p_{t,2}^2}}{\sqrt{p_{t,2}^2/p_{t,1}^2}+\sqrt{p_{t,1}^2/p_{t,2}^2}}}\left\{\mathrm{sin}^2\mathrm{\Delta }\varphi _{SS}+2\mathrm{sin}^2\varphi _1\mathrm{sin}^2\varphi _2\right\}`$ (1)
$`{\displaystyle \frac{(m_\pi /2)^2\mathrm{exp}(y_{t\mathrm{\Sigma }})}{2\mathrm{cosh}(y_{t\mathrm{\Delta }})}}\left\{\mathrm{sin}^2(\varphi _\mathrm{\Delta }/\sqrt{2})_{SS}+2\mathrm{sin}^2\left(\varphi _\mathrm{\Delta }/2\sqrt{2}\right)^2\right\}`$
for $`\varphi `$, with a similar expression for $`\eta `$.
The inter-jet azimuth angular correlation of particle pairs from opposed jets (dijets) is wider than the intra-jet azimuth correlation. The excess width or ‘$`k_t`$ broadening’ is attributed to the initial (intrinsic) transverse momenta $`k_t`$ of partons prior to scattering, which produces an acoplanarity of the two jets with respect to the collision axis. The symmetrized relation between angular correlation width and inferred $`k_t`$ is
$`\overline{z^2k_{t\varphi }^2}_{12}{\displaystyle \frac{\sqrt{p_{t,1}^2p_{t,2}^2}}{\sqrt{p_{t,2}^2/p_{t,1}^2}+\sqrt{p_{t,1}^2/p_{t,2}^2}}}\left\{{\displaystyle \frac{\mathrm{sin}^2\mathrm{\Delta }\varphi _{AS}\mathrm{sin}^2\mathrm{\Delta }\varphi _{SS}}{12\mathrm{sin}^2\mathrm{\Delta }\varphi _{SS}}}\right\}`$ (2)
$`{\displaystyle \frac{(m_\pi /2)^2\mathrm{exp}(y_{t\mathrm{\Sigma }})}{2\mathrm{cosh}(y_{t\mathrm{\Delta }})}}\left\{{\displaystyle \frac{\mathrm{sin}^2(\varphi _\mathrm{\Delta }/\sqrt{2})_{AS}\mathrm{sin}^2(\varphi _\mathrm{\Delta }/\sqrt{2})_{SS}}{12\mathrm{sin}^2(\varphi _\mathrm{\Delta }/\sqrt{2})_{SS}}}\right\}.`$
Those symmetrized expressions are used to extract fragment transverse momenta relative to thrust axis from angular correlations from p-p collisions. Fragment angular correlations are thus obtained for previously inaccessible low-$`Q^2`$ partons.
## 10 Angular correlation measurements
Fig. 8 (first panel) shows two-particle correlations on transverse-rapidity space $`(y_{t1},y_{t2})`$ described on sum and difference variables $`y_{t\mathrm{\Sigma }}y_{t1}+y_{t2}`$ and $`y_{t\mathrm{\Delta }}y_{t1}y_{t2}`$. Hard-component fractions for angular correlation measurements are defined by the grid of rectangles or bins (the narrow yellow lines) along $`y_{t\mathrm{\Sigma }}`$ numbered $`1,\mathrm{},12`$. The solid green boxes in the upper-right corner represent regions typically explored in leading-particle analyses based on high-$`p_t`$ ‘trigger’ particles . The dashed extensions represent cuts for extended associated-particle conditions recently applied to heavy ion collisions . The first two panels of Fig. 7 show angular autocorrelations for the second and eleventh bins on $`y_{t\mathrm{\Sigma }}`$ in Fig. 8 as examples. In the first panel the most probable combination is two particles each with $`p_t0.6`$ GeV/c. The same-side peak (jet cone) is broad but well-defined. The away-side ridge (particle pairs from dijets) is uniform on $`\eta _\mathrm{\Delta }`$ as expected. This jet correlation is a remarkable result for particles with such low $`p_t`$ and illustrates the power of the autocorrelation technique. The second panel shows correlations for the 11<sup>th</sup> bin, corresponding to $`p_t2.5`$ GeV/c for each particle. The same-side cone is much narrower, as is the away-side peak on azimuth $`\varphi `$, a result more typical of a high-$`p_t`$ leading-particle analysis.
The general trend is monotonic reduction of peak widths on $`(\eta _\mathrm{\Delta },\varphi _\mathrm{\Delta })`$ with increasing $`y_{t\mathrm{\Sigma }}`$. The angular widths are determined by the mean fragment momentum $`\widehat{ȷ}_t`$ perpendicular to the parton momentum (jet thrust) axis, relative to the particle $`p_t`$ in the form $`\widehat{ȷ}_t/p_t`$. We extract from those angular correlations near-side widths on $`\eta _\mathrm{\Delta }`$ and $`\varphi _\mathrm{\Delta }`$, and an away-side width on $`\varphi _\mathrm{\Delta }`$. We fit the same-side ($`|\varphi _\mathrm{\Delta }|<\pi /2`$) angular correlations (jet cone) with a 2D Gaussian and extract widths on pseudorapidity and azimuth for like-sign (LS), unlike-sign (US) and total (CI) pairs. Those widths are plotted in the last two panels of Fig. 7 vs $`\mathrm{exp}\{y_{t\mathrm{\Sigma }}/2\}1/\sqrt{p_{t1}p_{t2}}`$, that is, proportional to the inverse geometric mean of the pair momenta. The corresponding mean particle $`p_t`$ is given at the top of the third panel.
The dashed diagonal lines in those panels represent ‘$`\widehat{ȷ}_t`$ scaling’: the assumption that $`\widehat{ȷ}_t`$ (r.m.s. $`j_t`$) is independent of hadron momentum as a limiting case for larger $`p_t`$. We observe a dramatic departure from that expectation for these low fragment momenta ($`p_t<2.5`$ GeV/c). However, with the autocorrelation technique we observe that jet structure is very well defined, even down to hadron pairs with particle $`p_t0.35`$ GeV/c. The hatched regions to the left in those panels show the values measured for minimum-bias hard-component hadrons, which are weighted toward lower $`p_t`$. In the following section we examine the systematics of $`\widehat{ȷ}_t`$ vs $`y_{t\mathrm{\Sigma }}`$ separately on $`\eta `$ and $`\varphi `$, and for values of parton $`Q^2`$ previously inaccessible.
## 11 Inferred $`\widehat{ȷ}_t`$ and $`\widehat{k}_t`$ distributions
By applying Eq. (1) with the values of $`(y_{t\mathrm{\Sigma }},y_{t\mathrm{\Delta }})`$ defined on the $`(y_{t1},y_{t2})`$ cut space as in Fig. 6 (first panel) we obtain $`\widehat{ȷ}_t`$ values for $`x=\eta ,\varphi `$ from the corresponding angular widths. In Fig. 8 (second and third panels) we observe a strong decrease in both $`\widehat{ȷ}_t`$ components with decreasing $`y_{t\mathrm{\Sigma }}`$ ($``$ parton $`Q^2`$). We also observe substantial differences in the $`\widehat{ȷ}_t`$ $`\eta `$ and $`\varphi `$ components, converging however to the ‘$`\widehat{ȷ}_t`$ scaling’ expectation at larger $`y_{t\mathrm{\Sigma }}`$ (hatched bands). It is notable that the orientation of the $`\widehat{ȷ}_t`$ asymmetry (major axis lying in the $`(\varphi ,p_t)`$ plane) is parallel to the ‘contact plane’ of the parton collision. That asymmetry orientation is orthogonal to what is observed for low-$`Q^2`$ jet fragments in central heavy ion collisions .
We also measured the angular width of the away-side azimuth correlation of hadron pairs ($`|\varphi _1\varphi _2|>\pi /2`$) and compared it with the same-side azimuth correlation to infer the mean transverse momentum $`\widehat{k}_t`$ of partons prior to scattering. We converted the angular widths to $`\widehat{k}_t`$ values according to Eq. (2), where $`z`$ estimates the mean value of the ratio of ‘trigger particle’ momentum to parton momentum. At larger $`p_t`$ that value is approximately 0.75, but is a falling function of $`p_t`$ for smaller fragment $`p_t`$, as in this case. Below $`y_{t\mathrm{\Sigma }}/23.3`$ (hadron $`p_t2`$ GeV/c) we observe that $`\widehat{k}_t`$ values fall below the $`k_t`$-scaling trend (hatched band) by an increasing amount with smaller $`y_{t\mathrm{\Sigma }}`$ (smaller $`Q^2`$). The trend of Pythia $`\widehat{k}_t`$ values (dotted line) is substantially displaced below the data.
## 12 Summary
We have presented a broad survey of newly-obtained two-particle correlations from 200 GeV p-p collisions at RHIC. Correlations from longitudinal string fragmentation and transverse scattered parton fragmentation are clearly distinguished. The former is dominated by local momentum and charge conservation. The latter provides new access to scattering and fragmentation of minimum-bias partons. Using newly-devised analysis techniques we find that low-$`Q^2`$ parton fragmentation is accessible down to hadron $`p_t`$ = 0.35 GeV/c for both hadrons of a correlated pair. Jet morphology for low-$`Q^2`$ partons requires a more general treatment of fragment $`p_t`$ distributions and angular correlations. Fragment distributions on transverse rapidity $`y_t`$ are ‘infrared safe’ and exhibit interesting systematic behaviors which can be compared with the MLLA. Jet angular correlations show a dramatic asymmetry about the thrust axis at low-$`Q^2`$ (up to 2:1 aspect ratio, with larger width in the azimuth direction), possibly related to nonperturbative details of soft parton collisions. These measurements of p-p correlations provide an important reference for the study of in-medium modification of parton scattering and fragmentation in heavy ion collisions.
This work was supported in part by the Office of Science of the U.S. DoE under grant DE-FG03-97ER41020.
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# Eigenvalue problems of Nordhaus-Gaddum type
## 1 Introduction
Our notation is standard (e.g., see , , and ); in particular, all graphs are defined on the vertex set $`\{1,2,\mathrm{},n\}=\left[n\right]`$ and $`G(n,m)`$ stands for a graph with $`n`$ vertices and $`m`$ edges. We write $`\mathrm{\Gamma }\left(u\right)`$ for the set of neighbors of the vertex $`u`$ and set $`d\left(u\right)=\left|\mathrm{\Gamma }\left(u\right)\right|.`$ Given a graph $`G`$ of order $`n,`$ we assume that the eigenvalues of the adjacency matrix of $`G`$ are ordered as $`\mu \left(G\right)=\mu _1\left(G\right)\mathrm{}\mu _n\left(G\right)`$. As usual, $`\overline{G}`$ denotes the complement of a graph $`G`$ and $`\omega (G)`$ stands for the clique number of $`G.`$
Nosal showed that for every graph $`G`$ of order $`n,`$
$$n1\mu \left(G\right)+\mu \left(\overline{G}\right)<\sqrt{2}n.$$
(1)
Quite of attention has been given to second of these inequalities. In it was shown that
$$\mu \left(G\right)+\mu \left(\overline{G}\right)\sqrt{\left(2\frac{1}{\omega (G)}\frac{1}{\omega (\overline{G})}\right)n\left(n1\right)},$$
(2)
improving earlier results in , , , and . Unfortunately inequality (2) is not much better then (1) when both $`\omega (G)`$ and $`\omega (\overline{G})`$ are large enough. Thus, it is natural to ask whether $`\sqrt{2}`$ in (1) can be replaced by a smaller absolute constant for $`n`$ sufficiently large. In this note we answer this question in the positive but first we state a more general problem.
###### Problem 1
For every $`1kn`$ find
$$f_k\left(n\right)=\underset{v\left(G\right)=n}{\mathrm{max}}\left|\mu _k\left(G\right)\right|+\left|\mu _k\left(\overline{G}\right)\right|.$$
It is difficult to determine precisely $`f_k\left(n\right)`$ for every $`n`$ and $`k,`$ so at this stage it seems more practical to estimate it asymptotically. In this note we show that
$$\frac{4}{3}n2f_1\left(n\right)<\left(\sqrt{2}c\right)n$$
(3)
for some $`c>8\times 10^7`$ independent of $`n.`$ For $`f_2\left(n\right)`$ we give the following tight bounds
$$\frac{\sqrt{2}}{2}n3<f_2\left(n\right)<\frac{\sqrt{2}}{2}n.$$
(4)
We also show that
$$\frac{\sqrt{2}}{2}n3<f_n\left(n\right)\frac{\sqrt{3}}{2}n$$
(5)
Finally for fixed $`k,`$ $`2<k<n,`$ and $`n`$ large, we prove that
$`{\displaystyle \frac{n}{k}}1`$ $`f_k\left(n\right)\sqrt{{\displaystyle \frac{2}{k}}}n,`$
$`{\displaystyle \frac{n}{k}}+1`$ $`f_{nk}\left(n\right)\sqrt{{\displaystyle \frac{2}{k}}}n.`$
## 2 Bounds on $`f_1\left(n\right)`$
Before stating the main result of this section, we shall recall two auxiliary results whose proofs can be found in . Given a graph $`G=G(n,m),`$ let
$$s\left(G\right)=\underset{uV\left(G\right)}{}\left|d\left(u\right)\frac{2m}{n}\right|.$$
###### Proposition 2
For every graph $`G=G(n,m),`$
$$\frac{s^2\left(G\right)}{2n^2\sqrt{2m}}\mu _1\left(G\right)\frac{2m}{n}\sqrt{s\left(G\right)},$$
(6)
and
$$\mu _n\left(G\right)+\mu _n\left(\overline{G}\right)1\frac{s^2\left(G\right)}{n^3}.$$
(7)
Decreasing the constant $`\sqrt{2}`$ in (1) happened to be a surprisingly challenging task for the author. The little progress that has been made is given in the following theorem.
###### Theorem 3
There exists $`c8\times 10^7`$ such that
$$\mu _1\left(G\right)+\mu _1\left(\overline{G}\right)\left(\sqrt{2}c\right)n.$$
for every graph $`G`$ of order $`n`$.
Proof Assume the opposite: let $`\epsilon =8\times 10^7`$ and let there exist a graph $`G`$ of order $`n`$ such that
$$\mu _1\left(G\right)+\mu _1\left(\overline{G}\right)>\left(\sqrt{2}\epsilon \right)n.$$
Writing $`A\left(G\right)`$ for the adjacency matrix of $`G,`$ we have
$$\underset{i=1}{\overset{n}{}}\mu _i^2\left(G\right)=tr\left(A^2\left(G\right)\right)=2e\left(G\right),$$
(8)
implying that
$$\mu _1^2\left(G\right)+\mu _n^2\left(G\right)+\mu _1^2\left(\overline{G}\right)+\mu _n^2\left(\overline{G}\right)2e\left(G\right)+2e\left(\overline{G}\right)<n^2.$$
From
$$\mu _1^2\left(G\right)+\mu _1^2\left(\overline{G}\right)\frac{1}{2}\left(\mu _1\left(G\right)+\mu _1\left(\overline{G}\right)\right)^2>\left(1\frac{\epsilon }{\sqrt{2}}\right)^2n^2>\left(12\epsilon \right)n^2$$
we find that
$$\left|\mu _n\left(G\right)\right|+\left|\mu _n\left(\overline{G}\right)\right|\sqrt{2\left(\mu _n^2\left(G\right)+\mu _n^2\left(\overline{G}\right)\right)}<\sqrt{4\epsilon }n^2,$$
(9)
and so, $`\mu _n\left(G\right)+\mu _n\left(\overline{G}\right)>\sqrt{4\epsilon }n.`$ We thus have $`\sqrt{4\epsilon }n^4s^2\left(G\right).`$ On the other hand, by (6) and in view of $`s\left(G\right)=s\left(\overline{G}\right),`$ we see that
$$\mu _1\left(G\right)+\mu _1\left(\overline{G}\right)n1+2\sqrt{s\left(G\right)}<n+2\sqrt{s\left(G\right)},$$
and, by (9), it follows that
$$\left(\sqrt{2}\epsilon \right)n<n+2\left(4\epsilon \right)^{1/8}n.$$
Dividing by $`n`$, we obtain $`\left(\sqrt{2}1\right)<\epsilon +2^{5/4}\epsilon ^{1/8}`$, a contradiction for $`\epsilon =8\times 10^7.`$ $`\mathrm{}`$
It is certain that the upper bound given by Theorem 3 is far from the best one. We shall give below a lower bound on $`f_1\left(n\right)`$ which seems to tight.
For every $`1r<n`$ the graph $`G=K_r+\overline{K_{nr}}`$ (see, e.g. ) satisfies
$$\mu _1\left(G\right)+\mu _1\left(\overline{G}\right)=\frac{r1}{2}+\sqrt{nr\frac{3r^2+2r1}{4}}+nr1=n\frac{r+3}{2}+\sqrt{nr\frac{3r^2+2r1}{4}}.$$
The right-hand side of this inequality is increasing in $`r`$ for $`0r\left(n1\right)/3`$ and we find that
$$f_1\left(n\right)>\frac{4n}{3}2.$$
This gives some evidence for the following conjecture.
###### Conjecture 4
$$f_1\left(n\right)=\frac{4n}{3}+O\left(1\right).$$
We conclude this section with an improvement of the lower bound in (1). Using the first of inequalities (6) we obtain
$`\mu _1\left(G\right)+\mu _1\left(\overline{G}\right)`$ $`n1+{\displaystyle \frac{s^2\left(G\right)}{2n^2}}\left({\displaystyle \frac{1}{\sqrt{2e\left(G\right)}}}+{\displaystyle \frac{1}{\sqrt{2e\left(\overline{G}\right)}}}\right)`$
$`n1+\sqrt{2}{\displaystyle \frac{s^2\left(G\right)}{n^3}}.`$
## 3 A class of graphs
In this section we shall describe a class of graphs that give the right order of $`f_2\left(G\right)`$ and, we believe, also of $`f_n\left(G\right).`$
Let $`n4`$. Partition $`\left[n\right]`$ in $`4`$ classes $`A,B,C,D`$ so that $`\left|A\right|\left|B\right|\left|C\right|\left|D\right|\left|A\right|1.`$ Join every two vertices inside $`A`$ and $`D,`$ join each vertex in $`B`$ to each vertex in $`AC,`$ join each vertex in $`D`$ to each vertex in $`C.`$ Write $`G\left(n\right)`$ for the resulting graph.
Note that if $`n`$ is divisible by $`4,`$ the sets $`A,B,C,D`$ have equal cardinality and we see that $`G\left(n\right)`$ is isomorphic to its complement.
Our main goal to the end of this section is to estimate the eigenvalues of $`G\left(n\right).`$ Write $`ch\left(A\right)`$ for the characteristic polynomial of a matrix $`A.`$ The following general theorem holds.
###### Theorem 5
Suppose $`G`$ is a graph and $`V\left(G\right)=_{i=1}^kV_i`$ is a partition in sets of size $`n`$ such that
(i) for all $`1ik,`$ either $`e\left(V_i\right)=\left(\genfrac{}{}{0pt}{}{n}{2}\right)`$ or $`e\left(V_i\right)=0`$;
(ii) for all $`1i<jk,`$ either $`e(V_i,V_j)=n^2`$ or $`e(V_i,V_j)=0.`$
Let the sets $`V_1,\mathrm{},V_p`$ be independent and $`V_{p+1},\mathrm{},V_k`$ induce a complete graph. Then for the characteristic polynomial of the adjacency matrix of $`G`$ we have
$$ch\left(A\left(G\right)\right)=x^{pnp}\left(1x\right)^{\left(kp\right)n\left(kp\right)}ch\left(R\right),$$
where $`R=\left(r_{ij}\right)`$ is a $`k\times k`$ matrix such that
$$r_{ij}=\{\begin{array}{ccc}0\hfill & \text{if }ij\hfill & \text{and }e(V_i,V_j)=0\hfill \\ n\hfill & \text{if }ij\hfill & \text{and }e(V_i,V_j)=n^2\text{ }\hfill \\ 0\hfill & \text{if }i=j\hfill & \text{and }e\left(V_i\right)=0\hfill \\ n1\hfill & \text{if }i=j\hfill & \text{and }e\left(V_i\right)=\left(\genfrac{}{}{0pt}{}{n}{2}\right)\hfill \end{array}.$$
The proof of this theorem is a straight exercise in determinants, so we shall omit it.
If $`n`$ is divisible by $`4,`$ say $`n=4k,`$ by Theorem 5, for the characteristic polynomial of $`A\left(G\left(n\right)\right)`$ we have
$$ch\left(A\left(G\left(n\right)\right)\right)=x^{2k2}\left(1x\right)^{2k2}\left[\begin{array}{cccc}k1x& k& 0& 0\\ k& x& k& 0\\ 0& k& x& k\\ 0& 0& k& k1x\end{array}\right].$$
By straightforward calculations, setting $`a=11/k`$ and $`y=x/k`$, we see that
$`ch\left(A\left(G\left(n\right)\right)\right)`$ $`=x^{2k2}\left(1x\right)^{2k2}\left[\left(ay\right)\left(y^2\left(ay\right)+2ya\right)\left(y^2ay1\right)\right]`$
$`=x^{2k2}\left(1x\right)^{2k2}\left(y^2\left(1+a\right)y\left(1a\right)\right)\left(y^2+\left(1a\right)y\left(a+1\right)\right).`$
Hence, we find that
$`\mu _2\left(G\right)`$ $`={\displaystyle \frac{1}{2}}+\sqrt{{\displaystyle \frac{1}{4}}+2{\displaystyle \frac{n}{4}}^2{\displaystyle \frac{n}{4}}}`$
$`\mu _n\left(G\right)`$ $`={\displaystyle \frac{1}{2}}\sqrt{{\displaystyle \frac{1}{4}}+2{\displaystyle \frac{n}{4}}^2{\displaystyle \frac{n}{4}}.}`$
If $`n`$ is not divisible by $`4,`$ we will give some tight estimates of $`\mu _2\left(G\right)`$ and $`\mu _n\left(G\right).`$ Notice first that $`G\left(4n/4\right)`$ is an induced graph of $`G\left(n\right)`$ which in turn is an induced graph of $`G\left(4n/4\right).`$ Thus the adjacency matrix of $`G\left(4n/4\right)`$ is a principal submatrix of the adjacency matrix of $`G\left(n\right)`$ which in turn is a principal submatrix of the adjacency matrix of $`G\left(4n/4\right).`$ Since the eigenvalues of a matrix and its principal matrices are interlaced (, Theorem 4.3.15), we obtain
$`{\displaystyle \frac{1}{2}}+\sqrt{{\displaystyle \frac{1}{4}}+2{\displaystyle \frac{n}{4}}^2{\displaystyle \frac{n}{4}}}`$ $`\mu _2\left(G\right){\displaystyle \frac{1}{2}}+\sqrt{{\displaystyle \frac{1}{4}}+2{\displaystyle \frac{n}{4}}^2{\displaystyle \frac{n}{4}}},`$ (10)
$`{\displaystyle \frac{1}{2}}\sqrt{{\displaystyle \frac{1}{4}}+2{\displaystyle \frac{n}{4}}^2{\displaystyle \frac{n}{4}}}`$ $`\mu _n\left(G\right){\displaystyle \frac{1}{2}}\sqrt{{\displaystyle \frac{1}{4}}+2{\displaystyle \frac{n}{4}}^2{\displaystyle \frac{n}{4}}}.`$ (11)
## 4 The asymptotics of $`f_2\left(n\right)`$
In this section we shall prove inequalities (4). From (10) we readily have
$$f_2\left(n\right)\frac{1}{2}+\sqrt{\frac{1}{4}+2\frac{n}{4}^2\frac{n}{4}}>\frac{\sqrt{2}}{2}n3,$$
so all we need to prove is that $`f_2\left(n\right)n/\sqrt{2}`$.
By (8) we have
$$\mu _1^2\left(G\right)+\mu _2^2\left(G\right)+\mu _n^2\left(G\right)+\mu _1^2\left(\overline{G}\right)+\mu _2^2\left(\overline{G}\right)+\mu _n^2\left(\overline{G}\right)n\left(n1\right).$$
(12)
By Weyl’s inequalities (, p. 181), for every graph $`G`$ of order $`n,`$ we have
$$\mu _2\left(G\right)+\mu _n\left(\overline{G}\right)\mu _2\left(K_n\right)=1.$$
Hence, using $`\mu _20`$ and $`\mu _n1`$ we obtain
$$\mu _2^2\left(G\right)\mu _n^2\left(\overline{G}\right)+2\mu _n\left(\overline{G}\right)+1<\mu _n^2\left(\overline{G}\right).$$
Hence, from (12) and $`\mu _1\left(G\right)+\mu _1\left(\overline{G}\right)n1,`$ we find that
$$\frac{\left(n1\right)^2}{2}+2\mu _2^2\left(G\right)+2\mu _2^2\left(\overline{G}\right)\mu _1^2\left(G\right)+\mu _2^2\left(G\right)+\mu _n^2\left(G\right)+\mu _1^2\left(\overline{G}\right)+\mu _2^2\left(\overline{G}\right)+\mu _n^2\left(\overline{G}\right)n\left(n1\right).$$
After some algebra, we deduce that
$$\mu _2\left(G\right)+\mu _2\left(\overline{G}\right)\frac{\sqrt{2}}{2}n,$$
completing the proof of inequalities (4).
## 5 Bounds on $`f_n\left(n\right)`$
In this section we shall prove inequalities (5). From (11), as above, we have
$$f_n\left(n\right)>\frac{\sqrt{2}}{2}n3.$$
We believe that, in fact, the following conjecture is true.
###### Conjecture 6
$$f_n\left(G\right)=\frac{\sqrt{2}n}{2}+O\left(1\right).$$
However we can only prove that $`f_n\left(G\right)<\left(\sqrt{3}/2\right)n`$ which is implied by the following theorem.
###### Theorem 7
For every graph $`G`$ of order $`n,`$
$$\mu _n^2\left(G\right)+\mu _n^2\left(\overline{G}\right)\frac{3}{8}n^2.$$
Proof Indeed, suppose $`(u_1,\mathrm{},u_n)`$ and $`(w_1,\mathrm{},w_n)`$ are eigenvectors to $`\mu _n\left(G\right)`$ and $`\mu _n\left(\overline{G}\right).`$ Let
$$U=\{i:u_i>0\},\text{ }W=\{i:w_i>0\}.$$
Setting $`V=\left[n\right],`$ we clearly have $`\mu _n^2\left(G\right)E_G(U,V\backslash U)`$ and $`\mu _n^2\left(\overline{G}\right)E_{\overline{G}}(W,V\backslash W)`$. Since $`E_G(U,V\backslash U)E_{\overline{G}}(W,V\backslash W)=\mathrm{},`$ we see that the graph
$$G^{}=(V,E_G(U,V\backslash U)E_{\overline{G}}(W,V\backslash W))$$
is at most $`4`$-colorable and hence $`G^{}`$ contains no $`4`$-cliques. By Turán’s theorem (e.g., see ), we obtain $`e\left(G^{}\right)\left(3/8\right)n^2,`$ completing the proof. $`\mathrm{}`$
## 6 Bounds on $`f_k\left(n\right),`$ $`2<k<n`$
In this section we shall give simple bounds on $`f_k\left(n\right)`$ for $`2<k<n`$. Write for the Turán graph of order $`n`$ with $`k`$ classes. Recall that $`T_k\left(n\right)`$ is a complete $`k`$-partite graph whose vertex classes differ by at most $`1`$ in size. We assume that $`k`$ is fixed and $`n`$ is large enough. Since $`\mu _k\left(T_k\left(n\right)\right)=0`$ and $`\mu _{nk}\left(T_k\left(n\right)\right)n/k`$ for $`n`$ large, we immediately have
$`f_k\left(n\right)`$ $`n/k1,`$
$`f_{nk}\left(n\right)`$ $`n/k+1.`$
We next turn to upper bounds on $`f_k\left(n\right).`$
###### Theorem 8
For any fixed $`k`$ and any graph $`G`$ of sufficiently large order $`n,`$
$$\left|\mu _k\left(G\right)\right|+\left|\mu _k\left(\overline{G}\right)\right|<\sqrt{\frac{2}{k}}n.$$
(13)
and
$$\left|\mu _{nk}\left(G\right)\right|+\left|\mu _{nk}\left(\overline{G}\right)\right|<\sqrt{\frac{2}{k}}n.$$
(14)
Proof Set $`e\left(G\right)=m.`$ Our first goal is to prove that $`\left|\mu _k\left(G\right)\right|\sqrt{2e\left(G\right)/k}.`$ If $`\mu _k\left(G\right)0,`$ we have in view of (8)
$$k\mu _k^2\left(G\right)\underset{i=1}{\overset{n}{}}\mu _i^2\left(G\right)=2m.$$
If $`\mu _k\left(G\right)<0`$ and $`\left|\mu _k\left(G\right)\right|>\sqrt{2m/k}`$ then
$$\underset{i=1}{\overset{n}{}}\mu _i^2\left(G\right)\left(nk\right)\mu _k^2\left(G\right)>2m\frac{nk}{k}>2m,$$
a contradiction. Hence, $`\left|\mu _k\left(G\right)\right|\sqrt{2e\left(G\right)/k},`$ and, by symmetry, $`\left|\mu _k\left(\overline{G}\right)\right|\sqrt{2e\left(\overline{G}\right)/k}.`$ Now
$$\left|\mu _k\left(G\right)\right|+\left|\mu _k\left(\overline{G}\right)\right|\sqrt{2e\left(G\right)/k}+\sqrt{2e\left(\overline{G}\right)/k}\sqrt{\frac{2}{k}n\left(n1\right)}<\sqrt{\frac{2}{k}}n,$$
proving inequality (13). The proof of inequality (14) goes along the same lines, so we will omit it. $`\mathrm{}`$
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# Submm line imaging of Orion-KL with the Submillimeter Array
## 1 Introduction
Orion-KL is the most studied region of massive star formation. At a distance of $``$450 pc, its molecular spectral line emission is very strong (e.g., Schilke et al. 1997a). Most molecular line studies toward Orion-KL have been carried out with single-dish instruments (e.g., Sutton et al. 1985; Blake et al. 1987; Schilke et al. 1997a, 2001; Comito et al. 2005) and thus do not resolve the spatial distribution of the molecular gas. Molecular line studies with interferometers have been conducted only at wavelengths longer than 1 mm (e.g., Wilner et al. 1994; Wright et al. 1996; Blake et al. 1996; Wilson et al. 2000; Liu et al. 2002). While many such studies were single-molecule observations, the first dedicated interferometric line-surveys were conducted by Wright et al. (1996) and Blake et al. (1996). The interferometric investigations found that the single-dish spectral features known as the hot core, compact ridge and plateau are all associated with the KL region whereas the dust continuum source CS1 about $`25^{\prime \prime }`$ north-east of source I exhibits narrower line widths, lower temperatures and only weak indications of star formation (e.g., Wright et al. 1996).
The hot core as traced by dust continuum emission and nitrogen-bearing molecules such as CH<sub>3</sub>CN or NH<sub>3</sub> appears as a chain of dense clumps offset by $`1^{\prime \prime }`$ from the radio source I with estimated temperatures between 130 and 335 K (e.g., Wilner et al. 1994; Wright et al. 1996; Wilson et al. 2000). The compact ridge a few arcseconds south-west of source I is reported to have lower temperatures (of the order 100 K) and to be especially prominent in oxygen-bearing molecules (e.g., Wright et al. 1996; Liu et al. 2002). Furthermore, the region exhibits a complex cluster of infrared sources studied from near- to mid-infrared wavelengths (Dougados et al., 1993; Greenhill et al., 2004; Shuping et al., 2004). At least two outflows are driven from the region on scales $`>10^4`$ AU, one high-velocity outflow in the south-east north-west direction observed in molecular lines and in the optical and near-infrared (e.g., Allen & Burton 1993; Wright et al. 1995; Chernin & Wright 1996; Schultz et al. 1999), and one lower velocity outflow in the north-east south-west direction best traced by thermal SiO and H<sub>2</sub>O maser emission as well as some H<sub>2</sub> bow shocks (e.g., Genzel & Stutzki 1989; Blake et al. 1996; Chrysostomou et al. 1997; Stolovy et al. 1998). The driving source(s) of the outflows are uncertain: initial claims that it might be IRc2 are outdated now, and possible culprits are the radio sources I and/or the infrared source n, also known as radio source L (Menten & Reid, 1995). Radio source I lies close to the center of the north-east south-west outflow (Gezari, 1992; Menten & Reid, 1995; Gezari et al., 1998; Greenhill et al., 2003) and has not been detected in the near- to mid-infrared (Greenhill et al., 2004).
A study of the spatial distribution of the different molecular lines can be used to identify and separate different centers of activity via their different excitation conditions. The Submillimeter Array (SMA<sup>1</sup><sup>1</sup>1The Submillimeter Array is a joint project between the Smithsonian Astrophysical Observatory and the Academia Sinica Institute of Astronomy and Astrophysics, and is funded by the Smithsonian Institution and the Academia Sinica.) on Mauna Kea was dedicated in November 2003, and Orion was one of its early science targets. Therefore, we present the first high-spatial-resolution ($`450`$ AU) submm wavelength molecular line investigation at 865 $`\mu `$m. Since the dataset is incredibly rich, a complete analysis and interpretation of all existing features is out of the scope of this paper. Here, we present the overall content of the dataset and focus on a few aspects in more detail. Follow-up investigations are currently undertaken, and the data are publicly available (raw data via the SMA web-site http://cfa-www.harvard.edu/rtdc/index-sma.html, and calibrated data from the lead author of this paper).
The submm continuum emission was presented in a separate paper (Beuther et al., 2004). The continuum data resolved source I from the hot core and detected, for the first time, source n at a wavelength shorter than 7 mm. Furthermore, a new continuum peak between the sources I and n is found, SMA1, which is either an independent deeply embedded protostellar object or part of the more extended hot core. These submm continuum sources, reproduced in Figure 1, will be used as a reference frame for the discussion of the line emission throughout this paper.
## 2 Observations
Orion-KL was observed with the Submillimeter Array (SMA) in December 2003 and February 2004 at 348 GHz ($`865\mu `$m) with 7 antennas in two configurations and with projected baselines between 15 and 230 k$`\lambda `$. The phase center was the nominal position of source I as given by Plambeck et al. (1995): R.A. \[J2000\] $`5^h35^m14.50^s`$ and Dec. \[J2000\] $`5^{}22^{}30^{\prime \prime }.45`$. For bandpass calibration we used Jupiter and Mars. The flux scale was derived from observations of Callisto and is estimated to be accurate to within 15%. Phase and amplitude calibration was done via frequent observations of the quasar 0420-014 about 17 from the phase center. The zenith opacities, measured with the NRAO tipping radiometer located at the Caltech Submillimeter Observatory, were excellent during both tracks with $`\tau (348\mathrm{G}\mathrm{H}\mathrm{z})0.1`$ & 0.14 (scaled from the 225 GHz measurement via $`\tau (348\mathrm{G}\mathrm{H}\mathrm{z})2.8\times \tau (225\mathrm{G}\mathrm{H}\mathrm{z})`$). The receiver operated in a double-sideband mode with an IF band of 4-6 GHz so that the upper and lower sideband were separated by 10 GHz. The correlator had a bandwidth of 2 GHz and the spectral resolution was 0.8125 MHz. Measured double-sideband system temperatures corrected to the top of the atmosphere were between 250 and 600 K, mainly depending on the elevation of the source.
The uv-coverage with no projected baselines below 15 k$`\lambda `$ implies that source structure on spatial scales $`13^{\prime \prime }`$ is filtered out by the observations. Interferometric observations of sources close to Declination $`0^{}`$ like Orion-KL result in a rather poor uv-coverage and dirty beam (Fig. 2), further complicating the image deconvolution. Since there is significant large-scale emission associated with most spectral lines, we deconvolved the Point Spread Function with the CLEAN-algorithm over a large spatial area ($`26^{\prime \prime }\times 26^{\prime \prime }`$, about 2/3 of the primary beam size) for each spectral line to derive the final images. It is possible that, for some spectral lines, specific box-cleaning could improve the final images but this is beyond the scope of this paper. We applied different uv-weightings for various lines depending mainly on the line strength, and obtained average synthesized beams of $`1.1^{\prime \prime }\times 0.9^{\prime \prime }`$ (P.A. $`44^{}`$).
To get a feeling for the amount of missing flux, we convolved the data to a $`20^{\prime \prime }`$ beam and compared them with the single-dish line survey by Schilke et al. (1997a). Many strong and extended lines have missing flux values up to 90%. For one of our main target lines, <sup>28</sup>SiO(8–7), the measured SMA flux is even lower, $`10\%`$ of that measured with the CSO (see also the discussion at the end of §4.1). Since this missing flux has to be distributed over scales $`>13^{\prime \prime }`$, possibly as large as the $`20^{\prime \prime }`$ beam of the CSO, the missing <sup>28</sup>SiO(8–7) flux scaled to the $`1^{\prime \prime }`$ synthesized SMA beam can be estimated to approximately 3 Jy beam<sup>-1</sup>. For additional discussion on missing flux in particular lines see §4.1 & 4.2.
The sensitivity was dynamic-range limited by the side-lobes of the strongest emission peaks and thus varied between the line maps of different molecules and molecular transitions. This limitation was due to the incomplete sampling of short uv-spacings and the presence of extended structures. The theoretical $`1\sigma `$ rms per 1 km s<sup>-1</sup> channel was $`100`$ mJy whereas the measured $`1\sigma `$ rms in a 1 km s<sup>-1</sup> channel image was of the order 500 mJy because of the dynamic-range problem and thus insufficient cleaning. The $`1\sigma `$ rms for the velocity-integrated molecular line maps (the velocity ranges for the integrations were chosen for each line separately depending on the line-widths and intensities) ranged between 70 and 320 mJy with a mean $`1\sigma `$ value of $`170`$ mJy. We calibrated the data within the IDL superset MIR developed for the Owens Valley Radio Observatory and adapted for the SMA; the imaging was performed in MIRIAD. For more details on the array and its capabilities see Ho et al. (2004).
## 3 Results
Figure 3 presents the whole observed bandpass with its extremely rich line forest, a closer zoom into each sideband is shown in Figures 4 & 5. The line identifications are not always unambiguous because different velocity components and large line-widths are present in the region. At first, the lines were identified by comparing our spectra with those of the previous single-dish study of Schilke et al. (1997a). Then, in selected cases we used the JPL, CDMS and LOVAS line catalogs (Poynter & Pickett, 1985; Müller et al., 2001; Lovas, 2004) for confirmations and further identifications. Altogether, we find about 145 lines above an approximate flux level of 8 Jy on the short baseline of 21 m (Figs. 3, 4 & 5). We note that this is not the sensitivity limit of the data but rather a cutoff below which confusion becomes an even larger problem and line identifications get considerably more difficult. Of the 145 lines listed in Table 1, we are able to identify approximately 90%. Tentative identifications are marked in Table 1 and Figures 3, 4, 5, & 6 with “?”.
In addition to the listed lines, especially the upper sideband data (Fig. 5) show many features at low intensities, of which a majority is likely not noise but consists of weak additional spectral lines. At the given noise level, we do not investigate these features any further, but it is likely that finding real line-free spectral regions to study the pure continuum emission in such hot core type sources is difficult. The produced continuum emission will mostly be “polluted” by underlying weak line emission. As outlined below and in Beuther et al. (2004), imaging of the lines helps significantly, e.g., source I is nearly line-free (except for SiO) and thus the continuum emission is reliable, whereas this is less likely the case for the hot core region. To get an estimate of the line contamination, we produced a second pseudo-continuum image averaging the whole upper sideband data (in contrast to the continuum image presented in Beuther et al. (2004) and Fig. 1, where we excluded all strong lines from the upper sideband). We did not use the lower sideband data because the line contamination there is significantly higher due to the strong CH<sub>3</sub>OH bands. Deriving the source peak flux values from this new image for each continuum sub-source, the values are only 5-10% higher than for the original image. While the original image still covers a considerable bandpass of the upper sideband with low-level emission, it is unlikely that the low-level emission exceeds the contribution of the strong lines in the bandpass. Therefore, we estimate the continuum flux uncertainties due to line contamination in the image presented in Beuther et al. (2004) and Fig. 1 to be $`<10\%`$, independent of the sub-sources. This is well within our calibration uncertainties of 15% (§2). We note that this line contamination estimate is only valid for this particular dataset and cannot be extrapolated to single-dish studies, because the spatial filtering for the line and continuum data with the interferometer varies and single-dish bolometers cover a significantly larger bandwidth (e.g., $`80`$ GHz for the MAMBO array, Kreysa et al. 1998)
Fifteen percent of the lines listed in Table 1 were not reported in the previous single-dish survey (Schilke et al., 1997a), not to mention the additional low-level emission. In contrast, we do not identify any clear interferometric non-detection of a previously detected single-dish line, which could occur for very broadly distributed molecular gas, implying that most molecular emission must have some compact components.
Many of the lines listed in Table 1 can be attributed to a few molecules, and we identify 13 different species, 6 isotopologues, and 5 vibrational excited states within the dataset (Table 2). Most of these have clearly separated lines without much blending and thus can be imaged. Choosing one of the strongest lines of each species, we produced velocity integrated images of each. Figure 6 presents a compilation of line images from representative species in the dataset. A first inspection of these line images highlights the following interesting characteristics:
* The main molecular species detected toward source I is SiO. It clearly delineates the known low-velocity north-east south-west outflow. Other species are either not detected toward source I (e.g., CH<sub>3</sub>OH) or exhibit only weak emission (e.g., SO<sub>2</sub>).
* Most emission from nitrogen-bearing species is strong toward the hot core and has a similar morphology to NH<sub>3</sub> (Wilson et al. 2000). This confirms previous results at lower frequencies by Wright et al. (1996) and Blake et al. (1996).
* Most oxygen-bearing molecules are weaker toward the hot core but stronger in the south-west toward the so-called compact ridge<sup>2</sup><sup>2</sup>2The “compact ridge” has no clearly defined peak positions but extends $`5^{\prime \prime }10^{\prime \prime }`$ in east-west direction with peak positions varying with the molecular line transitions (see Fig. 6). (see also Wright et al. 1996).
* Sulphur-bearing species show emission toward the hot core and the compact ridge<sup>2</sup> (see also Wright et al. 1996). Vibrationally excited lines are stronger toward the hot core and SMA1.
* Both, oxygen- and sulphur-bearing molecules show additional emission toward the north-west, spatially associated with IRc6. Millimeter continuum emission from this region has been previously reported by Blake et al. (1996).
* The imaging of the molecular lines allows a better identification of the spectral lines. For example, the previous identification of the HCOOH$`(15_{4,12}14_{4,11})`$ line at 338.144 GHz (Schilke et al., 1997a) is called into question because its spatial structure resembles rather that of an N-bearing than of an O-bearing species. Its spatial distribution appears also different to its 1 mm counterpart presented in Liu et al. (2002). Therefore, this line is more likely CH<sub>3</sub>CH<sub>2</sub>CN$`(37_{3,34}36_{3,33})`$. Some of the tentative identifications in Table 1 and Figures 3, 4, 5, & 6 are uncertain due to this as well.
Another way to study the spatial differences of the molecular line emission is via extracting spectra from the whole spectral data-cube at selected positions within the field of view. For this purpose, we imaged the large data-cube with a lower spectral resolution of 3 km s <sup>-1</sup> and extracted spectra toward the positions of source I, source n, the hot core, SMA1, three CH<sub>3</sub>CN peak positions, the southern CH<sub>3</sub>OH peak position, the prominent western gas peak and the most northern peak which is strong in HC<sub>3</sub>N (Fig. 7).
The most striking difference appears to be between source I and most of the other positions. Toward source I, only <sup>28</sup>SiO and its rarer isotopologue <sup>30</sup>SiO are strong, whereas toward the other positions many more lines are detected but SiO is weaker. The hot core and SMA1 as well as the CH<sub>3</sub>CN peak positions are all strong line emitters. Regarding CH<sub>3</sub>OH, the strongest and most prominent position is the southern peak associated with the compact ridge. It shows the strongest CH<sub>3</sub>OH line bands in the ground state as well as the vibrationally excited states. The north-western peak is still rather strong in the ground state CH<sub>3</sub>OH lines but the vibrationally excited bands are significantly weaker than toward the southern peak.
A comparison of these different spectra shows the necessity of high spatial resolution in such complex regions. At lower angular resolution, the complexity of spatial variations and different excitation conditions blends together and results in integrated spectra as seen in many single-dish studies and also in Fig. 3.
## 4 Analysis and Discussion
### 4.1 <sup>28</sup>SiO(8–7) and <sup>30</sup>SiO(8–7)
The main <sup>28</sup>SiO(8–7) line was one of our prime target lines within this survey. Previous studies of the lower excitation ground state <sup>28</sup>SiO lines showed that in Orion-KL SiO traces the collimated north-east south-west outflow as well as the less collimated structures in the north-west south-east direction (Chandler & de Pree, 1995; Wright et al., 1995; Blake et al., 1996). A peculiarity of the <sup>28</sup>SiO v=0 emission is that the (1–0) and the (2–1) lines show maser emission in a bow-tie-like structure close to source I (Chandler & de Pree, 1995; Wright et al., 1995). The orientation of the bow-tie is along the axis of the north-east south-west outflow, its two peak positions are approximately $`0.5^{\prime \prime }`$ offset from source I in both directions. The <sup>28</sup>SiO(5–4) observations of Blake et al. (1996) lacked the spatial resolution to disentangle the potential maser positions from the other SiO emission, but at their given spatial resolution they do not find extremely high brightness temperatures which would imply obvious <sup>28</sup>SiO(5–4) maser emission. We now resolve the potential maser positions better in the <sup>28</sup>SiO(8–7) line but we also do not find any indication for <sup>28</sup>SiO(8–7) maser emission toward these positions.
However, there is thermal <sup>28</sup>SiO(8–7) and <sup>30</sup>SiO(8–7) emission on larger spatial scales tracing the molecular outflows. Figures 8 & 9 show channel maps of the <sup>28</sup>SiO and <sup>30</sup>SiO(8–7) data. While the velocity range from $`5`$ to 20 km s<sup>-1</sup> is dominated by the collimated north-east south-west outflow, the higher velocity channels, especially of <sup>28</sup>SiO(8–7), also show more extended emission. We stress that the <sup>28</sup>SiO(8–7) emission suffers strongly from missing short spacings and thus side-lobe problems. Especially, the north-east south-western stripes at the edges of the field presented in the <sup>28</sup>SiO(8–7) channel map (Fig. 8, the middle row) are side-lobe emission and not real. Comparing the channel maps with a vector-averaged <sup>28</sup>SiO(8–7) spectrum (Fig. 10), one finds that the collimated and the extended emission features are clearly separated in the spectrum as well. The more extended features are represented by the secondary peaks at $`<5`$ and $`>20`$ km s<sup>-1</sup>, whereas the collimated north-east south-west structure is confined to the lower velocities in between.
To investigate the different distributions, we averaged the <sup>28</sup>SiO(8–7) emission over the red and blue secondary peaks, the red and blue collimated outflow, and the central velocities. The final images are presented in Figure 11. We used the main isotopologue for this purpose because, even considering the short spacings problems, it is more sensitive to the larger-scale emission. The collimated structures at low to intermediate velocities are similar for both isotopologues. Obviously, the central velocities are dominated by the collimated outflow (Fig. 11 left). Going to slightly higher velocities we find blue and red emission toward the south-west and the north-east of source I (Fig. 11 middle). Since we resolve the outflow spatially well in north-east south-west direction and still find significant overlap of blue and red emission on both sides of source I, this indicates that the outflow is likely close to the plane of the sky (For an outflow orientation which is close to the plane of the sky, expanding outflow lobes along the line of sight produce the observed red and blue signature.). Going to the higher velocity secondary peaks (Fig. 11 right), the emission is more extended and spatially consistent with previous <sup>28</sup>SiO(2–1) observations by Wright et al. (1995). The blue emission might be part of the larger-scale north-west south-east outflow but the signature of the red higher velocity <sup>28</sup>SiO emission is less clear. It shows emission toward the south-east as expected from the large-scale outflow, but it has additional emission toward the south-west as well. As already mentioned, the <sup>28</sup>SiO emission suffers from missing flux problems due to the missing short spacings (Figs. 8 & 11), and it is difficult to interpret the large-scale emission in more detail.
Source n may be the driver of one of the outflows in the region (e.g., Menten & Reid 1995). However, source n is a strong near-infrared source and thus not deeply embedded and not particularly young. In addition, source n is weak in the usually outflow-tracing SiO emission (Figs. 6 & 7), which might be indicative of no recent molecular outflow activity. In contrast to this, there is indicative evidence for a potential outflow from source n based on radio, H<sub>2</sub>O maser and infrared emission signatures (e.g., Menten & Reid 1995; Stolovy et al. 1998; Greenhill et al. 2004; Shuping et al. 2004). While the detection of SiO emission often traces molecular outflows (e.g., Schilke et al. 1997b; Gueth & Guilloteau 1999; Cesaroni et al. 1999), the SiO non-detection close to source n does not necessarily imply the opposite. For example, dissociative shocks with velocities $`>60`$ km s<sup>-1</sup> do not produce significant SiO emission (e.g., Flower et al. 1996), and SiO can also be found in distinct bullets further down from the outflow center. However, it remains interesting to note that the SiO emission from the sources I and n is so different. Since source n is less deeply embedded than source I, it may either have driven an outflow in the past and we might still observe the shocked remnants in H<sub>2</sub>O maser emission (Menten & Reid, 1995), or the outflow is potentially of a slightly different, maybe more evolved nature than the ones usually observed in SiO.
Figure 12 shows a comparison of the <sup>28</sup>SiO and <sup>30</sup>SiO spectrum integrated over the central $`4^{\prime \prime }\times 4^{\prime \prime }`$ region. Both spectra cover about the same velocity range, and the <sup>30</sup>SiO spectrum is only $`20\%`$ weaker than the main isotopologue <sup>28</sup>SiO. This is different from the line ratio observed for these lines with the CSO 10.4 m single-dish telescope. Schilke et al. (1997a) find in their survey a <sup>28</sup>SiO/<sup>30</sup>SiO ratio of about 8, consistent with a <sup>28</sup>SiO opacity of 2. It appears that the observed line ratio decreases at higher spatial resolution. However, missing short spacings are a severe problem in comparing fluxes from different species and isotopologues because the spatial filtering affects various species differently. Trying to quantify this effect, we convolved the <sup>28</sup>SiO and <sup>30</sup>SiO maps to the $`20^{\prime \prime }`$ resolution of the CSO observations and compared the resulting peak fluxes. Additionally, we did the same again but this time masking out all negative features. The first approach gives an extreme lower limit to the recovered line intensities because the negative flux caused by the missing large-scale emission is included (Table 4). With these two approaches, we recover between 0.3 and 13% of the <sup>28</sup>SiO line flux and between 15 and 49% of the <sup>30</sup>SiO line flux. An additional effect caused by the missing short spacings is that the outflow emission is embedded in a larger-scale “bowl” of negative emission which to first order lowers the observed peak fluxes as well. Trying to estimate the distortion of the measured fluxes by this “bowl”, we find that the main <sup>28</sup>SiO line intensity is reduced by approximately $`15\%`$ whereas the rarer <sup>30</sup>SiO is affected by the “bowl” only to about $`5\%`$. While it is difficult to give exact numbers how much emission is filtered out in each line, these estimates show that the spatial filtering affects the main isotopologue <sup>28</sup>SiO significantly stronger than its rarer species <sup>30</sup>SiO. As outlined in §2, the missing flux has to be distributed over large spatial scales, and compared to the approximate missing flux of 3 Jy per synthesized SMA beam for <sup>28</sup>SiO (§2), this value is considerably lower for <sup>30</sup>SiO, between 0.21 and 0.35 Jy per synthesized SMA beam (depending on the SMA <sup>30</sup>SiO flux measurement in Table 4). In addition to this, it is likely that the line opacities toward the central $`(4^{\prime \prime })^2`$ are higher than in the larger region traced by the CSO ($`20^{\prime \prime }`$ beam). This would also decrease the observed <sup>28</sup>SiO/<sup>30</sup>SiO ratio. Hence, the observed <sup>28</sup>SiO/<sup>30</sup>SiO ratio of approximately unity is explicable by a combination of spatial filtering and opacity effects.
### 4.2 CH<sub>3</sub>OH emission
The methanol molecule shows by far the most lines within our spectral coverage. We identify 49 transitions from the $`v_t=0,1,2`$ states of CH<sub>3</sub>OH as well as its isotopologue <sup>13</sup>CH<sub>3</sub>OH (46 and 3, respectively). As shown in Figure 6, the methanol spatial distribution is significantly different from N-bearing hot core molecules with a strong additional peak in the south-west toward the compact ridge. We also find that the CH<sub>3</sub>OH emission is stronger toward the SMA1 submm continuum peak than toward the strongest submm continuum peak, which usually determines the hot core peak position.
The rich CH<sub>3</sub>OH spectra not only provide spatial and kinematic information but also allow possible estimates of the excitation condition of the gas throughout the region. The observed CH<sub>3</sub>OH transitions have upper state energy levels spanning from $``$ 75K to $``$ 700 K, readily indicating the high-temperature nature of the gas in this region. Assuming the CH<sub>3</sub>OH line emission originates from the same parcels of gas along each line of sight in local thermodynamic equilibrium (LTE, thus having the same size, velocity, and a single temperature), rotational temperatures can be derived by the rotation or population diagram analysis (e.g., Goldsmith & Langer 1999). However, line blending among the CH<sub>3</sub>OH transitions as well as with other species complicates the line intensity estimates and subsequent analysis. To remedy the situation, we estimate the temperature by optimizing synthetic spectra to observed spectra, an approach similar to those employed by, for example, Nummelin et al. (1998) and Comito et al. (2005). First, the observed CH<sub>3</sub>OH lines toward a particular line-of-sight are considered independent of their intensities, associated only by their relative rest frequency differences. By using the minimum $`\chi ^2`$ algorithm between the synthetic and observed spectra, we fitted all CH<sub>3</sub>OH lines simultaneously with Gaussian line profiles to derive common center velocities and linewidths. With the fixed velocities and linewidths, the rotational temperatures are then derived via a second $`\chi ^2`$ minimization, matching the synthetic to the observed intensities. In this derivation, the line opacity is assumed to have Gaussian profiles, a beam filling factor is employed, and the background radiation has been ignored.
Figure 13 presents observed and synthetic (fitted) spectra toward the hot core position and a location in the vicinity of the compact ridge but offset from the torsionally excited lines (see also Fig. 14, the position is approximately the center of the HCOOH emission presented in Liu et al. 2002). While the CH<sub>3</sub>OH lines from both torsionally ground and excited states are apparent toward the hot core, only transitions from torsionally ground states are noticeable at the latter position. Given that the torsionally excited states require either high densities and temperatures for collisional excitation, or strong infrared radiation fields for radiative pumping, the above behavior readily demonstrates variations in excitation conditions at different locations.
Figure 14 presents the derived rotational temperature structure toward the Orion-KL region. To increase the S/N ratio, we have smoothed/binned the original $`0.1^{\prime \prime }`$ pixel images to a larger pixel size of $`0.9^{\prime \prime }`$, of the order of the beam size. Still, this map gives one of the most detailed spatial pictures of the temperature distribution in this region. The derived rotational temperatures are between 50 and 350 K, with a mean value of $``$ 150 K over the region. Formal statistical errors from the fitting mostly range between 10 and 50 K. The temperatures we obtained are consistent with previous studies based on (sub)mm molecular line emission (e.g., Blake et al. 1987; Wilner et al. 1994; Schilke et al. 1997a; Wilson et al. 2000; Comito et al. 2005) and mid-infrared observations (Gezari et al., 1998). The highest temperature locations are mostly found to be close to the hot core positions such as the main submm continuum peak. In some regions, such as toward the south-west of SMA1 near the compact ridge, lines from the torsionally ground and excited states have similar intensities and thus appear to be optically thick (Fig. 13). In such optically thick cases, the fitting becomes insensitive to the rotational temperature along the particular line of sight. Toward locations where more than half of the CH<sub>3</sub>OH lines show fitted opacities $`>1`$, we refrained from temperature estimates and blanked the pixels in Figure 14. However, in the optical thick case the apparent intensity (brightness temperature) is directly related to true brightness temperature, differing only by the filling factor. For example, using the averaged temperature of 150 K and typical observed peak intensity of 4 Jy beam<sup>-1</sup> toward such optically thick directions, we find filling factors of 0.3 or less in the smoothed data cube. For positions where the spectral fitting is successful, the filling factors are even smaller, 0.14 on average. Such small filling factors may appear somewhat surprising, given past studies. Comito et al. (2005) for example, adopted a source size of 10<sup>′′</sup> for CH<sub>3</sub>OH emission but also obtained a similar rotational temperature. One expects that, when imaged at a higher angular resolution, the CH<sub>3</sub>OH emission would be resolved and thus fill the smaller beam. The fact that we find rotational temperatures comparable to past studies but with filling factor less than unity even for the optically thick lines suggests that a significant amount of CH<sub>3</sub>OH emission is missing and the observed hot CH<sub>3</sub>OH gas is actually clumpy. Indeed, smoothing the data to the $`20^{\prime \prime }`$ CSO beam and comparing these spectra with the single-dish observations by Schilke et al. (1997a), as much as 90$`\%`$ of the integrated CH<sub>3</sub>OH flux is resolved out by the interferometer. Since the spatial filtering affects various CH<sub>3</sub>OH lines differently, this has to be taken into account as well. As discussed in §2 & §4.1, the missing flux has to be distributed over large-spatial scales, and the more interesting parameter is the missing flux per synthesized beam. Estimating the latter values, we find that the CH<sub>3</sub>OH missing flux per synthesized SMA beam is in the worst case at a 10-20% level, within the calibration uncertainty. Therefore, it should not affect the temperature determination significantly.
We note that there are a few caveats in our synthetic spectral fitting process which does not always result in satisfactory fitted spectra. As mentioned above, the fitting process fails toward optically thick regions. Moreover, at some positions the CH<sub>3</sub>OH lines appear to be double-peaked, or a velocity offset between the torsionally ground state and excited transitions is noticeable (Fig. 13). If various kinematic components are present our fitting process often fails to converge. Additionally, (minor) disagreement between the observed and fitted spectra is not limited to the above cases but also present in the two top spectra of Fig. 13. Since all the transitions were observed simultaneously, calibration errors due to poor pointing or amplitude scaling are likely to be canceled out. Such deviations therefore suggest that additional non-LTE effects are at work. Furthermore, radiative pumping of the torsionally excited as well as the ground state transitions also needs to be considered, although collisional excitation is likely to dominate in high density regions like the Orion KL (Menten et al., 1986). A full statistical equilibrium analysis such as suggested by Leurini et al. (2004), which is outside the scope of this paper, and the inclusion of torsionally excited states would be an interesting next step to investigate these effects.
### 4.3 Vibrational excited emission
SMA1: The three vibrationally/torsionally excited lines in the survey (CH<sub>3</sub>OH, SO<sub>2</sub>, HC<sub>3</sub>CN) all show strong emission toward the submm continuum source SMA1 (Fig. 6). Since these lines are usually excited by infrared radiation, this emission indicates a deeply embedded infrared source at the position of SMA1. Similar conclusions were drawn by de Vicente et al. (2002) from lower spatial resolution Plateau de Bure observations of vibrational excited HC<sub>3</sub>N(10–9). Just based on the submm continuum study, we could not judge well whether SMA1 is an independent protostellar source, or whether it is an extension of the hot core (Beuther et al., 2004). These vibrationally excited line studies support the independent protostellar object interpretation for SMA1. Because SMA1 is detected neither at infrared nor at cm wavelength, it is likely one of the youngest sources of the evolving cluster.
The compact ridge: The torsionally excited CH<sub>3</sub>OH is also observed toward the compact ridge. We find only weak vibrationally excited SO<sub>2</sub> and no HC<sub>3</sub>N emission in this region. It is possible that an additional embedded infrared source exists in the region of the compact ridge, but we do not detect any emission above the noise in the 865 $`\mu `$m continuum there (Beuther et al., 2004). In contrast to this, Blake et al. (1996) detected weak mm continuum emission in that region. However, there mm continuum map does not show a strong peak but is more broadly distributed. The compact ridge is believed to be the interface region between one (or more?) of the molecular outflows with the dense ambient gas (e.g., Blake et al. 1987; Liu et al. 2002). Sine torsionally excited CH<sub>3</sub>OH can also be produced in regions of extremely high densities (critical densities of the order $`10^{10}`$ cm<sup>-3</sup>), it might also be possible that the CH<sub>3</sub>OH $`v_t=1,2`$ lines are excited in this interface region without an embedded infrared source.
## 5 Conclusion
The presented observations comprise a large set of molecular line data ($`145`$ lines from $`24`$ species, isotopologues and vibrational excited states) taken simultaneously during only two observing nights. The large number of observed CH<sub>3</sub>OH lines (46) with $`v_t`$=0,1,2 are a powerful tool to derive physical parameters (e.g., temperatures) of the molecular gas, and we find temperatures between 50 and 350 K throughout the region.
The SiO lines trace the collimated low-velocity molecular jet emanating from source I as well as larger-scale emission likely associated with a different outflow. In contrast to previously reported ground state <sup>28</sup>SiO(1–0) & (2–1) maser emission close to source I, we do not find detectable <sup>28</sup>SiO(8–7) maser emission there. Since source n is weak in the outflow-tracing SiO emission, this could indicate that it may currently drive no molecular outflow anymore.
Comparing different molecular species, for example, nitrogen-, oxygen- and sulphur-bearing molecules, we find strong chemical gradients over the observed region. For example, source I exhibits mainly SiO emission whereas the hot core is especially strong in nitrogen-bearing molecules like CH<sub>3</sub>CN. The strongest emission features of most oxygen-bearing molecules are south of the hot core toward the so-called compact ridge. These differences have important implications for studies lower-spatial-resolution studies and/or of sources at larger distances. For example, analyzing data of a typical massive star-forming region at approximately 4 kpc distance and deriving spatially averaged temperatures from the CH<sub>3</sub>OH data does not necessarily imply that the same temperatures can be attributed to the (sub)mm dust continuum cores. Furthermore, the imaging of spectral lines can be used as an additional tool for line identifications. Since new interstellar molecular line identifications get more complicated for large and complex molecules (see, e.g., the controversy about glycine, Kuan et al. 2003; Snyder et al. 2005), imaging each molecular line and comparing their spatial distributions is an important consistency check for rigorous line identifications.
All vibrationally/torsionally excited lines show strong emission toward the submm continuum peak SMA1. Because this emission is usually excited by infrared emission, it supports the notion that the recently identified source SMA1 is likely of protostellar nature.
This dataset is particularly rich and has potential for follow-up investigations. To mention just a couple: Imaging and potentially identifying the low-level emission peaks, or better constraining the underlying physical processes, which are responsible for the observed chemical gradients, is of great interest.
We like to thank all people at the SMA for their large efforts over many years to get the instrument going! An additional Thank You to the referee Peter Schilke for detailed comments improving the paper. H.B. acknowledges financial support by the Emmy-Noether-Program of the Deutsche Forschungsgemeinschaft (DFG, grant BE2578/1).
All figure captions are directly below the figures. Just for Fig. 6, where the caption did not properly fit on the page, we set it here.
Caption to Fig. 6.— Continuum-subtracted images of molecular species: the thick contours present the molecular emission (labeled in each panel) from 10 to 90% (step 10%) of each species’ peak emission. Additional parameters for each line image are given in Table 3. The axis are offsets in the directions of Right Ascension and Declination. The dashed contours show the negative emission from 10 to 90% (step 20%) which is due to missing short spacings. Source I, the hot core (HC), SMA1, source n, IRc6 and the compact ridge are marked in each panel and labeled in the <sup>30</sup>SiO panel (2nd at top). Two additional CH<sub>3</sub>CN positions, where spectra are shown from in Fig. 7, are marked with two additional stars in the CH<sub>3</sub>CN panel. The top-two rows show the Si- and and S-bearing species, the middle-two rows present the OH-bearing molecules, and the bottom-two rows focus on the N-bearing species. Lines with tentative or uncertain identifications are marked with “?” (see Tables 1 & 3).
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# Iwahori-Hecke Algebras of 𝑆𝐿₂ over 2-dimensional Local Fields
## Introduction
Hecke algebras play important roles in the representation theory of $`p`$-adic groups. There are two important classes of Hecke algebras. One is the spherical Hecke algebra attached to a maximal compact open subgroup, and the other is the Iwahori-Hecke algebra attached to a Iwahori subgroup. The spherical Hecke algebra is isomorphic to the center of the corresponding Iwahori-Hecke algebra. In the theory of higher dimensional local fields , a $`p`$-adic field is a $`1`$-dimensional local field. So the theory of $`p`$-adic groups and their Hecke algebras is over $`1`$-dimensional local fields.
Recently, the representation theory of algebraic groups over $`2`$-dimensional local fields was initiated by the works of Kapranov, Kazhdan and Gaitsgory . In their development, Cherednik’s double affine Hecke algebras () appear as an analogue of Iwahori-Hecke algebras. The common feature of the works mentioned above is the use of rank one integral structure of a $`2`$-dimensional local field. But there is also a rank two integral structure in a $`2`$-dimensional local field, and it is important in the arithmetic theory to use the rank two integral structure ; we refer the reader to Fesenko’s article and to the references there for recent developments.
In their paper , Kim and Lee constructed an analogue of spherical Hecke algebras of $`SL_2`$ over $`2`$\- dimensional local fields using rank two integral structure. They also established a Satake isomorphism using Fesenko’s $`((X))`$-valued measure defined in (See also ). In particular, the algebra is proved to be commutative. A similar result is expected in the case of $`GL_n`$ and, eventually, in the case of reductive algebraic groups.
In this paper, we construct an analogue of Iwahori-Hecke algebras of $`SL_2`$ over $`2`$-dimensional local fields coming from rank two integral structure. Our basic approach will be similar to that of . More precisely, an element of the algebra is an infinite linear combination of characteristic functions of double cosets satisfying certain conditions. We define a convolution product of two characteristic functions, using coset decompositions of double cosets of an Iwahori subgroup and imposing a restriction on the support of the product. Since a double coset of the Iwahori subgroup is an uncountable union of cosets in general, it is necessary to prove that the convolution product is well-defined. This will be done in the first part of the paper.
In the second part, we study the structure of the Hecke algebra. It has a natural $``$-grading and contains the affine Hecke algebra of $`SL_2`$ as a subalgebra. We will find a big commutative subalgebra, and determine its structure completely. Surprisingly, it is different from the group algebra of double cocharacters. We will also calculate the center of the Hecke algebra, and it turns out to be the same as the center of the affine Hecke algebra of $`SL_2`$. The classical Iwahori-Hecke algebra has a well-known presentation due to Iwahori and Matsumoto . The relations can be understood as deformations of Coxeter relations of the affine Weyl group. The corresponding Weyl group of a reductive algebraic group over a $`2`$-dimensional local field is not a Coxeter group. In $`SL_2`$ case, A.N. Parshin obtained an explicit presentation of the (double affine) Weyl group . We will investigate Iwahori-Matsumoto type relations of the Hecke algebra in light of Parshin’s presentation.
Now that we have spherical Hecke algebras and Iwahori-Hecke algebras of $`SL_2`$ over $`2`$-dimensional local fields with respect to rank two integral structure, next natural steps would be considering representations of $`SL_2`$ over $`2`$-dimensional local fields, constructing the Hecke algebras for more general reductive algebraic groups, and understanding these algebras in connection with the works of Kapranov, Kazhdan and Gaitsgory mentioned earlier. It seems that another interesting approach to the representation theory over two-dimensional local fields could be obtained from the work of Hrushovski and Kazhdan , in which they developed a theory of motivic integration. Actually, in the appendix of the paper , given by Avni, an Iwahori-Hecke algebra of $`SL_2`$ is constructed using motivic integration. It would be nice if one can find any connection of it to the constructions of this paper.
There are three sections and an appendix in this paper. In the first section, we fix notations and recall the Bruhat decomposition. The next section is devoted to the construction of the Hecke algebra. We will show that the convolution product is well-defined. In the third section, we study the structure of the Hecke algebra, determining the center of the algebra and finding Iwahori-Matsumoto type relations. In the appendix, we present a complete set of formulas for convolution products of characteristic functions. These formulas will be essentially used for many calculations in this paper.
### Acknowledgments
I would like to thank Henry Kim and Ivan Fesenko for their encouragement and useful comments during preparation of this paper.
## 1. Bruhat Decomposition
In this section, we fix notations and collect some results on double cosets and coset decompositions we will use later. We assume that the reader is familiar with basic definitions in the theory of $`2`$-dimensional local fields, which can be found in .
Let $`F(=F_2)`$ be a two dimensional local field with the first residue field $`F_1`$ and the last residue field $`F_0(=𝔽_q)`$ of $`q`$ elements. We fix a discrete valuation $`v:F^\times ^2`$ of rank two. Recall that $`^2`$ is endowed with the lexicographic ordering from the right. Let $`t_1`$ and $`t_2`$ be local parameters with respect to the valuation $`v`$. We have the ring $`O`$ of integers of $`F`$ with respect to the rank-two valuation $`v`$. There is a natural projection $`p:OO/t_1O=F_0`$. Note that the ring $`O`$ is different from the ring $`O_{21}`$ of integers of $`F`$ with respect to the rank-one valuation $`v_{21}:F^\times `$.
Let $`G`$ be a connected split semisimple algebraic group defined over $``$. We fix a maximal torus $`T`$ and a Borel subgroup $`B`$ such that $`TBG`$, and we have $`W_0=N_G(T)/T`$, the Weyl group of $`G`$. We write $`G=G(F)`$ and consider $`I=\{xG(O):p(x)B(F_0)\}`$, the double Iwahori subgroup, and $`W=N_G(T)/T(O)`$, the double affine Weyl group, and obtain the following decomposition.
###### Proposition 1.1.
We have
$$G=\underset{wW}{}IwI,$$
and the resulting identification $`I\backslash G/IW`$ is independent of the choice of representatives of elements of $`W`$.
From now on, in the rest of this paper, we assume that $`G=SL_2`$ and $`B`$ is the subgroup of upper triangular matirices. The following lemma gives explicit formulas for the decomposition in Proposition 1.1. The proof is straightforward, so we omit it.
###### Lemma 1.2.
Assume that $`x=\left(\begin{array}{cc}a& b\\ c& d\end{array}\right)G`$.
1. If $`v(a)v(b)`$ and $`v(a)<v(c)`$, then $`xI\left(\begin{array}{cc}a& 0\\ 0& a^1\end{array}\right)I`$.
2. If $`v(b)<v(a)`$ and $`v(b)<v(d)`$, then $`xI\left(\begin{array}{cc}0& b\\ b^1& 0\end{array}\right)I`$.
3. If $`v(c)v(a)`$ and $`v(c)v(d)`$, then $`xI\left(\begin{array}{cc}0& c^1\\ c& 0\end{array}\right)I`$.
4. If $`v(d)v(b)`$ and $`v(d)<v(c)`$, then $`xI\left(\begin{array}{cc}d^1& 0\\ 0& d\end{array}\right)I`$.
We denote by $`C_{i,j}^{(1)}`$ and $`C_{i,j}^{(2)}`$, $`(i,j)^2`$, the double cosets
$$I\left(\begin{array}{cc}t_1^it_2^j& 0\\ 0& t_1^it_2^j\end{array}\right)I\text{and}I\left(\begin{array}{cc}0& t_1^it_2^j\\ t_1^it_2^j& 0\end{array}\right)I,\text{respectively}.$$
In the following lemma, we obtain complete sets of coset representatives in the decomposition of double cosets of the subgroup $`I`$ into right cosets.
###### Lemma 1.3.
We have $`C_{i,j}^{(a)}=_zIz`$ where the disjoint union is over $`z`$ in the following list.
1. If $`a=1`$ and $`(i,j)(0,0)`$, then
$$z=\left(\begin{array}{cc}t_1^it_2^j& 0\\ 0& t_1^it_2^j\end{array}\right),\left(\begin{array}{cc}t_1^it_2^j& 0\\ t_1^{i+k}t_2^{j+l}u& t_1^it_2^j\end{array}\right)\text{ for }(1,0)(k,l)<(2i+1,2j),$$
where $`uO^\times `$ are units belonging to a fixed set of representatives of $`O/t_1^{2ik+1}t_2^{2jl}O`$.
2. If $`a=1`$ and $`(i,j)<(0,0)`$, then
$$z=\left(\begin{array}{cc}t_1^it_2^j& 0\\ 0& t_1^it_2^j\end{array}\right),\left(\begin{array}{cc}t_1^it_2^j& t_1^{i+k}t_2^{j+l}u\\ 0& t_1^it_2^j\end{array}\right)\text{ for }(0,0)(k,l)<(2i,2j),$$
where $`uO^\times `$ are units belonging to a fixed set of representatives of $`O/t_1^{2ik}t_2^{2jl}O`$.
3. If $`a=2`$ and $`(i,j)(0,0)`$, then
$$z=\left(\begin{array}{cc}0& t_1^it_2^j\\ t_1^it_2^j& 0\end{array}\right),\left(\begin{array}{cc}0& t_1^it_2^j\\ t_1^it_2^j& t_1^{i+k}t_2^{j+l}u\end{array}\right)\text{ for }(0,0)(k,l)<(2i+1,2j),$$
where $`uO^\times `$ are units belonging to a fixed set of representatives of $`O/t_1^{2ik+1}t_2^{2jl}O`$.
4. If $`a=2`$ and $`(i,j)<(0,0)`$, then
$$z=\left(\begin{array}{cc}0& t_1^it_2^j\\ t_1^it_2^j& 0\end{array}\right),\left(\begin{array}{cc}t_1^{i+k}t_2^{j+l}u& t_1^it_2^j\\ t_1^it_2^j& 0\end{array}\right)\text{ for }(1,0)(k,l)<(2i,2j),$$
where $`uO^\times `$ are units belonging to a fixed set of representatives of $`O/t_1^{2ik}t_2^{2jl}O`$.
###### Proof.
Since the other cases are similar, we only prove the part (1). Consider
$$z=\left(\begin{array}{cc}t_1^it_2^j& 0\\ 0& t_1^it_2^j\end{array}\right)\left(\begin{array}{cc}a& b\\ c& d\end{array}\right),z^{}=\left(\begin{array}{cc}t_1^it_2^j& 0\\ 0& t_1^it_2^j\end{array}\right)\left(\begin{array}{cc}a^{}& b^{}\\ c^{}& d^{}\end{array}\right)$$
where $`\left(\begin{array}{cc}a& b\\ c& d\end{array}\right),\left(\begin{array}{cc}a^{}& b^{}\\ c^{}& d^{}\end{array}\right)I`$. We see that the condition $`Iz=Iz^{}`$ is equivalent to
$$c^{}dcd^{}t_1^{2i+1}t_2^{2j}O.$$
We write $`(c,d)(c^{},d^{})`$ if $`c^{}dcd^{}t_1^{2i+1}t_2^{2j}O`$. Note that if $`\left(\begin{array}{cc}a& b\\ c& d\end{array}\right)I`$ then $`ct_1O`$ and $`d`$ is a unit. Let $`C`$ be the set of pairs $`(c,d)O^2`$ such that $`ct_1O`$ and $`d`$ is a unit. Then $``$ is an equivalence relation on $`C`$. In order to determine different cosets, we need only to determine a set of representatives of the equivalence relation $``$, which turn out to be
$$(0,1)\text{and}(t_1^kt_2^lu,1)\text{ for }(1,0)(k,l)<(2i+1,2j),$$
where $`uO^\times `$ are units belonging to a fixed set of representatives of $`O/t_1^{2ik+1}t_2^{2jl}O`$. These yield the elements $`z`$ in the part (1). ∎
###### Remark 1.4.
The disjoint union $`C_{i,j}^{(a)}=_zIz`$ is an uncountable union unless $`j=0`$. The same is true for the double cosets of $`K=SL_2(O)`$; see .
## 2. Iwahori-Hecke algebras
In this section, we define the convolution product of two characteristic functions of double cosets of the subgroup $`I`$, and we show that the product is well-defined. Then we construct an analogue of the Iwahori-Hecke algebra of $`SL_2`$.
We fix a set of representatives $`R`$ of the double affine Weyl group $`W`$ to be
(2.1)
$$R=\{\eta _{i,j}^{(1)}:=\left(\begin{array}{cc}t_1^it_2^j& 0\\ 0& t_1^it_2^j\end{array}\right),\eta _{i,j}^{(2)}:=\left(\begin{array}{cc}0& t_1^it_2^j\\ t_1^it_2^j& 0\end{array}\right):(i,j)^2\}.$$
We define a map $`\eta :GR`$ so that $`xI\eta (x)I`$ for each $`xG`$. That is, we assign to an element $`x`$ of $`G`$ its representative $`\eta (x)R`$ in the decomposition $`G=_{wW}IwI`$ of Proposition 1.1.
We put
(2.2)
$$C_j=\underset{m}{}C_{m,j}^{(1)}C_{m,j}^{(2)},j.$$
We denote by $`\chi _{i,j}^{(1)}`$ and $`\chi _{i,j}^{(2)}`$ the characteristic functions of the double cosets $`C_{i,j}^{(1)}`$ and $`C_{i,j}^{(2)}`$, respectively. We will consider the following types of functions
(2.3)
$$\underset{ri}{}c_r\chi _{r,j}^{(a)}(j>0),\underset{ri}{}c_r\chi _{r,j}^{(a)}(j<0),\text{and}\underset{iri^{}}{}c_r\chi _{r,0}^{(a)}$$
for $`i,i^{},j`$, $`a=1,2`$ and $`c_r`$. Now we define the convolution product of two characteristic functions.
###### Definition 2.4.
For $`a,b=1,2`$ and $`(i,j),(k,l)^2`$, we define
(2.5)
$$\left(\chi _{i,j}^{(a)}\chi _{k,l}^{(b)}\right)(x)=\{\begin{array}{cc}q^1_z\chi _{i,j}^{(a)}\left(\eta (x)z^1\right)\hfill & \text{if }xC_{j+l}\text{ and }jl0,\hfill \\ 0\hfill & \text{otherwise},\hfill \end{array}$$
where the sum is over the representatives $`z`$ of the decomposition $`C_{k,l}^{(b)}=_zIz`$.
###### Remark 2.6.
One can construct a certain invariant $`((X))`$-valued measure $`d\gamma `$ on $`G`$. Then we could define the convolution product of two functions $`f`$ and $`g`$ on $`G`$ by
$$(fg)(x)=_Gf(xy^1)g(y)𝑑\gamma (y).$$
The definition of the convolution product given above is derived from this formula.
Since the cardinality of the set of the representatives $`z`$ in the union is uncountable (Remark 1.4) in general, we need to prove the following.
###### Theorem 2.7.
The convolution product $`\chi _{i,j}^{(a)}\chi _{k,l}^{(b)}`$ yields a well-defined function of one of the types in (2.3) for any $`a,b=1,2`$ and $`(i,j),(k,l)^2`$.
###### Proof.
We write
$$\chi _{i,j}^{(a)}\chi _{k,l}^{(b)}=\underset{m}{}c_m^{(1)}\chi _{m,j+l}^{(1)}+\underset{m}{}c_m^{(2)}\chi _{m,j+l}^{(2)},c_m^{(1)},c_m^{(2)},m.$$
1. Assume that $`b=1`$, $`j0`$ and $`(k,l)(0,0)`$. By Lemma 1.3, we need only to consider $`z=\left(\begin{array}{cc}t_1^kt_2^l& 0\\ t_1^{k+k^{}}t_2^{l+l^{}}u& t_1^kt_2^l\end{array}\right)\text{ for }(1,0)(k^{},l^{})<(2k+1,2l)`$, where $`uO^\times `$ are units belonging to a fixed set of representatives of $`O/t_1^{2kk^{}+1}t_2^{2ll^{}}O`$.
1. If $`\eta (x)=\eta _{m,j+l}^{(1)}`$, then
$$\eta (x)z^1=\left(\begin{array}{cc}t_1^{mk}t_2^j& 0\\ t_1^{mk+k^{}}t_2^{j2l+l^{}}u& t_1^{m+k}t_2^j\end{array}\right).$$
1. If $`a=1`$, then we have $`mk=i`$, and either $`(mk+k^{},j2l+l^{})>(mk,j)`$ or $`(mk+k^{},j2l+l^{})>(m+k,j)`$ by Lemma 1.2. The first case gives $`(k^{},l^{})>(2m,2j+2l)=(2i+2k,2j+2l)`$, and so $`j=0`$, $`l^{}=2l`$ and $`2i+2k<k^{}<2k+1`$. Counting the number of possible representative units in $`O/t_1^{2kk^{}+1}O`$, we see that $`c_{i+k}^{(1)}`$ is finite and $`c_m^{(1)}=0`$ for all $`mi+k`$. The second case gives $`(k^{},l^{})>(2k,2l)`$, but $`(k^{},l^{})<(2k+1,2l)`$ from the assumption, and so $`c_m^{(1)}=0`$ for all $`m`$.
2. If $`a=2`$, then we must have $`l^{}=2l`$ and $`mk+k^{}=i`$. Then $`k^{}<2k+1`$ implies $`m<i+k+1`$, and counting the number of possible representative units in $`O/t_1^{2kk^{}+1}O`$, we see that $`c_m^{(2)}`$ is finite for each $`m`$.
2. If $`\eta (x)=\eta _{m,j+l}^{(2)}`$, then
$$\eta (x)z^1=\left(\begin{array}{cc}t_1^{mk+k^{}}t_2^{j+l^{}}u& t_1^{m+k}t_2^{j+2l}\\ t_1^{mk}t_2^{j2l}& 0\end{array}\right).$$
1. If $`a=1`$, then $`mk+k^{}=i`$ and $`l^{}=0`$, and it follows from Lemma 1.2 that $`jl`$. Since $`j0`$ and $`l0`$, we have $`j=l=l^{}=0`$. The condition $`1k^{}<2k+1`$ gives $`ik1<mi+k1`$ and $`c_m^{(2)}`$ is finite for each $`m`$.
2. If $`a=2`$, then $`l=l^{}=0`$, $`m+k=i`$ and $`1k^{}<2k+1`$. Thus $`c_{ik}^{(2)}`$ is finite and $`c_m^{(2)}=0`$ for $`mik`$.
2. Assume that $`b=1`$, $`j0`$ and $`(k,l)<(0,0)`$. By Lemma 1.3, we need only to consider $`z=\left(\begin{array}{cc}t_1^kt_2^l& t_1^{k+k^{}}t_2^{l+l^{}}u\\ 0& t_1^kt_2^l\end{array}\right)\text{ for }(0,0)(k^{},l^{})<(2k,2l)`$, where $`uO^\times `$ are units belonging to a fixed set of representatives of $`O/t_1^{2kk^{}}t_2^{2ll^{}}O`$.
1. If $`\eta (x)=\eta _{m,j+l}^{(1)}`$, then
$$\eta (x)z^1=\left(\begin{array}{cc}t_1^{mk}t_2^j& t_1^{m+k+k^{}}t_2^{j+2l+l^{}}u\\ 0& t_1^{m+k}t_2^j\end{array}\right).$$
1. If $`a=1`$, then we have $`mk=i`$, and either $`(mk,j)(m+k+k^{},j+2l+l^{})`$ or $`(m+k,j)(m+k+k^{},j+2l+l^{})`$ by Lemma 1.2. The first case yields $`(k^{},l^{})(2k,2l)`$, but $`(k^{},l^{})<(2k,2l)`$ from the assumption. Thus $`c_m^{(1)}=0`$ for all $`m`$. The second case gives $`(k^{},l^{})(2m,2j2l)=(2i2k,2j2l)`$, and we must have $`j=0`$, $`l^{}=2l`$ and $`2i2kk^{}<2k`$. Thus $`c_{i+k}^{(1)}`$ is finite and $`c_m^{(1)}=0`$ for all $`mi+k`$.
2. If $`a=2`$, then we have $`l^{}=2l`$ and $`m+k+k^{}=i`$. The condition $`k^{}<2k`$ leads to $`i+k<m`$, and $`c_m^{(2)}`$ is finite for each $`m`$.
2. If $`\eta (x)=\eta _{m,j+l}^{(2)}`$, then
$$\eta (x)z^1=\left(\begin{array}{cc}0& t_1^{m+k}t_2^{j+2l}\\ t_1^{mk}t_2^{j2l}& t_1^{m+k+k^{}}t_2^{j+l^{}}u\end{array}\right).$$
1. If $`a=1`$, then $`m+k+k^{}=i`$ and $`l^{}=0`$, and we have $`jl`$ by Lemma 1.2. Since $`j0`$ and $`l0`$, we have $`j=l=l^{}=0`$. The condition $`0k^{}<2k`$ yields $`i+km<ik`$ and $`c_m^{(2)}`$ is finite for each $`m`$.
2. If $`a=2`$, then $`l=l^{}=0`$, $`m+k=i`$ and $`0k^{}<2k`$. Thus $`c_{ik}^{(2)}`$ is finite and $`c_m^{(2)}=0`$ for $`mik`$.
3. Assume that $`b=2`$, $`j0`$ and $`(k,l)(0,0)`$. By Lemma 1.3, we need only to consider $`z=\left(\begin{array}{cc}0& t_1^kt_2^l\\ t_1^kt_2^l& t_1^{k+k^{}}t_2^{l+l^{}}u\end{array}\right)\text{ for }(0,0)(k^{},l^{})<(2k+1,2l)`$, where $`uO^\times `$ are units belonging to a fixed set of representatives of $`O/t_1^{2kk^{}+1}t_2^{2ll^{}}O`$.
1. If $`\eta (x)=\eta _{m,j+l}^{(1)}`$, then
$$\eta (x)z^1=\left(\begin{array}{cc}t_1^{mk+k^{}}t_2^{j+l^{}}u& t_1^{m+k}t_2^{j+2l}\\ t_1^{mk}t_2^{j2l}& 0\end{array}\right).$$
1. If $`a=1`$, then it is similar to (1)(b)(i).
2. If $`a=2`$, then it is similar to (1)(b)(ii).
2. If $`\eta (x)=\eta _{m,j+l}^{(2)}`$, then
$$\eta (x)z^1=\left(\begin{array}{cc}t_1^{mk}t_2^j& 0\\ t_1^{mk+k^{}}t_2^{j2l+l^{}}u& t_1^{m+k}t_2^j\end{array}\right).$$
1. If $`a=1`$, then it is similar to (1)(a)(i).
2. If $`a=2`$, then it is similar to (1)(a)(ii).
4. Assume that $`b=2`$, $`j0`$ and $`(k,l)<(0,0)`$. By Lemma 1.3, we need only to consider $`z=\left(\begin{array}{cc}t_1^{k+k^{}}t_2^{l+l^{}}u& t_1^kt_2^l\\ t_1^kt_2^l& 0\end{array}\right)\text{ for }(1,0)(k^{},l^{})<(2k,2l)`$, where $`uO^\times `$ are units belonging to a fixed set of representatives of $`O/t_1^{2kk^{}}t_2^{2ll^{}}O`$.
1. If $`\eta (x)=\eta _{m,j+l}^{(1)}`$, then
$$\eta (x)z^1=\left(\begin{array}{cc}0& t_1^{m+k}t_2^{j+2l}\\ t_1^{mk}t_2^{j2l}& t_1^{m+k+k^{}}t_2^{j+l^{}}u\end{array}\right).$$
1. If $`a=1`$, then it is similar to (2)(b)(i).
2. If $`a=2`$, then it is similar to (2)(b)(ii).
2. If $`\eta (x)=\eta _{m,j+l}^{(2)}`$, then
$$\eta (x)z^1=\left(\begin{array}{cc}t_1^{mk}t_2^j& t_1^{m+k+k^{}}t_2^{j+2l+l^{}}u\\ 0& t_1^{m+k}t_2^j\end{array}\right).$$
1. If $`a=1`$, then it is similar to (2)(a)(i).
2. If $`a=2`$, then it is similar to (2)(a)(ii).
A complete set of formulas for the convolution product $`\chi _{i,j}^{(a)}\chi _{k,l}^{(b)}`$ can be found in Appendix, and we obtain the following result.
###### Corollary 2.8.
(2.9)
$$\chi _{i,j}^{(a)}\chi _{k,l}^{(b)}=\{\begin{array}{cc}\underset{mi+k}{}c_m\chi _{m,j+l}^{(b)},c_m,\hfill & \text{ if }a=2,j>0\text{ and }l>0,\hfill \\ \underset{m>i+k}{}c_m\chi _{m,j+l}^{(b)},c_m,\hfill & \text{ if }a=2,j<0\text{ and }l<0,\hfill \\ \text{a finite sum}\hfill & \text{ otherwise.}\hfill \end{array}$$
We denote by $`(G,I)`$ the $``$-vector space generated by the functions of types in (2.3). We linearly extend the convolution product $``$ defined in (2.5) to the whole space $`(G,I)`$. It follows from Corollary 2.8 that it is well-defined. Thus we have obtained a $``$-algebra structure on the space $`(G,I)`$.
###### Definition 2.10.
The $``$-algebra $`(G,I)`$ will be called the Iwahori-Hecke algebra of $`G(=SL_2)`$.
It can be easily checked that $`\iota :=q\chi _{0,0}^{(1)}`$ is the identity element of the algebra $`(G,I)`$. Furthermore, we have:
###### Proposition 2.11.
The algebra $`(G,I)`$ is an associative algebra.
###### Proof.
Assume that $`\alpha =\chi _{i,j}^{(a)}`$, $`\beta =\chi _{k,l}^{(b)}`$ and $`\gamma =\chi _{m,n}^{(c)}`$. We fix the sets of representatives $`\{z_1\}`$, $`\{z_2\}`$ and $`\{z_3\}`$ in the decompositions $`C_{i,j}^{(a)}=_{z_1}Iz_1`$, $`C_{k,l}^{(b)}=_{z_2}Iz_2`$ and $`C_{m,n}^{(c)}=_{z_3}Iz_3`$, respectively. We write
$$\alpha \beta =\underset{\sigma }{}c(\alpha ,\beta ;\sigma )\sigma ,\text{ where }\sigma =\chi _{p,j+l}^{(d)},d=1,2,p.$$
Then
$$c(\alpha ,\beta ;\sigma )=\mathrm{Card}\{z_2:\eta _{p,j+1}^{(d)}z_2^1C_{i,j}^{(a)}\}=\mathrm{Card}\{(z_1,z_2):I\eta _{p,j+l}^{(d)}=Iz_1z_2\}.$$
The coefficient $`c(\alpha ,\beta ;\sigma )`$ is finite for any $`\sigma `$ by Theorem 2.7. Similarly, we write
$$\sigma \gamma =\underset{\tau }{}c(\sigma ,\gamma ;\tau )\tau ,\text{ where }\tau =\chi _{r,j+l+n}^{(e)},e=1,2,r.$$
We define
$$c(\alpha ,\beta ,\gamma ;\tau )=\mathrm{Card}\{(z_1,z_2,z_3):I\eta _{r,j+l+n}^{(e)}=Iz_1z_2z_3\}.$$
Since $`(\alpha \beta )\gamma `$ is defined, the number $`_\sigma c(\alpha ,\beta ;\sigma )c(\sigma ,\gamma ;\tau )`$ is finite. Now it is not difficult to see that
$$\underset{\sigma }{}c(\alpha ,\beta ;\sigma )c(\sigma ,\gamma ;\tau )=c(\alpha ,\beta ,\gamma ;\tau ).$$
Similarly, one can show that
$$\underset{\sigma ^{}}{}c(\alpha ,\sigma ^{};\tau )c(\beta ,\gamma ;\sigma ^{})=c(\alpha ,\beta ,\gamma ;\tau ),$$
where $`\sigma ^{}`$ is defined with regard to $`\alpha (\beta \gamma )`$. It proves the assertion of the proposition. ∎
###### Remark 2.12.
The argument is essentially the same as in the case of Hecke operators on the space of modular forms; see .
## 3. The structure of $`(G,I)`$
In this section, we investigate the structure of the Hecke algebra $`(G,I)`$. We will see that it has a big commutative subalgebra. After that, the center will be determined and Iwahori-Matsumoto type relations will be found.
For each $`j`$, we let $`_j`$ be the subspace of $`(G,I)`$, consisting of the functions with their supports contained in $`C_j`$, where the set $`C_j`$ is defined in (2.2). We put
$$_{}=\underset{j<0}{}_j\text{ and }_+=\underset{j>0}{}_j.$$
Then, clearly, we have
$$(G,I)=_{}_0_+.$$
It is easy to see that $`_0`$ is isomorphic to the (usual) affine Hecke algebra of $`SL_2`$ and that each $`_j`$ $`(j)`$ is a right and left $`_0`$-module.
We define
$$\begin{array}{c}\mathrm{\Theta }_{1,0}=\chi _{1,0}^{(1)},\hfill \\ \mathrm{\Theta }_{1,0}=\chi _{1,0}^{(1)}(q1)\chi _{1,0}^{(2)}(q1)\chi _{0,0}^{(2)}+q(q+q^12)\chi _{0,0}^{(1)},\hfill \\ \mathrm{\Theta }_{0,1}=\chi _{0,1}^{(1)},\hfill \\ \mathrm{\Theta }_{0,1}=\chi _{0,1}^{(1)}(q1)\underset{i0}{}q^i\chi _{i,1}^{(2)}.\hfill \end{array}$$
One can check $`\mathrm{\Theta }_{1,0}^1=\mathrm{\Theta }_{1,0}`$. (Recall that $`\iota =q\chi _{0,0}^{(1)}`$ is the identity.) The elements $`\mathrm{\Theta }_{1,0}`$ and $`\mathrm{\Theta }_{1,0}`$ are the same as appear in Bernstein’s presentation of the affine Hecke algebra $`_0`$.
###### Lemma 3.1.
The elements $`\mathrm{\Theta }_{1,0}`$, $`\mathrm{\Theta }_{1,0}`$, $`\mathrm{\Theta }_{0,1}`$and $`\mathrm{\Theta }_{0,1}`$ commute with each other.
Moreover, we have:
1. $`\mathrm{\Theta }_{1,0}^i\mathrm{\Theta }_{0,1}^j=q^{(i+j1)}\chi _{i,j}^{(1)}`$ for $`i`$ and $`j0`$,
2. $`\mathrm{\Theta }_{1,0}^i\mathrm{\Theta }_{0,1}^j=q^{(i+j1)}\chi _{i,j}^{(1)}(q1)q^{(i+j1)}\chi _{i,j}^{(2)}+{\displaystyle \underset{\genfrac{}{}{0pt}{}{m>i}{a=1,2}}{}}c_m^{(a)}\chi _{m,j}^{(a)}`$
for $`i`$ and $`j>0`$ and for $`c_m^{(a)}`$.
###### Proof.
We use the formulas in Appendix and inductions on $`i`$ and $`j`$ to obtain (1) and (2). ∎
Let us denote by $`𝒜`$ the commutative subalgebra of $`(G,I)`$ generated by the elements $`\mathrm{\Theta }_{1,0}`$, $`\mathrm{\Theta }_{1,0}`$, $`\mathrm{\Theta }_{0,1}`$and $`\mathrm{\Theta }_{0,1}`$. The structure of $`𝒜`$ is described by the following proposition.
###### Proposition 3.2.
The algebra $`𝒜`$ is isomorphic to the quotient of the algebra $`[X,X^1,Y,Z]`$ by the relation $`YZ=0`$.
###### Proof.
We have the surjective homomorphism $`\varphi :[X,X^1,Y,Z]/(YZ)𝒜`$ defined by
$$\varphi (X)=\mathrm{\Theta }_{1,0},\varphi (X^1)=\mathrm{\Theta }_{1,0},\varphi (Y)=\mathrm{\Theta }_{0,1}\text{and}\varphi (Z)=\mathrm{\Theta }_{0,1}.$$
It follows from (1) and (2) of Lemma 3.1 that $`\varphi `$ is also injective. ∎
Our next task is to determine the center of $`(G,I)`$.
###### Theorem 3.3.
The center of $`(G,I)`$ is the same as the center of $`_0`$ generated by the element $`\mathrm{\Theta }_{1,0}+\mathrm{\Theta }_{1,0}`$.
###### Proof.
It is well known that the element $`\mathrm{\Theta }_{1,0}+\mathrm{\Theta }_{1,0}`$ generates the center of $`_0`$ (). Let us check if it commutes with $`\chi _{i,j}^{(a)}`$, $`j0`$. First, we assume $`j>0`$. Since $`\chi _{i,j}^{(1)}𝒜`$, it commutes with $`\mathrm{\Theta }_{1,0}+\mathrm{\Theta }_{1,0}`$. We have $`\chi _{i,j}^{(2)}=q\chi _{i,j}^{(1)}\chi _{0,0}^{(2)}`$. Since $`\chi _{0,0}^{(2)}_0`$, we get $`\chi _{i,j}^{(2)}\left(\mathrm{\Theta }_{1,0}+\mathrm{\Theta }_{1,0}\right)=\left(\mathrm{\Theta }_{1,0}+\mathrm{\Theta }_{1,0}\right)\chi _{i,j}^{(2)}`$. Next, we assume $`j<0`$. We obtain, using the formulas in Appendix,
$$\chi _{i,j}^{(1)}\left(\mathrm{\Theta }_{1,0}+\mathrm{\Theta }_{1,0}\right)=q\chi _{i+1,j}^{(1)}+q^1\chi _{i1,j}^{(1)}=\left(\mathrm{\Theta }_{1,0}+\mathrm{\Theta }_{1,0}\right)\chi _{i,j}^{(1)}.$$
Since $`\chi _{i,j}^{(2)}=\chi _{i,j}^{(1)}\chi _{0,0}^{(2)}(1q^1)\chi _{i,j}^{(1)}`$, the element $`\chi _{i,j}^{(2)}`$ also commutes with $`\mathrm{\Theta }_{1,0}+\mathrm{\Theta }_{1,0}`$. Therefore, the element $`\mathrm{\Theta }_{1,0}+\mathrm{\Theta }_{1,0}`$ is in the center of the algebra $`(G,I)`$.
Suppose that $`\zeta `$ is an element in the center of $`(G,I)`$. We write $`\zeta =_j\zeta _j`$, $`\zeta _j_j`$. Since $`\chi _{0,0}^{(2)}_j_j`$ and $`_j\chi _{0,0}^{(2)}_j`$ for each $`j`$, the equality $`\zeta \chi _{0,0}^{(2)}=\chi _{0,0}^{(2)}\zeta `$ yields $`\zeta _j\chi _{0,0}^{(2)}=\chi _{0,0}^{(2)}\zeta _j`$ for each $`j`$. First, we assume $`j>0`$. Suppose that we choose the largest $`i`$ so that
$$\zeta _j=c_1\chi _{i,j}^{(1)}+c_2\chi _{i,j}^{(2)}+\underset{\genfrac{}{}{0pt}{}{m<i}{a=1,2}}{}c_m^{(a)}\chi _{m,j}^{(a)},c_10\text{ or }c_20.$$
We get
$$\chi _{0,0}^{(2)}\zeta _j=c_1(1q^1)\chi _{i,j}^{(1)}+c_2(1q^1)\chi _{i,j}^{(2)}+\underset{\genfrac{}{}{0pt}{}{m<i}{a=1,2}}{}c_{}^{}{}_{m}{}^{(a)}\chi _{m,j}^{(a)}.$$
On the other hand,
$$\zeta _j\chi _{0,0}^{(2)}=c_1q^1\chi _{i,j}^{(2)}+c_2\chi _{i,j}^{(1)}+c_2(1q^1)\chi _{i,j}^{(2)}+\underset{\genfrac{}{}{0pt}{}{m<i}{a=1,2}}{}c_{}^{\prime \prime }{}_{m}{}^{(a)}\chi _{m,j}^{(a)}.$$
Thus we have $`c_1=c_2=0`$, a contradiction. It implies that $`\zeta _j=0`$.
A similar argument also works for the case $`j<0`$, and we have $`\zeta _j=0`$ in this case, too. Thus $`\zeta =\zeta _0_0`$. It completes the proof. ∎
The double affine Weyl group $`W`$ is not a Coxeter group, but it has a similar presentation as one can see in the following proposition due to A. N. Parshin. It is also related to Kyoji Saito’s elliptic Weyl groups .
###### Proposition 3.4.
The group $`W`$ has a presentation given by
(3.5)
$$W<s_0,s_1,s_2|s_0^2=s_1^2=s_2^2=e,(s_0s_1s_2)^2=e>.$$
We can easily determine elements of $`W`$ corresponding to the generators in the presentation. For example, we can take, using the same notation,
$$s_0=\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right),s_1=\left(\begin{array}{cc}0& t_1^1\\ t_1& 0\end{array}\right),\text{and}s_2=\left(\begin{array}{cc}0& t_2^1\\ t_2& 0\end{array}\right).$$
We define
(3.6)
$$\varphi _0=q^{\frac{1}{2}}\chi _{0,0}^{(2)}=q^{\frac{1}{2}}\chi _{Is_0I},\varphi _1=q^{\frac{1}{2}}\chi _{1,0}^{(2)}=q^{\frac{1}{2}}\chi _{Is_1I},\text{and}\varphi _2=q^{\frac{1}{2}}\chi _{0,1}^{(2)}=q^{\frac{1}{2}}\chi _{Is_2I}.$$
The elements $`\varphi _0`$ and $`\varphi _1`$ have the special property
$$\varphi _0_{}=0\text{and}\varphi _1_+=0,$$
which follows from the formulas (2)(f) in Appendix.
###### Proposition 3.7.
The following identities hold in $`(G,I)`$:
(3.8)
$$\begin{array}{cc}& \varphi _0\varphi _0=(q^{\frac{1}{2}}q^{\frac{1}{2}})\varphi _0+\iota ,\hfill \\ & \varphi _1\varphi _1=(q^{\frac{1}{2}}q^{\frac{1}{2}})\varphi _1+\iota ,\hfill \\ \text{and }\hfill & \varphi _0\varphi _1\varphi _2\varphi _0\varphi _1=\varphi _2.\hfill \end{array}$$
###### Proof.
We check all the relations using the formulas in Appendix. ∎
###### Remark 3.9.
We can consider the relations in (3.8) as Iwahori-Matsumoto type relations. The first two relations in (3.8) are the usual deformation. The last one in (3.8) reflects the structure of the group algebra of $`W`$. However, we have
$$\varphi _2\varphi _2=(q^{\frac{1}{2}}q^{\frac{1}{2}})\underset{m>0}{}q^{m\frac{1}{2}}\chi _{m,2}^{(2)},$$
which reveals a new feature of the Hecke algebra $`(G,I)`$.
## Appendix
1. 1. If $`(i,j)(0,0)`$ and $`(k,l)(0,0)`$, or if $`(i,j)<(0,0)`$ and $`(k,l)<(0,0)`$, then
$$\chi _{i,j}^{(1)}\chi _{k,l}^{(1)}=q^1\chi _{i+k,j+l}^{(1)}\text{and}\chi _{i,j}^{(1)}\chi _{k,l}^{(2)}=q^1\chi _{i+k,j+l}^{(2)}.$$
2. If $`i0`$, $`j=0`$ and $`l<0`$, or if $`i<0`$, $`j=0`$ and $`l>0`$, then
$$\chi _{i,0}^{(1)}\chi _{k,l}^{(1)}=q^{2|i|1}\chi _{i+k,l}^{(1)}\text{and}\chi _{i,0}^{(1)}\chi _{k,l}^{(2)}=q^{2|i|1}\chi _{i+k,l}^{(2)}.$$
3. If $`j>0`$, $`k<0`$ and $`l=0`$, then
$$\begin{array}{cc}& \chi _{i,j}^{(1)}\chi _{k,0}^{(1)}=q^{2k1}\chi _{i+k,j}^{(1)}+(1q^1)\underset{m=i+k}{\overset{ik1}{}}q^{ikm1}\chi _{m,j}^{(2)}\hfill \\ \text{and}\hfill & \chi _{i,j}^{(1)}\chi _{k,0}^{(2)}=(1q^1)\underset{m=i+k+1}{\overset{ik1}{}}q^{ikm1}\chi _{m,j}^{(1)}+q^{2k2}\chi _{i+k,j}^{(2)}.\hfill \end{array}$$
4. If $`j<0`$, $`k0`$ and $`l=0`$, then
$$\begin{array}{cc}& \chi _{i,j}^{(1)}\chi _{k,0}^{(1)}=q^{2k1}\chi _{i+k,j}^{(1)}+(1q^1)\underset{m=ik}{\overset{i+k1}{}}q^{i+k+m}\chi _{m,j}^{(2)}\hfill \\ \text{and}\hfill & \chi _{i,j}^{(1)}\chi _{k,0}^{(2)}=(1q^1)\underset{m=ik}{\overset{i+k}{}}q^{i+k+m}\chi _{m,j}^{(1)}+q^{2k}\chi _{i+k,j}^{(2)}.\hfill \end{array}$$
5. If $`i0`$, $`j=0`$, $`k<0`$ and $`l=0`$, then
$$\begin{array}{cc}& \chi _{i,0}^{(1)}\chi _{k,0}^{(1)}=q^{\mathrm{min}\{2i1,2k1\}}\chi _{i+k,0}^{(1)}+(1q^1)\underset{m=\mathrm{max}\{i+k,ik\}}{\overset{ik1}{}}q^{ikm1}\chi _{m,0}^{(2)}\hfill \\ \text{and}\hfill & \chi _{i,0}^{(1)}\chi _{k,0}^{(2)}=(1q^1)\underset{m=\mathrm{max}\{i+k+1,ik\}}{\overset{ik1}{}}q^{ikm1}\chi _{m,0}^{(1)}+q^{\mathrm{min}\{2i1,2k2\}}\chi _{i+k,0}^{(2)}.\hfill \end{array}$$
6. If $`i<0`$, $`j=0`$, $`k0`$ and $`l=0`$, then
$$\begin{array}{cc}& \chi _{i,0}^{(1)}\chi _{k,0}^{(1)}=q^{\mathrm{min}\{2i1,2k1\}}\chi _{i+k,0}^{(1)}+(1q^1)\underset{m=ik}{\overset{\mathrm{min}\{i+k1,ik1\}}{}}q^{i+k+m}\chi _{m,0}^{(2)}\hfill \\ \text{and}\hfill & \chi _{i,0}^{(1)}\chi _{k,0}^{(2)}=(1q^1)\underset{m=ik}{\overset{\mathrm{min}\{i+k,ik1\}}{}}q^{i+k+m}\chi _{m,0}^{(1)}+q^{\mathrm{min}\{2i1,2k\}}\chi _{i+k,0}^{(2)}.\hfill \end{array}$$
2. 1. If $`j>0`$ and $`l>0`$, then
$$\begin{array}{cc}& \chi _{i,j}^{(2)}\chi _{k,l}^{(1)}=(1q^1)\underset{mi+k}{}q^{i+km}\chi _{m,j+l}^{(1)}\hfill \\ \text{and}\hfill & \chi _{i,j}^{(2)}\chi _{k,l}^{(2)}=(1q^1)\underset{mi+k}{}q^{i+km}\chi _{m,j+l}^{(2)}.\hfill \end{array}$$
2. If $`j<0`$ and $`l<0`$, then
$$\begin{array}{cc}& \chi _{i,j}^{(2)}\chi _{k,l}^{(1)}=(1q^1)\underset{m>i+k}{}q^{ik+m1}\chi _{m,j+l}^{(1)}\hfill \\ \text{and}\hfill & \chi _{i,j}^{(2)}\chi _{k,l}^{(2)}=(1q^1)\underset{m>i+k}{}q^{ik+m1}\chi _{m,j+l}^{(2)}.\hfill \end{array}$$
3. If $`(i,j)(0,0)`$, $`k<0`$ and $`l=0`$, or if $`(i,j)<(0,0)`$, $`k0`$ and $`l=0`$, then
$$\chi _{i,j}^{(2)}\chi _{k,0}^{(1)}=q^1\chi _{ik,j}^{(2)}\text{and}\chi _{i,j}^{(2)}\chi _{k,0}^{(2)}=q^1\chi _{ik,j}^{(1)}.$$
4. IF $`j>0`$, $`k0`$ and $`l=0`$, then
$$\begin{array}{cc}& \chi _{i,j}^{(2)}\chi _{k,0}^{(1)}=(1q^1)\underset{m=ik+1}{\overset{i+k}{}}q^{i+km}\chi _{m,j}^{(1)}+q^{2k1}\chi _{ik,j}^{(2)}\hfill \\ \text{and}\hfill & \chi _{i,j}^{(2)}\chi _{k,0}^{(2)}=q^{2k}\chi _{ik,j}^{(1)}+(1q^1)\underset{m=ik}{\overset{i+k}{}}q^{i+km}\chi _{m,j}^{(2)}.\hfill \end{array}$$
5. If $`j<0`$, $`k<0`$ and $`l=0`$, then
$$\begin{array}{cc}& \chi _{i,j}^{(2)}\chi _{k,0}^{(1)}=(1q^1)\underset{m=i+k+1}{\overset{ik}{}}q^{ik+m1}\chi _{m,j}^{(1)}+q^{2k1}\chi _{ik,j}^{(2)}\hfill \\ \text{and}\hfill & \chi _{i,j}^{(2)}\chi _{k,0}^{(2)}=q^{2k2}\chi _{ik,j}^{(1)}+(1q^1)\underset{m=i+k+1}{\overset{ik1}{}}q^{ik+m1}\chi _{m,j}^{(2)}.\hfill \end{array}$$
6. If $`i0`$, $`j=0`$ and $`l<0`$, or if $`i<0`$, $`j=0`$ and $`l>0`$, then
$$\chi _{i,0}^{(2)}\chi _{k,l}^{(1)}=0\text{and}\chi _{i,0}^{(2)}\chi _{k,l}^{(2)}=0.$$
7. If $`i0`$, $`j=0`$ and $`l>0`$, then
$$\begin{array}{cc}& \chi _{i,0}^{(2)}\chi _{k,l}^{(1)}=(1q^1)\underset{m=i+k}{\overset{i+k}{}}q^{i+km}\chi _{m,l}^{(1)}\hfill \\ \text{and}\hfill & \chi _{i,0}^{(2)}\chi _{k,l}^{(2)}=(1q^1)\underset{m=i+k}{\overset{i+k}{}}q^{i+km}\chi _{m,l}^{(2)}.\hfill \end{array}$$
8. If $`i<0`$, $`j=0`$ and $`l<0`$, then
$$\begin{array}{cc}& \chi _{i,0}^{(2)}\chi _{k,l}^{(1)}=(1q^1)\underset{m=i+k+1}{\overset{i+k1}{}}q^{ik+m1}\chi _{m,l}^{(1)}\hfill \\ \text{and}\hfill & \chi _{i,0}^{(2)}\chi _{k,l}^{(2)}=(1q^1)\underset{m=i+k+1}{\overset{i+k1}{}}q^{ik+m1}\chi _{m,l}^{(2)}.\hfill \end{array}$$
9. If $`i0`$, $`j=0`$, $`k0`$ and $`l=0`$, then
$$\begin{array}{cc}& \chi _{i,0}^{(2)}\chi _{k,0}^{(1)}=(1q^1)\underset{m=\mathrm{max}\{ik+1,i+k\}}{\overset{i+k}{}}q^{i+km}\chi _{m,0}^{(1)}+q^{\mathrm{min}\{2i,2k1\}}\chi _{ik,0}^{(2)}\hfill \\ \text{and}\hfill & \chi _{i,0}^{(2)}\chi _{k,0}^{(2)}=q^{\mathrm{min}\{2i,2k\}}\chi _{ik,0}^{(1)}+(1q^1)\underset{m=\mathrm{max}\{ik,i+k\}}{\overset{i+k}{}}q^{i+km}\chi _{m,0}^{(2)}.\hfill \end{array}$$
10. If $`i<0`$, $`j=0`$, $`k<0`$ and $`l=0`$, then
$$\begin{array}{cc}& \chi _{i,0}^{(2)}\chi _{k,0}^{(1)}=(1q^1)\underset{m=i+k+1}{\overset{\mathrm{min}\{ik,i+k1\}}{}}q^{ik+m1}\chi _{m,0}^{(1)}+q^{\mathrm{min}\{2i2,2k1\}}\chi _{ik,0}^{(2)}\hfill \\ \text{and}\hfill & \chi _{i,0}^{(2)}\chi _{k,0}^{(2)}=q^{\mathrm{min}\{2i2,2k2\}}\chi _{ik,0}^{(1)}+(1q^1)\underset{m=i+k+1}{\overset{\mathrm{min}\{i+k1,ik1\}}{}}q^{ik+m1}\chi _{m,0}^{(2)}.\hfill \end{array}$$
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# On making predictions in a multiverse: conundrums, dangers, and coincidences
## 1 Introduction
The standard model of particle physics and the standard model of cosmology are both rife with numerical parameters that must have values fixed by hand to explain the observed world. The world would be a radically different place if some of these constants took a different value. In particular, it has been argued that if any one of six (or perhaps a few more) numbers did not have rather particular values, then life as we know it would not be possible : atoms would not exist, or no gravitationally bound structures would form in the universe, or some other calamity would occur that would appear to make the (alter)-universe a very dull and lifeless place. How, then, did we get so lucky as to be here?
This question is an interesting one because all of the possible answers to it that I have encountered or devised entail very interesting conclusions. An essentially exhaustive list of such answers is:
1. We just got very lucky: all of the numbers could have been very different, in which case the universe would have been barren – but they just happened by pure chance to take values in the tiny part of parameter space that would allow life. We owe our existence to one very, very, very lucky roll of the dice.<sup>1</sup><sup>1</sup>1If the parameters were explained – by some deeper theory – in terms of fewer or no parameters, this does not change much: it would explain the origin of the parameters, but their hospitality to life would still be dumb luck.
2. We weren’t particularly lucky: almost any set of parameters would have been fine, because life would find a way to arise in nearly any type of universe. This is quite interesting because it implies (at least theoretically) the existence of life forms radically different from our own, existing for example in universes with no atoms or with no bound structure, or overrun with black holes, etc.
3. The universe was specifically designed for life. The choice of constants only happened once, but their values were determined in some way by the need for us to arise. This might be divine agency, or some radical form of Wheeler’s “self-creating universe”, or super-advanced beings that travel back in time to set the constants at the beginning of the universe, etc. However the reader feels about this possibility, they must admit that it would be interesting if true.
4. We did not have to get lucky, because there are many universes with different sets of constants – i.e., the dice were rolled many, many times. We are necessarily in one of the universes that allows life, just as we necessarily reside on a planet that supports life, even when most others may not. This is interesting because it means that there are other very different universes coexisting with ours in a “multiverse”.
These four answers – luck, elan vital, design, and multiverse – will appeal at different levels to different readers. But I think it is hard to argue that the multiverse is necessarily less reasonable than the alternatives. Moreover, as is discussed at length elsewhere in this volume, there are quite independent reasons to believe, on the basis of inflation, quantum cosmology, and string/M theory, that there might quite naturally be many regions larger than our observable universe, governed by different sets of low-energy physics. I am not aware of any independent scientific argument for the other three possible explanations.
Whether they are contemplated as an answer to the “why are we lucky” question, or because they are forced upon us from other considerations, multiverses come at a high price. Even if we have in hand a physical theory and cosmological model that lead to a multiverse, how do we test it? If there are many sets of constants, which ones do we compare to those we observe? In the next section of this chapter I will outline what I think a sound prediction in a multiverse would look like. As will become clear, this requires many ingredients, and there are some quite serious difficulties in generating some of these ingredients, even with a full theory in hand. For this reason, many short-cuts have been devised to try to make predictions more easily. In the third section I will describe a number of these, and show the cost that this convenience entails. Finally, in Section 4 I will focus on the interesting question of whether the anthropic approach to cosmology might lead to any general conclusions about how the study of cosmology will look in coming years.
## 2 Making predictions in a multiverse
Imagine that we have a candidate physical theory and set of cosmological boundary conditions (hereafter denoted $`𝒯`$) that predicts an ensemble of physically realized systems, each of which is approximately homogeneous in some coordinates and can be characterized by a set of parameters (i.e. the constants appearing in the standard models of particle physics and cosmology; I assume here that the laws of physics themselves retain the same form). Let us denote each such system a “universe” and the ensemble a “multiverse”. Given that we can observe only one of these universes, what conclusions can we draw regarding the correctness of $`𝒯`$, and how?
One possibility would be if there were a parameter for which none of the universes in the ensemble had the value we observe. In this case $`𝒯`$ would be ruled out. (Note that any $`𝒯`$ in which at least one parameter has a range of values that it does not take in any universe is thus rigorously falsifiable, which is a nice thing for a theory to be). Or perhaps some parameter takes only one value in all universes, and this value matches the observed one. This would obviously be a significant accomplishment of the theory. Both possibilities are good as far as they go, and seem completely uncontroversial. But they do not go far enough. What if our observed parameter values appear in some but not all of the universes? Could we still rule out the theory if those values are incredibly rare, or gain confidence if they are extremely common?
I find it hard to see why not. If some theory predicts outcome A of some experiment with $`p=0.99999999`$ probability, and outcome B with probability $`1p`$, I think we would be reluctant to accept the theory if a single experiment were performed and showed outcome B, even if we did not get to repeat the experiment. In fact, it seems consistent with all normal scientific methodology to rule out the theory at $`99.999999\%`$ confidence – the problem is just that without repeating our measurements we will not be able to increase this confidence. This seems to be exactly analogous to the multiverse if we can compute, given our $`𝒯`$, the probability that we should observe a given value for some observable.
Can we compute this probability distribution in a multiverse? Perhaps. I will argue that to do so in a sensible way, we would need seven successive ingredients.
1. First, of course, we require a multiverse: an ensemble of regions, each of which would be considered a universe to observers inside it (i.e. its properties would be uniform for as far as those observers could see), but each of which may have different properties.
2. Next we need to isolate the set of parameters characterizing the different universes. This might be the set of 20-odd free parameters in the standard model of particle physics (see, e.g., Ref. and references therein), plus a dozen or so cosmological parameters . There might be additional parameters that become important in other universes, or differences (such as different forms of the physical laws) that cannot be characterized by differences in a finite set of parameters. But for simplicity let us assume that some set of $`N`$ numbers $`\alpha _i`$ (where $`i=1..N`$) fully specify each universe.
3. Given our parameters, we need some measure with which to calculate the multi-dimensional probability distribution $`P(\alpha _i)`$ for the parameters. We might, for example, “count each universe equally” to obtain the probability $`P_U(\alpha _i)`$, defined to be the chance that a randomly chosen universe from the ensemble would have the parameter values $`\alpha _i`$.<sup>2</sup><sup>2</sup>2Note that this is really shorthand for $`\frac{d𝒫}{d\alpha _1..d\alpha _N}d\alpha _1..d\alpha _N`$, the probability that $`\alpha _i`$ are all within the interval $`[\alpha _i,\alpha _i+d\alpha _i]`$, where $`𝒫(\alpha _i)`$ is a cumulative probability distribution. This can be a bit tricky, however, because it depends on how we delineate the universes: suppose that $`\alpha _1=a`$ universes happen to be $`10^{10}`$ times larger than $`\alpha _1=b`$ universes. What would then prevent us from “splitting” each $`\alpha _1=a`$ universe into 10, or 100, or $`10^{10}`$ universes, thus radically changing the relative probability of $`\alpha _1=a`$ vs. $`\alpha _1=b`$? These considerations might lead us to take a different measure such as volume, e.g. to define $`P_V(\alpha _i)`$, the chance that a randomly point in space would reside in a universe with parameter values $`\alpha _i`$. But in an expanding universe volume increases, so this would depend on the time at which we choose to evaluate the volume in each universe. We might then consider some “counting” object that endures, say a baryon (which is relatively stable), and define $`P_B(\alpha _i)`$, the chance that a randomly chosen baryon would reside in a universe with parameter values $`\alpha _i`$. But now we have excluded from consideration universes with no baryons. Do we want to do that? This will be addressed in step (v). For now, note only that it is not entirely clear, even in principle, which measure we should place over our multiverse. We can call this the “measure problem.”
4. Once we choose a measure object $`M`$, we still need to actually compute $`P_M(\alpha _i)`$, and this may be far from easy. For example, in computing $`P_V`$, some universes may have infinite volume. In this case values of $`\alpha _i`$ leading to universes with finite volume will have zero probability. How, though, do we compare two infinite volumes? The difficulty can be seen by considering how we would count the fraction of red vs. blue marbles in an infinite box. We could pick one red, then one blue, and find a 50-50 split. But we could also repeatedly pick 1 red, then 2 blue, or 5 red, then 1 blue. We could do this forever and so obtain any ratio we like. What we would like to do is just “grab a bunch of marbles at random” and count the ratio. But in the multiverse case it is not so clear how to perform this random ordering of marbles to pick. This difficulty, which might be termed the “ordering problem” , has been discussed a number of times in the context of eternal inflation and a number of plausible prescriptions have been proposed. But there does not seem to be any generic solution, or convincing way to prove that one method is correct.
5. If we have managed to calculate $`P_M(\alpha _i)`$, do we have a prediction? Sortof. We have an answer to the question: “given that I am (or can associate myself with) a randomly chosen $`M`$-object, which sort of universe am I in?” But this is not necessarily the same as the more general question: “what sort of universe am I in?” First of all, different $`M`$-objects will generally give different probabilities, and they cannot all be the answer to the same question. Second, we may not be all that closely associated with our $`M`$-object (which was chosen mainly to provide some way to compute probabilities) because it does not take into account important requirements for our existence. For example, if $`M`$ were volume, I would be asking what I should observe given that I am at a random point in space; but we are not at a random point in space (which would on average have a density of $`10^{29}`$g/cc), but rather at one of the very rare points with density $`1`$ g/cc. The reason for this improbable situation is obviously “anthropic” – we just do not worry about it because we can observe many other regions at the proper density (if we could not see such regions, we might be more reluctant to accept a cosmological model with such a low average density.) Finally, it might be argued that the question we have answered through our calculation is not nearly as specific a question as we could ask, because we know a lot more about the universe than that it contains volume, or baryons. We might, instead, ask “given that I am in a universe with the properties we have already observed, what should I observe in the future?”
As discussed at length in , these different specific questions can be usefully thought of as arising from different choices of conditionalization. The probabilities $`P_M(\alpha _i)`$ are conditioned on as little as possible, whereas the anthropic question of “given that I am a randomly chosen observer, which should I measure” specifies probabilities conditioned on the existence of an “observer”, while the approach of “given what I know now, what will I see” specifies probabilities conditioned on being in a universe with all of the properties that we have already observed. These are three genuinely different approaches to making predictions in a multiverse that may be termed, respectively, “bottom-up”, “anthropic”, and “top down”.
Let us denote by $`O`$ the conditionalization object used to specify these conditional probabilities. In bottom-up reasoning, it would be the same as the $`M`$-object; in the anthropic approach it would be an “observer”, and in the top-down approach it could be a universe with the currently-known properties of our universe. It can be seen that they inhabit a spectrum, from the weakest conditionalization (bottom-up) to the most stringent (top-down). Like our initial $`M`$-object, choosing a conditionalization is unavoidable and important, and there is no obviously correct choice to make. (See Refs. for similar conditionalization schemas.)
6. Having decided on a conditionalization object $`O`$, the next step is to compute the number of $`N_{O,M}(\alpha _i)`$ of $`O`$-objects per $`M`$-object, for each set of values of the parameters $`\alpha _i`$. For example, if we have chosen to condition on observers, but have used baryons to define our probabilities, then we need to calculate the number of observers per baryon as a function of cosmological parameters. We can then calculate $`P_O(\alpha _i)P_M(\alpha _i)N_{O,M}(\alpha _i)`$, i.e. the probability that a randomly chosen $`O`$-object (observer) resides in a universe with parameters $`\alpha _i`$. There are a few possible pitfalls in doing this. First, if $`N_{O,M}`$ is infinite, then the procedure clearly breaks because $`P_O`$ then becomes undefined. This is why the $`M`$-object should be chosen to requires as little as possible for its existence (and hence be associated with the minimal-conditionalization bottom-up approach). This difficulty will generically occur if the existence of an $`O`$-object does not necessarily entail the existence of an $`M`$-object. For example, if the $`M`$-object were a baryon but the $`O`$-object were a bit of volume, then $`N`$ would be infinite for $`\alpha _i`$ corresponding to universes with no baryons. The problem arises because baryons require volume to be in, but volume does not require a baryon to be in it. This seems straightforward, but gets much murkier when we consider the second difficulty when calculating $`N`$, which is that we may not be able to precisely define what an $`O`$-object is, or what it takes to make one. If we say that the $`O`$-object is an observer, what exactly does that mean? A human? A carbon-based life form? Can observers exist without water? Without heavy elements? Without baryons? Without volume? It seems quite hard to say. We are forced, then, to choose some proxy for an observer, e.g. a galaxy, or a star with possible planets, etc. But our probabilities will perforce depend on the chosen proxy and this must be kept in mind.
It is worth noting a small bit of good news here. If we do manage to consistently compute $`N_{O,M_1}`$ for some measure object $`M_1`$, then insofar as we want to condition our probabilities on $`O`$-objects, we have solved the measure problem: if we could consistently calculate $`N_{O,M_2}`$ for a different measure object $`M_2`$, then we should obtain the same result for $`N_O`$, i.e. $`N_{O,M_1}P_{M_1}=N_{O,M_2}P_{M_2}`$. Thus our choice of $`M_1`$ (rather than $`M_2`$) becomes unimportant.
7. The final step in making predictions is to make the assumption that the probability that we will measure some set of $`\alpha _i`$ is given by the probability that a randomly chosen $`O`$-object will. This assumption really entails two others: first, that we are some how directly associated with $`O`$-objects, and second that we have not, simply by bad luck, observed highly improbable values of the parameters. The assumption that we are typical observers has been termed the “principle of mediocrity” . One may argue about this assumption, but some assumption is necessary if we are to connect our computed probabilities to observations, and it is difficult to see what alternative assumption would be more reasonable.
The result of all this work would be the probability $`P_O(\alpha _i)`$ that a randomly selected $`O`$-object (out of all of the $`O`$-objects that exist in multiverse) would reside in a universe governed by parameters $`\alpha _i`$, along with a reason to believe that this same probability distribution should govern what we will observe. We can then make the observations (or consider some already-made ones). If the observations are highly improbable according to our predictions, we can rule out the candidate $`𝒯`$ at some confidence that depends on how improbable our observations were. Apart from the manifest and grave difficulties involved in actually completing the seven listed steps in a convincing way, I think the only real criticism that can be leveled at this approach is that unless $`P=0`$ for our observed paramters, there will always be the chance that the $`𝒯`$ was correct and we measured an unlikely result. Usually, we can rid ourselves of this problem by repeating our experiments to make $`P`$ as small as we like (at least in principle), while here we do not have that option – once we have “used up” the measurement of all of the paramters required to describe our universe (which appears to be rather surprisingly few, at least according to current theories), we are done.
## 3 Making predictions in a multiverse more easily: a bestiary of shortcuts
Although the idea of a multiverse has been around for quite a while, no one has ever really come close to making the sort of calculation outlined in the previous section.<sup>3</sup><sup>3</sup>3The most ambitious attempt is probably the recent one by Tegmark . Instead, those wishing to make predictions in a multiverse context have made strong assumptions about which parameters $`\alpha _i`$ actually vary across the ensemble, about the choice of $`O`$-object, and about the quantities $`P_M(\alpha _i)`$ and $`N_{O,M}(\alpha _i)`$ that go into predicting their probabilities for measurement. Some of these shortcuts aim simply to make a calculation tractable; others are efforts to avoid anthropic considerations, or alternatively to use anthropic considerations to avoid other difficulties.
I would not have listed any ingredients that I thought could be omitted from a really sound calculation, thus all of these shortcuts are necessarily incomplete (some, in my opinion, disastrously so). But by listing and discussing them, I hope to give the reader both a flavor for what sort of anthropic (or anthropic-esque) arguments have been made in the literature, and where they may potentially go astray.
### 3.1 The “Maybe anthropic considerations only allow one set of parameters” hope
This assumption underlies a sort of anthropic reasoning that has earned the anthropic principle a lot of ill will. It goes something like: “Let’s assume that lots of universes governed by lots of different parameter values exist. Then since only universes with parameter values almost exactly the same as ours allow life, we must be in one of those, and we should not find it strange if our parameter values seem special.” In the conventions I have described, this is essentially equivalent setting $`O`$-objects to be observers, then hoping that the “hospitality factor” $`N_{\mathrm{obs},\mathrm{M}}(\alpha _i)`$ is very narrowly peaked around one particular set of parameters. In this case, the a priori probabilities $`P_M`$ are pretty much irrelevant because the shape of $`N_{\mathrm{obs},\mathrm{M}}`$ will pick out just one set of parameters. Because our observed values $`\alpha _i^{\mathrm{obs}}`$ definitely allow observers, the allowed set must then be very near $`\alpha _i^{\mathrm{obs}}`$.
Three problems with this type of reasoning are as follows. First, it is rather circular: it entails picking the $`O`$-object to be an observer, but then quickly substituting a “universe just like ours” for the $`O`$-object, with the reasoning that such universes will definitely support life.<sup>4</sup><sup>4</sup>4One can also argue that if there were other, more common, universes that supported life, we ought to be in them; since we are, not, we should assume that almost all life-supporting universes are like ours. But this is also circular in assuming that the whole anthropic argument works, in formulating the argument. Thus we have arrived at: the universe we observe should be pretty much like the observed universe. The way to avoid this silliness is to allow at least the possibility that there are life-supporting universes with $`\alpha _i\alpha _i^{\mathrm{obs}}`$, i.e. to discard the unproven assumption that $`N_{\mathrm{obs},\mathrm{M}}`$ has a single, dominant, narrow peak.
The second problem is that if $`N_{\mathrm{obs},\mathrm{M}}`$ were really so narrowly peaked as to render $`P_M`$ irrelevant, then we would be in serious trouble as theorists, because we would lose any ability to distinguish between candidates for our fundamental theory: unless our observed universe is impossible in the theory, then the anthropic factor would force the predictions of the theory to match our observations. As discussed in the next section, this is not good.
The third problem with $`N_{\mathrm{obs},\mathrm{M}}`$ being an extremely peaked function is that it does not appear to be true! As discussed in below, it appears that for any reasonable surrogate for observers (e.g. galaxies like ours, or stars with heavy elements, etc.), calculations done using our current understanding of galaxy and structure formation indicate that the region of parameter space in which there can be many of those objects may be small compared to the full parameter space, but it is much larger than the region compatible with our observations.
### 3.2 The “Just look for zero probability regions ($`P=0`$) in parameter space” approach
As mentioned in Section 2, there is a (relatively!) easy thing to do with a multiverse theory $`𝒯`$: work out which parameters combinations cannot occur in any universe. If the combination we actually observe is one of these, then the theory is ruled out. This is unobjectionable, but a rather weak way to test a theory because given two theories that are not ruled out, we have no way whatsoever of judging one to be better, even if the parameter values we observe are in some sense generic in one and absurdly rare in the other.<sup>5</sup><sup>5</sup>5Amusingly, in terms of testing $`𝒯`$, this approach which makes no assumptions about $`N_{O,M}`$ is equivalent to the approach just described of making the very strong assumption that $`N_{O,M}`$ allows only one specific set of parameter values, because in either case a theory can only be ruled out if our observed values are impossible in that theory.
This is not how science usually works. For example, suppose our theory is that a certain coin-tossing process is unbiased. If our only way to test this theory was to look for experimental outcomes that are impossible, then the theory would unfalsifiable: we would have to accept it theory for any coin we are confronted with, because no sequence of tosses would be impossible in it! Even if 10,000 tosses in a row all came up heads, we would have no grounds for doubting our theory because while getting heads 10,000 times in a row on a fair coin is absurdly improbably, it is not impossible. Nor would we have reason to prefer the (seemingly much better) “nearly every toss comes out heads” theory. Clearly this is a situation we would like to improve on, as much in universes as in coin tosses.
### 3.3 The “Let’s look for overwhelmingly more probable values” suggestion
One possible improvement would be to assume that we will observe a “typical” set of parameters in the ensemble, i.e. that we will employ “bottom-up” reasoning as described in the first section, by using the a priori (or “prior”) probabilities $`P_M`$ for some choice of measure-object such as universes, and just ignore the conditionalization factor $`N_{O,M}`$. There are two possible justifications for this. First, we might simply want to avoid any sort of anthropic issues on principle. Second, we might hope that some parameter values are much, much more common than others, to the extent that the $`N_{O,M}`$-factor becomes irrelevant – in other words that $`P_M`$ (rather than $`N_{O,M}`$) is a very strongly peaked around some particular parameters.
The problem with this approach is the “measure problem” discussed above: there is an implicit choice of basing probabilities on universes (say) rather than on (say) volume elements or baryons. Each of these measures has problems – for example, it seems that probabilities based on “universes” depends on how the universes are delineated, which can be ambiguous. Moreover there seems to be no reason to believe that predictions made using any two measures should agree particularly well. For example, as discussed elsewhere in this volume, in the string theory “landcape” there are many possible parameter sets, depending on which metastable minimum one chooses in a potential that depends in turn on a number of fluxes that can take a large range of discrete values. Imagine that exponentially many more minima lead to $`\alpha _1=a`$ than lead to $`\alpha _1=b`$. Should we expect to observe $`\alpha _1=a`$? Not necessarily, because the relative number of $`a`$-universes vs. $`b`$-universes that actually come into existence may easily differ exponentially from the relative number of $`a`$-minima vs. $`b`$-minima. (This seems likely to me in an eternal-inflation context, where the relative number universes could depend on exponentially-suppressed tunnelings between vacua.) Worse yet, these may in turn differ exponentially (or even by an infinite factor) from the relative numbers of baryons, or relative volumes.
In short, while we are free to use bottom-up reasoning with any choice of measure object we like, we are not free to assert that other choices would give similar predictions, or that conditionalization can be rendered irrelevant. So we had better have a pretty good reason for the choice we make.
### 3.4 The “Let’s fix some parameters to the observed values and predict others” shortcut
Another way in which one might hope to circumvent anthropic issues is to condition the probabilities on some or all observations that have already been made. In this “top-down” (or perhaps “pragmatic”) approach we ask: given everything that has been observed so far, what will we observe in some future measurement? It has a certain appeal, as this is often what is done in experimental science: we do not try to predict what our laboratory will look like, just what will happen given that the lab is in a particular state at a given time. In the conventions of Section 2, the approach could consist of choosing the $`O`$-object to be universes with parameters agreeing with the measured values.
While appealing, this approach suffers some deficiencies:
* It still does not completely avoid the measure problem, because even once we have limited our consideration to universes that match our current observations, we must still choose a measure with which to calculate the probabilities for the remaining ones.
* Through our conditioning, we may accept theories for which our parameter values are wildly improbable, without supplying any justification as to why we observe such improbable values. This is rather strange. Imagine that I have a theory in which the cosmological constant $`\mathrm{\Lambda }`$ is (with very high probability) much higher than we observe, and the dark matter particle mass $`m_{\mathrm{DM}}`$ is almost certainly $`>1000`$GeV. I condition on our observed $`\mathrm{\Lambda }`$, simply accepting that I am in an unusual universe. Now say I measure $`m_{\mathrm{DM}}=1`$GeV. I would like to say my theory is ruled out. Fine, but here is where it gets odd: according to top-down reasoning, I should also have already ruled it out if I had done my calculation in 1997, before $`\mathrm{\Lambda }`$ was measured. And someone who invented the very same theory next week – but had not been told that I have already ruled it out – would not rule it out, but instead just take the low value of $`m_{\mathrm{DM}}`$ (along with the observed $`\mathrm{\Lambda }`$) as part of the conditionalization!
* If we condition on everything we have observed, we obviously give up the possibility of explaining anything we have observed (which at this point is quite a lot in cosmology) through our theory.
The last two issues motivate variations on the top-down approach in which only some current observations are conditioned on. Two of which I am aware are:
1. We might start by conditioning on all observations, then progressively condition on less and less and try to “predict” the things we have decided not to condition on (as well, of course, as any new observations). The more we can predict, the better our theory is. The problem is that either (a) we will get to the point where we are conditioning on as little as possible (the bottom-up approach), and hence the whole conditionalization process will have been a waste of time, or (b) we will still have to condition on some things, and admit either that these have an anthropic explanation, or that we just choose to condition on them (leading to the funny issues discussed above).
2. We might choose at the outset to condition on things that we think may be fixed anthropically (without trying to actually generate this explanation), then try to predict the others . This is nice in being relatively easy, and in providing a justification for the conditionalization. It suffers from the problems of (a) guessing which parameters are anthropically important and which are not, (b) even if a parameter is anthropically unimportant, it may be strongly correlated in $`P_M`$ with one that is, and (c) we still have to face the measure problem, which we cannot avoid by counting conditioning on observers, because we are avoiding anthropic considerations.
### 3.5 The “Let’s assume just one parameter varies” simplification
Most of the “shortcuts” discussed so far have been attempts to avoid anthropic considerations. But we may, instead, consider how me might try to formulate an anthropic prediction (or explanation) for some observable, without going through the full calculation outlined in Section 2. The way of doing this that has been employed in the literature (largely in the efforts of Vilenkin and collaborators) is as follows.
First, one fixes all but one (or perhaps two) of the parameters to the observed values. This is done for tractability and/or because one hopes that they will have non-anthropic explanations. Let us call the parameter that is allowed to vary across the ensemble $`\alpha `$.
Next, an $`O`$-object is chosen such that given that only $`\alpha `$ varies, it is hoped that (a) the number $`N_{O,M}`$ of these objects (per baryon, or per comoving volume element) in a given universe is calculable, and (b) this number is arguably proportional to the number of observers. For example, if only $`\mathrm{\Lambda }`$ varies across the ensemble, galaxies might make reasonable $`O`$-objects because a moderately different $`\mathrm{\Lambda }`$ will probably not change the number of observers per galaxy, but will change the number of galaxies in a way that can be computed using fairly well-understood theories of galaxy and structure formation (for examples see ).
Third, it is assumed that $`P_M(\alpha )`$ is either flat or a simple power-law, without any complicated structure. This can be done just for simplicity, but it is often argued to be natural . The flavor of this argument is as follows. If $`P_M`$ is to have interesting structure over the relatively small range in which observers are abundant, there must be a parameter of order the observed $`\alpha `$ in the expression for $`P_M`$. But precisely this absence is what motivated the anthropic approach. For example, if the expression for $`P_M(\mathrm{\Lambda })`$ contained the energy scale $`0.01`$eV corresponding to the observed $`\mathrm{\Lambda }`$, the origin of that energy scale would probably be more interesting than our anthropic argument, as it would provide the basis for a (non-anthropic) solution to the cosmological constant problem!
Under these (fairly strong) assumptions we can then actually calculate $`P_O(\alpha )`$ and see whether or not the observed value is reasonably probable given this predicted distribution. For example when $`\mathrm{\Lambda }`$ alone is varied, a randomly chosen galaxy is predicted to lie in a universe with $`\mathrm{\Lambda }`$ comparable to (but somewhat larger than) the value we see .<sup>6</sup><sup>6</sup>6This is for a “flat” probability distribution $`dP_M/d\lambda \lambda ^\alpha `$ with $`\alpha =0`$. For $`\alpha >0`$, higher values would be predicted, and with $`\lambda <0`$ lower values would be favored.
I actually think this sort of reasoning is pretty respectable, given the assumptions made. In particular, the anthropic argument in which only $`\mathrm{\Lambda }`$ varies is a relatively clean one. But there are a number of pitfalls when it is applied to parameters other than $`\mathrm{\Lambda }`$, or when one allows multiple parameters to vary simultaneously.
* Assuming that the abundance of observers is strictly proportional to that of galaxies only makes sense if the number of galaxies – and not their properties – changes as $`\alpha `$ varies. However, changing nearly any cosmological parameter will change the properties of typical galaxies. For example, increasing $`\mathrm{\Lambda }`$ will decrease galaxy numbers, but also make galaxies smaller on average, because a high $`\mathrm{\Lambda }`$ squelches structure formation at late times when massive galaxies form. Increasing the amplitude of primordial perturbations would similarly lead to smaller, denser – but more numerous – galaxies, as would increasing ratio of dark matter to baryons. In these cases, we must specify in more detail what properties an observer-supporting galaxy should have, and this is very difficult to do without falling into the circular-argument trap of assuming that only galaxies like ours support life. Finally, this sort of strategy seems unlikely to work if we try to change non-cosmological parameters, as this could lead to radically different physics and the necessity of thinking very hard about what sort of observers there might be.
* The predicted probability distribution clearly depends on $`P_M`$, and the assumption that $`P_M`$ is flat, or a simple power law, can break down. This can happen even for $`\mathrm{\Lambda }`$ but perhaps more naturally for other parameters such as the dark matter density for which particle physics models can already yield sensible values. Moreover, this breakdown is much more probable if (as discussed below and contrary to the assumption made above) the hospitality factor $`N_{O,M}(\alpha )`$ is significant over many orders of magnitude in $`\alpha `$.
* Calculations of the hospitality factor $`N_{O,M}(\alpha )`$ can go awry if $`\alpha `$ is changed more than a little. For example, a neutrino mass slightly larger than we observe would suppress galaxy formation by erasing small-scale structure. But neutrinos with a large ($`>100`$eV) mass would act as dark matter and lead to strong halo formation. Whether these galaxies would be hospitable is questionable (they would be very baryon-poor), but the point is that the physics becomes qualitatively different. As another example, a lower photon/baryon ratio $`n_\gamma /n_b`$ would lead to earlier-forming, denser galaxies. But a much smaller value would lead to qualitatively different structure formation, as well as the primordial generation of heavy elements . As discussed at length in ref. , these changes are very dangerous because over orders of magnitude in $`\alpha `$, $`P_M(\alpha )`$ will tend to change by many orders of magnitude. Thus even if these alter-universes only have a few observers in them, they may dominate $`P_O`$ and hence qualitatively change the predictions.
* Along the same lines, but perhaps even more pernicious, when multiple parameters are varied simultaneously, the effects of some variations can offset the effect of others so that universes quite different from ours can support many of our chosen $`O`$-objects. For example, increasing $`\mathrm{\Lambda }`$ cuts off galaxy formation at a given cosmic density, but raising the perturbation amplitude $`Q`$ causes galaxies to form earlier (thus nullifying the effect of $`\mathrm{\Lambda }`$). This can be seen in the calculations of , and is discussed explicitly in . Many such deneneracies exist, because rasing $`\mathrm{\Lambda }`$, $`n_\gamma /n_b`$, or the neutrino mass all decrease the efficiency of structure formation, while raising $`\mathrm{\Omega }_{\mathrm{DM}}/\mathrm{\Omega }_b`$ or $`Q`$ increase the efficiency. As an extreme case, it was shown in that if $`Q`$ and $`n_\gamma /n_b`$ are allowed to vary with $`\mathrm{\Lambda }`$, then universes with $`\mathrm{\Lambda }`$ of $`10^{17}`$ times our observed value could arguably support observers! Including more cosmological parameters, or non-cosmological parameters, can only make this problem worse.
These problems indicate that while anthropic arguments concerning $`\mathrm{\Lambda }`$ in the literature are relatively “clean”, it is unclear whether other parameters (taken individually) will work as nicely. More importantly, a number of issues arise when several parameters are allowed to vary at once, and there does not seem to be any reason to believe that success in explaining one parameter anthropically will persist when additional parameters are allowed to vary. In some cases, it may: for example, allowing neutrino masses to vary in addition to $`\mathrm{\Lambda }`$ does not appear to spoil the anthropic explanation of a small but nonzero cosmological constant . On the other hand, allowing $`Q`$ to vary does, unless $`P_M(Q)`$ is strongly peaked at small values of $`Q`$ . I suspect that allowing $`\mathrm{\Omega }_{\mathrm{DM}}/\mathrm{\Omega }_b`$ or $`n_\gamma /n_b`$ to vary along with $`\mathrm{\Lambda }`$ would have a similar effect.
### 3.6 So what should we do?
For those serious about making predictions in a multiverse, I would propose that rather than working to generate additional incomplete anthropic arguments by taking shortcuts, a much better job must be done in each of the individual ingredients. For one example, our understanding of galaxy formation is sufficiently strong that the multi-dimensional hospitality factor $`N_{O,M}(\alpha _i)`$ could probably be computed for $`\alpha _i`$ within a few orders of magnitude of the observed values, for $`O`$objects of galaxies with properties within some range. Second, despite some nice previous work, I think the problem of how to compute $`P_M`$ in eternal inflation is a pretty open one. Finally, the string/M theory landscape (which is generating a lot of interest in the present topic right now) cannot hope to say much of anything about $`P_M`$ until its place in cosmology is understood – in particular, we need both a better understanding of the statistical distribution of field values that result from evolution in a given potential, and also an understanding of how transitions between vacua with different flux values occur, and exactly what is transitioning.
## 4 Cosmic coincidences and living dangerously: are there general predictions of anthropic reasoning?
The preceding sections should have suggested to the reader that it will be a huge project to compute a sound prediction of cosmological and physical parameters from a multiverse theory in which they vary. It may be so hard that it will be a very long time before any such calculation is at all believable. It is worth asking then: is there any way nature might give us an indication as to whether the anthropic approach is a sensible one, i.e. does the anthropic approach make any sort of general predictions even without the full calculation of $`P_O`$? Interestingly, I think the answer might be yes: I am aware of two such general (though somewhat vague) predictions of the anthropic approach.
To understand the first, assume that only one parameter, $`\alpha `$, varies, and consider $`p(\mathrm{log}\alpha )=\alpha P_M(\alpha )`$, the probability distribution in $`\mathrm{log}\alpha `$, given by some theory $`𝒯`$. For $`\mathrm{log}\alpha `$ near the observed value $`\mathrm{log}\alpha ^{\mathrm{obs}}`$, $`p`$ can basically only be doing one of three things: it can rise with $`\mathrm{log}\alpha `$, fall with $`\mathrm{log}\alpha `$, or be approximately constant. In the first two cases, the theory $`𝒯`$ would predict that we should see a value of $`\alpha `$ that is, respectively, higher or lower than we actually do if no anthropic conditionalization $`N_{\mathrm{obs},\mathrm{M}}`$ is applied. Now suppose we somehow compute $`N_{\mathrm{obs},\mathrm{M}}(\alpha )`$ and find that it falls off quickly for values of $`\alpha `$ much smaller or larger than we observe, i.e. that only a range $`\alpha _{\mathrm{min}}<\alpha <\alpha _{\mathrm{max}}`$ is “anthropically acceptable”. Then we have an anthropic argument explaining $`\alpha ^{\mathrm{obs}}`$, because this falloff means that $`P_{\mathrm{obs}}`$ will only be significant near $`\alpha ^{\mathrm{obs}}`$. But now note that within the anthropically acceptable range, $`P_{\mathrm{obs}}`$ will be peaked near $`\alpha _{\mathrm{max}}`$ if $`p`$ is increasing with $`\alpha `$, or near $`\alpha _{\mathrm{min}}`$ if $`p`$ is decreasing with $`\alpha `$. That is, we should expect $`\alpha ^{\mathrm{obs}}`$ at one edge of the anthropically acceptable range. This idea has been called the “principle of living dangerously” . It asserts that for a parameter that is anthropically determined, we should expect that a calculation of $`N_{\mathrm{obs},\mathrm{M}}`$ would reveal that observers would be strongly suppressed either for $`\alpha `$ slightly larger or slightly smaller than $`\alpha _{\mathrm{obs}}`$, depending on whether $`p`$ is rising or falling.
Now, this is not a very specific prediction: exactly where we would expect $`\alpha ^{\mathrm{obs}}`$ to lie depends both on how steep $`p(\mathrm{log}\alpha )`$ is, and how sharp the cutoff in $`N_{\mathrm{obs},\mathrm{M}}`$ is for $`\alpha `$ outside of the anthropically acceptable range. And it would not apply to anthropically-determined parameters in all possible cases. (For example, if $`p`$ were flat near $`\alpha ^{\mathrm{obs}}`$, but also very high at $`\alpha \alpha ^{\mathrm{obs}}`$, anthropic effects would be required to explain why we do not observe the very high value; but any region within the anthropically acceptable range would be equally probable, so we would not expect to be, so to speak, living on the edge.) Despite these caveats, this is a prediction of sorts, because the naive expectation would probably be for our observation to place us somewhere in the interior of the region of parameter space that is hospitable to life, rather than at the edge.
A second sort of general prediction of anthropic reasoning is connected to what might be called “cosmic coincidences.” For example, many cosmologists have asked themselves (and each other) why the current density in vacuum energy, dark matter, baryons, and neutrinos are all within a couple of orders of magnitude of each other – making the universe a much more complicated place than it might be. Conventionally, it has been assumed that these coincidences are just that, and follow directly from fundamental physics that we do not understand. But if the anthropic approach to cosmology is really correct (that is, if it is the real answer to the question of why these densities take the particular values they do), then the explanation is quite different: the densities are bound together by the necessity of observers’ existence, because only certain combinations will do.
More explicitly, suppose several cosmological parameters are governed by completely unrelated physics, so that their individual prior probabilities $`P_M`$ simply multiply to yield the multidimensional probability distribution. For example, we might have $`P_M(\mathrm{\Lambda },\mathrm{\Omega }_{\mathrm{DM}}/\mathrm{\Omega }_b,Q)=P_M(\mathrm{\Lambda })P_M(\mathrm{\Omega }_{\mathrm{DM}}/\mathrm{\Omega }_b)P_M(Q)`$. But even if $`P`$ factors, the hospitality factor $`N_{\mathrm{obs},\mathrm{M}}`$ will almost certainly not: if galaxies are $`O`$-objects, the number of galaxies formed at a given $`\mathrm{\Lambda }`$ will depend on both other parameters, and only certain combinations will give a significant number of observers. Thus $`P_O=N_{\mathrm{obs},\mathrm{M}}P_M`$ will likewise have correlations between the different parameters that lead to only particular combinations (for example those with $`\mathrm{\Omega }_{\mathrm{DM}}/\mathrm{\Omega }_b110`$ for a given $`Q`$ and $`\mathrm{\Lambda }`$) having high probabilities. The cosmic coincidences would be explained in this way.
This anthropic explanation of coincidences, however, should not only apply to things that we have already observed. If it is correct, then it should apply also to future observations; that is, we should expect to uncover yet more bizarre coincidences between quantities that seem to follow from quite unrelated physics.
How might this actually happen? Consider dark matter. We know fairly precisely how much dark matter there is in the universe, and what its basic properties are. But we have no real idea what it actually is, and there are many, many possible candidates that have been proposed in the literature. In fact, we have no observational reason to believe that dark matter is one substance at all: in principle it could be equal parts axions, supersymmetric particles, and primordial black holes. The reason most cosmologists do not expect this is that it would be a strange coincidence if three substances involving quite independent physics all wound up with essentially the same density in our universe. But of course this would be just like the suprising-but-true coincidences that hold in already-observed cosmology.
In the anthropic approach, these comparable densities could be quite natural . To see why, imagine that there are two completely independent types of dark matter permeating the ensemble: in each universe, they have some particular densities $`\rho _1`$ and $`\rho _2`$ out of a wide range of possibilities, so that the densities in a randomly chosen universe (or around a randomly chosen baryon, etc.) will be given probabilistically by $`P_M(\rho _1)P_M(\rho _2)`$. Under these assumptions there is no reason to expect that we should observe $`\rho _1\rho _2`$ based just on these a priori probabilities. Now suppose, though, that $`N_{O,M}`$ picks out a particular narrow range of total dark matter density as anthropically acceptable. That is, $`N_{O,M}(\rho _1+\rho _2)`$ is narrowly peaked about some $`\rho _{\mathrm{anth}}`$. In this case, the peak of the probability distribution $`P_O(\rho _1,\rho _2)`$, which indicates what values a randomly chosen observer should see, will occur where $`P_M(\rho _1)P_M(\rho _2)`$ is maximized subject to the condition that $`\rho _1+\rho _2\rho _{\mathrm{anth}}`$. For simplicity let both prior probabilities be power laws: $`P_M(\rho _1)\rho _1^\alpha `$ and $`P_M(\rho _2)\rho _2^\beta `$. Now the coincidence: it is not hard to show that if $`\alpha 0`$ and $`\beta 0`$, then the maximum will occur when $`\rho _1/\rho _2=\alpha /\beta `$. That is, the two components are likely to have similar densities unless the power law indices of their probability distributions differ by orders of magnitude.<sup>7</sup><sup>7</sup>7Extremely high power law indices are uncomfortable in the anthropic approach because they would lead to $`P_O`$ being peaked where $`N_{O,M}`$ is declining, i.e. we should be living outside the anthropically comfortable range, not just dangerously but downright recklessly. Of course, there are many ways in which this coincidence could fail to occur (e.g. negative power-law indices, or correlated probabilities), but the point is that there is a quite natural set of circumstances in which the components are coincident, even though the fundamental physics is completely unrelated.
## 5 Conclusions
The preceding sections should have convinced the reader that there are good reasons for scientists to be very worried if we live in a multiverse: in order to test a multiverse theory in a sound manner, we must perform a fiendishly difficult calculation of $`P_O(\alpha _i)`$, the probability that an $`O`$-object will reside in a universe characterized by parameters $`\alpha _i`$. And because of the shortcomings of the shortcuts one may (and presently must) take in doing this, almost any particular multiverse prediction is going to be easy to criticize; only a quite good calculation is going to be at all convincing. Much worse, we face an unavoidable and important choice in what $`O`$ should be: a possible universe, or an existing universe, or a universe matching current observations, or a bit of volume, or a baryon, or a galaxy, or an “observer”, etc.
I find it disturbingly plausible that “observers” really are the correct conditionalization object, that their use as such is the correct answer to the measure problem, and that anthropic effects are the real explanation for the values of some parameters (just as for the local density that we observe). Many cosmologiest appear to believe that taking the necessity of observers into account is shoddy thinking, and is employed only because it is the easy way out of solving problems the “right” way. But the arguments of this paper suggest that the truth may well be exactly the opposite: the anthropic approach may be the right thing to do in principle, but nearly impossible in practice.
Nonetheless, we cannot do away with multiverses just by wishing them away: we may in fact live in one, whatever the inconvenience to cosmologists. The productive strategy then seems to be one of accepting multiverses as a possibility, and working toward understanding how to calculate the various ingredients necessary to make predictions in one. Whether really performing such a calculation will turn out to be possible, but it is certainly impossible if not attempted.
Even if we cannot calculate $`P_O`$ in the foreseeable future, however, cosmology in a multiverse may not be completely devoid of predictive power. For example, of anthropic effects are at work, they should leave certain clues. First, if we could determine the region of parameter space hospitable to observers, we should find that we are living in the outskirts of the livable region, rather than somewere in its midst. Second, if the anthropic effects are the explanation of the parameter values – and coincidences between them – that we see, then it ought to predict that new coincidences will be observed in future observations.
If in the next several decades dark matter is resolved into several equally important components, dark energy is found to be three independent substances, and several other “cosmic coincidences” are observed, even someone the most die-hard skeptics might accede that the anthropic approach may have validity – why else would be universe be so very baroque? On the other hand, if we are essentially finished in defining the basic cosmological constituents, and the defining parameters are in the midst of a relatively large region of parameter space that might arguably support observers, then I think the anthropic approach would lose almost all appeal it has; we would be forced to ask: why isn’t the universe much wierder?
## References
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# On the uniqueness of the moonshine vertex operator algebra
## 1. Introduction
The moonshine vertex operator algebra $`V^{\mathrm{}}`$ constructed by Frenkel-Lepowsky-Meurman \[FLM1\],\[FLM2\] not only proves a conjecture by McKay-Thompson but also plays a fundamental role in shaping the theory of vertex operator algebra. In the introduction of \[FLM2\], Frenkel-Lepowsky-Meurman conjectured that the $`V^{\mathrm{}}`$ can be characterized by the following three conditions:
(a) the VOA $`V^{\mathrm{}}`$ is the only irreducible ordinary module for itself;
(b) the central charge of $`V^{\mathrm{}}`$ is 24;
(c) $`V_1^{\mathrm{}}=0.`$
We call their conjecture the Frenkel-Lepowsky-Meurman conjecture. These conditions are natural analogues of conditions which characterize the binary Golay code and the Leech lattice.
Conditions (b) and (c) are clear from the construction. Condition (a) is proved in \[D\] by using the 48 commuting Virasoro elements of central charge $`\frac{1}{2}`$ discovered in \[DMZ\]. Furthermore, $`V^{\mathrm{}}`$ is rational \[DLM2\],\[DGH\]. Although the theory of vertex operator algebra has developed a lot since \[FLM2\], including some uniqueness results for certain VOAs \[LX\], \[DM2\], \[DM3\], there has been no real progress in proving their conjecture.
In this paper we prove two weak versions of the Frenkel-Lepowsky-Meurman conjecture:
###### Theorem 1.
Let $`V`$ be a $`C_2`$-cofinite vertex operator algebra satisfying (a)-(c). We also assume that $`V_2`$ is isomorphic to the Griess algebra. Then $`V`$ is isomorphic to $`V^{\mathrm{}}.`$
In the second main theorem, we replace condition (a) by the assumption that $`dimV_ndimV_n^{\mathrm{}}`$ for $`n3.`$
###### Theorem 2.
Let $`V`$ be a simple vertex operator algebra satisfying (b)-(c). We also assume that $`V_2`$ is isomorphic to the Griess algebra and $`dimV_ndimV_n^{\mathrm{}}`$ for $`n3.`$ Then $`V`$ is isomorphic to $`V^{\mathrm{}}.`$
We now discuss the theorems and background. The weight two subspace $`V_2^{\mathrm{}}`$ of $`V^{\mathrm{}}`$ with the product which takes the pair $`u,v`$ to $`u_1v`$, where $`u_1`$ is the component operator of the vertex operator $`Y(u,z)=_nu_nz^{n1}`$ \[FLM2\], is the Griess algebra \[G\], which is a commutative nonassociative algebra of dimension 196884. Moreover, $`V^{\mathrm{}}`$ is generated by $`V_2^{\mathrm{}}`$ and $`V^{\mathrm{}}`$ is an irreducible module for the affinization of the Griess algebra \[FLM2\]. So, in order to understand the moonshine vertex operator algebra, one must know the Griess algebra and its affinization very well. It seems that a complete proof of FLM’s uniqueness conjecture needs a better understanding of the Griess algebra. Unfortunately, there does not yet exist a characterization of the Griess algebra (independent of its connection to the monster simple group). Also, the affinization of the Griess algebra is not a Lie algebra and lacks a highest weight module theory. From this point of view, $`V^{\mathrm{}}`$ is a very difficult vertex operator algebra.
The study of the moonshine vertex operator algebras in terms of minimal series of the Virasoro algebras was initiated in \[DMZ\]. This is equivalent to the study the maximal associative subalgebra of the Griess algebra. In \[DMZ\], we find 48 mutually commutative Virasoro algebras with central charge $`\frac{1}{2}.`$ As a result, a tensor product $`T_{48}`$ of 48 vertex operator algebras, associated to the highest weight unitary representations of the Virasoro algebra with central charge $`\frac{1}{2}`$ is a subalgebra of $`V^{\mathrm{}}`$ and $`V^{\mathrm{}}`$ decomposes into a direct sum of finitely many irreducible modules for $`T_{48}`$ as $`T_{48}`$ is rational and the homogeneous summands for $`V`$ are finite dimensional. A lot of progress on the study of the moonshine vertex operator algebra has been made by using the subalgebra $`T_{48}`$ and vertex operator subalgebras associated to the other minimal unitary series for the Virasoro algebras \[DLMN\], \[DGH\], \[KLY\], \[M3\]. The discovery of the $`T_{48}`$ inside $`V^{\mathrm{}}`$ also inspired the study of code vertex operator algebras and framed vertex operator algebras \[M2\], \[DGH\].
A frame in $`V^{\mathrm{}}`$ is a set of 48 mutually orthogonal Virasoro elements with central charge $`\frac{1}{2}.`$ The subalgebra $`T_{48}`$ depends on a frame as studied in \[DGH\]. It is proved in \[DGH\] that for any choice of 48 commuting Virasoro algebras there are two codes $`C`$ and $`D`$ associated to the decomposition of $`V^{\mathrm{}}`$ into irreducible $`T_{48}`$-modules. Each irreducible $`T_{48}`$-module is a tensor product of 48 unitary highest weight modules $`L(\frac{1}{2},h)`$ for the Virasoro algebra with central charge $`\frac{1}{2}`$ where $`h`$ can take only three values $`0,\frac{1}{2},\frac{1}{16}.`$ The code $`C`$ tells us the irreducible $`T_{48}`$-modules occurring in $`V^{\mathrm{}}`$ which are a tensor product of $`L(\frac{1}{2},h)`$ for $`h=0`$ or $`\frac{1}{2}.`$ Similarly, the code $`D`$ indicates the appearance of irreducible $`T_{48}`$-modules whose tensor factors have at least one $`L(\frac{1}{2},\frac{1}{16}).`$ The fusion rules for the vertex operator algebra $`L(\frac{1}{2},0)`$ indicate that we should consider a frame so that $`C`$ has maximal possible dimension and $`D`$ minimal possible dimension. These respective dimensions are 41 and 7. The reason for using our particular frame is that, for the code VOA which arises, all irreducible modules are simple currents (see Theorem 6.10 in this paper). The uniqueness of $`V^{\mathrm{}}`$ then follows from known uniqueness results for certain smaller VOAs, those which are simple current extensions of code VOAs.
The main strategy in proving the theorem is to use this particular frame. Since we assume that the weight 2 subspace of the abstract vertex operator algebra in the theorem is isomorphic to the Griess algebra, we can use the theory of framed vertex operator algebra developed in \[DGH\] and \[M2\] to investigate the structure of such vertex operator algebras.
Although we assume that $`V_2V_2^{\mathrm{}}`$ (as algebras), we can not claim automatically that any VF in $`V^{\mathrm{}}`$ corresponds to a VF in $`V`$. The difficult point is to prove that a Virasoro vector in $`V_2^{\mathrm{}}`$ generates a subVOA which is simple, i.e. an irreducible highest weight module. This is where we make use of the other assumptions in our main theorems. The proof involves both character theory for the Virasoro algebra with central charge $`\frac{1}{2}`$ and an explicit expression for the $`J`$-function.
It seems that there is still a long way to go to settle the FLM conjecture. The main difficulty is that we do not have much theory of finite dimensional commutative nonassociative algebras which could be applicable to a 196884-dimensional degree 2 summand of a VOA satisfying our conditions (a,b,c) (see \[G1\]). In a sense, this paper reduces the uniqueness of the moonshine vertex operator algebra to the uniqueness of the Griess algebra.
## 2. Notations
Most of our notations are fairly standard in the VOA literature. For the reader, we note a few below.
codes $`C=C(F),D=D(F)`$: see Section 4;
codes $`𝒞,𝒟`$: see Section 6;
$`j(q),J(q)`$ : the elliptic modular function and the elliptic modular function with constant term set equal to 0, i.e., $`J(q)=j(q)744`$;
$`\omega _i`$ : the subVOA generated by $`\omega _i`$;
VF : Virasoro frame, see Section 4;
$`Vir(\omega _i)`$ : the Virasoro algebra spanned by the modes of the Virasoro element $`\omega _i`$ and the scalars;
$`V^I`$ : see Section 4;
$`V^0`$ or $`V^{\mathrm{}}`$ : the case of $`V^I`$ for $`I=0`$ or $`\mathrm{}`$;
$`V^{\mathrm{}}`$ : the moonshine VOA, constructed in \[FLM2\];
$`(V^{\mathrm{}})^0`$ : this is $`V^0`$ for $`V=V^{\mathrm{}}`$.
## 3. Various modules for vertex operator algebras
Let $`(V,Y,\mathrm{𝟏},\omega )`$ be a vertex operator algebra. We recall various notion of modules (cf. \[FLM2\], \[DLM1\]).
A weak $`V`$ module is a vector space $`M`$ with a linear map $`Y_M:VEnd(M)[[z,z^1]]`$ where $`vY_M(v,z)=_nv_nz^{n1}`$, $`v_nEnd(M)`$. In addition $`Y_M`$ satisfies the following:
1) $`v_nw=0`$ for $`n>>0`$ where $`vV`$ and $`wM`$;
2) $`Y_M(\text{1},z)=Id_M`$;
3) The Jacobi Identity
$`z_0^1\delta \left({\displaystyle \frac{z_1z_2}{z_0}}\right)Y_M(u,z_1)Y_M(v,z_2)z_0^1\delta \left({\displaystyle \frac{z_2z_1}{z_0}}\right)Y_M(v,z_2)Y_M(u,z_1)`$
$`=z_2^1\delta \left({\displaystyle \frac{z_1z_0}{z_2}}\right)Y_M(Y(u,z_0)v,z_2)`$
holds.
An admissible $`V`$ module is a weak $`V`$ module which carries a $`_+`$-grading, $`M=_{n_+}M(n)`$, such that $`v_mM(n)M(n+\mathrm{wt}vm1)`$
An ordinary $`V`$ module is a weak $`V`$ module which carries a $``$-grading, $`M=_\lambda M_\lambda `$, such that:
1) $`dim(M_\lambda )<\mathrm{}`$ for all $`\lambda `$;
2) $`M_{\lambda +n}=0`$ for fixed $`\lambda `$ and $`n<<0`$ (depending on $`\lambda `$);
3) $`L(0)w=\lambda w=\mathrm{wt}(w)w`$, for $`wM_\lambda `$.
It is easy to prove that an ordinary module is admissible.
A vertex operator algebra is called rational if every admissible module is a direct sum of simple admissible modules. That is, a VOA is rational if there is complete reducibility of the category of admissible modules. It is proved in \[DLM2\] that if $`V`$ is rational there are only finitely many irreducible admissible modules up to isomorphism and each irreducible admissible module is ordinary.
A vertex operator algebra $`V`$ is called holomorphic if it is rational and the only irreducible ordinary module is itself. In this case $`V`$ is also the only irreducible admissible module.
A vertex operator algebra is called regular if every weak module is a direct sum of simple ordinary modules. So, regularity implies rationality.
A vertex operator algebra $`V`$ is called $`C_2`$-cofinite if $`V/C_2(V)`$ is finite dimensional where $`C_2(V)=u_2v|u,vV.`$
## 4. Framed vertex operator algebras
In this section we review the framed vertex operator algebras and related results from \[DMZ\] and \[DGH\].
Let $`L(c,h)`$ be the irreducible highest weight module for the Virasoro algebra with central charge $`c`$ and highest weight $`h.`$ The $`L(\frac{1}{2},0)`$-module $`L(\frac{1}{2},h)`$ is unitary if and only if $`h=0,\frac{1}{2},\frac{1}{16}`$ \[FQS\], \[GKO\]. Moreover, $`L(\frac{1}{2},0)`$ is a rational vertex operator algebra and $`L(\frac{1}{2},h)`$ for $`h=0,\frac{1}{2},\frac{1}{16}`$ gives a complete list of inequivalent irreducible $`L(\frac{1}{2},h)`$-modules.
We first recall the notion of framed vertex operator algebra. Let $`r`$ be a nonnegative integer. A framed vertex operator algebra (FVOA) is a simple vertex operator algebra $`(V,Y,\mathrm{𝟏},\omega )`$ satisfying the following conditions: there exist $`\omega _iV`$ for $`i=1`$$`\mathrm{}`$$`r`$ such that (a) each $`\omega _i`$ generates a copy of the simple Virasoro vertex operator algebra $`L(\frac{1}{2},0)`$ of central charge $`\frac{1}{2}`$ and the component operators $`L^i(n)`$ of $`Y(\omega _i,z)=_nL^i(n)z^{n2}`$ satisfy $`[L^i(m),L^i(n)]=(mn)L^i(m+n)+\frac{m^3m}{24}\delta _{m,n};`$ (b) The $`r`$ Virasoro algebras $`Vir(\omega _i)`$, spanned by the modes of $`Y(\omega _i,z)`$ and the identity, are mutually commutative; and (c) $`\omega =\omega _1+\mathrm{}+\omega _r`$. The set $`\{\omega _1,\mathrm{},\omega _r\}`$ is called a Virasoro frame (VF).
¿From now on we assume that $`V`$ is a FVOA of central charge $`\frac{r}{2}`$ with frame $`F:=\{\omega _1,\mathrm{},\omega _r\}`$. Let $`T_r`$ be the vertex operator algebra generated by $`\omega _i`$ for $`i=1,\mathrm{},r.`$ Then $`T_r`$ is isomorphic to $`L(\frac{1}{2},0)^r`$ and its irreducible modules are the $`L(h_1,\mathrm{},h_r):=L(\frac{1}{2},h_1)\mathrm{}L(\frac{1}{2},h_r)`$ for $`h_i=0,\frac{1}{2},\frac{1}{16}.`$ Since $`T_r`$ is a rational vertex operator algebra, $`V`$ is a completely reducible $`T_r`$-module. That is,
$$V\underset{h_i\{0,\frac{1}{2},\frac{1}{16}\}}{}m_{h_1,\mathrm{},h_r}L(h_1,\mathrm{},h_r)$$
(4.1)
where the nonnegative integer $`m_{h_1,\mathrm{},h_r}`$ is the multiplicity of $`L(h_1,\mathrm{},h_r)`$ in $`V`$. In particular, all the multiplicities are finite and $`m_{h_1,\mathrm{},h_r}`$ is at most $`1`$ if all $`h_i`$ are different from $`\frac{1}{16}`$.
There are two binary codes $`C=C(F)`$ and $`D=D(F)`$ associated to the decomposition (4.1). In order to define the code $`D`$ we identify a subset $`I`$ of $`\{1,\mathrm{},r\}`$ with a codeword $`d=(d_1,\mathrm{},d_r)𝔽_2^r`$ where $`d_i=1`$ if $`iI`$ and $`d_0=0`$ elsewhere. Let $`I`$ be a subset of $`\{1,\mathrm{},r\}`$. Define $`V^I`$ as the sum of all irreducible submodules isomorphic to one of the irreducibles $`L(h_1,\mathrm{},h_r)`$ such that $`h_i=\frac{1}{16}`$ if and only if $`iI`$. Then
$$V=\underset{I\{1,\mathrm{},r\}}{}V^I.$$
Set
$$D=D(F):=\{I𝔽_2^rV^I0\}.$$
(4.2)
For $`c=(c_1,\mathrm{},c_r)𝔽_2^r`$, we define $`V(c)=m_{h_1,\mathrm{},h_r}L(h_1,\mathrm{},h_r)`$ where $`h_i=\frac{1}{2}`$ if $`c_i=1`$ and $`h_i=0`$ elsewhere. Set
$$C=C(F):=\{c𝔽_2^rV(c)0\}.$$
(4.3)
Then $`V^{\mathrm{}}=V^0=_{cC}V(c)`$.
Here we summarize the main result about FVOAs from \[DGH\]
###### Theorem 4.1.
Let $`V`$ be a FVOA. Then
(a) $`V=_{n0}V_n`$ with $`V_0=\mathrm{𝟏}.`$
(b) $`V`$ is rational.
(c) $`C`$ and $`D`$ are binary codes and
$$CD^{}=\{x=(x_1,\mathrm{},x_r)𝔽_2^r|xd=0dD\}.$$
Moreover, $`V`$ is holomorphic if and only if $`C=D^{}.`$
(d) $`V^0`$ is a simple vertex operator algebra and the $`V^I`$ are irreducible $`V^0`$-modules. Moreover $`V^I`$ and $`V^J`$ are inequivalent if $`IJ`$.
(e) For any $`I,JD`$ and $`0vV^J`$ we have $`V^{I+J}=span\{u_nv|uV^I,n\}.`$
(f) Let $`I\{1,\mathrm{},r\}`$ be given and suppose that $`(h_1,\mathrm{},h_r)`$ and $`(h_1^{},\mathrm{},h_r^{})`$ are $`r`$-tuples with $`h_i`$, $`h_i^{}\{0,\frac{1}{2},\frac{1}{16}\}`$ such that $`h_i=\frac{1}{16}`$ (resp. $`h_i^{}=\frac{1}{16}`$) if and only if $`iI`$. If both $`m_{h_1,\mathrm{},h_r}`$ and $`m_{h_1^{},\mathrm{},h_r^{}}`$ are nonzero then $`m_{h_1,\mathrm{},h_r}=m_{h_1^{},\mathrm{},h_r^{}}`$. That is, all irreducible modules inside $`V^I`$ for $`T_r`$ have the same multiplicities.
(g) For any $`c,dC`$ and $`0vV(d)`$ we have $`V(c+d)=span\{u_nv|uV(c),n\}.`$
## 5. Code VOA $`M_C`$
In this section we review and extend results on code VOAs and their modules, following \[M1\]-\[M3\] and \[La\].
We shall sometimes consider an integer modulo 2 as its Euclidean lift, i.e., its representative 0 or 1 in $``$, so that when $`\alpha _2`$, $`\frac{1}{2}\alpha `$ makes sense as the rational number 0 or $`\frac{1}{2}`$.
Let $`C`$ be an even binary code. For any $`\alpha =(\alpha _1,\mathrm{},\alpha _n)C`$, denote
$$M_\alpha =L(\frac{1}{2},\frac{\alpha _1}{2})\mathrm{}L(\frac{1}{2},\frac{\alpha _n}{2})\text{ and }M_C=\underset{\alpha C}{}M_\alpha .$$
Note that $`M_C`$ is a simple current extension of $`T_n=L(\frac{1}{2},0)^n`$ and it has a unique VOA structure over $``$ (cf. \[DM2\], \[M2\]). This will be used to deduce the uniqueness of $`V^{\mathrm{}}`$.
###### Remark 5.1.
We use $`M_C`$ for a code VOA instead of $`M_D`$ given in \[M2\] in this paper. This is consistent with our code $`C`$ defined in Section 3. In fact, $`M_C`$ is a framed VOA with frame $`F`$ satisfying $`C(F)=C`$ and $`D(F)=0.`$
###### Remark 5.2.
For any $`\beta _2^n`$, one can define an automorphism $`\sigma _\beta :M_CM_C`$ by
$$\sigma _\beta (u)=(1)^{\alpha ,\beta }u\text{ for }uM_\alpha .$$
This automorphism is called a coordinate automorphism. Note that $`\sigma _\beta =\sigma _\beta ^{}`$ if and only if $`\beta +\beta ^{}C^{}`$ and the subgroup $`P`$ generated by $`\{\sigma _\beta |\beta _2^n\}`$ is isomorphic to $`_2^n/C^{}`$. Moreover, the fixed subalgebra $`M_{C}^{}{}_{}{}^{P}`$ is $`T_n`$ (cf. \[M1\]).
We first study the representations of the code VOA $`M_C.`$ Let $`W`$ be an irreducible $`M_C`$-module. Then $`W`$ can be written as a direct sum of irreducible $`T:=T_n`$-modules,
$$W\underset{h_i\{0,\frac{1}{2},\frac{1}{16}\}}{}m_{h_1,\mathrm{},h_n}L(h_1,\mathrm{},h_n).$$
###### Definition 5.3.
Define $`\tau (L(h_1,\mathrm{},h_n))=(a_1,\mathrm{},a_n)`$ $`_2^n`$ such that
$$a_i=\{\begin{array}{ccc}0\hfill & \text{ if }\hfill & h_i=0\text{ or }\frac{1}{2}\hfill \\ 1\hfill & \text{ if }\hfill & h_i=\frac{1}{16}\hfill \end{array}.$$
This binary word is called the $`\tau `$-word of $`L(h_1,\mathrm{},h_n).`$
By the fusion rules for $`L(\frac{1}{2},0)`$, the $`\tau `$-words for all irreducible $`T`$-submodules of $`W`$ are the same. Thus, we can also define the $`\tau `$-word of $`W`$ by
$$\tau \left(W\right)=\tau (L(h_1,\mathrm{},h_n)),$$
where $`L(h_1,\mathrm{},h_n)`$ is any irreducible $`T`$-submodule of $`W`$.
The following proposition is an easy consequence of the fusion rules (cf. \[DGH\] and \[M2\]).
###### Proposition 5.4.
Let $`C`$ be an even code and let $`W`$ be an irreducible module of $`M_C`$. Then $`\tau \left(W\right)`$ is orthogonal to $`C`$.
Now we shall give more details about the structure of the irreducible module $`W.`$ The details can be found in \[M2\]. Let $`\beta C^{}:=\{\alpha _2^n|\alpha ,\gamma =0\text{ for all }\gamma C\}`$ and $`C_\beta :=\{\alpha C|\mathrm{supp}\alpha \mathrm{supp}\beta \}`$.
Let the group $`\widehat{C}=\left\{\pm e^k\right|kC\}`$ be a central extension of $`C`$ by $`\{\pm 1\}`$ such that
$$e^he^k=(1)^{h,k}e^ke^h$$
for any $`h,kC`$ and denote $`\widehat{C}_\beta :=\left\{\pm e^k\right|kC_\beta \}\widehat{C}`$. Let $`H`$ be a maximal self-orthogonal subcode of $`C_\beta `$. Then $`\widehat{H}=\{\pm e^\alpha |\alpha H\}`$ is a maximal abelian subgroup of $`\widehat{C}_\beta `$ (it is automatically normal since it contains the commutator subgroup of $`\widehat{C}_\beta `$). Take a linear character $`\chi :\widehat{H}\{\pm 1\}`$ with $`\chi \left(e^0\right)=1`$ and define a 1-dimensional $`\widehat{H}`$-module $`F_\chi `$ by the action
$$e^\alpha p=\chi \left(e^\alpha \right)p\text{ for }pF_\chi ,\text{ }\alpha H.$$
We use “$`h_1\times h_2`$” to abbreviate a few of the well-known fusion rules involving $`L(\frac{1}{2},h_1)`$ and $`L(\frac{1}{2},h_2)`$, i.e., $`0\times h=h\times 0=h`$ for $`h\{0,\frac{1}{2},\frac{1}{16}\}`$, $`\frac{1}{2}\times \frac{1}{2}=0`$ and $`\frac{1}{2}\times \frac{1}{16}=\frac{1}{16}\times \frac{1}{2}=\frac{1}{16}`$.
For any $`h^i\{0,\frac{1}{2},\frac{1}{16}\},i=1,\mathrm{},n`$, with $`\tau \left(_{i=1}^nL(\frac{1}{2},h^i)\right)=\beta `$, we define
$$U=\left(_{i=1}^nL(\frac{1}{2},h^i)\right)F_\chi .$$
Then $`U`$ becomes an $`M_H`$-module with the vertex operator defined by
$$Y(\left(_{i=1}^nu^i\right)e^\alpha ,z)=\left(_{i=1}^nI^{\frac{a_i}{2},h^i}(u^i,z)\right)\chi \left(e^\alpha \right),$$
where $`u^iL(\frac{1}{2},\frac{a_i}{2})`$, $`(a_1,\mathrm{},a_n)H`$, and $`I^{\frac{a_i}{2},h^i}`$ is a nonzero intertwining operator of type
$$\left(\begin{array}{ccc}& L(\frac{1}{2},\frac{a_i}{2}\times h^i)& \\ L(\frac{1}{2},\frac{a_i}{2})& & L(\frac{1}{2},h^i)\end{array}\right).$$
We shall denote this $`M_H`$-module by $`U(\left(h^i\right),\chi )`$ or $`U\left(h^i\right)F_\chi `$.
Let $`\{\beta _j=\left(b_j^i\right)\}_{j=1}^s`$ be a transversal of $`H`$ in $`C`$ and
$$X=\underset{\beta _jC/H}{}\left\{U\left(h^i\times \frac{b_j^i}{2}\right)\left(e^{\beta _j}_{\widehat{H}}F_\chi \right)\right\},$$
Note that $`X`$ does not depend on the choice of the transversal of $`H`$ in $`C`$ and $`X`$ is an $`M_H`$-module.
The following results can be found in Miyamoto \[M2\].
###### Theorem 5.5.
$`X`$ is an $`M_C`$-module with
$$Y(u^\gamma e^\gamma ,z)=\left(_{i=1}^nI(u^i,z)\right)e^\gamma $$
for any $`\gamma C`$ and $`u^\gamma =_{i=1}^nu^iM_\gamma `$. We shall denote $`X`$ by $`\mathrm{Ind}_H^CU((h^i),\chi )`$.
###### Theorem 5.6.
For any irreducible $`M_C`$-module $`W`$, there is a pair $`(\left(h^i\right),\chi )`$ such that
$$W\mathrm{Ind}_H^C\left(U(\left(h^i\right),\chi )\right),$$
where $`\tau \left(W\right)=\tau \left(L(h^1,\mathrm{},h^n)\right)=\beta `$, $`H`$ is a maximal self-orthogonal subcode of $`C_\beta =\{\alpha C|\mathrm{supp}\alpha \mathrm{supp}\beta \}`$ and $`\chi `$ is a linear character of $`\widehat{H}`$. Moreover, the structure of the $`M_C`$-module $`W`$ is uniquely determined by an irreducible $`M_H`$-submodule of $`W`$.
Next we shall give a description of all irreducible $`M_C`$-modules by using some binary words. Let $`C`$ be an even code of length $`n.`$ For a given $`\beta C^{}`$ and $`\gamma _2^n`$, we define
$$h_{\beta ,\gamma }=(h_{\beta ,\gamma }^1,\mathrm{},h_{\beta ,\gamma }^n)\{0,\frac{1}{2},\frac{1}{16}\}^n$$
such that
$$h_{\beta ,\gamma }^i=\{\begin{array}{cc}\frac{1}{16}\hfill & \text{ if }\beta _i=1,\hfill \\ \frac{\gamma _i}{2}\hfill & \text{ if }\beta _i=0.\hfill \end{array}$$
Denote $`U(h_{_{\beta ,\gamma }})=U(h_{\beta ,\gamma }^1,\mathrm{},h_{\beta ,\gamma }^n)=L(h_{\beta ,\gamma }^1,\mathrm{},h_{\beta ,\gamma }^n)`$. Fix a maximal self-orthogonal subcode $`H_\beta `$ of the code $`C_\beta =\{\alpha C|\mathrm{supp}\alpha \mathrm{supp}\beta \}`$ and define a character $`\chi _\gamma :\widehat{H_\beta }`$ of the abelian group $`\widehat{H}_\beta `$ by
$$\chi _\gamma (e^0)=1\text{ and }\chi _\gamma (e^\alpha )=(1)^{\alpha ,\gamma }\text{ for }\alpha H_\beta .$$
Then $`(\beta ,\gamma )`$ determines an irreducible $`M_C`$-module
$$M_C(\beta ,\gamma )=\mathrm{Ind}_{H_\beta }^CU(h_{\beta ,\gamma }^1,\mathrm{},h_{\beta ,\gamma }^n)F_{\chi _\gamma }.$$
When there is no confusion, we shall simply denote $`M_C(\beta ,\gamma )`$ by $`M(\beta ,\gamma )`$.
###### Lemma 5.7.
The definition of $`M(\beta ,\gamma )`$ is independent of the choice of the self-orthogonal subcode $`H_\beta `$ of $`C_\beta `$.
Proof: Let $`H`$ be another maximal self-orthogonal subcode of $`C_\beta `$ and let $`\psi _\gamma :\widehat{H}`$ be a character of $`\widehat{H}`$ such that $`\psi _\gamma (e^\xi )=(1)^{\xi ,\gamma }`$ and $`\psi _\gamma (e^0)=1`$. Then we can construct another $`M_C`$-module
$$\mathrm{Ind}_H^CU(h_{\beta ,\gamma }^1,\mathrm{},h_{\beta ,\gamma }^n)F_{\psi _\gamma }.$$
By Miyamoto’s Theorem (Theorem 5.6), the structure of this module is uniquely determined by the structure of the $`M_H`$ submodule $`U(h_{\beta ,\gamma }^1,\mathrm{},h_{\beta ,\gamma }^n)F_{\psi _\gamma }`$. Thus,
$$\mathrm{Ind}_H^CU(h_{\beta ,\gamma }^1,\mathrm{},h_{\beta ,\gamma }^n)F_{\psi _\gamma }\mathrm{Ind}_{H_\beta }^CU(h_{\beta ,\gamma }^1,\mathrm{},h_{\beta ,\gamma }^n)F_{\chi _\gamma }$$
if and only if $`\mathrm{Ind}_{H_\beta }^CU(h_{\beta ,\gamma }^1,\mathrm{},h_{\beta ,\gamma }^n)F_{\chi _\gamma }`$ contains an $`M_H`$-submodule isomorphic to $`U(h_{\beta ,\gamma }^1,\mathrm{},h_{\beta ,\gamma }^n)F_{\psi _\gamma }.`$ It is equivalent to the fact that $`\mathrm{Res}_{\widehat{H}}\mathrm{Ind}_{\widehat{H}_\beta }^{\widehat{C_\beta }}\chi _\gamma ,\psi _\gamma 0,`$ where $`\mathrm{Res}_{\widehat{H}}\mathrm{Ind}_{\widehat{H}_\beta }^{\widehat{C_\beta }}\chi _\gamma ,\psi _\gamma `$ denotes the multiplicity of the character $`\psi _\gamma `$ in $`\mathrm{Res}_{\widehat{H}}\mathrm{Ind}_{\widehat{H}_\beta }^{\widehat{C_\beta }}\chi _\gamma `$.
On the other hand,
$$\mathrm{Res}_{\widehat{H}}\mathrm{Ind}_{\widehat{H}_\beta }^{\widehat{H_\beta +H}}F_{\chi _\gamma }\underset{\alpha (H+H_\beta )/H_\beta }{}e^\alpha F_{\chi _\gamma }\underset{\alpha H/HH_\beta }{}e^\alpha F_{\chi _\gamma }.$$
Let
$$w_\psi =\frac{1}{\left|HH_\beta \right|}\underset{\alpha H}{}\psi _\gamma (\alpha )e^\alpha v,$$
where $`vF_{\chi _\gamma }`$. Then $`w_\psi \mathrm{Ind}_{\widehat{H}_\beta }^{\widehat{H_\beta +H}}F_{\chi _\gamma }`$ and for any $`xH`$,
$`e^xw_\psi `$ $`={\displaystyle \frac{1}{\left|HH_\beta \right|}}{\displaystyle \underset{\alpha H}{}}(1)^{\gamma ,\alpha }e^xe^\alpha v`$
$`=(1)^{\gamma ,x}{\displaystyle \frac{1}{\left|HH_\beta \right|}}{\displaystyle \underset{\alpha H}{}}(1)^{\gamma ,\alpha +x}e^{x+\alpha }v`$
$`=(1)^{\gamma ,x}w_\psi =\psi _\gamma (e^x)w_\psi .`$
Hence $`w_\psi `$ affords the $`\widehat{H}`$-character $`\psi _\gamma `$ inside $`\mathrm{Ind}_{\widehat{H}_\beta }^{\widehat{H_\beta +H}}F_{\chi _\gamma }\mathrm{Ind}_{\widehat{H}_\beta }^{\widehat{C_\beta }}F_{\chi _\gamma }`$ and
$`\mathrm{Res}_{\widehat{H}}\mathrm{Ind}_{\widehat{H}_\beta }^{\widehat{C_\beta }}\chi _\gamma ,\psi _\gamma 0`$
as desired. ∎
###### Lemma 5.8.
Let $`\beta _1,\beta _2C^{}`$ and $`\gamma _1,\gamma _2_2^n`$. Let $`H_\beta `$ be a maximal self-orthogonal subcode of $`C_\beta `$ and let
$$(H_\beta )^_\beta :=\{\alpha _2^n|\mathrm{supp}\alpha \mathrm{supp}\beta \text{ and }\alpha ,\xi =0\text{ for all }\xi H_\beta \}.$$
Then the irreducible $`M_C`$-modules $`M(\beta _1,\gamma _1)`$ and $`M(\beta _2,\gamma _2)`$ are isomorphic if and only if
$$\beta _1=\beta _2\text{ and }\gamma _1+\gamma _2C+(H_\beta )^_\beta .$$
Proof: By the definition of $`M(\beta ,\gamma )`$, it is easy to see that $`M(\beta _1,\gamma _1)M(\beta _2,\gamma _2)`$ if $`\beta _1=\beta _2`$ and $`\gamma _1+\gamma _2C`$. Moreover, if $`\beta _1=\beta _2=\beta `$ and $`\gamma _1+\gamma _2(H_\beta )^_\beta `$, then $`h_{\beta _1,\gamma _1}=h_{\beta _2,\gamma _2}`$ and $`\chi _{\gamma _1}=\chi _{\gamma _2}`$ for any choice of $`H_\beta `$. Thus, $`M(\beta _1,\gamma _1)M(\beta _2,\gamma _2)`$ if
$$\beta _1=\beta _2\text{ and }\gamma _1+\gamma _2C+(H_\beta )^_\beta ,$$
Now suppose that $`M(\beta _1,\gamma _1)M(\beta _2,\gamma _2)`$. Then they have the same $`\tau `$-word and $`\beta _1=\beta _2`$. Let $`\beta :=\beta _1=\beta _2`$. Let $`H_\beta `$ be a maximal self-orthogonal subcode of $`C_\beta `$. Since $`M(\beta _1,\gamma _1)M(\beta _2,\gamma _2)`$, $`M(\beta _1,\gamma _1)`$ contains the $`M_{H_\beta }`$-module $`U(h_{\beta ,\gamma _2})\chi _{\gamma _2}`$. Thus, there exists an element $`\delta C`$ such that
$$h_{\beta ,\gamma _1}\times \frac{\delta }{2}=h_{\beta ,\gamma _2}\text{ and }e^\delta F_{\chi _{\gamma _1}}F_{\chi _{\gamma _2}}.$$
Since $`h_{\beta ,\gamma _1}\times \frac{\delta }{2}=h_{\beta ,\gamma _2}`$, $`\delta +\gamma _1+\gamma _2_2^\beta `$, where $`_2^\beta =\{\alpha _{2}^{}{}_{}{}^{n}|\mathrm{supp}\alpha \mathrm{supp}\beta \}`$. Moreover, $`e^\delta F_{\chi _{\gamma _1}}F_{\chi _{\gamma _2}}`$ implies that
$$(1)^{\delta +\gamma _1,\alpha }=(1)^{\gamma _2,\alpha }\text{ for all }\alpha H_\beta .$$
Therefore, $`\delta +\gamma _1+\gamma _2H_{\beta }^{}{}_{}{}^{}`$ and we have $`\delta +\gamma _1+\gamma _2H_{\beta }^{}{}_{}{}^{}_2^\beta =(H_\beta )^_\beta `$ and $`\gamma _1+\gamma _2C+(H_\beta )^_\beta .`$
###### Lemma 5.9.
The code $`C+H_\beta ^_\beta `$ is independent of the choice of the self-orthogonal subcode $`H_\beta `$.
###### Proof.
Let $`H`$ be another maximal self-orthogonal subcode of $`C_\beta `$. Then we have $`|H|=|H_\beta |`$. First we will consider the intersection $`HH_\beta `$ of $`H`$ and $`H_\beta `$ and its orthogonal complement in $`_2^\beta `$.
Claim: $`(HH_\beta )^_\beta =H_\beta ^_\beta +H`$.
It is easy to see that $`H_\beta ^_\beta `$ and $`H`$ are both contained in $`(HH_\beta )^_\beta `$. Hence we have $`H_\beta ^_\beta +H(HH_\beta )^_\beta `$. Now note that $`H_\beta ^_\beta H=H_\beta ^_\beta (C_\beta H)=(H_\beta ^_\beta C_\beta )H=H_\beta H`$ and $`dimH=dimH_\beta `$. By computing the dimensions, we have
$$\begin{array}{cc}\hfill dim(H_\beta ^_\beta +H)& =dimH_\beta ^_\beta +dimHdim(H_\beta ^_\beta H)\hfill \\ & =(|\beta |dimH_\beta )+dimHdim(H_\beta H)\hfill \\ & =|\beta |dim(H_\beta H)\hfill \\ & =dim(HH_\beta )^_\beta .\hfill \end{array}$$
Hence we have $`(HH_\beta )^_\beta =H_\beta ^_\beta +H`$.
By the claim, we have
$$(HH_\beta )^_\beta =H+H_\beta ^_\beta C+H_\beta ^_\beta .$$
Therefore, $`C+(HH_\beta )^_\beta C+H_\beta ^_\beta `$. On the other hand, $`C+H_\beta ^_\beta `$ is clearly contained in $`C+(HH_\beta )^_\beta `$ and thus $`C+H_\beta ^_\beta =C+(HH_\beta )^_\beta `$. Similarly, we also have $`C+H^_\beta =C+(HH_\beta )^_\beta `$ and hence $`C+H_\beta ^_\beta =C+H^_\beta `$ as desired. ∎
Next we shall compute the fusion rules among some irreducible $`M_C`$-modules. We recall a theorem proved by Miyamoto \[M2, M3\]. Let $`C`$ be an even linear code.
###### Theorem 5.10.
For any $`\alpha _2^n`$, the $`M_C`$-module $`M(0,\alpha )=M_{\alpha +C}=_{\delta \alpha +C}M_\delta `$ is a simple current module. Moreover,
$$M_{\alpha +C}\times M(\beta ,\gamma )=M(\beta ,\alpha +\gamma )$$
for any irreducible $`M_C`$-module $`M(\beta ,\gamma )`$.
Now by using the associativity and commutativity of the fusion rules, we also have the following Lemma.
###### Lemma 5.11.
Let $`\beta _1,\beta _2C^{}`$ and $`\gamma _2^n`$. Then
$$dimI_{M_C}\left(\genfrac{}{}{0pt}{}{M(\beta _1+\beta _2,\gamma )}{M(\beta _1,0)M(\beta _1,0)}\right)=dimI_{M_C}\left(\genfrac{}{}{0pt}{}{M(\beta _1+\beta _2,\alpha _1+\alpha _2+\gamma )}{M(\beta _1,\alpha _1)M(\beta _1,\alpha _2)}\right)$$
for any $`\alpha _1,\alpha _2_2^n`$.
Proof: For any $`\gamma _2^n`$, let
$$m_\gamma =dimI_{M_C}\left(\genfrac{}{}{0pt}{}{M(\beta _1+\beta _2,\gamma )}{M(\beta _1,0)M(\beta _1,0)}\right).$$
Then we have
$$M(\beta _1,0)\times M(\beta _2,0)=\underset{\gamma _2^n/K}{}m_\gamma M(\beta _1+\beta _2,\gamma ),$$
where $`K=C+(H_{\beta _1+\beta _2})^{_{\beta _1+\beta _2}}`$.
Since the fusion product is associative and commutative, we have
$$\begin{array}{cc}& M(\beta _1,\alpha _1)\times M(\beta _2,\alpha _2)\hfill \\ \hfill =& \left[M(0,\alpha _1)\times M(\beta _1,0)\right]\times \left[M(0,\alpha _2)\times M(\beta _2,0)\right]\hfill \\ \hfill =& \left[M(0,\alpha _1)\times M(0,\alpha _2)\right]\times \left[M(\beta _1,0)\times M(\beta _2,0)\right]\hfill \\ \hfill =& M(0,\alpha _1+\alpha _2)\times \left[M(\beta _1,0)\times M(\beta _2,0)\right]\hfill \\ \hfill =& M(0,\alpha _1+\alpha _2)\times \left(\underset{\gamma _2^n/K}{}m_\gamma M(\beta _1+\beta _2,\gamma )\right)\hfill \\ \hfill =& \underset{\gamma _2^n/K}{}m_\gamma M(\beta _1+\beta _2,\alpha _1+\alpha _2+\gamma ).\hfill \end{array}$$
Hence, we also have
$$m_\gamma =dimI_{M_C}\left(\genfrac{}{}{0pt}{}{M(\beta _1+\beta _2,\alpha _1+\alpha _2+\gamma )}{M(\beta _1,\alpha _1)M(\beta _1,\alpha _2)}\right)$$
as desired. ∎
For the later purpose we also need some facts about the Hamming code VOA $`M_{H_8}`$ from \[M2, M3\] (see also \[La\]).
Let $`H_8`$ be the Hamming code $`[8,4,4]`$ code, i.e., the code generated by the rows of
$$\left(\begin{array}{cc}1111& 1111\\ 1111& 0000\\ 1100& 1100\\ 1010& 1010\end{array}\right).$$
Let $`\{e^1,\mathrm{},e^8\}`$ be the standard frame for $`M_{H_8}`$. Let $`q^0=\mathrm{𝟏}`$ be the vacuum element of $`L(\frac{1}{2},0)`$ and let $`q^1`$ be a highest weight vector of $`L(\frac{1}{2},\frac{1}{2})`$ such that $`q_0^1q^1=\mathrm{𝟏}`$. For any $`\alpha =(\alpha _1,\mathrm{},\alpha _8)H_8`$, let
$$q^\alpha =q^{\alpha _1}\mathrm{}q^{\alpha _8}M_\alpha ,$$
where $`q^{\alpha _k}`$ is a norm 1 highest weight vector for the $`k`$-th tensor factor with respect to the action of our $`T_8`$. Then $`q^\alpha `$ is a highest weight vector in $`M_\alpha `$. Moreover, we have
$$q_{}^{\alpha }{}_{1}{}^{}q^\beta =\{\begin{array}{cc}2_1^8\alpha _ie^i\hfill & \text{ if }\alpha =\beta ,\hfill \\ q^{\alpha +\beta }\hfill & \text{ if }|\alpha \beta |=2,\hfill \\ 0\hfill & \text{ otherwise,}\hfill \end{array}$$
for any $`\alpha ,\beta H_8`$ with $`|\alpha |=|\beta |=4`$.
The following results are obtained in \[M2\].
###### Lemma 5.12.
Let $`\nu _i`$ be the binary word whose $`i`$-th entry is $`1`$ and all other entries are $`0`$. Define $`\alpha _i:=\nu _1+\nu _i`$.
In the Hamming code VOA $`M_{H_8}`$, there exist exactly three Virasoro frames, namely,
$$\{e^1,\mathrm{},e^8\},\{d^1,\mathrm{},d^8\},\text{ and }\{f^{1,}\mathrm{},f^8\}$$
where
$$d^i=S^{\alpha ^i}=\frac{1}{8}(e^1+\mathrm{}+e^8)+\frac{1}{8}\underset{\beta H_8,\left|\beta \right|=4}{}(1)^{\alpha _i,\beta }q^\beta e^\beta ,$$
$$f^i=S^{\nu _i}=\frac{1}{8}(e^1+\mathrm{}+e^8)+\frac{1}{8}\underset{\beta H_8,\left|\beta \right|=4}{}(1)^{\nu _i,\beta }q^\beta e^\beta $$
.
###### Theorem 5.13.
Let $`L`$ be an irreducible $`M_{H_8}`$-module with half-integral or integral weight. Then, $`L`$ is isomorphic to one of the following:
1. $`M_{\nu _1+\nu _i+H_8}`$ with respect to $`\{e^1,\mathrm{},e^8\}`$ for all $`i=1,\mathrm{},8.`$
2. $`M_{\nu _i+H_8}`$ with respect to $`\{e^1,\mathrm{},e^8\}`$ for all $`i=1,\mathrm{},8.`$
3. $`M_{\nu _i+H_8}`$ with respect to $`\{d^1,\mathrm{},d^8\}`$ for all $`i=1,\mathrm{},8.`$
4. $`M_{\nu _i+H_8}`$ with respect to $`\{f^{1,}\mathrm{},f^8\}`$ for all $`i=1,\mathrm{},8.`$
Moreover, all modules in (3) and (4) are isomorphic to $`_{i=1}^8L(\frac{1}{2},\frac{1}{16})`$ as $`T_8`$-modules.
As a corollary, we have the following theorem. The proof can be found in \[La\] (see also \[M2, M3\]).
###### Theorem 5.14.
For any $`\beta _1=(0^8)`$ or $`(1^8)`$ and $`\beta _2H_8`$, we have
$$M(\beta _1,\alpha _1)\times _{_{M_{H_8}}}M(\beta _2,\alpha _2)=M(\beta _1+\beta _2,\alpha _1+\alpha _2)$$
Consequently, all irreducible $`M_{H_8}`$-modules with half-integral or integral weight are simple current modules.
###### Remark 5.15.
For any $`\alpha _2^8/H_8`$, $`\alpha `$ uniquely determines a character $`\chi _\alpha \mathrm{Irr}H_8`$ such that $`\chi _\alpha (\gamma )=(1)^{\alpha ,\gamma }`$ for any $`\gamma H_8`$. By using this identification, our module $`M(\beta ,\alpha )`$ actually corresponds to the class $`[\beta ,\chi _\alpha ]`$ defined in Section 5 of \[La\].
## 6. The moonshine vertex operator algebra $`V^{\mathrm{}}`$
Let $`V^{\mathrm{}}`$ be the moonshine vertex operator algebra \[FLM1\]-\[FLM2\]. The following theorem can be found in \[DGH\].
###### Theorem 6.1.
There exists a VF in $`V^{\mathrm{}}`$, called $`F:=\{\omega _1,\mathrm{},\omega _{48}\}`$, so that the code $`𝒞:=C(F)`$ associated to this VF has length $`48`$ and dimension $`41`$. The code $`𝒟:=D(F)=𝒞^{}`$ has generator matrix
$$\left(\begin{array}{ccc}1111111111111111& 0000000000000000& 0000000000000000\\ 0000000000000000& 1111111111111111& 0000000000000000\\ 0000000000000000& 0000000000000000& 1111111111111111\\ 0000000011111111& 0000000011111111& 0000000011111111\\ 0000111100001111& 0000111100001111& 0000111100001111\\ 0011001100110011& 0011001100110011& 0011001100110011\\ 0101010101010101& 0101010101010101& 0101010101010101\end{array}\right)_.$$
###### Remark 6.2.
The weight enumerator of $`𝒟`$ is given by
$$X^{48}+3X^{32}+120X^{24}+3X^{16}+1$$
and the minimal weight of $`𝒞`$ is $`4`$. Moreover, $`𝒟`$ is self-orthogonal and hence $`𝒟𝒞.`$ (The codes $`𝒟`$ and $`𝒞`$ are denoted by $`S^{\mathrm{}}`$ and $`D^{\mathrm{}}`$, respectively in \[M3\].)
###### Lemma 6.3.
The code $`𝒞`$ in Theorem 6.1 is generated by the weight 4 codewords.
Proof: First we note that the code $`𝒟=D(F)`$ is generated by the elements of the form
$$(1^{16},0^{16},0^{16}),(0^{16},1^{16},0^{16}),(0^{16},0^{16},1^{16})\text{ and }(\alpha ,\alpha ,\alpha ),\alpha \mathrm{RM}(1,4),$$
where $`\mathrm{RM}(r,m)`$ denote the $`r`$-th order Reed-Muller code of length $`2^m`$ (cf. \[CS\]).
Since $`\mathrm{RM}(1,4)^{}=\mathrm{RM}(2,4)`$, we have
$$𝒞=𝒟^{}=\{(\alpha ,\beta ,\gamma )|\alpha +\beta +\gamma \mathrm{RM}(2,4),\alpha ,\beta ,\gamma \text{ even }\}.$$
Hence the code $`𝒞`$ can be generated by the elements
$$(\alpha ,0,0),(0,\beta ,0),(0,0,\gamma ),\alpha ,\beta ,\gamma \text{ are generators of }\mathrm{RM}(2,4)$$
and
$$(\alpha ,\beta ,0),(\alpha ,0,\beta ),(0,\alpha ,\beta ),\alpha ,\beta \text{ are even and }\alpha +\beta \text{ is a generator of }\mathrm{RM}(2,4).$$
Note that the Reed Muller code $`\mathrm{RM}(2,4)`$ is of dimension $`11`$ and is generated by the elements of the form
$$(\alpha ,0),(0,\alpha ),\alpha H_8,$$
and
$$(\mathrm{1100\hspace{0.17em}0000\hspace{0.17em}1100\hspace{0.17em}0000}),(\mathrm{1010\hspace{0.17em}0000\hspace{0.17em}1010\hspace{0.17em}0000}),(\mathrm{1000\hspace{0.17em}1000\hspace{0.17em}1000\hspace{0.17em}10000}).$$
Since the Hamming code $`H_8`$ is generated by its weight $`4`$ elements, the codes $`\mathrm{RM}(2,4)`$ and $`𝒞`$ are generated by the weight $`4`$ codewords also. ∎
In the next theorem, see 4.1(d) for the meaning of $`(V^{\mathrm{}})^0`$.
###### Lemma 6.4.
The vertex operator subalgebra $`(V^{\mathrm{}})^0`$ is isomorphic to $`M_𝒞`$ and is uniquely determined by the set of weight 4 codewords of $`𝒞.`$
Proof: By the uniqueness of the code VOA, $`(V^{\mathrm{}})^0`$ and $`M_𝒞`$ are isomorphic. Since $`𝒞`$ is generated by the weight 4 codewords of $`𝒞,`$ the vertex operator algebra structure of $`(V^{\mathrm{}})^0`$ is uniquely determined by the generators of the group $`𝒞.`$
We now determine the irreducible modules and the fusion rules for the code VOA $`M_𝒞`$.
###### Remark 6.5.
In the next result, $`\mathrm{RM}(r,m)`$ denote the $`r`$-th order Reed-Muller code of length $`2^m`$ (cf. \[CS\]). Note that the Reed Muller codes are nested in the sense that $`\mathrm{RM}(r+1,m)\mathrm{RM}(r,m)`$ and $`\mathrm{RM}(r+1,m+1)\mathrm{RM}(r,m)\mathrm{RM}(r,m)`$, where the direct sum corresponds to a partition of indices by an affine hyperplane and its complement.
The following properties of the code $`𝒞`$ can be derived easily from the definition.
###### Proposition 6.6.
Let $`𝒟`$ and $`𝒞`$ be defined as above. For any $`\beta 𝒟`$, denote
$$𝒞_\beta :=\{\alpha 𝒞|\mathrm{supp}\alpha \mathrm{supp}\beta \}.$$
1. If $`|\beta |=16`$, then $`𝒞_\beta \mathrm{RM}(2,4)`$.
2. If $`|\beta |=24`$, then $`𝒞_\beta \{(\alpha ,\gamma ,\delta )|\alpha +\gamma +\delta H_8\text{ and }\alpha ,\gamma ,\delta \text{ even}\}.`$
3. If $`|\beta |=32`$, then $`𝒞_\beta \mathrm{RM}(3,5).`$
4. If $`|\beta |=48`$, then $`𝒞_\beta =𝒞`$.
Note that the Hamming code $`H_8\mathrm{RM}(1,3)`$. Hence, for $`\beta 0`$, $`𝒞_\beta `$ contains a self-dual subcode which is isomorphic to a direct sum of $`|\beta |/8`$ copies of the Hamming code $`H_8`$.
Proof: Let $`\beta 𝒟`$ and $`n=|\beta |`$, the weight of $`\beta `$. Let $`p_\beta :_2^{48}_2^n`$ be the natural projection of $`_2^{48}`$ to the support of $`\beta `$. Since $`𝒞=𝒟^{}`$, it is easy to see that $`𝒞_\beta p_\beta (𝒟)^{}`$.
Case 1. $`|\beta |=16`$. In this case, $`p_\beta (𝒟)`$ is generated by the codewords $`(1^{16})`$, $`(0^8\mathrm{\hspace{0.17em}1}^8),(0^41^4)^2`$, $`(0^21^2)^4`$ and $`(\mathrm{0\; 1})^8`$ and is isomorphic to the Reed Muller code $`\mathrm{RM}(1,4)`$. Since
$$\mathrm{RM}(r,m)^{}\mathrm{RM}(mr1,m)$$
for any $`1rm`$ we have $`𝒞_\beta \mathrm{RM}(2,4)`$ as desired.
Case 2. $`|\beta |=24`$. In this case, $`p_\beta (𝒟)`$ is of dimension $`6`$ and is isomorphic to a code generated by
$$(1^8\mathrm{\hspace{0.17em}0}^8\mathrm{\hspace{0.17em}0}^8),(0^8\mathrm{\hspace{0.17em}1}^8\mathrm{\hspace{0.17em}0}^8),(0^8,0^8,1^8)\text{ and }(\alpha ,\alpha ,\alpha ),\alpha H_8.$$
Hence $`𝒞_\beta \{(\alpha ,\gamma ,\delta )|\alpha +\gamma +\delta H_8\text{ and }\alpha ,\gamma ,\delta \text{ even}\}`$.
Case 3. $`|\beta |=32`$. $`p_\beta (𝒟)\mathrm{RM}(1,5)`$ and hence $`𝒞_\beta \mathrm{RM}(3,5)`$.
Case 4. $`|\beta |=48`$. It is clear that $`𝒞_\beta =𝒞`$ in this case. ∎
Now by using Lemma 5.8, we have the following theorem.
###### Theorem 6.7.
Let $`𝒟`$ and $`𝒞`$ be defined as above. Then
$$\{M_𝒞(\beta ,\gamma )|\beta 𝒟\text{ and }\gamma _2^{48}/𝒞\}$$
is the set of all inequivalent irreducible modules for $`M_𝒞`$.
Proof: By the previous proposition, we can choose $`H_\beta `$ such that it is a direct sum of $`|\beta |/8`$ copies of the Hamming code $`H_8`$. In this case, $`(H_\beta )^_\beta =H_\beta `$ and we have $`𝒞=𝒞+(H_\beta )^_\beta `$. Hence $`\{M_𝒞(\beta ,\gamma )|\beta 𝒟\text{ and }\gamma _2^{48}/𝒞\}`$ is the set of all inequivalent irreducible modules for $`M_𝒞`$ by Lemma 5.8 . ∎
Next we shall compute the fusion rules among irreducible $`M_𝒞`$-modules. The main tool is the representation theory of the Hamming code VOA $`M_{H_8}`$ given in Section 4. First we recall the following theorem from \[DL\].
###### Theorem 6.8.
Let $`W^1,W^2`$ and $`W^3`$ be $`V`$-modules and let $`I`$ be an intertwining operator of type
$$\left(\begin{array}{ccc}& W^3& \\ W^1& & W^2\end{array}\right).$$
Assume that $`W^1`$ and $`W^2`$ have no proper submodules containing $`v^1`$ and $`v^2`$, respectively. Then $`I(v^1,z)v^2=0\text{ implies }I(,z)=0.`$
###### Lemma 6.9.
For any $`\beta _1,\beta _2,\beta _3𝒟`$ and $`\alpha _1,\alpha _2,\alpha _3_2^{48}`$, we have
$$dimI_{M_𝒞}\left(\genfrac{}{}{0pt}{}{M(\beta _3,\alpha _3)}{M(\beta _1,\alpha _1)M(\beta _2,\alpha _2)}\right)1$$
and
$$dimI_{M_𝒞}\left(\genfrac{}{}{0pt}{}{M(\beta _3,\alpha _3)}{M(\beta _1,\alpha _1)M(\beta _2,\alpha _2)}\right)=0$$
unless $`\beta _3=\beta _1+\beta _2`$ and $`\alpha _3\alpha _1+\alpha _2\mathrm{mod}𝒞`$.
Proof: Recall that $`dim𝒟=7`$ and the weight enumerator of $`𝒟`$ is $`X^{48}+3X^{32}+120X^{24}+3X^{16}+1.`$
Without loss, we may assume that $`\beta _3=\beta _1+\beta _2`$; otherwise,
$$dimI_{M_𝒞}\left(\genfrac{}{}{0pt}{}{M(\beta _3,\alpha _3)}{M(\beta _1,\alpha _1)M(\beta _2,\alpha _2)}\right)=0.$$
Let $`\overline{\beta _1}=(1^{48})+\beta _1`$. Then $`\overline{\beta _1}`$ is also in $`𝒟`$. Thus, there exist self-orthogonal codes $`H_{\beta _1}`$ and $`H_{\overline{\beta _1}}`$ of $`𝒞`$ such that both $`H_{\beta _1}`$ and $`H_{\overline{\beta _1}}`$ are direct sums of Hamming $`[8,4,4]`$ codes. Let $`E=H_{\beta _1}H_{\overline{\beta _1}}H_{8}^{}{}_{}{}^{6}`$. If the weight of $`\beta _2`$ is a multiple of $`16`$ (i.e., $`0,16,32,`$ or $`48`$), then $`|\mathrm{supp}\beta _1\mathrm{supp}\beta _2|`$ is a multiple of $`8`$. In this case, it is possible to find maximal self-orthogonal subcodes $`H_{\beta _2}`$ of $`𝒞_{\beta _2}`$ and $`H_{\beta _1+\beta _2}`$ of $`𝒞_{\beta _1+\beta _2}`$ such that $`H_{\beta _2}`$ and $`H_{\beta _1+\beta _2}`$ are isomorphic to direct sums of Hamming codes and are both contained in $`E`$. Then as an $`M_E`$-module,
$$M_𝒞(\beta _i,\alpha _i)=\underset{\delta 𝒞/E}{}M_E(\beta _i,\alpha _i+\delta ).$$
Note that $`H_{\beta _i}E`$ for any $`i=1,2,3`$ and hence $`M_𝒞(\beta _i,\alpha _i)`$ is a direct sum of inequivalent irreducible $`M_E`$-modules. Thus by Theorem 5.14 and 6.8, we have
$$dimI_{M_𝒞}\left(\genfrac{}{}{0pt}{}{M(\beta _1+\beta _2,\alpha _3)}{M(\beta _1,\alpha _1)M(\beta _2,\alpha _2)}\right)dimI_{M_E}\left(\genfrac{}{}{0pt}{}{M_E(\beta _1+\beta _2,\alpha _3)}{M_E(\beta _1,\alpha _1)M_E(\beta _2,\alpha _2)}\right)1$$
and
$$dimI_{M_𝒞}\left(\genfrac{}{}{0pt}{}{M(\beta _1+\beta _2,\alpha _3)}{M(\beta _1,\alpha _1)M(\beta _2,\alpha _2)}\right)=0$$
unless $`\alpha _3=\alpha _1+\alpha _2`$.
Finally, we shall treat the case for which all $`\beta _1`$, $`\beta _2`$ and $`\beta _1+\beta _2`$ are of weight $`24`$. For simplicity, we may assume that $`\beta _1=(1^8\mathrm{\hspace{0.17em}0}^8\mathrm{\hspace{0.17em}1}^8\mathrm{\hspace{0.17em}0}^8\mathrm{\hspace{0.17em}1}^8\mathrm{\hspace{0.17em}0}^8)`$ and $`\beta _2=(1^4\mathrm{\hspace{0.17em}0}^4\mathrm{}\mathrm{\hspace{0.17em}1}^4\mathrm{\hspace{0.17em}0}^4)`$. Then $`\beta _3=\beta _1+\beta _2=(0^4\mathrm{\hspace{0.17em}1}^8\mathrm{\hspace{0.17em}0}^8\mathrm{\hspace{0.17em}1}^8\mathrm{\hspace{0.17em}0}^8\mathrm{\hspace{0.17em}1}^8\mathrm{\hspace{0.17em}0}^4)`$. In this case, we have $`E=H_{\beta _1}H_{\overline{\beta }_1}H_{8}^{}{}_{}{}^{6}`$, $`H_{\beta _2}H_8H_8H_8`$ and $`H_{\beta _1+\beta _2}H_8H_8H_8`$.
Note that $`E+H_{\beta _1+\beta _2}=E+H_{\beta _2}`$ in this case. Moreover, we have
$$E_{\beta _2}=\{\alpha E|\mathrm{supp}\alpha \mathrm{supp}\beta _2\}=EH_{\beta _2}\text{ and }E_{\beta _3}=E_{\beta _1+\beta _2}=EH_{\beta _1+\beta _2}.$$
Let $`:=E+H_{\beta _2}=E+H_{\beta _1+\beta _2}`$. Then the $`M_𝒞`$-module $`M_𝒞(\beta _i,\alpha _i),i=2,3,`$ can be decomposed as
$$M_𝒞(\beta _i,\alpha _i)=\underset{\delta 𝒞/}{}M_{}(\beta _i,\alpha _i+\delta ).$$
Claim: $`M_{}(\beta _i,\alpha _i+\delta )`$ is irreducible as an $`M_E`$-module for any $`\delta 𝒞/`$.
Proof. Let $`W=M_{}(\beta _i,\alpha _i+\delta )`$ and $`_{\beta _i}=\{\alpha |\mathrm{supp}\alpha \mathrm{supp}\beta _i\}`$. Then $`H_{\beta _i}`$ is a maximal self-orthogonal subcode of $`_{\beta _i}`$. Let $`U(𝐡)F_\chi `$ be an irreducible $`M_{H_{\beta _i}}`$-submodule of $`W`$. Then
$$W=\mathrm{Ind}_{H_{\beta _i}}^{}U(𝐡)F_\chi =\underset{\delta /H_{\beta _i}}{}U(𝐡\times \frac{\delta }{2})(e^\delta F_\chi ).$$
Since $`E_{\beta _i}=EH_{\beta _i}H_{\beta _i}`$, $`U(𝐡)F_\chi `$ is also an irreducible $`M_{E_{\beta _i}}`$-module. Hence,
$$W^{}=\mathrm{Ind}_{E_{\beta _i}}^EU(𝐡)F_\chi =\underset{\delta E/E_{\beta _i}}{}U(𝐡\times \frac{\delta }{2})(e^\delta F_\chi )$$
is an irreducible $`M_E`$-submodule of $`W`$. Note that
$$|/H_{\beta _i}|=|(E+H_{\beta _i})/H_{\beta _i}|=|E/(EH_{\beta _i})|=|E/E_{\beta _i}|.$$
Therefore, we have $`W=W^{}`$ and $`W`$ is an irreducible $`M_E`$-module.
Now, by Theorem 5.14 and 6.8, we have
$$\begin{array}{cc}& dimI_{M_𝒞}\left(\genfrac{}{}{0pt}{}{M(\beta _1+\beta _2,\alpha _3)}{M(\beta _1,\alpha _1)M(\beta _2,\alpha _2)}\right)\hfill \\ \hfill & dimI_M_{}\left(\genfrac{}{}{0pt}{}{M_{}(\beta _1+\beta _2,\alpha _3)}{M_{}(\beta _1,\alpha _1)M_{}(\beta _2,\alpha _2)}\right)\hfill \\ \hfill & dimI_{M_E}\left(\genfrac{}{}{0pt}{}{M_E(\beta _1+\beta _2,\alpha _3)}{M_E(\beta _1,\alpha _1)M_E(\beta _2,\alpha _2)}\right)1\hfill \end{array}$$
and
$$dimI_{M_𝒞}\left(\genfrac{}{}{0pt}{}{M(\beta _1+\beta _2,\alpha _3)}{M(\beta _1,\alpha _1)M(\beta _2,\alpha _2)}\right)=0$$
unless $`\alpha _3=\alpha _1+\alpha _2`$. Note that $`M_{}(\beta _2,\alpha _2)=M_E(\beta _2,\alpha _2)`$ and $`M_{}(\beta _1+\beta _2,\alpha _1+\alpha _2)=M_E(\beta _1+\beta _2,\alpha _1+\alpha _2)`$ as $`M_E`$-modules. ∎
###### Theorem 6.10.
The fusion rules among irreducible $`M_𝒞`$ modules are given by
$$M(\beta _1,\alpha _1)\times M(\beta _2,\alpha _2)=M(\beta _1+\beta _2,\alpha _1+\alpha _2),$$
where $`\beta _1,\beta _2𝒟`$ and $`\alpha _1,\alpha _2_2^{48}/𝒞`$. In particular, each irreducible $`M_𝒞`$-module is a simple current.
Proof: By Lemma 5.11 and 6.9, it remains to show that
$$I_{M_𝒞}\left(\genfrac{}{}{0pt}{}{M(\beta _1+\beta _2,\alpha _1+\alpha _2)}{M(\beta _1,\alpha _1)M(\beta _2,\alpha _2)}\right)0,$$
for some $`\alpha _1,\alpha _2_2^{48}`$. Nevertheless, such kind of intertwining operators does exist and can be realized inside the Leech lattice VOA $`V_\mathrm{\Lambda }`$. In fact, there exists a Virasoro frame of $`V_\mathrm{\Lambda }`$ such that $`V_\mathrm{\Lambda }`$ can be decomposed as
$$V_\mathrm{\Lambda }\underset{\beta 𝒟}{}M_𝒞(\beta ,\gamma _\beta ),\text{ for some }\gamma _\beta _2^{48}/𝒞.$$
We shall refer to \[DGH\] or \[M3\] for details. ∎
## 7. Proof of the main theorems
We first prove Theorem 1. So we assume that (1) $`V`$ is a vertex operator algebra satisfying conditions (a)-(c), (2) $`V_2`$ is isomorphic to the Griess algebra, (3) $`V`$ is $`C_2`$-cofinite.
###### Lemma 7.1.
$`V`$ is truncated below 0 and $`V_0=\mathrm{𝟏}.`$
Proof: First we prove that $`V_n=0`$ if $`n`$ is negative. If this is not true, take the smallest $`n`$ such that $`V_n0.`$ Then each $`0vV_n`$ generates a highest weight module for the Lie algebra $`L(1)L(0)L(1)`$ (which is isomorphic to $`sl(2,).`$) According to the structure of the highest weight modules for $`sl(2,)`$ we know that $`L(1)^iv0`$ for $`i=0,\mathrm{},2n.`$ Since $`n`$ is less than or equal to $`1`$ we see that $`L(1)^{n+1}v0.`$ Since the weight of $`L(1)^{n+1}v`$ is 1, we immediately have a contradiction as $`V_1=0`$ by assumption.
We now prove that $`V_0`$ is one dimensional. Note that $`L(1)V_0=0.`$ So each nonzero vector $`vV_0`$ is a vacuum-like vector \[Li\]. As a result, we have a $`V`$-module isomorphism $`f_v:VV`$ by sending $`u`$ to $`u_1v`$ for $`uV`$ \[Li\]. By Schur’s lemma, $`f_v`$ must be a multiple of the identity map. As a result, $`f_v(\mathrm{𝟏})=v`$ is a multiple of the vacuum. This shows that $`V_0`$ is spanned by the vacuum. ∎
###### Lemma 7.2.
$`V`$ is a holomorphic vertex operator algebra.
Proof: It is proved in \[DLM2\] that if $`U`$ is a vertex operator algebra such that $`U=_{n0}U_n`$ with $`U_0`$ being 1-dimensional and $`U_1=0`$ and that $`U`$ is the only irreducible ordinary module for itself then any ordinary module is completely reducible. So any ordinary $`V`$-module is a direct sum of copies of $`V.`$ Since $`V`$ is $`C_2`$-cofinite, any submodule generated by a single vector in any admissible module is an ordinary module (see \[ABD\]). This shows that any admissible $`V`$-module is completely reducible. That is, $`V`$ is rational. This together with condition (a) gives conclusion that $`V`$ is holomorphic. ∎
###### Lemma 7.3.
The $`q`$-dimension $`ch_qV=q^1_{n0}(dimV_n)q^n`$ of $`V`$ is $`J(q).`$
Proof: Since $`V`$ is holomorphic and $`C_2`$-cofinite, by the modular invariance result in \[Z\], $`ch_qV`$ is a modular function on the full modular group, and thus equal to $`J(q)`$ by noting that $`V_0=\mathrm{𝟏}`$ and $`V_1=0.`$
Since $`V`$ is irreducible and $`V_0/L(1)V_1`$ is one dimensional, there is a unique nondegenerate symmetric invariant bilinear form $`(,)`$ on $`V`$ such that $`(\mathrm{𝟏},\mathrm{𝟏})=1`$ (see \[Li\]). That is,
$$(Y(u,z)v,w)=(z^2)^{\mathrm{wt}u}(v,Y(e^{zL(1)}u,z^1)w)$$
for homogeneous $`uV.`$ In particular, the restriction of $`(,)`$ to each $`V_n`$ is nondegenerate. As a result, $`(,)`$ defines a nondegenerate symmetric invariant bilinear form on the Griess algebra $`V_2`$ such that $`(u,v)=u_3v`$ for $`u,vV_2.`$
From now on we will fix the vectors $`\{\omega _1,\mathrm{},\omega _{48}\}`$ of $`V_2`$ given in Theorem 6.1. Since we only assume that $`V_2`$ is isomorphic to the Griess algebra we do not know if the bilinear form $`(u,v)=u_3v`$ defined on $`V_2`$ is the same as the bilinear form defined on $`V_2^{\mathrm{}}`$ using the same formula. So it is not clear that $`\{\omega _1,\mathrm{},\omega _{48}\}`$ forms a VF in $`V.`$
Since $`V_2`$ is a simple commutative nonassociative algebra, we need a result on the bilinear forms over a finite dimensional simple commutative nonassociative algebra $`B.`$ A bilinear form $`(,)`$ on $`B`$ is called invariant if $`(ab,c)=(b,ac),`$ for all $`a,b,cB`$. The next result applies to any finite dimensional simple algebra.
###### Lemma 7.4.
The space of nondegenerate symmetric invariant bilinear forms on $`B`$ is at most one-dimensional.
Proof: Let $`(,)`$ and $`,`$ be two nondegenerate symmetric invariant bilinear forms on $`B.`$ Then there is a linear isomorphism $`f:BB`$ such that $`(u,v)=f(u),v`$ for all $`u,vV.`$ For any $`aB`$ we have
$$f(au),v=(au,v)=(u,av)=f(u),av=af(u),v.$$
That is, $`f(au)=af(u).`$ Let $`B_\lambda `$ be the eigenspace of $`f`$ with eigenvalue $`\lambda 0`$ Then $`B_\lambda `$ is an ideal of $`B.`$ This shows that $`B=B_\lambda .`$ So $`f=\lambda id_B.`$ As a result, $`(,)=\lambda ,,`$ as desired. ∎
###### Lemma 7.5.
Each $`\omega _i`$ is a Virasoro vector with central charge $`\frac{1}{2}`$ and for all $`m,n`$,
$$[L^i(m),L^j(n)]=0$$
if $`ij`$ where $`Y(\omega _i,z)=_nL^i(n)z^{n2}.`$
Proof: We first prove that each $`\omega _i`$ is a Virasoro vector of central charge $`\frac{1}{2}.`$ That is, the component operators $`L^i(n)`$ of $`Y(\omega _i,z)=_nL^i(n)z^{n2}`$ satisfies the Virasoro algebra relation with central charge $`\frac{1}{2}.`$
Clearly $`\omega _i\omega _i=L^i(0)\omega _i=2\omega _i`$ by the product in $`B.`$ So, $`\omega _i`$ is a Virasoro vector with central charge $`c_i`$ defined by $`c_i\mathrm{𝟏}=2L^i(2)\omega _i.`$ Note that $`L^i(0)`$ is semisimple on $`V_2`$ and the eigenvalues of $`L^i(0)`$ are $`2,0,\frac{1}{2}`$ and $`\frac{1}{16}`$ (see \[DGH\]). Since the bilinear form is invariant, we see that the eigenspaces with different eigenvalues are orthogonal. So the restriction of the bilinear form to each eigenspace is nondegenerate. It is known from \[DGH\] that the eigenspace with eigenvalue $`2`$ is one dimensional and is spanned by $`\omega _i.`$ As a result, $`L^i(0)\omega _i`$ is nonzero and $`c_i0.`$ We must prove that $`c_i=\frac{1}{2}.`$
Recall from \[DM3\] that the Griess algebra is a simple commutative nonassociative algebra. Let $`,`$ be the bilinear from defined on $`V_2^{\mathrm{}}`$ and $`(,)`$ be the bilinear form defined on $`V_2.`$ By Lemma 7.4, $`(,)`$ is a multiple of $`,.`$ Note that $`\omega ,\omega =(\omega ,\omega )=12.`$ We conclude that these two bilinear forms are exactly the same. So $`(\omega _i,\omega _i)=\omega _i,\omega _i=\frac{1}{4}.`$ That is $`c_i=\frac{1}{2}.`$
Let $`ij.`$ Since $`(\omega _i,\omega _j)=0`$ and $`L^i(0)\omega _j=0`$, we see immediately that $`[L^i(m),L^j(n)]=0`$ for all $`m,n.`$
###### Theorem 7.6.
The $`\{\omega _1,\mathrm{},\omega _{48}\}`$ forms a VF in $`V`$ and $`V`$ is a FVOA.
Proof: We only need to prove that vertex operator subalgebra $`\omega _i`$ generated by $`\omega _i`$ is isomorphic to $`L(\frac{1}{2},0)`$ for the Virasoro algebra $`Vir_i`$ generated by $`L^i(m)`$ for $`m.`$ It is clear that $`\omega _i`$ is a highest weight module with highest weight $`0`$ for $`Vir_i.`$ Then there are two possibilities. Either $`\omega _i`$ is the Verma module modulo the submodule generated by $`L^i(1)\mathrm{𝟏}`$ or $`\omega _i`$ is isomorphic to $`L(\frac{1}{2},0)`$, according to the structure theory of highest weight modules for the Virasoro algebra with central charge $`\frac{1}{2}`$ \[FF\]. We now assume that the first possibility happens.
In this case the $`q`$-character of $`\omega _i`$ is equal to
$$ch_q\omega _i=q^{1/48}\frac{1}{_{n2}(1q^n)}.$$
Let $`U`$ be the vertex operator subalgebra of $`V`$ generated by $`\omega _j`$ for $`j=1,\mathrm{},48.`$ Then we have
$$U=\omega _1\mathrm{}\omega _{48}$$
is a tensor product. Let $`f(q^{\frac{1}{n}}),g(q^{\frac{1}{n}})[[q^{1/n},q^{1/n}]]`$ for some positive integer $`n.`$ We write $`f(q^{\frac{1}{n}})g(q^{\frac{1}{n}})`$ if the coefficient of $`q^m`$ in $`f(q^{\frac{1}{n}})`$ is less than or equal to that in $`g(q^{\frac{1}{n}})`$ for all $`m.`$ It is well known that the $`q`$-character of $`L(\frac{1}{2},0)`$ is equal to
$$\frac{1}{2}q^{1/48}\left(\underset{n0}{}(1+q^{n+\frac{1}{2}})+\underset{n0}{}(1q^{n+\frac{1}{2}})\right)$$
(cf. \[KR\]). Thus we have
$$ch_qUf(q^{\frac{1}{2}})$$
(7.1)
where
$$f(q^{\frac{1}{2}}):=q^1\frac{1}{2^{47}}\left(\underset{n0}{}(1+q^{n+\frac{1}{2}})+\underset{n0}{}(1q^{n+\frac{1}{2}})\right)^{47}\underset{n2}{}\frac{1}{(1q^n)}.$$
Clearly, both $`ch_qU`$ and $`f(q^{\frac{1}{2}})`$ are convergent for $`0<|q|<1`$, when $`q`$ is regarded as a complex number. So, we can and do treat both $`ch_qU`$ and $`f(q^{\frac{1}{2}})`$ as functions for $`0<q<1`$ and the inequality (7.1) still holds as functions.
We have already proved in Lemma 7.3 that the graded dimension of $`V`$ is $`J(q)`$ which of course also converges for $`0<|q|<1.`$ In the following we will take $`q`$ to be a real number in the domain $`(0,1).`$ Since $`U`$ is a subspace of $`V,`$ we have
$$\frac{ch_qU}{J(q)}1.$$
Let $`L`$ be the Niemeier lattice of type $`D_{24}.`$ Then the lattice vertex operator algebra $`V_L`$ is a module for the affine Lie algebra $`D_{24}^{(1)}.`$ Denote the irreducible highest weight module for $`D_{24}^{(1)}`$ of level $`k`$ by $`L_k(\lambda )`$ where $`\lambda `$ is a dominant weight of the finite dimensional Lie algebra of type $`D_{24}.`$ Let $`\lambda _i`$ be the fundamental weights of Lie algebra of type $`D_{24}`$ for $`i=1,\mathrm{},24`$ so that $`\lambda _{23}`$ and $`\lambda _{24}`$ are the half spin weights. (We are using the labelling of simple roots given in \[H\].) Then as a module for $`D_{24}^{(1)}`$ $`V_L`$ is a direct sum
$$V_L=L_1(0)L_1(\lambda _{23})$$
following from the structure of lattice $`L.`$ It is well-known that
$$ch_qV_L=\frac{\theta _L(q)}{\eta (q)^{24}}=J(q)+2\times (24)^224$$
where $`2\times (24)^224=1128`$ is the dimension of the Lie algebra of type $`D_{24},`$
$$\theta (q)=\underset{\alpha L}{}q^{(\alpha ,\alpha )/2}$$
is the theta function of the lattice $`L`$ and
$$\eta (q)=q^{1/24}\underset{n1}{}(1q^n).$$
So we have
$$J(q)<ch_qV_L$$
as a function in $`q(0,1).`$
On the other hand, using the Boson-Fermion correspondence given in \[F\], we see that the characters of the fermion realizations of $`L_1(0)`$ and $`L_1(\lambda _{23})`$ satisfy the following relations
$$ch_qL_1(0)ch_qL_1(0)+ch_qL_1(\lambda _1)=q^1\underset{n0}{}(1+q^{n+\frac{1}{2}})^{48}$$
$$ch_qL_1(\lambda _{23})=q^1\underset{n>0}{}(1+q^n)^{48}<2q^1\underset{n0}{}(1+q^{n+\frac{1}{2}})^{48}$$
As a result we have
$$J(q)ch_qV_L3q^1\underset{n0}{}(1+q^{n+\frac{1}{2}})^{48}.$$
Note that
$$f(q^{\frac{1}{2}})q^1\frac{1}{2^{47}}\underset{n0}{}(1+q^{n+\frac{1}{2}})^{47}\underset{n2}{}\frac{1}{(1q^n)}$$
So finally we have
$$\frac{ch_qU}{ch_qV}\frac{1}{2^{47}3}\underset{n0}{}\frac{1}{(1+q^{n+\frac{1}{2}})}\underset{n2}{}\frac{1}{(1q^n)}.$$
(7.2)
Clearly, the right hand side of (7.2) goes to infinity as $`q`$ goes to 1. This is a contradiction to $`\frac{ch_qU}{ch_qV}1.`$
###### Remark 7.7.
From the proof Theorem 7.6 we see that we in fact prove a stronger result: If $`\{u_1,\mathrm{},u_{48}\}`$ are 48 mutually commutative Virasoro elements of central charge $`\frac{1}{2}`$ then $`\{u_1,\mathrm{},u_{48}\}`$ is a VF.
###### Remark 7.8.
In the proof of Theorem 7.6 we only use the fact that $`ch_qV=J(q).`$ In fact, the proof goes through if we assume that $`dimV_nV_n^{\mathrm{}}`$ for $`n3.`$ So Theorem 7.6 holds with the assumptions given in Theorem 2.
Proof of Theorem 1: By Theorem 7.6, $`V`$ is an FVOA with VF $`F:=\{\omega _1,\mathrm{},\omega _{48}\}.`$ Let $`U`$ be the vertex operator subalgebra generated by $`V_2.`$ Then $`U`$ is also a FVOA with the same VF. Since $`F`$ is a VF in both $`U`$ and $`V`$, we use a subscript $`U`$ to indicate dependence of the associated binary codes on $`U`$. We have that $`𝒞`$ is a subcode of $`C_U(F)`$ and $`𝒟`$ is a subcode of $`D_U(F).`$ Since $`D_U(F)C_U(F)^{},`$ and $`𝒟=𝒞^{}`$ we immediately see that $`𝒞=C_U(F)`$ and $`𝒟=D_U(F).`$
Note that $`𝒞`$ is a subgroup of $`C(V)`$ and $`𝒟`$ is a subgroup of $`D(F).`$ Since $`V`$ is holomorphic by Lemma 7.3, $`C(F)=D(F)^{}`$ (see Theorem 4.1). This implies that $`C_U(F)=C(F)=𝒞`$ and $`D_U(F)=D(F)=𝒟.`$
Now $`M_𝒞=U^0`$ is a vertex operator subalgebra of $`V.`$ Then by Theorem 4.1, $`V`$ is a direct sum of inequivalent irreducible $`M_𝒞`$-modules. By Theorem 6.7, for each $`\delta 𝒟`$ there exists a unique $`\gamma _\delta _2^{48}/𝒞`$ such that $`M(\delta ,\gamma _\delta )`$ is isomorphic to a submodule of $`V.`$ Then
$$V\underset{\delta D}{}M(\delta ,\gamma _\delta )$$
as $`M_𝒞`$-module. Similarly, $`V^{\mathrm{}}`$ has a decomposition
$$V^{\mathrm{}}\underset{\delta D}{}M(\delta ,\beta _\delta )$$
where $`\beta _\delta _2^{48}/𝒞`$. In the case that the lowest weight of $`M(\delta ,\beta _\delta )`$ is 0 or 2, we have $`\beta _\delta =\gamma _\delta .`$ Since every module for $`M_𝒞`$ is a simple current by Theorem 6.10, by the uniqueness of simple current extension theorem in \[DM2\], it is sufficient to show that $`M(\delta ,\gamma _\delta )`$ and $`M(\delta ,\beta _\delta )`$ are isomorphic $`M_𝒞`$-modules.
For $`\delta 𝒟`$ we denote the lowest weight of $`M(\delta ,\beta _\delta )`$ by $`w(\delta ).`$ Set
$$X=\{(\delta ,\beta _\delta )|\delta 𝒟,w(\delta )=0,2\}.$$
Since $`V^{\mathrm{}}`$ is generated by $`V_2^{\mathrm{}}`$ (see \[FLM2\]), the group $`G:=\{(\delta ,\beta _\delta )|\delta 𝒟\}`$ is a subgroup of $`D\times _2^{48}/𝒞`$ generated by $`X.`$ So, the group $`H:=\{(\delta ,\gamma _\delta )|\delta 𝒟\}`$ is a subgroup of $`D\times _2^{48}/𝒞`$ and contains $`G`$ as a subgroup. As a result, $`G=H.`$ By Theorem 6.7, $`M(\delta ,\gamma _\delta )`$ and $`M(\delta ,\beta _\delta )`$ are indeed isomorphic $`M_𝒞`$-modules. ∎
Proof of Theorem 2: In this case, the conclusions of Lemmas 7.1, 7.4, 7.5 and Theorem 7.6 still hold (see Remark 7.8).
Let $`U`$ be as in the proof of Theorem 1. Since $`U`$ is generated by the Griess algebra, and $`𝒟𝒞,`$ $`C(U)=𝒞`$ and $`M_𝒞`$ is a subalgebra of $`U.`$ From the proof of Theorem 1 we see that
$$U\underset{\delta D}{}M(\delta ,\gamma _\delta )$$
as $`M_𝒞`$-modules. The same argument used in the proof of Theorem 1 shows that $`U`$ and $`V^{\mathrm{}}`$ are isomorphic. So we have
$$J(q)=ch_qV^{\mathrm{}}ch_qVch_qU=J(q).$$
As a result, $`U=V.`$ This completes the proof. ∎
We give an application of Theorem 2. Let $`U`$ be the $`_3`$ orbifold construction given in \[DM1\]. It has been expected for a long time that $`U`$ and $`V^{\mathrm{}}`$ are isomorphic vertex operator algebras. The isomorphism follows from Theorem 2 easily now.
###### Corollary 7.9.
$`V^{\mathrm{}}`$ and $`U`$ are isomorphic.
Proof: $`U`$ satisfies the conditions in Theorem 2. In particular, $`\mathrm{ch}_qU=J(q).`$
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# 1 Introduction
## 1 Introduction
Last year some Labs. reported that they observed a new resonance state $`\mathrm{\Theta }`$ with positive strangeness $`𝒮=+1`$ . The mass of this $`\mathrm{\Theta }`$ particle is around $`M_\mathrm{\Theta }=1540`$ MeV and the upper limit of the width is about $`\mathrm{\Gamma }_\mathrm{\Theta }<25`$ MeV. People suggested this exotic particle as a pentaquark state because its strangeness quantum number is +1. Although there are also several negative reports from some other Labs. , it has motivated an enormous amount of experimental and theoretical studies in the baryon physics area, because if the existence of this exotic particle can be confirmed, it will be the first multi-quark state people discovered. Actually there have been a lots of theoretical works on the study of the pentaquark baryons with various quark models, e.g. , and other approaches , but its structure is still a challenging problem.
It is well known that the non-perturbative quantum chromodynamics (NPQCD) effect is very important in the light quark system, but up to now there is no serious approach to really solve the NPQCD problem. In this sense, people still need QCD-inspired models to help in the low energy region. The constituent quark model is quite successful in explaining the baryon spectrum, especially when the non-perturbative QCD effect is considered by introducing the chiral field coupling , the Roper resonance and the low excited states of $`\mathrm{\Lambda }`$ can be explained simultaneously. At the same time, the chiral SU(3) quark model can reasonably reproduce the binding energy of deuteron, the nucleon-nucleon ($`NN`$) and the kaon-nucleon ($`KN`$) scattering phase shifts of different partial waves, and the hyperon-nucleon ($`YN`$) cross sections by the resonating group method (RGM) calculations . Inspired by these achievements, we try to extend this model to study the $`5`$ quark system.
In our previous work , we calculated the energies of eight low configurations of the $`5q`$ system, four lowest configurations of $`J^\pi =\frac{1}{2}^{}`$ with $`4q`$ partition $`[4]_{orb}(0s^4)[31]^{\sigma f}`$ and four of $`J^\pi =\frac{1}{2}^+`$ with $`4q`$ partition $`[31]_{orb}(0s^30p)[4]^{\sigma f}`$. But the results of the adiabatic approximation calculation show that the mass of the lowest $`5q`$ cluster state is about $`150300`$ MeV higher than the observed one of the $`\mathrm{\Theta }^+`$ particle when the model parameters are taken in the reasonable region.
The purpose of this paper is to do a further study on the $`5q`$ system with strangeness $`𝒮=+1`$ based on Ref.. Some modifications are made: (1) The orbital wave function is extended as an expansion of the harmonic oscillator form with 4 different sizes $`b_i(i=14)`$ to improve the adiabatic approximation. (2) Three various forms (quadratic, linear and error function form) of the color confinement potential are considered to examine the effects from the different confinement potential. (3) Seven low lying $`J^\pi =\frac{1}{2}^+`$ states of partition $`[4]_{orb}(0s^30p)[31]^{\sigma f}`$ are added and the mixing between configurations $`[31]_{orb}(0s^30p)[4]^{\sigma f}`$ and $`[4]_{orb}(0s^30p)[31]^{\sigma f}`$ is also studied. Meanwhile, the parameters are chosen three different groups, one is fitted by the $`NN`$ and $`YN`$ scattering experimental data and the other two are fitted by the $`KN`$ scattering phase shifts . The results show that with various groups of parameters and different color confinement potentials, the $`T=0`$ state is still always the lowest one for both $`J^\pi =\frac{1}{2}^{}`$ and $`J^\pi =\frac{1}{2}^+`$ states, and $`J^\pi =\frac{1}{2}^{}`$, $`T=0`$ state is always lower than that of $`J^\pi =\frac{1}{2}^+`$. In addition, the modification of the adiabatic approximation and the confinement potential of error function form can improve the calculated energy by several tens MeV to the lowest state, respectively. When the parameters are taken for the case by fitting the $`KN`$ phase shifts, the energies of the system are much higher than those for the case by fitting $`NN`$ and $`YN`$ data, because the $`S`$ wave interaction between $`K`$ and $`N`$ is repulsive. Some of the $`J^\pi =\frac{1}{2}^+`$ states with 4q partition $`[4]_{orb}(0s^30p)[31]^{\sigma f}`$ are lower than those with partition of $`[31]_{orb}(0s^30p)[4]^{\sigma f}`$ and the mixing between these two configurations can make the results $`4090`$ MeV lower. As a consequence, all of these modifications can offer several tens to hundred MeV effect, and the calculated energy of the lowest state is still about $`245`$ MeV higher than the experimental mass of $`\mathrm{\Theta }`$. It seems that when the model space is chosen as $`5q`$ cluster, it is difficult to get the calculated mass close to the observed one by using the chiral quark model with the reasonable parameters.
The paper is arranged as follows. The theoretical framework of the chiral SU(3) quark model and the determination of parameters are briefly introduced in Section 2. The calculation results of different confinement potentials with three groups of parameters are listed and discussed in Section 3. Finally conclusions are drawn in Section 4.
## 2 Theoretical framework
### 2.1 The model
As mentioned in Ref , the Hamiltonian of the $`5q`$ system is written as
$$H=\underset{i}{}T_iT_G+\underset{i<j=1}{\overset{4}{}}V_{ij}+\underset{i=1}{\overset{4}{}}V_{i5},$$
$`(1)`$
where $`_iT_iT_G`$ is the kinetic energy of the system, $`V_{ij}(i,j=14)`$ and $`V_{i5}(i=14)`$ represent the interactions between quark-quark ($`qq`$) and quark-antiquark ($`q\overline{q}`$) respectively. In the chiral SU(3) quark model, the quark-quark interaction includes three parts: color confinement potential $`V_{ij}^{conf}`$, one gluon exchange (OGE) interaction $`V_{ij}^{OGE}`$ and chiral field coupling induced interaction $`V_{ij}^{ch}`$,
$$V_{ij}=V_{ij}^{conf}+V_{ij}^{OGE}+V_{ij}^{ch},$$
$`(2)`$
where the confinement potential $`V_{ij}^{conf}`$, which provides the non-perpurbative QCD effect in the long distance, is taken as three different forms in this work, i.e. quadratic, linear, and error function form. And the expression of $`V_{ij}^{OGE}`$ is
$$V_{ij}^{OGE}=\frac{1}{4}g_ig_j\left(\lambda _i^c\lambda _j^c\right)\left\{\frac{1}{r_{ij}}\frac{\pi }{2}\delta \left(\stackrel{}{r}_{ij}\right)\left(\frac{1}{m_{qi}^2}+\frac{1}{m_{qj}^2}+\frac{3}{4}\frac{1}{m_{qi}m_{qj}}\left(\stackrel{}{\sigma }_i\stackrel{}{\sigma }_j\right)\right)\right\}$$
$$+V_{\text{tensor}}^{OGE}+V_{\stackrel{}{\text{l}}\stackrel{}{\text{s}}}^{OGE},$$
$`(3)`$
which governs the short-range perturbative QCD behavior. $`V_{ij}^{ch}`$ represents the interactions from chiral field couplings and describes the nonperturbative QCD effect of the low-momentum medium-distance range, which can be derived from the chiral-quark coupling interaction Lagrangian
$$_I^{ch}=g_{ch}F(q^2)\overline{\psi }\left(\underset{a=0}{\overset{8}{}}\sigma _a\lambda _a+i\underset{a=0}{\overset{8}{}}\pi _a\lambda _a\gamma _5\right)\psi ,$$
$`(4)`$
where $`\lambda _0`$ is a unitary matrix, $`\sigma _0`$,….,$`\sigma _8`$ are the scalar nonet fields and $`\pi _0`$,….,$`\pi _8`$ the pseudoscalar nonet fields. $`_I^{ch}`$ is invariant under the infinitesimal chiral SU(3)$`{}_{L}{}^{}\times `$ SU(3)<sub>R</sub> transformation. Thus $`V_{ij}^{ch}`$ can be expressed as
$$V_{ij}^{ch}=\underset{a=0}{\overset{8}{}}V_{s_a}(\stackrel{}{r}_{ij})+\underset{a=0}{\overset{8}{}}V_{ps_a}(\stackrel{}{r}_{ij}).$$
$`(5)`$
The expressions of $`V_{s_a}(\stackrel{}{r}_{ij})`$ and $`V_{ps_a}(\stackrel{}{r}_{ij})`$ can be found in Refs .
The interaction between $`q`$ and $`\overline{q}`$ includes two parts: direct interaction and annihilation part,
$$V_{q\overline{q}}=V_{q\overline{q}}^{dir}+V_{q\overline{q}}^{ann},$$
$`(6)`$
$$V_{q\overline{q}}^{dir}=V_{q\overline{q}}^{conf}+V_{q\overline{q}}^{OGE}+V_{q\overline{q}}^{ch},$$
$`(7)`$
with
$$V_{q\overline{q}}^{ch}(\stackrel{}{r})=\underset{i}{}(1)^{G_i}V_{qq}^{ch,i}(\stackrel{}{r}).$$
$`(8)`$
Here $`(1)^{G_i}`$ describes the $`G`$ parity of the $`i`$th meson. For the $`\mathrm{\Theta }`$ particle case without vector meson exchanges, $`q\overline{q}`$ can only annihilation into a $`K`$ meson, thus $`V_{i5}^{ann}`$ can be expressed as
$$V_{q\overline{q}}^{ann}=V_{ann}^K,$$
$`(9)`$
with
$$V_{ann}^K=C_{ann}^K\left(\frac{1\stackrel{}{\sigma }_q\stackrel{}{\sigma }_{\overline{q}}}{2}\right)_{spin}\left(\frac{2+3\lambda _q\lambda _{\overline{q}}^{}}{6}\right)_{color}\left(\frac{19}{9}+\frac{1}{6}\lambda _q\lambda _{\overline{q}}^{}\right)_{flavor}\delta \left(\stackrel{}{r}_q\stackrel{}{r}_{\overline{q}}\right),$$
$`(10)`$
where we treat $`C_{ann}^K`$ as a parameter and adjust it to fit the mass of the $`K`$ meson.
In this work, calculations are carried on with three groups of parameters. First, we take the parameters which can reasonably reproduce the experimental data of $`NN`$ and $`NY`$ scattering (case I) . By some special constraints, the model parameters are fixed in the following way: the chiral coupling constant $`g_{ch}`$ is fixed by
$$\frac{g_{ch}^2}{4\pi }=\left(\frac{3}{5}\right)^2\frac{g_{NN_\pi }^2}{4\pi }\frac{m_u^2}{M_N^2},$$
$`(11)`$
with $`g_{NN_\pi }^2/4\pi =13.67`$ taken as the experimental value. The mass of the mesons are also adopted the experimental values, except for the $`\sigma `$ meson, whose mass is treated as an adjustable parameter. $`\eta `$, $`\eta ^{}`$ mesons are mixed by $`\eta _0`$, $`\eta _8`$,
$$\eta ^{}=\eta _8\mathrm{sin}\theta ^{PS}+\eta _0\mathrm{cos}\theta ^{PS},$$
$$\eta =\eta _8\mathrm{cos}\theta ^{PS}\eta _0\mathrm{sin}\theta ^{PS},$$
$`(12)`$
with the mixing angle $`\theta ^{PS}`$ taken to be the usual value $`23^{}`$.
The one gluon exchange coupling constants $`g_u`$ and $`g_s`$ can be determined by the mass splits between $`N`$, $`\mathrm{\Delta }`$ and $`\mathrm{\Sigma }`$, $`\mathrm{\Lambda }`$ respectively. The confinement strengths $`a_{uu}^c`$, $`a_{us}^c`$, and $`a_{ss}^c`$ are fixed by the stability conditions of $`N`$, $`\mathrm{\Lambda }`$, and $`\mathrm{\Xi }`$, and the zero point energies $`a_{uu}^{c0}`$, $`a_{us}^{c0}`$, and $`a_{ss}^{c0}`$ by fitting the masses of $`N`$, $`\mathrm{\Sigma }`$, and $`\overline{\mathrm{\Xi }+\mathrm{\Omega }}`$, respectively.
Another two groups of parameters are fitted by $`KN`$ scattering , in which the scalar meson mixing between the flavor singlet and octet mesons is considered, i.e. $`\sigma `$ and $`ϵ`$ mesons are mixed by $`\sigma _0`$ and $`\sigma _8`$,
$$\sigma =\sigma _8\mathrm{sin}\theta ^S+\sigma _0\mathrm{cos}\theta ^S,$$
$$ϵ=\sigma _8\mathrm{cos}\theta ^S\sigma _0\mathrm{sin}\theta ^S.$$
$`(13)`$
The mixing angle $`\theta ^S`$ is an open problem because the structure of the $`\sigma `$ meson is unclear and controversial. We adopt two possible values by which we can get reasonable $`KN`$ phase shifts, one is ideally mixing $`\theta ^S=35.264^{}`$ (case II), which means that $`\sigma `$ only act on the $`u(d)`$ quark, and $`ϵ`$ on the $`s`$ quark, the other is $`\theta ^S=18^{}`$ (case III), which is provided by Dai and Wu based on their recent investigation .
The three sets of model parameters are tabulated in Table 1, the first column (case I) is for the case fitted by $`NN`$ and $`YN`$ scattering, and the second (case II) and third (case III) columns for the case fitted by $`KN`$ scattering.
### 2.2 The configurations
For the $`5q`$ system with $`𝒮=+1`$, there are a lot of different configurations. In the actually calculating process, we have to choose some lower configurations. Using this model, we considered 15 states of the $`5q`$ system with strangeness $`𝒮=+1`$: four lowest $`J^\pi =\frac{1}{2}^{}`$ states with $`4q`$ partition $`[4]_{orb}(0s)^4`$ $`[4]_{orb}[31]_{ts=01}^{\sigma f}\overline{s},LST=0\frac{1}{2}0,J^\pi =\frac{1}{2}^{}`$, $`[4]_{orb}[31]_{ts=10}^{\sigma f}\overline{s},LST=0\frac{1}{2}1,J^\pi =\frac{1}{2}^{}`$,
$`[4]_{orb}[31]_{ts=11}^{\sigma f}\overline{s},LST=0\frac{1}{2}1,J^\pi =\frac{1}{2}^{}`$, $`[4]_{orb}[31]_{ts=21}^{\sigma f}\overline{s},LST=0\frac{1}{2}2,J^\pi =\frac{1}{2}^{}`$,
and four configurations of $`J^\pi =\frac{1}{2}^+`$ with $`[31]_{orb}(0s)^3(0p)`$
$`[31]_{orb}[4]_{ts=00}^{\sigma f}\overline{s},LST=1\frac{1}{2}0,J^\pi =\frac{1}{2}^+`$, $`[31]_{orb}[4]_{ts=11}^{\sigma f}\overline{s},LST=1\frac{1}{2}1,J^\pi =\frac{1}{2}^+`$,$`[31]_{orb}[4]_{ts=11}^{\sigma f}\overline{s},LST=1\frac{3}{2}1,J^\pi =\frac{1}{2}^+`$, $`[31]_{orb}[4]_{ts=22}^{\sigma f}\overline{s},LST=1\frac{3}{2}2,J^\pi =\frac{1}{2}^+`$.
Here the symbols $`[f]_{orb}`$ and $`[f^{}]^{\sigma f}`$ are the partitions of orbital space and flavor-spin space respectively; $`[4]_{orb}`$ represents the total symmetric state in the orbital space, where four quarks are all in $`(0s)`$ state; and $`[31]_{orb}`$ is the orbital space partition of $`(0s)^3(0p)`$. Symbol $`t`$ and $`s`$ denote the isospin and spin of the four quark part; after coupling with the fifth quark $`\overline{s}`$, the total orbital angular momentum, spin and isospin of the five quark system are expressed as $`LST`$. The color part is $`(10)_c`$ of $`4q`$ and $`(01)_c`$ of the anti-quark $`\overline{s}`$ respectively, here we omitted them in the expressions.
Meanwhile, there are some other $`J^\pi =\frac{1}{2}^+`$ states (we did not consider them in our previous work ) in which the $`\overline{s}`$ can be in the $`(0p)`$ state. The wave functions of these states should be orthogonal to the excited states of the center of mass motion. The expression of the wave function with various $`ts`$ and $`LST`$, $`\mathrm{\Psi }_{ts}^{LST}(5q)`$, is given as following
$$\mathrm{\Psi }_{ts}^{LST}(5q)=\sqrt{\frac{m_s}{4m_u+m_s}}\left(\mathrm{\Psi }_{4q}((0s^30p)[4]_{orb}[31]_{ts}^{\sigma f}[211]_c)\mathrm{\Phi }_{0s}(\overline{s})\right)_{LST}$$
$$\sqrt{\frac{4m_u}{4m_u+m_s}}\left(\mathrm{\Psi }_{4q}((0s^4)[4]_{orb}[31]_{ts}^{\sigma f}[211]_c)\mathrm{\Phi }_{0p}(\overline{s})\right)_{LST}.$$
$`(14)`$
Where in the first term, the $`4q`$ orbital wave function is $`(0s)^3(0p)`$ with partition $`[4]_{orb}`$ and $`\overline{s}`$ is in $`(0s)`$, while in the second term, the $`4q`$ orbital wave function is $`(0s)^4`$ with partition $`[4]_{orb}`$ and $`\overline{s}`$ is in $`(0p)`$. $`[31]_{ts}^{\sigma f}`$ is the $`4q`$ partition of spin-flavor space. The total orbital angular momentum, spin and isospin of the five quark system are $`LST`$. We considered seven low states of $`\mathrm{\Psi }_{ts}^{LST}(5q)`$, they are
$`[4]_{orb}[31]_{ts=01}^{\sigma f}\overline{s},LST=1\frac{1}{2}0,J^\pi =\frac{1}{2}^+`$, $`[4]_{orb}[31]_{ts=10}^{\sigma f}\overline{s},LST=1\frac{1}{2}1,J^\pi =\frac{1}{2}^+`$, $`[4]_{orb}[31]_{ts=11}^{\sigma f}\overline{s},LST=1\frac{1}{2}1,J^\pi =\frac{1}{2}^+`$, $`[4]_{orb}[31]_{ts=21}^{\sigma f}\overline{s},LST=1\frac{1}{2}2,J^\pi =\frac{1}{2}^+`$,$`[4]_{orb}[31]_{ts=01}^{\sigma f}\overline{s},LST=1\frac{3}{2}0,J^\pi =\frac{1}{2}^+`$, $`[4]_{orb}[31]_{ts=11}^{\sigma f}\overline{s},LST=1\frac{3}{2}1,J^\pi =\frac{1}{2}^+`$,$`[4]_{orb}[31]_{ts=21}^{\sigma f}\overline{s},LST=1\frac{3}{2}2,J^\pi =\frac{1}{2}^+`$.
In order to improve the adiabatic approximation, in this work the so-called ”breath model” is taken to carry on the energy calculation. The trail wave function can be written as an expansion of the $`5q`$ states with several different harmonic oscillator frequency $`\omega _i`$,
$$\mathrm{\Psi }_{5q}=\underset{i}{\overset{n}{}}\alpha _i\mathrm{\Phi }_{5q}(b_i).$$
$`(15)`$
Where $`(b_i)^2=\frac{1}{m\omega _i}`$. Using the fractional parentage coefficient (f.p.) technique, we can easily write down the wave functions of the above 15 states and by solving the Schrödinger equation the energies of these states are obtained.
## 3 Results and discussions
We calculate energies of fifteen low configurations, four lowest configurations of $`J^\pi =\frac{1}{2}^{}`$ with $`4q`$ partition $`[4]_{orb}(0s^4)[31]^{\sigma f}`$, four of $`J^\pi =\frac{1}{2}^+`$ with $`4q`$ partition $`[31]_{orb}(0s^30p)[4]^{\sigma f}`$ and seven low configurations of $`J^\pi =\frac{1}{2}^+`$ with $`4q`$ partition $`[4]_{orb}(0s^30p)[31]^{\sigma f}`$. Comparing with our previous work , some modifications are made in this work. The trial wave function is taken as an expansion of harmonic oscillator wave functions with four different size parameter $`b_i(i=14)`$ to improve the adiabatic approximation. Three different forms (quadratic, linear and error function form) of the confinement potential are considered to study the effects from various confinement potentials. More low lying states of $`J^\pi =\frac{1}{2}^+`$ with $`4q`$ partition $`[4]_{orb}(0s^30p)`$ and the mixing between some configurations are considered. At the same time, parameters are chosen three different groups as listed in Table 1. The results are given in Table 2. In this table, case I means the parameters are fitted by $`NN`$ and $`YN`$ scattering, while case II by $`KN`$ scattering with $`\theta ^S=35^{}`$ and case III by $`KN`$ scattering with $`\theta ^S=18^{}`$. And $`r^2`$, $`r`$ and $`erf`$ represent the confinement potential is adopted as quadratic, linear and error function form respectively.
From Table 2, one can see that: (1) the isoscalar state $`T=0`$ is always the lowest state both in $`J^\pi =\frac{1}{2}^{}`$ and in $`J^\pi =\frac{1}{2}^+`$ cases, and $`[4]_{orb}[31]_{ts=01}^{\sigma f}\overline{s}`$, $`LST=0\frac{1}{2}0`$, $`J^\pi =\frac{1}{2}^{}`$ is the lowest one among all the states. For two sets of $`J^\pi =\frac{1}{2}^+`$ states, some configurations with $`4q`$ partition $`[4]_{orb}`$, especially ($`[4]_{orb}[31]_{ts=01}^{\sigma f}\overline{s}`$, $`LST=1\frac{1}{2}0`$, $`J^\pi =\frac{1}{2}^+`$), are about $`100200`$ MeV lower than those of $`[31]_{orb}`$ in various cases. It means that in $`J^\pi =\frac{1}{2}^+`$ state, the configurations with $`4q`$ partition $`[4]_{orb}`$ can not be neglected. (2) The effect of various confinement potentials is about several tens MeV. The energies obtained from the error function confinement potential are the lowest, and the influence is larger to the high energy configurations. (3) The energies of case I are lower than those of the other two cases, this is obvious, because in the $`NN`$ case the $`S`$ wave interaction is attraction, while in the $`KN`$ case it is repulsive.
The comparison between the results of adiabatic approximation (the wave function is taken as harmonic oscillator function with $`1b`$) and those of the ”breath model” (the wave function is treated as an expansion of the harmonic oscillator functions with 4 different frequency $`\omega `$) is shown in Table 3. Here the parameters are taken as case I, which is fitted by $`NN`$ and $`YN`$ scattering data. The results show that this modification can only reduce the energies about $`1030`$ MeV for different configurations. This means that even the trial wave function is modified, the energy of the $`5q`$ system can not be improved a lot.
In the case of $`J^\pi =\frac{1}{2}^+`$, $`([31]_{orb}[4]_{ts=11}^{\sigma f}\overline{s})_{1\frac{1}{2}1}`$ and $`([4]_{orb}[31]_{ts=11}^{\sigma f}\overline{s})_{1\frac{1}{2}1}`$ as well as $`([31]_{orb}[4]_{ts=11}^{\sigma f}\overline{s})_{1\frac{3}{2}1}`$ and $`([4]_{orb}[31]_{ts=11}^{\sigma f}\overline{s})_{1\frac{3}{2}1}`$ have the same quantum numbers. The configuration mixing has to be considered for these two cases. The results comparing with those without configuration mixing are given in Table 4. From this table, one can see that the configurations mixing effect is not very small, it can make energies about $`4090`$ MeV lower.
We also try to adjust the size parameter $`b_u`$ to be larger to see the influence. As an example, the results of lower configurations with $`b_u=0.6`$ fm in the chiral SU(3) quark model are given compared with the former results with $`b_u=0.5`$ fm in Table 5. In this case, the energies of all states become smaller, caused by the kinetic energy of the system is reduced for larger $`b_u`$. But the lowest energy $`1665`$ MeV is still about $`125`$ MeV higher than the experimental mass of the observed $`\mathrm{\Theta }`$.
In addition, we notice that recently Sachiko Takeuchi $`etal.`$ investigated $`uudd\overline{s}`$ pentaquarks by employing their quark models with the meson exchange and the effective gluon exchange as $`qq`$ and $`q\overline{q}`$ interactions, and dynamically solved the system by taking two quarks as a diquark-like $`qq`$ correlation. Their calculated value of the lowest configuration was $`19472144`$ MeV and the low mass close to the observed one could not be obtained, which was similar to our results. The main difference is that in Ref. there were parameter sets where the mass of the lowest positive-parity states became lower than that of the negative-parity states. However, in our work, where the model parameters are fitted by the $`NN`$,$`YN`$ and $`KN`$ scattering phase shifts and the $`q\overline{q}`$ $`s`$-channel interactions are fitted by the mass of the kaon meson, the results show that the $`J^\pi =\frac{1}{2}^{}`$, $`T=0`$ state is always the lowest one. As discussed in our previous work , if we omitted the interactions between $`4q`$ and $`\overline{q}`$, then the state with positive parity can be lower than that with negative parity, in agreement with what is claimed by Takeuchi and Shimizu , and Stancu and Riska . This means that how to treat the annihilation interactions reasonably is very important in the calculation.
Meanwhile, it is important to investigate the narrow width as well as the low mass of the $`\mathrm{\Theta }`$ particle. But in our present work, we concentrate our study on the masses of some low-lying $`5q`$ configurations with various quantum numbers. Since the coupling between $`5q`$ configuration states and continuum baryon-meson states such as $`KN`$ is not included, we can only get some qualitative information about the width from the wave functions. For example, by using the group theory method the lowest state of $`J^\pi =\frac{1}{2}^{}`$ and $`T=0`$, $`\mathrm{\Psi }([4]_{orb}(0s)^4[31]_{ts=01}^{\sigma f}\overline{s}(0s))_{LST=0\frac{1}{2}0,(00)_c}`$, can be expanded as:
$$\mathrm{\Psi }([4]_{orb}(0s)^4[31]_{ts=01}^{\sigma f}\overline{s}(0s))_{LST=0\frac{1}{2}0,(00)_c}=$$
$$\frac{1}{2}(\mathrm{\Phi }_{123})_{(s=\frac{1}{2})(t=\frac{1}{2})(11)_c}(\mathrm{\Phi }_{4\overline{5}})_{(s=0)(11)_c}$$
$$\frac{1}{2}\sqrt{\frac{1}{3}}(\mathrm{\Phi }_{123})_{(s=\frac{1}{2})(t=\frac{1}{2})(11)_c}(\mathrm{\Phi }_{4\overline{5}})_{(s=1)(11)_c}$$
$$+\sqrt{\frac{1}{3}}(\mathrm{\Phi }_{123})_{(s=\frac{3}{2})(t=\frac{1}{2})(11)_c}(\mathrm{\Phi }_{4\overline{5}})_{(s=1)(11)_c}$$
$$+\frac{1}{2}(\mathrm{\Phi }_{123})_{(s=\frac{1}{2})(t=\frac{1}{2})(00)_c}(\mathrm{\Phi }_{4\overline{5}})_{(s=0)(00)_c}$$
$$+\frac{1}{2}\sqrt{\frac{1}{3}}(\mathrm{\Phi }_{123})_{(s=\frac{1}{2})(t=\frac{1}{2})(00)_c}(\mathrm{\Phi }_{4\overline{5}})_{(s=1)(00)_c},$$
$`(16)`$
where the wave function is antisymmetrized and all terms are orthogonal and linearly independent each other. In the above expression (Eq.(16)), the fourth term is the component of $`S`$-wave $`KN`$, whose probability is $`(\frac{1}{2})^2=\frac{1}{4}`$. Such a large $`(25\%)`$ $`KN`$ component would consequently make the width of this state to be quite large. In this sense the $`5q`$ cluster $`\frac{1}{2}^+`$ state can be expected to have a smaller width. However in our present work, when the model parameters are fitted by the $`NN`$,$`YN`$ and $`KN`$ scattering phase shifts, the results show that the masses of the positive-parity states are higher than those of the negative-parity states, also much higher than the experimental value. In this framework, it is difficult to understand both the low mass and the narrow width of the $`\mathrm{\Theta }`$. But up to now, the existence of the $`\mathrm{\Theta }`$ particle is still a controversial problem, and some high-statistics experimental collaborations showed the negative results. From our results, we can just say that if the $`\mathrm{\Theta }`$ particle do exist, it can not be explained as a five-constituent-quark cluster, and its structure should be understood from other mechanism.
## 4 Conclusions
The structures of $`5q`$ cluster states with $`𝒮=+1`$ are further studied based on the $`5q`$ cluster configurations studied in the SU(3) chiral quark model. Fifteen low configurations and the mixing between some configurations are considered. In the calculation, the trial wave function is taken as an expansion of harmonic oscillator wave functions with four different size parameters. The effect of various color confinement potential is examined by taken three different forms (quadratic, linear and error function forms). With various groups of parameters and different color confinement potentials, the isoscalar state $`T=0`$ is always the lowest state both in $`J^\pi =\frac{1}{2}^{}`$ and in $`J^\pi =\frac{1}{2}^+`$ cases. And $`J^\pi =\frac{1}{2}^{}`$, $`T=0`$ is the lowest one among all of these configurations. The modification of the adiabatic approximation and the confinement potential of error function form can improve the calculated energy by several tens MeV to the lowest state, respectively. Since the $`S`$ wave interaction between $`K`$ and $`N`$ is repulsive, when the parameters are taken for the case by fitting the $`KN`$ phase shifts, the energies of the system are much higher than those for the case by fitting $`NN`$ and $`YN`$ data. The configurations mixing of $`J^\pi =\frac{1}{2}^+`$ states can make the calculated mass $`4090`$ MeV shifted. If we adjust the size parameter larger, $`b_u=0.6`$ fm, the energy of the lowest configuration will be $`1665`$ MeV. All of these modifications can only offer several tens to hundred MeV improvement, and the calculated value of the lowest state is still about $`125`$ MeV higher than the experimental mass of $`\mathrm{\Theta }`$. It seems that when the model space is chosen as $`5q`$ cluster, it is difficult to get the calculated mass close to the observed one in the framework of the chiral quark model with the reasonable parameters. Though the existence of the $`\mathrm{\Theta }`$ particle is still an open problem, this work just presents that it can not be regarded as a five-constituent-quark cluster, if it do exist, its structure should be explained by other mechanism.
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# Photoproduction of Baryons Decaying into N𝜋 and N𝜂
## 1 Introduction
The energy levels of bound systems and their decay properties provide valuable information about the constituents and their interactions Isgur:1995ei . In quark models, the dynamics of the three constituent quarks in baryons support a rich spectrum, much richer than the energy scheme experiments have established so far Capstick:bm ; Riska ; Metsch . This open issue is referred to as the problem of missing resonances. The intense discussion of the exotic baryon resonance $`\mathrm{\Theta }^+(1540)`$ Klempt:2004yz ; Dzierba:2004db ; Hicks:2005gp , of its existence and of its interpretation, has shown limits of the quark model and underlined the need for a deeper understanding of baryon spectroscopy. Here, the study of pentaquarks has played a pioneering role, but any new model has to be tested against the excitation spectrum of the nucleon as well. The properties of baryon resonances are presently under intense investigations at several facilities like ELSA (Bonn), GRAAL (Grenoble), JLab (Newport News), MAMI (Mainz), and SPring-8 (Hyogo). The aim is to identify the resonance spectrum, to determine spins, parities, and decay branching ratios and thus to provide constraints for models.
The largest part of our knowledge on baryons stems from pion induced reactions. In elastic $`\pi \mathrm{N}`$ scattering, the unitarity condition provides strong constraints for amplitudes close to the unitarity limit, since production couplings are related directly to the widths of resonances and to the cross section. If a resonance has however a large inelasticity, its production cross section in $`\pi \mathrm{N}`$ scattering is small and it contributes only weakly to the final state. Thus resonances may conceal themselves from observation in elastic scattering. This effect could be a reason why the number of observed states is much smaller than predicted by quark models Capstick:bm ; Riska ; Metsch . Information on resonances coupled weakly to the $`\pi \mathrm{N}`$ channel can be obtained from photoproduction experiments and the study of final states different from $`\pi \mathrm{N}`$ such as multibody final states or final states containing open strangeness.
The information from photoproduction experiments is complementary to experiments with hadronic beams and gives access to additional properties like helicity amplitudes. Experiments with polarised photons provide information which may be very sensitive to resonances having a small cross section. A clear example of such an effect is the observation of the $`\mathrm{N}(1520)\mathrm{D}_{13}`$ resonance in $`\eta `$ photoproduction. It contributes very little to the unpolarised cross section but its interference with $`\mathrm{N}(1535)\mathrm{S}_{11}`$ produces a strong effect in the beam asymmetry. Photoproduction can also provide a very strong selection tool: combining a circularly polarised photon beam and a longitudinally polarised target provides a tool to select states with helicity 1/2 or 3/2 depending on whether the target polarisation is parallel or antiparallel to the photon helicity.
Baryon resonances have large, overlapping widths rendering difficult the study of individual states, in particular of those only weakly excited. This problem can be overcome partly by looking at specific decay channels. The $`\eta `$ meson for example has isospin $`I=0`$ and consequently, the N$`\eta `$ final state can only be reached via formation of N resonances. Then even a small coupling of a resonance to N$`\eta `$ identifies it as N state. A key point in the identification of new baryon resonances is the combined analysis of data on photo- (and pion-) induced reactions with different final states. Resonances must have the same masses, total widths, and gamma-nucleon couplings, in all reactions under study. This imposes strong constraints for the analysis.
In the present paper we report results of a combined analysis of photoproduction experiments with $`\pi \mathrm{N}`$, $`\eta \mathrm{N}`$, $`\mathrm{K}\mathrm{\Lambda }`$, and $`\mathrm{K}\mathrm{\Sigma }`$ final states. This work is a first step of a forthcoming analysis of all reactions with production of baryon resonances in the intermediate state. This paper concentrates on the reactions $`\gamma \mathrm{p}\mathrm{N}\pi `$ and $`\mathrm{N}\eta `$, including available polarisation measurements. Results on photoproduction of open strangeness are presented in a subsequent paper sarantsev .
The outline of the paper is as follows: The fit method is described in section 2, data and fit are compared in section 3. In section 4 we present the main results of this analysis and discuss the statistical significance of new baryon resonances. Interpretations are offered for the newly found resonances. The paper ends with a short summary in section 5.
## 2 Fit method
### 2.1 Analytical properties of the amplitude and resonance–Reggeon duality
The choice of amplitudes used to describe the data is partly driven by experimental observations. In pion photoproduction, angular distributions exhibit strong variations indicating the presence of baryon resonances. On the other hand, all data on single–meson photoproduction have prominent forward peaks in the region above 2000 MeV which can be associated with $`t`$–channel exchange processes. Regge behaviour, extrapolated to the low–energy region, describes the cross section in the resonance region “on average”. This feature is known as Reggeon–resonance duality (see Duality1 and references therein). It gave hope for a self–consistent construction of hadron–hadron interactions in both, the low–energy and the high–energy region. However there is a problem with unitarity: The $`s`$–channel unitarity corrections destroy the one–Reggeon exchange picture, while the $`s`$–channel resonance amplitudes do not satisfy the $`t`$–, $`u`$–channel unitarity Duality2 . So it seems reasonable to extract the resonance structure of the amplitude together with phenomenological reggeized $`t`$– and $`u`$–channel exchange amplitudes.
The scattering amplitude has the following analytical properties. The partial–wave or multipole amplitudes contain singularities when the scattering particles can form a bound state with mass $`M`$. Unstable bound states with a finite width $`\mathrm{\Gamma }`$ have a pole singularity at $`s=M^2i\mathrm{\Gamma }M`$ in the complex plane. At the opening of thresholds, the amplitude acquires a square root singularity (right–hand singularity); $`t`$–exchange leads to left–hand singularities at $`t=\mu ^2`$ (one–particle exchange with mass $`\mu `$), $`t=4\mu ^2`$ (exchange of two of these particles) and so on. In three–body interactions the three–particle rescattering amplitude gives a triangle singularity which may contribute significantly to the cross section under some particular kinematical conditions triangle . Triangle singularities grow logarithmically and are thus weaker than a pole or a threshold singularity. In most cases, triangle singularities can be accounted for by introducing phases to resonance couplings.
In our present analysis, the primary goal is to get information about the leading (pole) singularities of the photoproduction amplitude. For this purpose, a representation of the amplitude as a sum of $`s`$–channel resonances together with some $`t`$– and $`u`$–exchange diagrams is an appropriate representation. Strongly overlapping resonances are parameterised as $`K`$–matrix. In many cases it is sufficient to use a relativistic Breit–Wigner parameterisation, though.
We emphasise that the amplitudes given below satisfy gauge invariance, analyticity and unitarity. However, when $`t`$–, $`u`$–, and $`s`$–channel amplitudes are added, unitarity is violated. In principle, this can be avoided by projecting the $`t`$– and $`u`$–channel amplitudes onto $`s`$–channel amplitudes of defined spins and parities. The projected amplitudes are however small, and the violation of unitarity is mild as long as $`t`$– and $`u`$–channel amplitudes contribute only a small fraction to the total cross section. In this analysis, amplitudes for photoproduction of baryon resonances and their decays are calculated in the framework of relativistic tensor operators. The formalism is fully described in Anisovich:2004zz ; here parameterisations of resonances used under different conditions are given.
### 2.2 Parameterisations of resonances
The differential cross section for production of two or more particles has the form:
$`\mathrm{d}\sigma ={\displaystyle \frac{(2\pi )^4|A|^2}{4\sqrt{(k_1k_2)^2m_1^2m_2^2}}}d\mathrm{\Phi }_n(k_1+k_2,q_1,\mathrm{},q_n)`$ (1)
where $`k_i`$ and $`m_i`$ are the four–momenta and masses of the initial particles (nucleon and $`\gamma `$ in the case of photoproduction) and $`q_i`$ are the four–momenta of final state particles. $`d\mathrm{\Phi }_n(k_1+k_2,q_1,\mathrm{},q_n)`$ is the n–body phase volume
$`d\mathrm{\Phi }_n(k_1+k_2,q_1,\mathrm{},q_n)=`$
$`\delta ^4(k_1+k_2{\displaystyle \underset{i=1}{\overset{n}{}}}q_i){\displaystyle \underset{i=1}{\overset{n}{}}}{\displaystyle \frac{d^3q_i}{(2\pi )^32q_{0i}}}`$ (2)
where $`q_{0i}`$ is time component (energy). The differential cross section for photoproduction of single mesons is given by
$`\mathrm{d}\sigma ={\displaystyle \frac{\sqrt{(s(m_\mu +m_B)^2))(s(m_\mu m_B)^2)}}{16\pi s(sm_N^2)}}|A|^2`$
where $`s=(k_1+k_2)^2=(q_1+q_2)^2`$ is the square of the total energy, $`m_\mu ,`$ ($`\mu =\pi ,\eta ,\mathrm{K}`$), $`m_B`$ ($`B=\mathrm{N},\mathrm{\Lambda },\mathrm{\Sigma }`$) the meson and baryon masses, respectively.
The $`\eta `$ photoproduction cross section is dominated by $`\mathrm{N}(1535)\mathrm{S}_{11}`$. It overlaps with $`\mathrm{N}(1650)\mathrm{S}_{11}`$ and the two $`\mathrm{S}_{11}`$ resonances are described as two–pole, four–channel $`K`$–matrix ($`\pi \mathrm{N}`$, $`\eta \mathrm{N}`$, $`\mathrm{K}\mathrm{\Lambda }`$ and $`\mathrm{K}\mathrm{\Sigma }`$). The photoproduction amplitude can be written in the $`P`$–vector approach since the $`\gamma \mathrm{N}`$ couplings are weak and do not contribute to rescattering. The amplitude is then given by
$`A_a=\widehat{P}_b(\widehat{I}i\widehat{\rho }\widehat{K})_{ba}^1.`$ (4)
The phase space $`\widehat{\rho }`$ is a diagonal matrix with
$`\rho _{ab}=\delta _{ab}\rho _a,a,b=\pi \mathrm{N},\eta \mathrm{N},\mathrm{K}\mathrm{\Lambda },\mathrm{K}\mathrm{\Sigma }.`$ (5)
and
$`\rho _a(s)={\displaystyle \frac{\sqrt{(s(m_\mu +m_B)^2))(s(m_\mu m_B)^2)}}{s}}.`$ (6)
The production vector $`\widehat{P}`$ and the $`K`$–matrix $`\widehat{K}`$ have the following parameterisation:
$`K_{ab}={\displaystyle \underset{\alpha }{}}{\displaystyle \frac{g_a^{(\alpha )}g_b^{(\alpha )}}{M_\alpha ^2s}}+f_{ab},P_b={\displaystyle \underset{\alpha }{}}{\displaystyle \frac{g_{\gamma \mathrm{N}}^{(\alpha )}g_b^{(\alpha )}}{M_\alpha ^2s}}+\stackrel{~}{f}_b`$
where $`M_\alpha `$, $`g_a^{(\alpha )}`$ and $`g_{\gamma \mathrm{N}}^{(\alpha )}`$ are the mass, coupling constant and production constant of the resonance $`\alpha `$; $`f_{ab}`$ and $`\stackrel{~}{f}_b`$ are non–resonant terms.
Other resonances were taken as Breit–Wigner amplitude:
$`A_a={\displaystyle \frac{g_{\gamma N}\stackrel{~}{g}_a(s)}{M^2siM\stackrel{~}{\mathrm{\Gamma }}_{tot}(s)}}`$ (8)
States with masses above 1700 MeV were parameterised with a constant width to fit exactly the pole position. For resonances below 1700 MeV, $`\stackrel{~}{\mathrm{\Gamma }}_{tot}(s)`$ was parameterised by
$`\stackrel{~}{\mathrm{\Gamma }}_{tot}(s)=\mathrm{\Gamma }_{tot}{\displaystyle \frac{\rho _{\pi N}(s)k_{\pi N}^{2L}(s)F^2(L,k_{\pi N}^2(M^2),r)}{\rho _{\pi N}(M^2)k_{\pi N}^{2L}(M^2)F^2(L,k_{\pi N}^2(s),r)}},`$
$`k_a^2(s)={\displaystyle \frac{(s(m_\mu +m_B)^2))(s(m_\mu m_B)^2)}{4s}}.`$ (9)
Here, $`L`$ is the orbital momentum and $`k`$ is the relative momentum for the decay into $`\pi N`$ ($`\mu =\pi `$, $`B=N`$). $`F(L,k^2,r)`$ are Blatt–Weiskopf form factors, taken with a radius $`r=0.8`$ fm. The exact form of these factors can be found e.g. in Anisovich:2004zz . $`g_{\gamma N}`$ is the production coupling and $`\stackrel{~}{g}_a`$ are decay couplings of the resonance into meson nucleon channels. These couplings are suppressed at large energies by a factor
$`\stackrel{~}{g}_a(s)=g_a\sqrt{{\displaystyle \frac{1.5\mathrm{GeV}^2}{1.0\mathrm{GeV}^2+k_a^2}}}.`$ (10)
The factor proved to be useful for two–meson photoproduction. For photoproduction of single mesons, it plays almost no role and is only introduced here for the sake of consistency.
The partial widths are related to the couplings as
$`M\mathrm{\Gamma }_a`$ $`=`$ $`\stackrel{~}{g}_a^2{\displaystyle \frac{\rho _a(M^2)k_{M^2}^{2L}}{F^2(L,k_{M^2}^2,r)}}{\displaystyle \frac{m_B+\sqrt{m_B^2+k_a^2}}{2m_B}}\beta _L,`$
$`\beta _L`$ $`=`$ $`{\displaystyle \frac{1}{L}}{\displaystyle \underset{l=1}{\overset{L}{}}}{\displaystyle \frac{2l1}{l}},J=L{\displaystyle \frac{1}{2}},`$
$`\beta _L`$ $`=`$ $`{\displaystyle \frac{1}{2L+1}}{\displaystyle \underset{l=1}{\overset{L}{}}}{\displaystyle \frac{2l1}{l}},J=L+{\displaystyle \frac{1}{2}}.`$ (11)
Here $`J`$ is the total momentum of the state.
### 2.3 $`𝐭`$– and $`𝐮`$–channel exchange parameterisations
At high energies, there are clear peaks in the forward direction of photoproduced mesons. The forward peaks are connected with meson exchanges in the $`t`$–channel. These contributions are parameterised as $`\pi `$, $`\rho (\omega )`$, $`\mathrm{K}`$, and $`\mathrm{K}^{}`$ exchanges.
These contributions are reggeized by using collins
$`T(s,t)=g_1(t)g_2(t){\displaystyle \frac{1+\xi exp(i\pi \alpha (t))}{\mathrm{sin}(\pi \alpha (t))}}\left({\displaystyle \frac{\nu }{\nu _0}}\right)^{\alpha (t)},`$
$`\nu ={\displaystyle \frac{1}{2}}(su).`$ (12)
Here, $`g_i`$ are vertex functions, $`\alpha (t)`$ is a function describing the trajectory, $`\nu _0`$ is a normalisation factor (which can be taken to be 1) and $`\xi `$ is the signature of the trajectory. Exchanges of $`\pi `$ and $`\mathrm{K}`$ have positive, $`\rho `$, $`\omega `$, and $`\mathrm{K}^{}`$ exchanges have negative signature.
For $`\rho `$($`\omega `$) exchange, $`\alpha (t)=0.50+0.85t`$. The pion trajectory is given by $`\alpha (t)=0.014+0.72t`$, the $`\mathrm{K}^{}`$ and $`\mathrm{K}`$ trajectories are represented by $`\alpha (t)=0.32+0.85t`$ and $`\alpha (t)=0.25+0.85t`$, respectively. The full expression for the $`t`$–channel amplitudes can be found in Anisovich:2004zz .
The $`u`$–channel exchanges were parameterised as nucleon, $`\mathrm{\Lambda }`$, or $`\mathrm{\Sigma }`$ exchanges.
## 3 Fits to the data
In this paper, we report results on baryon resonances and their coupling to $`\mathrm{N}\pi `$ and $`\mathrm{N}\eta `$. The results are based on a coupled–channel analysis of various data sets on photoproduction of different final states. The data comprise CB–ELSA $`\pi ^0`$ and $`\eta `$ photoproduction data Bartholomy:04 ; Crede:04 , the Mainz–TAPS data Krusche:nv on $`\eta `$ photoproduction, beam–asymmetry measurements of $`\pi ^0`$ and $`\eta `$ GRAAL1 ; SAID1 ; GRAAL2 , and data on $`\gamma \mathrm{p}\mathrm{n}\pi ^+`$ SAID2 . The high precision data from GRAAL GRAAL1 do not cover the low mass region; therefore we extract further data from the compilation of the SAID database SAID1 . This data allows us to define the ratio of helicity amplitudes for the $`\mathrm{\Delta }(1232)\mathrm{P}_{33}`$ resonance.
Data on photoproduction of $`\mathrm{K}^+\mathrm{\Lambda }`$, $`\mathrm{K}^+\mathrm{\Sigma }`$, and $`\mathrm{K}^0\mathrm{\Sigma }^+`$ from SAPHIR Glander:2003jw and CLAS McNabb:2003nf , and beam asymmetry data for $`\mathrm{K}^+\mathrm{\Lambda }`$, $`\mathrm{K}^+\mathrm{\Sigma }`$ from LEPS Zegers:2003ux are also included in the coupled–channel analysis. The results on couplings of baryon resonances to $`\mathrm{K}^+\mathrm{\Lambda }`$ and $`\mathrm{K}^+\mathrm{\Sigma }`$ are documented in a separate paper sarantsev .
The fit uses 14 N resonances coupling to N$`\pi `$, N$`\eta `$, $`\mathrm{K}\mathrm{\Lambda }`$, and $`\mathrm{K}\mathrm{\Sigma }`$ and 7 $`\mathrm{\Delta }`$ resonances coupling to N$`\pi `$ and $`\mathrm{K}\mathrm{\Sigma }`$. Most resonances are described by relativistic Breit–Wigner amplitudes. For the two S<sub>11</sub> resonances at 1535 and 1650 MeV, a four–channel $`K`$–matrix ($`\mathrm{N}\pi `$, $`\mathrm{N}\eta `$, $`\mathrm{K}\mathrm{\Lambda }`$, $`\mathrm{K}\mathrm{\Sigma }`$) is used. The background is described by reggeized $`t`$–channel $`\pi `$, $`\rho `$ ($`\omega `$), $`\mathrm{K}`$ and $`\mathrm{K}^{}`$ exchanges and by baryon exchanges in the $`s`$– and $`u`$–channels.
The $`\chi ^2`$ values for the final solution of the partial–wave analysis are given in Table 1. Weights are given to the different data sets included in this analysis with which they enter the fits. In the choice of weights, some judgement is needed. The CB-ELSA data on pion and $`\eta `$ photoproduction are the main source of the analysis and thus have large weights. The beam polarisation measurements for open strangeness production are also emphasized as discussed in sarantsev . Fits were performed with a variety of different weights; accepted solutions resulted not only in a good overall $`\chi ^2`$; emphasis was laid on having a good fit of all data sets. Changing the weights may result in pictures showing larger discrepancies; the changes of pole positions are only small.
The fit minimises a pseudo–chisquare function which we call $`\chi _{\mathrm{tot}}^2`$. It is given by
$`\chi _{\mathrm{tot}}^2={\displaystyle \frac{w_i\chi _i^2}{w_iN_i}}{\displaystyle N_i}`$ (13)
where the $`N_i`$ are given as $`N_{\mathrm{data}}`$ (per channel) in the $`2^{\mathrm{nd}}`$ column of Table 1 and the weights in the last column.
### 3.1 Fit to the $`𝒑𝝅^\mathrm{𝟎}`$ data
The differential cross sections for the CB–ELSA $`\gamma pp\pi ^0`$ data are shown in Fig. 1. The main fit is represented as solid line. The figure also shows the most important individual contributions. The contribution of $`\mathrm{\Delta }(1232)`$ (given as dashed line, on the left panel) dominates the low–energy region, for small photon energies it even exceeds the experimental cross section, thus underlining the importance of interference effects. Non–resonant background amplitudes, given by a pole at $`s1`$ GeV<sup>2</sup> and by a $`u`$–channel exchange diagram, are needed to describe the shape of the $`\mathrm{\Delta }(1232)`$. The pole at negative $`s`$ represents the left–hand cuts.
The two $`\mathrm{S}_{11}`$ resonances at 1535 and at 1650 MeV are described as $`K`$–matrix. Their sum is depicted as dotted line. The $`\mathrm{S}_{11}`$ contribution is flat in $`\mathrm{cos}\mathrm{\Theta }_{\mathrm{cm}}`$. The contribution of the $`\mathrm{D}_{13}(1520)`$ shown as dash–dotted line in Fig. 1 (left panel). It is strong in the $`14001600`$ MeV mass region. At higher energies (Fig. 1, right panel) the most significant contributions come from $`\mathrm{\Delta }(1700)\mathrm{D}_{33}`$ (dashed line) and from $`\mathrm{N}(1680)\mathrm{F}_{15}`$ (dotted line). For invariant $`\mathrm{p}\gamma `$ masses above 1800 MeV, the most forward point in Fig. 1 is not reproduced by the fit. If this point is given a very small error (to ensure that the fit describes these points), the overall agreement between data and fit becomes somewhat worse; resonance masses and widths change by a few MeV, at most.
Recent data from GRAAL GRAAL1 on the differential cross section for $`\gamma \mathrm{p}\mathrm{p}\pi ^0`$ and on the photon beam asymmetry $`\mathrm{\Sigma }`$ are compared to our fit in Figs. 2 and 4; older beam asymmetry data are shown in Fig. 4.
### 3.2 Fit to $`𝐧𝝅^\mathbf{+}`$ photoproduction data
It is important to include data on $`\mathrm{n}\pi ^+`$ photoproduction since the combination of the $`\mathrm{n}\pi ^+`$ and $`\mathrm{p}\pi ^\mathrm{o}`$ channels defines the isospin of $`s`$-channel baryons. Without this information, pairs of resonances like $`\mathrm{N}(1700)\mathrm{D}_{13}`$ and $`\mathrm{\Delta }(1700)\mathrm{D}_{33}`$ cannot be separated. A fit with both having large destructively interfering amplitudes may give a good $`\chi ^2`$ even though the fit is physically meaningless. For $`\gamma \mathrm{p}\mathrm{N}^{}\mathrm{n}\pi ^+`$ the isotopic coefficient is equal to $`\sqrt{2/3}`$, for $`\gamma \mathrm{p}\mathrm{N}^{}\mathrm{p}\pi ^0`$ it is equal to $`\sqrt{1/3}`$. In case of $`\mathrm{\Delta }`$ photoproduction, the respective isotopic coefficients are $`\sqrt{1/3}`$ for $`n\pi ^+`$ and $`\sqrt{2/3}`$ for $`\mathrm{p}\pi ^0`$.
Differential cross sections for $`\gamma \mathrm{p}\mathrm{n}\pi ^+`$ SAID2 and PWA result are compared in Fig. 5. In addition to resonances, a significant contribution stems from $`t`$–channel $`\pi `$ and $`\rho `$ exchanges (about 10% and 30%, respectively). This reaction has a large number of data points with small statistical errors but the largest ambiguities in its interpretation. Hence, a small weight is given to this channel to avoid that it has a significant impact on baryon masses, widths, or coupling constants. It was only used to stabilise the fits in case of isospin ambiguities.
### 3.3 Fit to the $`𝐩𝜼`$ channel
Differential cross section for $`\gamma \mathrm{p}\mathrm{p}\eta `$ in the threshold region were measured by the TAPS collaboration at MAINZ Krusche:nv . Data and fit are shown in Fig. 6. In the threshold region, the dominant contribution comes from the $`\mathrm{N}(1535)`$ $`\mathrm{S}_{11}`$ resonance which gives a flat angular distribution. This resonance strongly overlaps with $`\mathrm{N}(1650)\mathrm{S}_{11}`$, and a two–pole K-matrix parameterisation is used in the fit.
The CB–ELSA differential cross section Crede:04 is given in Fig. 7 and compared to the PWA results. The contribution of the two $`\mathrm{S}_{11}`$ resonances (dashed line, below 2 GeV) dominates the $`\eta `$ production region up to 1650 MeV. The most significant further contributions stem from production of $`\mathrm{N}(1720)\mathrm{P}_{13}`$ (dotted line, below 2 GeV), of $`\mathrm{N}(2070)\mathrm{D}_{15}`$ (dashed line, above 2 GeV) and $`\rho `$ ($`\omega `$) exchanges (dotted line, above 2 GeV).
Data on the photon beam asymmetry $`\mathrm{\Sigma }`$ for $`\gamma \mathrm{p}\mathrm{p}\eta `$, measured by GRAAL GRAAL1 are shown in Fig. 8. This data provides essential information on baryon resonances even if their (p$`\gamma `$)– and/or (p$`\eta `$)–couplings are weak. In addition, the beam asymmetry data are necessary to determine the ratio of helicity amplitudes.
## 4 Results
### 4.1 Total cross sections
From the differential cross sections presented in Figs. 1 and 7, absolute cross sections were determined by integration. The integration is performed by summation of the differential cross sections (dots with error bars) and using extrapolated values for bins with no data, and by integration of the fit curve.
In the total cross section for $`\pi ^0`$ photoproduction in Fig. 9, clear peaks are observed for the first, second, and third resonance region. With some good will, the fourth resonance region can be identified as broad enhancement at about 1900 MeV. The decomposition of the peaks into partial waves and their physical significance will be discussed below.
The $`\eta `$ photoproduction cross section (Fig. 10) shows the known strong peak at threshold due to the $`S_{11}(1535)`$. The cross section exhibits indications for one further resonance below 1800 MeV.
### 4.2 The best solution
The masses and widths of the observed states are presented in Table 2. Additionally, ratios of helicity amplitudes $`A_{1/2}/A_{3/2}`$ and fractional contributions normalised to the total cross section for the CB–ELSA $`\pi ^0`$– and $`\eta `$– photoproduction data are included.
A large number of fits (explorative fits plus more than 1000 documented fits) were performed to validate the solution. In these fits the number of resonances, their spin and parity, their parameterisation, and the relative weight of the different data sets were changed.
The errors are estimated from a sequence of fits in which one variable, e.g. a width of one resonance was changed to series of fixed values. All other variables were allowed to adjust freely; the $`\chi ^2`$ changes were monitored as a function of this variable. The errors given in Table 2 correspond to $`\chi ^2`$ changes of 9, hence to three standard deviations. However, the $`3\sigma `$ interval corresponds better to the systematic changes observed when changing the fit hypothesis.
The resonance properties are compared to PDG values PDG . Most resonance parameters converge in the fits to values compatible with previous findings within a 2$`\sigma `$ intervall of the combined error. The helicity ratios sometimes seem to be inconsistent, however they have large errors and the discrepancies are not really significant.
Three new resonances are necessary to describe the data, $`\mathrm{N}(1875)\mathrm{D}_{13}`$, $`\mathrm{N}(2070)\mathrm{D}_{15}`$ and $`\mathrm{N}(2200)`$ with uncertain spin and parity. The best fit is achieved for $`\mathrm{P}_{13}`$ quantum numbers. Two further resonances, $`\mathrm{N}(1840)\mathrm{P}_{11}`$ and $`\mathrm{N}(2170)\mathrm{D}_{13}`$, have masses which are not consistent with established resonances listed by the PDG. We list them also as new particles. Two resonances, $`\mathrm{N}(2000)\mathrm{F}_{15}`$ and $`\mathrm{\Delta }(1940)\mathrm{D}_{33}`$, are observed for the first time in photoproduction. PDG mass values for $`\mathrm{N}(2000)\mathrm{F}_{15}`$ range from 1882 to 2175 MeV. We find a mass of $`(1850\pm 25)`$ MeV. Our mass for $`\mathrm{\Delta }(1940)\mathrm{D}_{33}`$ is fully compatible with PDG. The $`\mathrm{\Delta }(1940)\mathrm{D}_{33}`$ contributes only at a marginal level. The $`\chi _{\mathrm{tot}}^2`$ changes by 143 units when this resonance is omitted. The $`\mathrm{\Delta }(1950)\mathrm{F}_{37}`$ is observed here at $`1893\pm 15`$ MeV instead of (PDG) $`1950\pm 10`$ MeV.
In this paper we concentrate on the $`\mathrm{N}(2070)\mathrm{D}_{15}`$ and $`\mathrm{N}(2200)`$. The $`\mathrm{N}(1840)\mathrm{P}_{11}`$, $`\mathrm{N}(1875)\mathrm{D}_{13}`$, and $`\mathrm{N}(2170)`$ $`\mathrm{D}_{13}`$ do not significantly contribute to $`\gamma \mathrm{p}\mathrm{p}\pi ^0`$, $`p\eta `$ and have large couplings to $`\mathrm{K}\mathrm{\Lambda }`$ and/or $`\mathrm{K}\mathrm{\Sigma }`$. They will be discussed in sarantsev .
Finally a comment is made on resonances with known photo-couplings but not seen in this analysis. $`\mathrm{N}(1990)\mathrm{F}_{17}`$, $`\mathrm{\Delta }(1600)\mathrm{P}_{33}`$, $`\mathrm{\Delta }(1910)\mathrm{P}_{33}`$, $`\mathrm{\Delta }(1930)\mathrm{D}_{35}`$, $`\mathrm{\Delta }(2420)\mathrm{H}_{\mathrm{3\hspace{0.17em}11}}`$, and $`\mathrm{N}(2190)\mathrm{G}_{17}`$ are not observed here. The latter resonance may however be misinterpreted as $`\mathrm{N}(2200)\mathrm{P}_{13}`$ (see Table 3). The photocouplings of most of these resonances are seen with weak evidence (one–star rating); only $`\mathrm{\Delta }(1600)`$ $`P_{33}`$ has a three–star photo–coupling, and the $`\mathrm{\Delta }(1930)\mathrm{D}_{35}`$ photocoupling has 2 stars. We have no explanation why these states are missing in this analysis. The $`\mathrm{\Delta }(1900)\mathrm{S}_{31}`$, $`\mathrm{\Delta }(1940)\mathrm{D}_{33}`$, and $`\mathrm{\Delta }(1930)\mathrm{D}_{35}`$ may form a spin triplet with intrinsic orbital angular momentum $`L=1`$ and total spin $`S=3/2`$ coupling to $`J=1/2,3/2`$, and $`5/2`$ as suggested in Klempt:2002cu . Two of these states are not observed in this analysis. Quark models do not reproduce these states predicting them to have masses above 2.1 GeV. Hence, the question remains open if these states exist at such a low mass.
### 4.3 Significance of resonance contributions
A systematic study of the significance of new resonances was carried out. For new resonances the quantum numbers were changed to any $`J^P`$ value with $`J9/2`$. In the new fits, all variables were left free for variations including masses, widths, and couplings of all resonances. The result of this study is summarised in Table 3. The Table illustrates the global deterioration of the fit and the $`\chi ^2`$ changes for the individual channels. Negative $`\chi ^2`$ changes indicate that the best quantum numbers are enforced by other data.
The $`\mathrm{N}(2070)\mathrm{D}_{15}`$ is the most significant new resonance. Omitting it changes $`\chi _{\mathrm{tot}}^2`$ by 1589, by 199 for the data on $`\eta `$ photoproduction and by 940 for the data on $`\pi ^0`$ photoproduction. Replacing the $`J^P`$ assignment from $`5/2^{}`$ to $`1/2^\pm `$, …, $`9/2^\pm `$, the $`\chi _{\mathrm{tot}}^2`$ deteriorates by more than 750. The deterioration of the fits is visible in the comparison of data and fit. One of the closest description for $`\eta `$ photoproduction was obtained fitting with a $`7/2^{}`$ state. In this case, Figs. 12 a,b show the fits of the differential cross section in the region of resonance mass and description of the beam asymmetry for highest energy bin. The shape of the differential cross section at small angles is close in both cases however the $`7/2^{}`$ state failed to describe the very forward two points. The beam asymmetry also clearly favours the $`5/2^{}`$ state. The $`\pi ^0`$ photoproduction cross sections measured by CB–ELSA are visually not too sensitive to $`5/2^{}`$ and $`7/2^{}`$ quantum numbers (see Fig. 12 c) but there is a clear difference between the two descriptions in the very backward region. The latest GRAAL results on the $`\mathrm{p}\pi ^0`$ differential cross section which were obtained after discovery of the $`\mathrm{N}(2070)\mathrm{D}_{15}`$ Bartholomy:04 confirmed $`5/2^{}`$ as favoured quantum numbers (see Fig. 12 d).
The mass scan of the $`\mathrm{D}_{15}(2070)`$ resonance ($`\chi ^2`$ as a function of the assumed $`\mathrm{D}_{15}`$ mass) is shown in Fig. 12. In the scan, the mass of the $`\mathrm{D}_{15}`$ was fixed at a number of values covering the region of interest while all other fit parameters were allowed to adjust newly. The sum of $`\chi ^2`$ for $`\pi ^0`$ photoproduction data (CB-ELSA, GRAAL 05) does not show any minimum in this region; the destributions are very flat. Fig. 12a shows separately the sum of $`\chi ^2`$ contributions from the CB–ELSA differential cross section plus the GRAAL 04 polarisation data, and the sum of the $`\chi ^2`$ for all $`\mathrm{\Lambda }\mathrm{K}^+`$ and all $`\mathrm{\Sigma }\mathrm{K}`$ reactions. A clear minimum is seen in all three data sets. The sum of $`\chi ^2`$ for all these reactions is given in Fig. 12b. The shaded area corresponds to the mass range assigned to this resonance, $`(2060\pm 30)`$ MeV. We conclude that the $`\mathrm{D}_{15}(2070)`$ is identified in its decays into $`\mathrm{N}\eta `$, $`\mathrm{\Lambda }\mathrm{K}^+`$ and $`\mathrm{\Sigma }\mathrm{K}`$. Its coupling to $`\mathrm{N}\pi `$ is weak, hence it is not surprising that it was not observed in pion induced reactions.
The $`\mathrm{N}(2200)`$ resonance is less significant. Omitting $`\mathrm{N}(2200)`$ from the analysis, changes $`\chi ^2`$ for the CB-ELSA data on $`\eta `$ photoproduction by 56, and by 20 for the $`\pi ^0`$–photoproduction data. Other quantum numbers than the preferred $`\mathrm{P}_{13}`$ lead to marginally larger $`\chi ^2`$ values. The mass scan for this state is shown in Fig. 13. The photoproduction data on $`\mathrm{d}\sigma /\mathrm{d}\mathrm{\Omega }`$ from CB-ELSA does not show any minimum, $`\eta `$ photoproduction data exhibit a shallow minimum slightly above 2200 MeV. The sum of all $`\mathrm{\Lambda }\mathrm{K}^+`$ and $`\mathrm{K}\mathrm{\Sigma }`$ reactions also have a minimum in this mass region. The sum of $`\chi ^2`$ for all these reactions is shown in Fig. 13 b and from this distribution the resonance mass can be well defined.
### 4.4 The four resonance regions
The first resonance region dominates pion photoproduction and is due to the excitation of the $`\mathrm{\Delta }(1232)\mathrm{P}_{33}`$. Its fractional contribution to $`\gamma \mathrm{p}\mathrm{p}\pi ^0`$ (Table 2) exceeds 1. There is strong destructive interference between $`\mathrm{\Delta }(1232)`$ $`\mathrm{P}_{33}`$, the $`\mathrm{P}_{33}`$ nonresonant amplitude and $`u`$–channel exchange. In the fit without latest GRAAL data on the cross section and beam asymmetry GRAAL1 the $`A_{1/2}/A_{3/2}`$ helicity ratio of excitation of the $`\mathrm{\Delta }(1232)\mathrm{P}_{33}`$ was found to be $`0.52\pm 0.06`$ which agrees favorably with the PDG average $`0.53\pm 0.04`$. With the new GRAAL05 data included, this value shifted to $`0.44\pm 0.06`$. The $`\mathrm{N}(1440)\mathrm{P}_{11}`$ Roper resonance provides a small contribution of about 1–3% compared to the $`\mathrm{\Delta }(1232)\mathrm{P}_{33}`$.
In the $`\mathrm{p}\pi ^0`$ final state $`\mathrm{N}(1520)\mathrm{D}_{13}`$ and the two $`\mathrm{S}_{11}`$ resonances yield contributions of similar strengths to the second resonance region. This is consistent with the known photocouplings and $`\mathrm{p}\pi `$ branching fractions of the three resonances.
The third bump in the $`\mathrm{p}\pi ^0`$ total cross section is due to three major contributions: the $`\mathrm{\Delta }(1700)\mathrm{D}_{33}`$ resonance provides the largest fraction ($`35`$%) of the peak, followed by $`\mathrm{N}(1680)\mathrm{F}_{15}`$ ($`25`$%) and $`\mathrm{N}(1650)\mathrm{S}_{11}`$ ($`20`$%) as extracted from the $`K`$–matrix parameterisation; observed as well are the $`\mathrm{\Delta }(1620)\mathrm{S}_{31}`$ ($`7`$%) and $`\mathrm{N}(1720)\mathrm{P}_{13}`$ ($`6`$%) resonances. The latter contributes to $`\mathrm{p}\eta `$ with a surprisingly large fraction; about 90% of the resonant intensity in this mass region is assigned to $`\mathrm{N}(1720)\mathrm{P}_{13}\mathrm{p}\eta `$ decays.
In the fourth resonance region we identify $`\mathrm{\Delta }(1950)\mathrm{F}_{37}`$ contributing $`41`$% to the enhancement and $`\mathrm{\Delta }(1920)\mathrm{P}_{33}`$ with $`35`$%. Additionally, the fit requires the presence of $`\mathrm{\Delta }(1905)\mathrm{F}_{35}`$ and $`\mathrm{\Delta }(1940)\mathrm{D}_{33}`$. The high–energy region is dominated by $`\rho `$($`\omega `$) exchange in the $`t`$ channel as can be seen by the forward peaking in the differential cross sections.
### 4.5 Discussion
Four new resonances are found in this analysis. The question arises of course why these resonances have not been found before. $`\mathrm{N}(2070)\mathrm{D}_{15}`$ has a large coupling to $`\mathrm{N}\eta `$ and may therefore have escaped discovery. The $`\mathrm{N}(1875)\mathrm{D}_{13}`$ and $`\mathrm{N}(2170)\mathrm{D}_{13}`$ states couple strongly to the $`\mathrm{K}\mathrm{\Lambda }`$ and $`\mathrm{K}\mathrm{\Sigma }`$ channels; the existence of the first state has already been suggested in Mart:1999ed from an analysis of older SAPHIR data on $`\gamma \mathrm{p}\mathrm{K}\mathrm{\Lambda }`$ Tran:1998qw . Cutkosky Cutkosky:1980rh reported two $`\mathrm{N}\mathrm{D}_{13}`$ resonances at $`(1880\pm 100)`$ and $`(2081\pm 80)`$ MeV with respective widths of $`(180\pm 60)`$ and $`(300\pm 100)`$ MeV. The $`\mathrm{N}(1840)\mathrm{P}_{11}`$ appears in all channels. The evidence for it is discussed in sarantsev . The $`\mathrm{N}(2200)`$ does not have such characteristic features. It improves the description of the data in a difficult mass range and further data will be required to establish or to disprove its existence. Its preferred quantum numbers are $`\mathrm{P}_{13}`$ but it seems not unlikely that $`\mathrm{N}(2200)`$ should be identified with $`\mathrm{N}(2190)\mathrm{G}_{17}`$ (which gives the second best PWA solution).
The three largest contributions to the $`\eta `$ photoproduction cross section stem from $`\mathrm{N}(1535)\mathrm{S}_{11}`$, $`\mathrm{N}(1720)\mathrm{P}_{13}`$, and
$`\mathrm{N}(2070)\mathrm{D}_{15}`$. We tentatively assign ($`J=1/2;L=1,S=1/2`$) quantum numbers to the first state; $`\mathrm{N}(1720)\mathrm{P}_{13}`$ and $`\mathrm{N}(1680)\mathrm{F}_{15}`$ form a spin doublet, hence the dominant quantum numbers of $`\mathrm{N}(1720)\mathrm{P}_{13}`$ must be ($`J=3/2;L=2,S=1/2`$). Thus it is tempting to assign ($`J=5/2;L=3,S=1/2`$) quantum numbers to $`\mathrm{N}(2070)\mathrm{D}_{15}`$. The three baryon resonances with strong contributions to the $`\mathrm{p}\eta `$ channel thus all have spin $`S=1/2`$ and orbital and spin angular momenta adding antiparallelly with $`J=L1/2`$. Fig. 14 depicts this scenario.
The large $`\mathrm{N}(1535)\mathrm{S}_{11}\mathrm{N}\eta `$ coupling has been a topic of a controversial discussion. In the quark model, this coupling arises naturally from a mixing of the two ($`J=1/2;L=1,S=1/2`$) and ($`J=1/2;L=1,S=3/2`$) harmonic-oscillator states Isgur:1978xj . However, $`\mathrm{N}(1535)\mathrm{S}_{11}`$ is very close to the $`\mathrm{K}\mathrm{\Lambda }`$ and $`\mathrm{K}\mathrm{\Sigma }`$ thresholds and the resonance can be understood as originating from coupled–channel meson–baryon chiral dynamics Kaiser:1995cy . Alternatively, the strong $`\mathrm{N}(1535)\mathrm{S}_{11}\mathrm{N}\eta `$ coupling can be explained as delicate interplay between confining and fine structure interactions Glozman:1995tb .
A consistent picture of the large $`\mathrm{N}(1535)\mathrm{S}_{11}\mathrm{N}\eta `$ coupling should explain the systematics of $`\mathrm{N}\eta `$ couplings. We note a kinematical similarity: The three resonances with large $`\mathrm{N}\eta `$ partial decay widths are those for which the dominant intrinsic orbital excitation $`L=1,2,3`$ and the decay orbital angular momenta $`\mathrm{}=0,1,2`$ are related by $`J=L1/2=\mathrm{}+1/2`$. The intrinsic quark spin configuration remains in a spin doublet.
## 5 Summary
We have presented a partial wave analysis of data on photoproduction of $`\pi \mathrm{N}`$, $`\eta \mathrm{N}`$, $`\mathrm{K}\mathrm{\Lambda }`$, and $`\mathrm{K}\mathrm{\Sigma }`$ final states. The data include total cross sections and angular distributions, beam asymmetry measurements as well as the recoil polarisation in case of hyperon production. A reasonable description of all data was achieved by introducing 14 $`\mathrm{N}^{}`$ and seven $`\mathrm{\Delta }^{}`$ resonances.
Most baryon resonances are found with masses, widths and ratios of helicity amplitudes which are fully compatible with previous findings. New resonances are required to fit the data, $`\mathrm{N}(1840)\mathrm{P}_{11}`$, $`\mathrm{N}(1875)\mathrm{D}_{13}`$, $`\mathrm{N}(2070)\mathrm{D}_{15}`$, $`\mathrm{N}(2170)\mathrm{D}_{13}`$, and $`\mathrm{N}(2200)`$. The $`\mathrm{N}(1840)\mathrm{P}_{11}`$ resonance could, however, be identical with $`\mathrm{N}(1710)\mathrm{P}_{11}`$ and $`\mathrm{N}(2170)`$ $`\mathrm{D}_{13}`$ with $`\mathrm{N}(2080)\mathrm{D}_{13}`$.
Three resonances are found to have very large couplings to $`\mathrm{N}\eta `$, $`\mathrm{N}(1535)\mathrm{S}_{11}`$, $`\mathrm{N}(1720)\mathrm{P}_{13}`$, and $`\mathrm{N}(2070)\mathrm{D}_{15}`$. The dynamical origin of this preference remains to be investigated.
### Acknowledgements
We would like to thank the CB-ELSA/TAPS Collaboration for numerous discussions on topics related to this work. We acknowledge financial support from the Deutsche Forschungsgemeinschaft within the SFB/TR16. The collaboration with St.Petersburg received funds from the DFG and the Russian Foundation for Basic Research. U. Thoma thanks for an Emmy Noether grant from the DFG. A.V. Anisovich and A.V. Sarantsev acknowledge support from the Alexander von Humboldt Foundation.
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# On Defect-Mediated Transitions in Bosonic Planar Lattices
## Abstract
We discuss the finite-temperature properties of Bose-Einstein condensates loaded on a $`2D`$ optical lattice. In an experimentally attainable range of parameters the system is described by the $`XY`$ model, which undergoes a Berezinskii-Kosterlitz-Thouless (BKT) transition driven by the vortex pair unbinding. The interference pattern of the expanding condensates provides the experimental signature of the BKT transition: near the critical temperature, the $`\stackrel{}{k}=0`$ component of the momentum distribution sharply decreases.
Bose-Einstein condensates (BECs) are nowadays routinely stored in $`1D`$ , $`2D`$ and $`3D`$ optical lattices. In this paper we review and further discuss the finite-temperature properties of Bose-Einstein condensates loaded on a $`2D`$ optical lattice : there we showed that at a critical temperature lower than the temperature $`T_{BEC}`$ at which condensation in each well occurs, $`2D`$ lattices of BECs may undergo a phase transition to a superfluid regime where the phases of the single-well condensates are coherently aligned allowing for the observation of a Berezinskii-Kosterlitz-Thouless (BKT) transition.
The BKT transition is the paradigmatic example of a defect-mediated transition which is exhibited by the two-dimensional $`XY`$ model , describing $`2`$-components spins on a two-dimensional lattice. In the low-temperature phase, characterized by the presence of bound vortex-antivortex pairs, the spatial correlations exhibit a power-law decay; above a critical temperature $`T_{BKT}`$, the decay is exponential and there is a proliferation of free vortices. The BKT transition has been observed in superconducting Josephson arrays and its predictions are well verified in the measurements of the superfluid density in $`{}_{}{}^{4}He`$ films . In finite magnetic systems with planar symmetry , the BKT transition point is signaled by the drop of a suitably defined magnetization . We shall discuss in the sequel the analogous of this magnetization for atomic systems.
For $`2D`$ optical lattices, when the polarization vectors of the two standing wave laser fields are orthogonal, the resulting periodic potential for the atoms is $`V_{opt}=V_0[\mathrm{sin}^2(kx)+\mathrm{sin}^2(ky)]`$ where $`k=2\pi /\lambda `$ is the wavevector of the laser beams. The potential maximum of a single standing wave $`V_0=sE_R`$ may be controlled by changing the intensity of the optical lattice beams, and it is conveniently measured in units of the recoil energy $`E_R=\mathrm{}^2k^2/2m`$ ($`m`$ is the atomic mass), while, typically, $`s`$ can vary from $`0`$ up to $`30`$. Gravity is assumed to act along the $`z`$-axis. Usually, superimposed to the optical potential there is an harmonic magnetic potential $`V_m=(m/2)[\omega _r^2(x^2+y^2)+\omega _z^2z^2]`$, where $`\omega _z`$ ($`\omega _r`$) is the axial (radial) trapping frequency. The minima of the lattice are located at the points $`\stackrel{}{j}=(j_1,j_2)\frac{\lambda }{2}`$ with $`j_1`$ and $`j_2`$ integers, and the potential around the minima is $`V_{opt}\frac{m}{2}\stackrel{~}{\omega }_r^2[(x\lambda j_1/2)^2+(y\lambda j_2/2)^2]`$ with $`\stackrel{~}{\omega }_r=\sqrt{2V_0k^2/m}`$. When $`\stackrel{~}{\omega }_r\omega _r,\omega _z`$, the system realizes a square array of tubes, i.e. an array of harmonic traps elongated along the $`z`$-axis . In the following we analyze only the situation in which the axial degrees of freedom are frozen out, which is realizable if $`\omega _z`$ is sufficiently large. We also assume that the harmonic oscillator length of the magnetic potential $`\sqrt{\mathrm{}/m\omega _r}`$ is enough larger than the size $`L`$ (in lattice units) of the optical lattice in order to reduce density inhomogeneity effects and to allow for a safe control of finite size effects.
When all the relevant energy scales are small compared to the excitation energies one can expand the field operator as $`\widehat{\mathrm{\Psi }}(\stackrel{}{r},t)=_\stackrel{}{j}\widehat{\psi }_\stackrel{}{j}(t)\mathrm{\Phi }_\stackrel{}{j}(\stackrel{}{r})`$ with $`\mathrm{\Phi }_\stackrel{}{j}(\stackrel{}{r})`$ the Wannier wavefunction localized in the $`\stackrel{}{j}`$-th well (normalized to $`1`$) and $`\widehat{N}_\stackrel{}{j}=\widehat{\psi }_\stackrel{}{j}^{}\widehat{\psi }_\stackrel{}{j}`$ the bosonic number operator. Substituting the expansion in the full quantum Hamiltonian describing the bosonic system, leads to the Bose-Hubbard model (BHM)
$$\widehat{H}=K\underset{<\stackrel{}{i},\stackrel{}{j}>}{}(\widehat{\psi }_\stackrel{}{i}^{}\widehat{\psi }_\stackrel{}{j}+h.c.)+\underset{\stackrel{}{j}}{}\left[\frac{U}{2}\widehat{N}_\stackrel{}{j}(\widehat{N}_\stackrel{}{j}1)\right]$$
(1)
where $`_{<\stackrel{}{i},\stackrel{}{j}>}`$ denotes a sum over nearest neighbours, $`U=(4\pi \mathrm{}^2a/m)𝑑\stackrel{}{r}\mathrm{\Phi }_\stackrel{}{j}^4`$ ($`a`$ is the $`s`$-wave scattering length) and $`K𝑑\stackrel{}{r}\left[\frac{\mathrm{}^2}{2m}\stackrel{}{}\mathrm{\Phi }_\stackrel{}{i}\stackrel{}{}\mathrm{\Phi }_\stackrel{}{j}+\mathrm{\Phi }_\stackrel{}{i}V_{ext}\mathrm{\Phi }_\stackrel{}{j}\right]`$.
Upon defining $`J2KN_0`$ (where $`N_0`$ is the average number of atoms per site), when $`N_01`$ and $`J/N_0^2U`$, the BHM reduces to $`\widehat{H}=H_{XY}\frac{U}{2}_\stackrel{}{j}\frac{^2}{\theta _\stackrel{}{j}^2}`$, which describes the so-called quantum phase model (see e.g. ): $`\theta _\stackrel{}{j}`$ is the phase of the $`j`$-th condensate and
$$H_{XY}=J\underset{<\stackrel{}{i},\stackrel{}{j}>}{}\mathrm{cos}(\theta _\stackrel{}{i}\theta _\stackrel{}{j})$$
(2)
stands for the Hamiltonian of the classical $`XY`$ model. When $`JU`$ and at temperatures $`TU/k_B`$, the quantum phase Hamiltonian may be well approximated by Eq.(2); therefore, the pertinent partition function describing the thermodynamic behaviour of the BECs stored in an optical lattice can be computed in these limits with the classical $`XY`$ model.
We remark that the BHM - for $`U=0`$ \- describes harmonic oscillators and thus cannot sustain any BKT transition. However, in the limits above specified ($`N_01`$ and $`J/N_0^2UJ`$) the Bose-Hubbard Hamiltonian reduces in the $`XY`$ model (2), which displays the BKT transition. We observe that in the chosen range of parameters, the condition $`UN_0^2/J=UN_0/2K1`$, satisfies the finite-temperature stability criterion recently derived in . In superconducting networks $`N_0`$ is the average number of excess Cooper pairs and thus, the requirements $`N_01`$ and $`J/N_0^2U`$ are always satisfied; at variance, in bosonic lattices $`N_0`$ varies usually between $`1`$ and $`10^3`$ and the validity of the mapping in the quantum phase Hamiltonian is not always guaranteed.
A simple estimate of the coefficients $`J`$ and $`U`$ may be obtained by approximating the Wannier functions with gaussians, whose widths are determined variationally . For a $`2D`$ lattice with $`V_0`$ between $`20`$ and $`25E_R`$ (and $`N_0170`$ as in ), the conditions $`JUJ/N_0^2`$ are rather well satisfied and the BKT critical temperature, $`T_{BKT}J/k_B`$, is between $`10`$ and $`30nK`$.
We observe that, using a coarse-graining approach to determine the finite-temperature phase-boundary line of the Bose-Hubbard model (1) , one gets - in $`2D`$ for $`N_01`$ and $`J/U1`$ \- a critical temperature $`T_{BKT}J/k_B`$. Indeed, within a coarse-graining approach, the phase-boundary line of the Bose-Hubbard model at a temperature $`T`$ is given by
$$1=Kd_0^\beta 𝑑\tau G(\tau )$$
(3)
where $`d`$ is the lattice dimension and $`\beta =1/k_BT`$. In Eq.(3) $`G(\tau )`$ is the correlation function, which, for $`N_01`$, is given by
$$G(\tau )N_0\frac{\vartheta _3(\pi (N_0+\tau /\beta ),q)}{\vartheta _3(\pi N_0,q)}e^{\tau U(1\tau /\beta )/2}$$
(4)
with $`\vartheta _3(z,q)=1+2_{n=1}^{\mathrm{}}\mathrm{cos}(2nz)q^{n^2}`$. Since $`d=2`$, and $`J=2KN_0`$, from Eq.(3) it follows $`1=JG(\omega =0)/N_0`$ where $`G(\omega =0)=_0^\beta 𝑑\tau G(\tau )`$; one has then
$$\frac{G(\omega =0)}{N_0}=\frac{_m\mathrm{exp}(\beta U\varphi _n^2/2)\frac{1\mathrm{exp}(\beta U(\varphi _n+1/2)}{U(\varphi _m+1/2)}}{_m\mathrm{exp}(\beta U\varphi _n^2/2)}:$$
(5)
where $`\varphi _m=N_0m`$. Eq.(5) yields, for $`Uk_BT`$, $`G(\omega =0)/N_0\beta `$, thus implying $`\beta J1`$, i.e. $`k_BT_{BKT}J`$, in fair agreement with $`XY`$ estimates.
We notice that the $`XY`$ Hamiltonian, in the limits $`N_01`$ and $`JUJ/N_0^2`$, can be retrieved extending the Gross-Pitaevskii (GP) Hamiltonian to finite temperature $`TU/k_B`$. In fact, replacing the tight-binding ansatz $`\mathrm{\Psi }(\stackrel{}{r},t)=_\stackrel{}{j}\psi _\stackrel{}{j}(t)\mathrm{\Phi }_\stackrel{}{j}(\stackrel{}{r})`$ (where $`\psi _\stackrel{}{j}`$ are classical fields and not operators) in the GP Hamiltonian one gets the lattice Hamiltonian $`H=K_{<\stackrel{}{i},\stackrel{}{j}>}(\psi _\stackrel{}{i}^{}\psi _\stackrel{}{j}+c.c.)+(U/2)_\stackrel{}{j}N_\stackrel{}{j}(N_\stackrel{}{j}1)`$, which is the classical version of the Bose-Hubbard model ($`N_\stackrel{}{j}\psi _\stackrel{}{j}^2`$). Writing $`N_\stackrel{}{j}=N_0+\delta N_\stackrel{}{j}`$, one may neglect the quadratic terms \[i.e. $`(\delta N_\stackrel{}{j})^2`$\] in the hopping part of the Hamiltonian $`K_{<\stackrel{}{i},\stackrel{}{j}>}(\psi _\stackrel{}{i}^{}\psi _\stackrel{}{j}+c.c.)`$, which then reduces to $`H_{XY}`$ (2).
Accurate Monte Carlo simulations yield - for the $`XY`$ model - the BKT critical temperature $`T_{BKT}=0.898J/k_B`$ . When $`UJ`$, a BKT transition still occurs at a slightly lower critical temperature $`T_{BKT}(U)`$. To evaluate more accurately the effects of the interaction term on the BKT critical temperature one may use a semiclassical approximation to evaluate a renormalized Josephson energy : the renormalization group equations yields then the following equation for $`𝒦\beta T_{BKT}(U)`$
$$2\pi 𝒦[1\frac{1}{4𝒦}\frac{1}{X_u𝒦^3}g(𝒦,X_u)]=0$$
(6)
where $`X_u=U/\pi J`$ and $`g(𝒦,X_u)_{n=1}^{\mathrm{}}[X_u𝒦^2/2(n^2/2)ln(1+X_u𝒦^2/n^2)]`$. One may see that for $`U/J<36/\pi `$, Eq.(6) has an unique solution; at variance, for $`U/J>36/\pi `$ Eq.(6) does not have any solution. The value $`U/J=36/\pi `$ is in reasonable agreement with the known mean-field prediction $`2zJ/U1`$ for the $`T=0`$ phase transition (see, e.g., Eq. (30) of Ref. for $`N_01`$). In the inset of Fig.1 we plot the value of $`T_{BKT}(U)/T_{BKT}`$ (where $`T_{BKT}`$ is the BKT critical temperature of the XY model with $`U=0`$) as a function of $`U/J`$ from the numerical solution of Eq.(6).
The emerging physical picture is the following: There are two relevant temperatures for the system, the temperature $`T_{BEC}`$, at which in each well there is a condensate, and the temperature $`T_{BKT}`$ at which the condensates phases start to coherently point in the same direction. Of course, in order to have well defined condensates phases one should have $`T_{BKT}<T_{BEC}`$. The critical temperature $`T_{BEC}0.94\mathrm{}N_0^{1/3}\overline{\omega }/k_B`$ where $`\overline{\omega }=(\stackrel{~}{\omega }_r^2\omega _z)^{1/3}`$. With the numerical values given in , $`T_{BEC}`$ turns out to be $`500nK`$ for $`s20`$. When $`T<T_{BEC}`$, the atoms in the well $`\stackrel{}{j}`$ of the $`2D`$ optical lattice may be described by a macroscopic wavefunction $`\psi _\stackrel{}{j}`$. Furthermore, when the fluctuations around the average number of atoms per site $`N_01`$ are strongly suppressed, one may put, apart from the factor $`\sqrt{N_0}`$ constant across the array, $`\psi _\stackrel{}{j}=e^{i\theta _\stackrel{}{j}}`$. The temperature $`T_{BKT}`$ is of order of $`J/k_B`$: with the experimental parameters of and $`V_0`$ between $`20`$ and $`25E_R`$, one has that $`T_{BKT}`$ is between $`10`$ and $`30nK`$, which is sensibly smaller than the condensation temperature $`T_{BEC}`$ of the single well. A similar picture describes also $`2D`$ arrays of superconducting Josephson junctions : they exhibit a temperature $`T_{BCS}`$ at which the metallic grains placed on the sites of the array become (separately) superconducting and the Cooper pairs may be described by macroscopic wavefunctions. At a temperature $`T_{BKT}<T_{BCS}`$, the array undergo a BKT transition and the system - as a whole - becomes superconducting.
The experimental signature of the BKT transition in bosonic planar lattices is obtained by measuring, as a function of the temperature, the central peak of the interference pattern obtained after turning off the confining potentials . In fact, the peak of the momentum distribution at $`\stackrel{}{k}=0`$ is the direct analog of the magnetization of a finite size $`2D`$ $`XY`$ magnet. For the $`XY`$ magnets, the spins can be written as $`\stackrel{}{S}_\stackrel{}{j}=(\mathrm{cos}\theta _\stackrel{}{j},\mathrm{sin}\theta _\stackrel{}{j})`$ and the magnetization is defined as $`M=(1/N)_\stackrel{}{j}\stackrel{}{S}_\stackrel{}{j}`$ where $`\mathrm{}`$ stands for the thermal average. A spin-wave analysis at low temperatures yields $`M=(2N)^{k_BT/8\pi J}`$ . With discrete BECs at $`T=0`$, all the phases $`\theta _\stackrel{}{j}`$ are equal at the equilibrium and the lattice Fourier transform of $`\psi _j`$, $`\stackrel{~}{\psi }_\stackrel{}{k}=\frac{1}{N}_\stackrel{}{j}\psi _\stackrel{}{j}e^{i\stackrel{}{k}\stackrel{}{j}}`$, exhibits a peak at $`\stackrel{}{k}=0`$ ($`\stackrel{}{k}`$ is in the first Brillouin zone of the $`2D`$ square lattice) and the magnetization is:
$$M=\stackrel{~}{\psi }_0.$$
(7)
The intuitive picture of the BKT transition is then the following: at $`T=0`$, all the spins point in the same direction. Increasing the temperature, bound vortex pairs are thermally induced. In Fig.2, left inset, a single free vortex is depicted: we plot a configuration of the phases such that $`\theta _\stackrel{}{j}`$ equals the polar angle of the site $`\stackrel{}{j}`$ with respect to the core vortex (notice that with $`\theta _\stackrel{}{j}`$ equals the polar angle plus $`\pi /2`$ one would have the usual plot of a vortex). As one can see, a single vortex modifies the phase distribution also much far from its core and the square modulus of its lattice Fourier transform $`\stackrel{~}{\psi }_\stackrel{}{k}^2`$ has a minimum at $`\stackrel{}{k}=0`$. At variance, a vortex-antivortex pair \[see Fig.2, right inset\] modifies the phase distribution only near the center of the pair (in this sense is a defect) and its lattice Fourier transform has a maximum at $`\stackrel{}{k}=0`$. Increasing the temperature, vortices are thermally induced. For $`T<T_{BKT}`$ only bound vortex pairs are present, and in average the spins continue to point in the same direction. When the condensates expand, a large peak (i.e., a magnetization) is observed in the central $`\stackrel{}{k}=0`$ momentum component. Rising further the temperature, due to the increasing number of vortex pairs, the central peak density decreases from the $`T=0`$ value. For $`TT_{BKT}`$ the pairs starts to unbind and free vortices begin to appear, determining a sharp decrease around $`T_{BKT}`$ of the magnetization. At high temperatures, only free vortices are present, leading to a randomization of the phases and to a vanishing magnetization. In Fig.1 we plot the intensity of the central peak of the momentum distribution (normalized to the value at $`T=0`$) in a 2D lattice as a function of the reduced temperature $`k_BT/J`$, evidencing the sharp decrease of the magnetization around the critical temperature .
At $`t=0`$ the amplitude of the $`\stackrel{}{k}=0`$ peak of the momentum distribution is simply given by the thermal average of $`\stackrel{~}{\psi }_0`$. By measuring the $`\stackrel{}{k}=0`$ peak (i.e., $`\stackrel{~}{\psi }_0^2`$) at different temperatures, one obtains the results plotted in Fig.1. The figure has been obtained using a Monte Carlo simulation of the $`XY`$ magnet for a $`40\times 40`$ array: we find $`k_BT_C1.07J`$. In Fig.1 we also plot the low-temperature spin wave prediction (solid line). At times different from $`t=0`$, the density profiles are well reproduced by the free expansion of the ideal gas. One obtains $`\stackrel{~}{\psi }(\stackrel{}{p},t)=\chi (p_z)e^{i[(p_z+mgt)^3p_z^3]/6m^2g\mathrm{}}\stackrel{~}{\phi }(p_x,p_y,t)`$ \[where $`\chi (p_x)e^{\sigma _x^2p_x^2/2\mathrm{}^2}`$ and similarly for $`\chi (p_y)`$ and $`\chi (p_z)`$\], i.e. a uniformly accelerating motion along $`z`$ and a free motion in the plan $`xy`$, with $`\stackrel{~}{\phi }(p_x,p_y,t)=\chi (p_x)\chi (p_y)\stackrel{~}{\psi }_\stackrel{}{k}e^{i(p_x^2+p_y^2)t/2\mathrm{}m}`$ giving the central and lateral peaks of the momentum distribution as a function of time for different temperatures.
An intense experimental work is now focusing on the Bose-Einstein condensation in two dimensions: at present a crossover to two-dimensional behaviour has been observed for $`Na`$ and $`Cs`$ atoms . Our analysis relies on two basic assumptions; namely, the validity of the tight-binding approximation for the Bose-Hubbard Hamiltonian and the requirement that the condensate in the optical lattice may be regarded as planar . It is easy to see that the first assumption is satisfied if, at $`T=0`$, $`V_0\mu `$ (where $`\mu `$ is the chemical potential), and, at finite temperature, $`\mathrm{}\stackrel{~}{\omega }_rk_BT`$. The second assumption is much more restrictive since it requires freezing the transverse excitations; for this to happen one should require a condition on the transverse trapping frequency $`\omega _z`$. Namely, one should have that, at $`T=0`$, $`\mathrm{}\omega _z8K`$ and that, at finite temperature, $`\mathrm{}\omega _zk_BT`$ (since $`\omega _z\stackrel{~}{\omega }_r`$, the latter condition also implies that $`\mathrm{}\stackrel{~}{\omega }_rk_BT`$). In it is $`V_0=sk_B0.15\mu K`$ and for $`s20`$ the tight-binding conditions are satisfied since $`10Hz8K/2\pi \mathrm{}`$; furthermore, if $`\omega _z=2\pi 1kHz`$, one may safely regard our finite temperature analysis to be valid at least up to $`T50nK`$. We notice that the experimental signature for the BKT transition for a continuous (i.e., without optical lattice) weakly interacting 2D Bose gas is also given by the central peak of the atomic density of the expanding condensates.
Our study evidences the possibility that Bose-Einstein condensates loaded on a $`2D`$ optical lattice may exhibit \- at finite temperature - a new coherent behaviour in which all the phases of the condensates located in each well of the lattice point in the same direction. The finite temperature transition, which is due to the thermal atoms in each well, is mediated by vortex defects and may be experimentally detectable by looking at the interference patterns of the expanding condensates. Our analysis strengthens - and extends at finite temperature - the striking and deep analogy of bosonic planar systems with $`2D`$ Josephson junction arrays .
Acknowledgements This work has been supported by MIUR through grant No. 2001028294 and by the DOE.
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# Rapid optimization of working parameters of microwave-driven multi-level qubits for minimal gate leakage
(June 29, 2005)
## Abstract
We propose an effective method to optimize the working parameters (WPs) of microwave-driven quantum logical gates implemented with multi-level physical qubits. We show that by treating transitions between each pair of levels independently, intrinsic gate errors due primarily to population leakage to undesired states can be estimated accurately from spectroscopic properties of the qubits and minimized by choosing appropriate WPs. The validity and efficiency of the approach are demonstrated by applying it to optimize the WPs of two coupled rf SQUID flux qubits for controlled-NOT (CNOT) operation. The result of this independent transition approximation (ITA) is in good agreement with that of dynamic method (DM). Furthermore, the ratio of the speed of ITA to that of DM scales exponentially as $`2^n`$ when the number of qubits $`n`$ increases.
A practical quantum computer would be comprised of a large number of coupled qubits and the coupled qubits must be kept in high-degree quantum coherence states for sufficiently long time. During the past decade, significant progress has been made on physical implementation of quantum computation. High-degree quantum coherence has been demonstrated experimentally in systems such as trapped ions Cirac and Zoller (1995); Monroe et al. (1996), nuclear spins Gershenfeld and Chuang (1997); Jones et al. (1998), atoms in optical resonators Turchette et al. (1995), and photons in microwave cavities Brune et al. (1996). However, it seems quite difficult to realize a large number of coupled qubits using these systems. Meanwhile, solid-state qubits are of particular interest because of their advantages of large-scale integration, flexibility in design, and easy connection to conventional electronic circuits Mooij et al. (1999). Of those, qubits based on superconducting devices have recently attracted much attention Makhlin et al. (2001) as manipulation of quantum coherent states being successfully demonstrated in a variety of single qubits Nakamura et al. (1999); Vion et al. (2002); Yu et al. (2002); Martinis et al. (2002); Chiorescu et al. (2003) and coupled two-qubit systems Pashkin et al. (2003); Yamamoto et al. (2003); Berkley et al. (2003); Chiorescu et al. (2004); Majer et al. (2005).
However, the solid-state qubits demonstrated in experiments so far all have relatively short coherence time and high probability of gate errors Nakamura et al. (1999); Vion et al. (2002); Yu et al. (2002); Martinis et al. (2002); Chiorescu et al. (2003); Pashkin et al. (2003); Yamamoto et al. (2003); Chiorescu et al. (2004). One of the causes of these problems is extrinsic gate error arising from interaction between the environment and qubits resulting in decoherence, such as dephasing and relaxation Makhlin et al. (2001). Another cause is intrinsic gate error resulting from population leakage to undesired states due to the typical multi-level structures of solid-state qubits Fazio et al. (1999); Zhou et al. (2002). The intrinsic gate error is crucial since it not only contributes to additional decoherence but also determines the ultimate performance of the quantum gates and cannot be eliminated by reducing the environment caused decoherence.
The leakage of a gate can be characterized quantitatively by summing up the maximum transition probabilities to all undesired states of a multi-level qubit Fazio et al. (1999); Zhou et al. (2002). It depends strongly on energy level structure and transition matrix elements, i.e., spectroscopic properties of the multi-level qubit determined completely by device parameters (DPs) and external control parameters of the qubit which we call working parameters (WPs) for simplicity. For instance, inductance (capacitance) of the superconducting flux (charge) qubit is a DP while external flux (gate voltage) is a WP. For the multi-level qubit with given DPs, the gate leakage is sensitive to their WPs and thus can be minimized by the use of appropriate WPs.
Conventionally, the leakage is calculated from transition probabilities to all undesired states by numerically solving the time-dependent Schrödinger equation (TDSE) Zhou et al. (2002, 2004). However, this kind of dynamic method (DM) not only needs complicated numerical algorithms but also a large amount of computational resources. For instance, optimizing the WPs of a many-qubit network needed for a practical quantum computer using DM may only be possible with ad hoc powerful quantum computers in the future. Thus a much faster approach is highly desirable.
In this Letter, we propose a very fast method to minimize the leakage of microwave-driven quantum logical gates by choosing appropriate WPs in a system of multi-level qubits with their DPs given in prior. We consider the case of weak microwave fields only since strong fields usually cause many additional types of intrinsic gate errors and thus should be avoided in general Zhou et al. (2002). Our method is based on an independent transition approximation (ITA) in which transitions in multi-level qubits are treated independently. The leakage is estimated using the spectroscopic properties obtained by solving the eigenvalue equation of and minimized by optimizing the WPs of the qubits. The method is applied to minimize leakage of controlled-NOT (CNOT) gate implemented with coupled rf superconducting quantum interference device (SQUID) flux qubits. The result is in good agreement with that obtained from DM. More importantly, the ITA is scalable as the number of qubit $`n`$ increases because the ratio of the speed of ITA to that of DM scales exponentially as $`2^n`$.
A microwave-driven gate is realized via coherent transitions between the computational states of qubits interacting with microwave fields. In general, correlation and interference between transitions cannot be ignored and transition probability can only be computed accurately by solving TDSE. However, in weak fields, which is the case considered here, only the transitions between levels with which the microwave field is resonant or nearly resonant are significant. For a given level, if the level spacings between it and all other levels are sufficiently different, the correlation and interference between the transitions from this level have negligible effect. Hence, each transition is expected to take place independently and the two levels involved can thus be treated as if they are isolated from the others. For each of the two-level sub-systems interacting with a rectangular pulse, the transition probability can be approximated by an analytical expression using the rotating-wave approximation (RWA). Assuming the rectangular pulse is $`ϵ\left(t\right)=ϵ_0\mathrm{cos}\left(\omega t\right)`$, the interaction between the qubit and the pulse is then $`V(t)=\mu \text{·}ϵ\left(t\right)`$, where $`ϵ_0`$ and $`\omega `$ are the amplitude and frequency of the pulse and $`\mu `$ is the dipole moment operator of the qubit. The maximum transition probability from state $`i`$ to $`j`$ in an $`N`$-photon process, $`𝒫_{ij}`$, is Kmetic et al. (1986)
$$𝒫_{ij}=\mathrm{\Omega }_{ij}^2/\left(𝒟_{ij}^2+\mathrm{\Omega }_{ij}^2\right),$$
(1)
where, $`\mathrm{\Omega }_{ij}=2\mu _{ij}`$·$`ϵ_0NJ_N(y_{ij})/y_{ij}`$ is the Rabi frequency for the $`N`$-photon resonance, $`𝒟_{ij}=\mathrm{\Delta }E_{ij}/\mathrm{}N\omega `$ is the detuning, $`\mu _{ij}=i\left|\mu \right|j`$ is the dipole transition matrix element, $`\mathrm{\Delta }E_{ij}=\left|E_jE_i\right|`$ is the level spacing between the states $`i`$ and $`j`$, $`J_N`$ is the Bessel function of integer order $`N`$, and $`y_{ij}=𝐝_{ij}`$·$`ϵ_0/\mathrm{}\omega `$ with $`𝐝_{ij}=\mu _{jj}\mu _{ii}`$. Note that $`𝒫_{ij}`$, $`\mathrm{\Omega }_{ij}`$, and $`𝒟_{ij}`$ depend on the number of photons $`N`$. In weak fields characterized by $`\mathrm{\Omega }_{ij}\omega `$ only transitions with small $`N`$ are important. Since the leakage is calculated from the maximum transition probabilities which are independent of pulse shape in weak fields Shore (1990) the optimized WPs obtained using rectangular pulse are also valid for other pulse shapes.
For the microwave-driven gate considered here, the microwave is resonant with a pair of computational states. According to Eq. (1), the leakage is suppressed if $`𝒟_{ij}\mathrm{\Omega }_{ij}`$ holds for all unintended transitions. This condition may be satisfied by using sufficiently weak fields and/or qubits with proper spectroscopic properties. However, the use of exceedingly weak fields will make the gate very slow. Even so the unintended transitions will still cause large leakage at or near resonance where $`𝒟_{ij}0`$ and $`𝒫_{ij}`$ $`1`$ regardless of intensity of the fields. Thus the leakage can be greatly reduced by setting the level spacings of the qubits substantially detuned from the microwave frequency and/or the dipole transition matrix elements sufficiently small for all unintended transitions.
To investigate the leakage in a multi-level qubit, we study what happens when the microwave acts on each of the computational states, which we call a component of the gate. For an $`n`$-bit gate, there are $`M=2^n`$ components. The leakage of the $`i`$th component is defined as $`\eta _i=_kP_{ik}`$, where $`P_{ik}`$ is the maximum probability of occupying an undesired state $`k`$ through all possible multi-photon transitions and the sum is over all the undesired states including all non-computational states as well as those computational states to which no transition is intended. The leakage of the $`n`$-bit gate, $`\eta `$, is then defined as $`\eta =\mathrm{max}(\eta _1,\eta _2,\mathrm{}\eta _M)`$ which is a function of the $`\nu `$ WPs $`q_1`$, $`q_2`$, $`\mathrm{}`$, and $`q_\nu `$. The set of WPs, $`q_1^{\text{opt}}`$, $`q_2^{\text{opt}}`$, $`\mathrm{}`$, and $`q_\nu ^{\text{opt}}`$ is obtained by minimizing the leakage of gate, namely, $`\eta (q_1^{\text{opt}},q_2^{\text{opt}},\mathrm{},q_\nu ^{\text{opt}})\mathrm{min}\left[\eta (q_1,q_2,\mathrm{},q_\nu )\right]`$.
As an example, we apply ITA to optimize the WPs of coupled rf SQUID flux qubits for CNOT gate. Each rf SQUID consists of a superconducting loop of inductance $`L`$ interrupted by a Josephson tunnel junction characterized by its critical current $`I_c`$ and shunt capacitance $`C`$ Danilov et al. (1983). A flux-biased rf SQUID with total magnetic flux $`\mathrm{\Phi }`$ enclosed in the loop is analogous to a ”flux” particle of mass $`m=C\mathrm{\Phi }_0^2`$, where $`\mathrm{\Phi }_0`$ $`=h/2e`$ is the flux quantum. The Hamiltonian is $`h\left(x\right)=p^2/2m+m\omega _{LC}^2\left(xx_e\right)^2/2E_J\mathrm{cos}\left(2\pi x\right)`$. Here, $`x=\mathrm{\Phi }/\mathrm{\Phi }_0`$ is the canonical coordinate of the “flux” particle, $`p=i\mathrm{}/x`$ is the canonical momentum conjugate to $`x`$, $`E_J=\mathrm{}I_c/2e=m\omega _{LC}^2\beta _L/4\pi ^2`$ is the Josephson coupling energy, $`\beta _L=2\pi LI_c/\mathrm{\Phi }_0`$ is the potential shape parameter, $`\omega _{LC}=1/\sqrt{LC}`$ is the characteristic frequency of the SQUID, and $`x_e=\mathrm{\Phi }_e/\mathrm{\Phi }_0`$ is the normalized external flux.
The coupled rf SQUID qubits comprise two rf SQUID qubits: a control qubit and a target qubit coupled via their mutual inductance $`M`$ Mooij et al. (1999); Zhou et al. (2005, in press). For simplicity, we assume that the two SQUIDs are identical: $`C_i=C`$, $`L_i=L`$, and $`I_{ci}=I_c`$ for $`i=1`$ and $`2`$. The Hamiltonian of the coupled SQUID qubits is $`H(x_1,x_2)=h\left(x_1\right)+h\left(x_2\right)+h_{12}(x_1,x_2)`$, where $`x_i`$, $`x_{ei}`$, and $`h\left(x_i\right)`$ ($`i=1`$ and $`2`$) are the canonical coordinate, normalized external flux, and Hamiltonian of the $`i`$th single SQUID qubit and $`h_{12}`$ is the interaction between the qubits given by $`h_{12}(x_1,x_2)=m\omega _{LC}^2\kappa \left(x_1x_{e1}\right)\left(x_2x_{e2}\right)`$. Here, $`\kappa =M/L`$ is the coupling constant. Note that each SQUID qubit is a multi-level system. The eigenstate $`|n)`$ and eigenenergy $`E_n`$ of the coupled qubits are computed by numerically solving the eigenvalue equation of $`H(x_1,x_2)`$ using the two-dimensional Fourier-grid Hamiltonian method Chu (1990). They are functions of the WPs $`x_{e1}`$, $`x_{e2}`$, and $`\kappa `$ for given device parameters $`L`$, $`C`$, and $`I_c`$ Filippov et al. (2003); van den Brink and Berkley (unpublished). For weak coupling $`\kappa 1`$, the eigenstate of the coupled qubits, denoted by $`|n)=|ij`$, can be well approximated by the product of the control qubit’s state $`|i`$ and the target qubit’s state $`|j`$, $`|n)=|i|j`$. When biased at $`x_{e1},`$ $`x_{e2}1/2`$ the potential of coupled SQUID qubits have four wells Mooij et al. (1999); Zhou et al. (2005, in press). The lowest eigenstates in each of the four wells, denoted as $`|1)=|00`$, $`|2)=|01`$, $`|3)=|10`$, and $`|4)=|11`$, are used as the computational states of the coupled SQUID qubits. For SQUIDs with $`L=100`$ pH, $`C=40`$ fF, and $`\beta _L=1.2`$, the energy levels and level spacings versus $`x_{e2}`$ and $`\kappa `$ are plotted in Fig. 1(a)$``$(d), respectively. It is shown that both the energy levels and level spacings are sensitive to the WPs $`x_{e2}`$ and $`\kappa `$. When $`x_{e2}<0.49877`$ or $`x_{e2}>0.50123`$ in Fig. 1(a) or $`\kappa >7.5\times 10^3`$ in Fig. 1(c) the energy level structures become quite complicated. Fig. 1(b) and Fig. 1(d) also show that in certain regions of the WP space the level spacings become crowded, meaning that they are degenerate or nearly degenerate, which may result in significant intrinsic gate errors.
Controlled two-bit gates can be realized by applying a resonant microwave pulse to the target qubit. The interaction between the microwave and the coupled qubits can be written as $`V(x_1,x_2,t)=m\omega _{LC}^2\left[\left(x_2x_{e2}\right)+\kappa \left(x_1x_{e1}\right)+x_m/2\right]x_m`$, where $`x_m\left(t\right)=x_{m0}\mathrm{cos}\left(\omega t\right)`$ is the magnetic flux (normalized to $`\mathrm{\Phi }_0`$) coupled to the target qubit from the microwave with amplitude $`x_{m0}`$ and frequency $`\omega `$. For the CNOT gate, a $`\pi `$-pulse with $`\omega =\mathrm{\Delta }E_{34}/\mathrm{}`$, where $`\mathrm{\Delta }E_{34}=\left|E_4E_3\right|`$, is used. Populations of the states $`|10`$ and $`|11`$ are exchanged after the $`\pi `$-pulse only if the initial state is $`|10`$ or $`|11`$ or a linear combination of them. Note that in this case in addition to the non-computational states, the undesired states also include the computational states $`|00`$ and $`|01`$. Using Eq. (1) and following the steps described above, the leakage of the CNOT gate $`\eta =\mathrm{max}(\eta _{|00},\eta _{|01},\eta _{|10},\eta _{|11})`$ is calculated. In the calculations, up to 3-photon processes are included. The $`N`$-photon processes for $`N>3`$ have negligible effects. For the SQUID qubits with the previously given DPs, the leakage is a function of WPs $`x_{e1}`$, $`x_{e2}`$, and $`\kappa `$. In Fig. 2(a) we plot the leakage of the CNOT gate versus $`x_{e2}`$ and $`\kappa `$ for $`x_{e1}=0.499`$ and $`x_{m0}=2\times 10^4`$. It is shown that the leakage of the CNOT gate is much smaller when the coupled SQUID qubits are operated around $`x_{e2}=0.4997`$ and $`\kappa =7.5\times 10^4`$ for $`x_{e1}=0.499`$.
To evaluate the results obtained from ITA, we performed an accurate dynamic calculation by solving the TDSE of the coupled qubits $`ic_n\left(\tau \right)/\tau =_n^{}H_{nn^{}}^R\left(\tau \right)c_n^{}\left(\tau \right)`$ for the probability amplitudes $`c_n`$ of the first 20 eigenstates $`|n)`$ using the split-operator method Hermann and Fleck, Jr. (1988), where $`\tau =\omega _{LC}t`$ and $`H_{nn^{}}^R=\left[E_n\delta _{nn^{}}+\left(n\left|V\right|n^{}\right)\right]/\mathrm{}\omega _{LC}`$ Zhou et al. (2005, in press). The probability of being in the state $`|n)`$ is $`\left|c_n\right|^2`$, from which the maximum probabilities on the undesired states and the leakage $`\eta (x_{e1},x_{e2},\kappa )`$ of the CNOT gate are calculated. The results are shown in Fig. 2(b) with the color scale identical to Fig. 2(a). To have a more quantitative comparison, we also plot, in Fig. 2(c) and (d), two line-cut figures at places with the richest structures in Fig. 2(a) and (b). The good agreement between the results of ITA and DM demonstrates the validity of using ITA to minimize the CNOT gate leakage.
One of the advantages of ITA is that it provides clear physical intuition and insight into the origin of intrinsic gate errors and thus how to reduce the problem by selecting appropriate WPs. Furthermore, ITA is orders of magnitude faster than DM. For instance, the result shown in Fig. 2(a) took less than 3 hours to compute on a dual-2.8 GHz processor Dell PRECISION 650 workstation while that shown in Fig. 2(b) took more than 300 hours on the same computer. Denoting $`\tau _S`$ and $`\tau _T`$ respectively the time needed to calculate the spectroscopic properties of an $`n`$-bit gate and that used to compute the probability evolution in DM for each component. The total times needed to compute the leakage by ITA and DM are approximately $`\tau _I\tau _S`$ and $`\tau _D\tau _S+2^n\tau _T`$, respectively, where the factor $`2^n`$ is the number of components for the $`n`$-bit gate. Hence, $`\tau _S<\tau _T`$ and the ratio $`\tau _D/\tau _I2^n\zeta `$ increases exponentially with $`n`$, where $`\zeta \tau _T/\tau _S`$. For weak fields, $`\zeta 1`$ and the ratio $`\tau _D/\tau _I`$ is very large. For example, for the CNOT gate above $`\zeta 25`$ and $`\tau _D/\tau _I100`$. Thus as the number of qubits increases the optimization of WPs of multi-bit quantum gates could become an intractable problem using DM.
To demonstrate quantitatively the effect of different WPs on intrinsic gate errors, we plot in Fig. 3(a) and (b) the population evolution of the computational states $`|10`$ and $`|11`$ for the CNOT gate operated with two different sets of WPs corresponding to points ”A” and ”B” marked in Fig. 2(a) \[also (b)\], respectively. It is clearly shown that inversion from the initial state $`|10`$ ($`|11`$) to the final state $`|11`$ ($`|10`$) after the $`\pi `$-pulse is almost complete when operated at the point ”A” but significantly incomplete at the point ”B”. We also computed the population evolution of the coupled qubits with the initial states $`|00`$ and $`|01`$ using the same microwave pulse. The system essentially stayed in the initial states for the gate operated at the point ”A” but leaked significantly to the non-computational states for the gate at the point ”B”. The quality of a gate can be described by gate fidelity $`F\overline{\text{Trace}\left[\rho _P\rho _I\right]}`$, where $`\rho _P`$ and $`\rho _I`$ are the physical and ideal density matrices after gate operation and the overline denotes averaging over all possible initial states Li et al. (2003). For the CNOT gates operated at the point ”A” we obtain $`F_A=0.9997`$ which is very close to the ideal CNOT gate. In contrast, the gate operated at the point ”B” has $`F_B=0.8061`$ which is too large to be tolerated. Furthermore, at the point ”A” the gate is about a factor of three faster than that at the point ”B”.
In summary, a very efficient method is proposed to optimize the WPs of microwave-driven quantum logical gates in multi-level qubits based on ITA. In this method, the leakage is estimated accurately from the spectroscopic properties of the qubits and minimized by choosing appropriate WPs. This method is exemplified by optimizing the WPs of coupled rf SQUID flux qubits for minimal leakage of the CNOT gate. The result is in good agreement with that obtained from dynamic calculations. Compared to the conventional dynamic method the ITA not only provides physical insight into the origin of gate leakage but also reduces the computational time by more than two orders of magnitude for the CNOT gate. Furthermore, since the ratio of the speed of ITA to that of DM scales exponentially as $`2^n`$ the ITA is scalable as the number of qubit $`n`$ increases. Our calculation also shows that high intrinsic fidelity CNOT gate can be achieved using microwave-driven rf SQUID qubits with properly selected WPs. Although the proposed ITA is only applied, as an example, to the rf SQUID qubits, it is also valid for other microwave-driven multi-level qubits. Therefore, ITA provides a much needed solution for the optimization of WPs of multi-bit gates implemented with multi-level physical qubits, for which DM is extremely time consuming or could even become intractable.
This work was supported in part by the NSF (DMR-0325551) and AFOSR, NSA and ARDA through DURINT grant (F49620-01-1-0439).
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# Pareto index induced from the scale of companies
## 1 Introduction
The pioneering discovery by Pareto is not only an important cue for research in fractal but also a significant issue in economics. The Pareto law states that a cumulative number distribution $`N_P(>x)`$ obeys power-law for income $`x`$ which is larger than a certain observational threshold $`x_0`$:
$`N_P(>x)x^\mu \mathrm{for}x>x_0,`$ (1)
where the exponent $`\mu `$ is called Pareto index. Recently Pareto’s law is checked with high accuracy by using digitalized data . Although the number of persons or companies in the range where Pareto’s law (1) holds is a few percent, the amount of the income occupies a large part of the total income. High income persons or high income companies can, therefore, influence economics. It is quite important to investigate the mechanism which governs them.
The research on the income distribution is well investigated in econophysics . Many models to explain the distribution are proposed (for instance recently Refs. ). Furthermore, Fujiwara et al. find that, without assuming any model, Pareto’s law can be derived kinematically from the law of detailed balance and Gibrat’s law which are observed in high quality digitalized data. They derive not only Pareto’s law but also the reflection law which correlates a positive growth rate distribution to negative one by the Pareto index $`\mu `$. This means that Pareto’s law and index are explained by the law of detailed balance and Gibrat’s law. We should note that these arguments are limited in the high income region ($`x>x_0`$).
The findings in Ref. are quite fascinating especially in the idea that the Pareto index can be understood by the growth rate distribution. In this paper, we first investigate the application as follows. The distribution of high income companies follows Pareto’s law and the Pareto index is about $`1`$. This is called Zipf’s law . It is also known that the distributions in most job categories follow Pareto’s law. Those Pareto indices, however, scatter around $`1`$ . We examine the reflection law that relates the change of the Pareto index to the change of the growth rate distribution in job categories.
There is, however, a problem in estimating the Pareto index by the growth rate distribution in each job category. The number of companies in several job classes is not enough to observe the reflection law. We clear this difficulty by employing profits data which have no threshold in contrast to high income data announced publicly.
By observing all positive profits data, we confirm the law of detailed balance. We also find that Gibrat’s law holds in the region where the only initial profits are conditioned to be larger than the threshold. In this case, Pareto’s law is also obtained and the reflection law can be observed in companies data classified into job categories. Using these arguments, we confirm the correlation between the Pareto index and the growth rate distribution.
In addition, we find that companies in a small scale job category have more possibilities of growing profits than those in a large scale job category. This kinematically explains that the Pareto index for the companies in the small job class is larger than that for the companies in the large scale job class, which is reported in Ref. .
## 2 Pareto’s law derived from the detailed balance and Gibrat’s law
In this section, we briefly review the proof that the detailed balance and Gibrat’s law lead to Pareto’s law and the reflection law .
Let the income at two successive points in time be denoted by $`x_1`$ and $`x_2`$. Growth rate $`R`$ is defined as the ratio $`R=x_2/x_1`$. The detailed balance and Gibrat’s law can be summarized as follows.
* Detailed balance (Time-reversal symmetry)
The joint probability distribution function (pdf) $`P_{12}(x_1,x_2)`$ is symmetric:
$`P_{12}(x_1,x_2)=P_{12}(x_2,x_1).`$ (2)
* Gibrat’s law
The conditional probability distribution of growth rate $`Q(R|x_1)`$ is independent of the initial value $`x_1`$:
$`Q(R|x_1)=Q(R),`$ (3)
where $`Q(R|x_1)`$ is defined by using the pdf $`P_1(x_1)`$ and the joint pdf $`P_{1R}(x_1,R)`$ as
$`Q(R|x_1)={\displaystyle \frac{P_{1R}(x_1,R)}{P_1(x_1)}}.`$ (4)
These phenomenological properties, for example, are observed in the data of high income Japanese companies in 2002 and 2003. In Japan, companies having annual income more than 40 million yen are announced publicly as “high income companies” every year. The database is published by Diamond Inc. In Fig. 1, all the companies are plotted, the income of which in 2002 ($`x_1`$) and 2003 ($`x_2`$) exceeded 40 million yen ($`x_0`$):
$`x_1>x_0\mathrm{and}x_2>x_0.`$ (5)
The number of companies which satisfy this condition is $`50,632`$. In Fig. 1, the time-reversal symmetry (2) is apparent. In Fig. 2, we confirm Pareto’s law (1) for the companies, the income of which satisfies the condition (5). The Pareto indices $`\mu `$ are very close to $`1`$, and the case is often referred to as Zipf’s law .
We also examine Gibrat’s law in the regime (5). Here we divide the range of $`x_1`$ into logarithmically equal bins as $`x_14\times [10^{4+0.2(n1)},10^{4+0.2n}]`$ with $`n=1,\mathrm{},5`$. In Fig. 3, the probability density for $`r`$ is expressed for each bin. In all the cases, all the probability density distributions for different $`n`$ collapse on a single curve approximately. This means that the distribution for the growth rate $`r`$ is approximately independent of the initial value $`x_1`$. Here the probability density for $`r`$ defined by $`q(r)`$ is related to that for $`R`$ by
$`\mathrm{log}_{10}Q(R=10^r)+r+\mathrm{log}_{10}(\mathrm{ln}10)=\mathrm{log}_{10}q(r).`$ (6)
Pareto’s law (1) can be derived from the detailed balance (2) and Gibrat’s law (3) without assuming any model . Due to the relation of $`P_{12}(x_1,x_2)dx_1dx_2=P_{1R}(x_1,R)dx_1dR`$ under the change of variables from $`(x_1,x_2)`$ to $`(x_1,R)`$, these two joint pdf’s are related to each other,
$`P_{1R}(x_1,R)=x_1P_{12}(x_1,x_2).`$ (7)
By the use of this relation, the detailed balance (2) is rewritten in terms of $`P_{1R}(x_1,R)`$ as follows:
$`P_{1R}(x_1,R)=R^1P_{1R}(x_2,R^1).`$ (8)
Substituting the joint pdf $`P_{1R}(x_1,R)`$ for the conditional probability $`Q(R|x_1)`$ defined in Eq. (4), the detailed balance is expressed as
$`{\displaystyle \frac{P_1(x_1)}{P_1(x_2)}}={\displaystyle \frac{1}{R}}{\displaystyle \frac{Q(R^1|x_2)}{Q(R|x_1)}}.`$ (9)
From the detailed balance and Gibrat’s law (3), one finds the following:
$`{\displaystyle \frac{P_1(x_1)}{P_1(x_2)}}={\displaystyle \frac{1}{R}}{\displaystyle \frac{Q(R^1)}{Q(R)}}.`$ (10)
By expanding this equation around $`R=1`$, the following differential equation is obtained
$`G^{}(1)P_1(x)+xP_1^{}(x)=0,`$ (11)
where $`x`$ denotes $`x_1`$ and $`G(R)`$ is defined as $`Q(R^1)/(RQ(R))`$. The solution is given by
$`P_1(x)=Cx^{\mu 1}`$ (12)
with $`G^{}(1)=\mu +1`$. This is Pareto’s law, which leads to the cumulative representation in Eq. (1) as follows: $`N_P(>x)=N_P(>x_0)P_1(>x)=N_P(>x_0)_x^{\mathrm{}}dyP_1(y)x^\mu `$.
From Eqs. (12) and (10), one finds the relation between positive growth rate ($`R>1`$) and negative one ($`R<1`$),
$`Q(R)=R^{\mu 2}Q(R^1).`$ (13)
In terms of $`q(r)`$ in Eq. (6), this is written as
$`\mathrm{log}_{10}q(r)=\mu r+\mathrm{log}_{10}q(r).`$ (14)
In Fig. 3, we examine this relation, which is called the ’Reflection law’ in Ref. . By the use of the reflection law to the best-fit line
$`\mathrm{log}_{10}q(r)=ct_+r\mathrm{for}r>0,`$ (15)
we obtain the line
$`\mathrm{log}_{10}q(r)=c+t_{}r\mathrm{for}r<0`$ (16)
with
$`\mu =t_+t_{},`$ (17)
where $`\mu 1`$. The reflected line (16) seems to match the data in Fig. 3 .
## 3 Pareto index induced from the growth rate distribution
From the argument in the preceding section, we expect the estimation of the Pareto index by using Eq. (17) under the detailed balance and Gibrat’s law. There is, however, a difficulty in measuring the slope $`t_{}`$ of the reflected line (16). The reason is that the data for large negative growth, $`r4+\mathrm{log}_{10}4\mathrm{log}_{10}x_1`$, are not available due to the limit (5) (Fig. 3). The reflected line (16) does not match the total data which are not classified into bins for instance (Fig. 4). Indeed, the Pareto index is estimated to be about $`0`$ from the best fit lines (15) for $`r>0`$ and (16) for $`r<0`$.
In order to avoid this problem, we must remove the threshold $`x_0`$ in the region (5). This is, however, impossible as long as we deal with the income data announced publicly in Japan. We employ profits as another economic quantity similar to the income. Although profits of companies are not announced publicly, the data are available on “CD Eyes” published by Tokyo Shoko Research, Ltd. for instance. This database contains financial information of over 250 thousand companies in Japan. We deal with the companies, the profits of which in 2002 ($`x_1`$) and 2003 ($`x_2`$) exceeded 0 yen:
$`x_1>0\mathrm{and}x_2>0.`$ (18)
In Fig. 5, all the companies which satisfy this condition are plotted, the number of which is $`132,499`$. The time-reversal symmetry (2) is apparent even in Fig. 5.
We desire to examine Gibrat’s law without the threshold. The condition is, therefore, employed:
$`x_1>x_0\mathrm{and}x_2>0,`$ (19)
where $`x_0`$ is 40 million yen. The number of the companies which satisfy this limit is $`28,644`$. In Fig. 6, Pareto’s law (1) for the companies is observed in the region (5), not in (19). Let us examine Gibrat’s law in the regime (19). We classify the data for the high profits companies, $`x_1>x_0`$, into five bins in the same manner in the preceding section. The result is shown in Fig. 7. We confirm Gibrat’s law without the forbid region caused by the threshold.
Under the extended limit (19), we can derive Pareto’s law. In the derivation, we should pay attention to the region where Gibrat’s law holds as follows:
$`Q(R|x_1)`$ $`=`$ $`Q(R)\mathrm{for}x_1>x_0\mathrm{and}x_2>0,`$ (20)
$`Q(R^1|x_2)`$ $`=`$ $`Q(R^1)\mathrm{for}x_1>0\mathrm{and}x_2>x_0.`$ (21)
In the region (5) where these two conditions are satisfying, Eq. (10) is relevant and Pareto’s law (12) holds as a result. This is observed in Fig. 6.
The reflection law (13) is relevant in the region (5) strictly. Here we assume that Eqs. (13) and (14) hold approximately near $`R=1`$ ($`r=0`$) even in the region (19), the validity of which should be checked against the results. Under this assumption, we evaluate the Pareto index by using Eq. (17). In Fig. 8, the Pareto index is estimated to be $`1.02\pm 0.02`$ from the best fit lines (15) for $`r>0`$ and (16) for $`r<0`$ for the total data which are not classified into bins. This is consistent with the assumption. We should comment that the Pareto index is estimated to be about $`0`$ for the companies, the profits of which in 2002 ($`x_1`$) and 2003 ($`x_2`$) exceeded $`x_0`$.
## 4 High profits companies in job categories
It is reported that income distributions in most job categories follow Pareto’s law, but those Pareto indices scatter around $`1`$ . The same phenomenon is observed in profits distributions. In this section, we classify the companies into job categories in order to verify the results in the preceding section.
In Japan, companies are categorized by Japan Standard Industrial Classification (JSIC). This industrial classification is composed of four stages, namely divisions, major groups, groups and details (industries). The composition includes 14 divisions, 99 major groups, 463 groups and 1,322 industries. The classification in the database “CD Eyes” follows JSIC.
We classify the $`28,644`$ companies into the 14 divisions, the profits of which satisfy the condition (19). Because the number of companies classified into 99 major groups is not enough to investigate the growth rate distribution. The name of each division and the number of companies are follows, A: Agriculture (84), B: Forestry (1), C: Fisheries (14), D: Mining (3,512), E: Construction (1,182), F: Manufacturing (6,428), G: Electricity, Gas, Heat Supply and Water (218), H: Transport and Communications (2,300), I: Wholesale, Retail Trade and Eating $`\&`$ Drinking Places (8,699), J: Finance and Insurance (1,994), K: Real Estate (293), L: Services (3,919), M: Government, N.E.C. (0) and N: Industries Unable to Classify (0).
In each job division, we measure $`t_+`$ and $`t_{}`$ of the best fit lines (15) and (16) of all the companies in the class. The Pareto index can be estimated from the relation (17) by assuming the detailed balance and Gibrat’s law. Here we have excluded A, B, C, M and N divisions, because the number of companies in each division is not enough to observe the growth rate distribution. By comparing the Pareto indices estimated by the growth rate with those directly measured in the power-law, we reconfirm the relation (17) in Fig. 9.
After identifying the Pareto index change as the change of the growth rate distribution in job categories, we consider an economic quantity which influences the growth rate. In balance-sheet, the sum of capital stock and retained earnings is equity capital, and the sum of equity capital and debt is assets. It is natural to consider that a company works based on the assets and gets profits ($``$ income). The positive correlation between assets and income is reported in Ref. . It is also pointed out that assets and sales correlate positively in the same reference. On the other hand, the positive correlation between sales and capital stock is observed in the database “CD Eyes”. Consequently there is a positive correlation between assets and capital stock<sup>2</sup><sup>2</sup>2 This correlation is probably caused by the Japanese commercial law. In Japan, capital stock is treated as an important quantity to protect creditors.. We employ, therefore, average “capital stock” as a quantity which statistically characterizes the scale of companies in each job category<sup>3</sup><sup>3</sup>3 Indeed, the database “CD Eyes” includes no information about equity capital or assets..
In Ref. , it is reported that there is a relation between the average capital stock of high income companies and the Pareto index in 99 major groups. The larger the average capital becomes, the smaller the Pareto index becomes. The same phenomenon is observed in profits distributions (Fig. 10). This means that the Pareto index is related to the scale of companies in job categories.
We can understand the kinematics by the relation between the average capital and the slope of growth rate distribution $`(t_+,t_{})`$. In Fig. 11, $`t_+`$ decreases with the average capital and $`t_{}`$ hardly responds to that. The companies in the small scale job category (Mining for instance) have more possibilities of growing profits than the companies in the large scale job category (Finance and Insurance for instance). On the other hand, the possibility of growing down is independent of the scale of the job category for high profits companies, $`x_1x_0`$. As a result, the relation (17) explains that the Pareto index for the companies in the small scale job category is larger than that for the companies in the large scale job category.
## 5 Conclusion
In this paper, we have first shown that the Pareto index can be estimated by profits growth rate distribution of companies in Japan. The point is that the data have no threshold ($`x_0`$) in contrast to the high income data. We find that the law of detailed balance and Gibrat’s law hold in the extended region $`(x_1>0,x_2>0)`$ and $`(x_1>x_0,x_2>0)`$, respectively.
By classifying companies into 14 job divisions, we have secondly confirmed the reflection law that relates the change of the Pareto index to the change of the growth rate distribution.
We thirdly find that the companies in the small scale job category have more possibilities of growing profits than the companies in the large scale job category by observing the relation between the average capital and the slope of growth rate distribution. The relation between the average capital and the Pareto index is also verified. Consequently it is kinematically explained that the Pareto index for the companies in the small scale job category is larger than that for the companies in the large scale job category.
We have concentrated our attention to the region ($`x_1>x_0`$, $`x_2>0`$) in order to examine Gibrat’s law. Beyond this regime, the breakdown of Gibrat’s law is studied, the recent works about that are done by Stanley’s group . In this paper, we find the law of detailed balance over the regime where Pareto’s law holds. We expect that the distribution in the region ($`x_1x_0`$, $`x_2>0`$) can be derived kinematically from the law of detailed balance and the results found by Stanley’s group. This argument will kinematically decide the distribution for the middle region where Pareto’s law does not hold.
## Acknowledgements
The author is grateful to the Yukawa Institute for Theoretical Physics at Kyoto University, where this work was initiated during the YITP-W-03-03 on “Econophysics - Physics-based approach to Economic and Social phenomena -”, and specially to Professor H. Aoyama for useful comments. Thanks are also due to Professor H. Kasama for careful reading of the manuscript.
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# Cepheid pulsation models at varying metallicity and Δ𝑌/Δ𝑍
## 1 Introduction
Classical Cepheids are the most reliable primary distance indicators for Local Group and (from the space) external galaxies, thanks to their characteristic period-luminosity-color (PLC) and period-luminosity (PL) relations. Moreover, through the calibration of secondary distance indicators, they allow us to reach cosmological distances (of the order of 100 Mpc), thus providing fundamental constraints on the Hubble constant (see e.g. Freedman et al. 2001, hereinafter F01; Saha et al. 2001). The problem of the dependence of the Cepheid PL relation on chemical composition has been widely debated in the recent literature but with quite different results, depending on the adopted method and the authors (see e.g. Kennicutt et al. 1998, Fiorentino et al. 2002, Storm et al. 2004, Groenewegen et al. 2004, Sakai et al. 2004, Romaniello et al. 2005). In the last few years we have studied the Cepheid pulsation properties through the computation of nonlinear, nonlocal and time-dependent convective pulsation models, which allow to predict all the relevant pulsational observables. In particular, the nonlinearity and the inclusion of a detailed treatment of the coupling between pulsation and convection allow these models to predict not only the periods and the blue boundary of the instability strip, but also the pulsation amplitudes, the detailed light and radial velocity curve morphology and the complete topology of the strip, including the red edge which is caused by the pulsation quenching due to convection (see Bono, Marconi & Stellingwerf 1999a, 2000a and references therein for details). On this basis, various sets of Cepheid models have been computed with varying chemical composition ($`0.004<`$Z$`<0.04`$, $`0.25<`$Y$`<0.33`$) and stellar mass from 2.8 M to 11 M (Bono et al. 1999b, 2000b, 2002b; Fiorentino et al. 2002, hereinafter F02). For each chemical composition and mass, an evolutionary mass-luminosity (ML) relation was adopted (see F02 for details) and a wide range of effective temperature was explored. As a result, we have found that the Cepheid properties, and in particular the location in the HR diagram of the instability strip and the coefficients of the multiband PL relations, depend on the pulsator metallicity with the amplitude of the effect decreasing from visual to near infrared magnitudes (Caputo, Marconi & Musella, 2000a, hereinafter C00). In particular, as the model metallicity increases from $`Z=0.004`$ to $`Z=0.03`$ the instability strip gets redder and the pulsator luminosity, at fixed period, gets fainter (Bono et al. 1999b, Caputo et al. 2000b, F02). This result is at variance with recent empirical evaluations of the metallicity effect (see e.g. Kennicutt et al. 1998; F01) and relies on the assumption of $`\mathrm{\Delta }Y/\mathrm{\Delta }Z=2.5`$. Specific computations for $`\mathrm{\Delta }Y/\mathrm{\Delta }Z=4.0`$ have shown that, at least for the higher metal contents ($`Z0.008`$), the location into the HR diagram of the Cepheid instability strip also depends on helium abundance, moving toward higher effective temperature as Y increases, at fixed Z (F02). On the basis of the above chemical composition dependent models, F02 have found that the adoption of LMC based $`V`$ and $`I`$ PL relations to get distance moduli with an uncertainty of $`\pm 0.1`$ mag is justified for variables with period shorter than 10 days. At longer periods, a correction to LMC based distances maybe needed, whose sign and amount depend on the helium and metal content of the Cepheids. In particular, model predictions were found to account for the empirical metallicity correction suggested by Kennicutt et al. (1998), provided that the adopted helium-to-metal enrichment ratio was about 3.5. Moreover, the above models provide a fairly good description of the data obtained by Romaniello et al. (2005, hereinafter R05), by relating the $`V`$ band residuals from the PL relation adopted by the HST Key Project (F01) to spectroscopic iron abundances measured for 37 Galactic and Magellanic Clouds Cepheids (see R05 for details). All the above results were essentially based on 2 values (at most) of helium content for each fixed metal abundance and therefore did not allow to properly investigate the helium effect on the whole metal content range of observed Cepheids. In order to perform a more accurate analysis of combined helium and metal effects, we extended our grid of models to other chemical compositions at varying $`\mathrm{\Delta }Y/\mathrm{\Delta }Z`$. In this paper we show the results of these new computations and further discuss the effect of chemical composition on Cepheid properties. In Sect. 2 we present the new model set and combine the results with the previous ones to investigate the dependence of the predicted pulsation observables on helium abundance and metallicity. In Sect. 3 the theoretical PL, PLC and Wesenheit relations are presented, whereas in Sect. 4 and 5 we discuss the comparison with recent observational data and the implication of the predictions presented in the previous sections for the Cepheid distance scale.
## 2 The new models
By adopting the same code and physical and numerical assumptions as in previous papers (Bono et al. 1999a, F02), we have computed new sequences of pulsation models for the input parameters listed in Table 1. In particular these new models correspond to different values of the $`\mathrm{\Delta }Y/\mathrm{\Delta }Z`$ parameter, ranging from 0.5 to 3.5. The upper limit is due to the evidence that, as noticed in F02, for $`\mathrm{\Delta }Y/\mathrm{\Delta }Z=4`$ no model is found to pulsate at the highest metallicities ($`Z0.04`$) covered by Cepheids in HST galaxies. On the other hand, $`\mathrm{\Delta }Y/\mathrm{\Delta }Z=0.5`$ is below the lower limit of the evaluations found in the recent literature (see e.g. Pagel & Portinari 1998, Pagel et al. 1992, Izotov, Thuan, & Lipovetsky 1997, Izotov & Thuan 2004). For each selected mass, the luminosity level is chosen on the basis of the same canonical evolutionary ML relation adopted in F02. For a detailed discussion of the effect of the ML selection on Cepheid properties, we refer the interested reader to a companion paper by Caputo et al. (2005, hereinafter C05). A wide range of effective temperatures is explored for each model mass and the modal stability is investigated for the fundamental mode. The first overtone mode pulsation is not studied in this paper because it is not expected to be significant at the selected metallicity and mass ranges and because it is known to be almost independent on chemical composition (see Bono et al. 2001, 2002b and references therein for details).
### 2.1 The instability strip
For each chemical composition, mass and luminosity level, model computations allowed us to derive the blue (FBE) and red (FRE) edges of the fundamental instability strip. The new strips are reported in Figs. 1 and 2 together with similar evaluations from our previous model sets. In Fig. 1 we show the location of the predicted instability strip in the HR diagram at fixed $`\mathrm{\Delta }Y/\mathrm{\Delta }Z`$ and for the labeled metal contents. This plot confirms our previous results (see F02 and references therein) concerning the shift of the instability strip toward lower effective temperatures, as the metal abundance increases from $`Z=0.004`$ to $`Z=0.03`$, and the narrowing of the strip when passing from $`Z=0.03`$ to $`Z=0.04`$. The latter occurrence is due to the reduced efficiency of pulsation associated to the low Hydrogen abundance for $`Z=0.04`$ and $`Y=0.29`$ or 0.33 (corresponding to $`\mathrm{\Delta }Y/\mathrm{\Delta }Z=1.5`$ or 2.5 respectively). In order to investigate the effect of varying the helium abundance at fixed metallicity, in the three panels of Fig. 2 we show the location of the instability strip in the HR diagram for $`Z=0.02`$ (top), $`Z=0.03`$ (middle) and $`Z=0.04`$ (bottom) and the labeled helium abundances. For $`Z=0.04`$ we note a narrowing of the instability strip when Y increases (and X decreases). The effect is less evident for the other two metal abundances. In particular, for $`Z=0.02`$ the FRE moves toward higher effective temperature as Y increases from $`Y=0.25`$ to $`Y=0.31`$, confirming the result found by F02, but the helium dependence of the FBE is much more complicated with a sort of to and fro behavior on the explored Y range. This occurrence is related to the competing pulsation driving role of H and He abundances in the associated ionization zones (see Bono et al. 1999a for details).
### 2.2 The light and radial velocity curves
One of the most important output of nonlinear pulsation codes is the predicted variation of relevant quantities (luminosity, radial velocity, effective temperature, surface gravity) along a model pulsation cycle. The bolometric light curves<sup>1</sup><sup>1</sup>1The radial velocity curves are also available upon request to the authors. for the models quoted in Table 1 are reported in Figs 3a-3g for each labeled mass and luminosity level. The model period and effective temperature is also reported in each panel. We notice that both the morphology and the amplitude of the curves vary with the position within the instability strip and depend on the adopted chemical composition. This behavior confirms our previous results (see Bono, Castellani & Marconi 2000b, Bono, Marconi & Stellingwerf 2000a, hereinafter BMS00). In particular, these plots support the empirical evidence for Galactic Cepheids originally found by Sandage & Tammann (1968, 1971) and Cogan (1980), and recently confirmed on the basis of a much larger sample of Cepheids by Sandage, Tammann & Reindl (2004) and Tammann, Sandage & Reindl (2003), that, in the period range $`\mathrm{log}P0.400.86`$ and for $`\mathrm{log}P>1.11.3`$ the largest luminosity amplitudes are attained close to the blue edge, while for $`0.85<\mathrm{log}P<1.11.3`$, the maximum is attained close to the red edge as a consequence of a phenomenon called Hertzsprung progression (HP) (Hertzsprung 1926; Ledoux & Walraven 1958). Classical Cepheids in the period range $`6<P<16`$ d show a secondary maximum (bump) along both the light and the radial velocity curves. The HP is the relationship between the phase of this bump and the pulsation period. In particular, for Galactic Cepheids the bump appears on the descending branch of the light curve for Cepheids with periods up to $`9`$ d, while it appears close to maximum light for $`9<P<12`$ d and moves at earlier phases for longer periods. On the basis of this observational evidence this group of variables was christened “Bump Cepheids”. As already obtained by BMS00 for $`0.004<Z<0.02`$, inspection of Figs. 3a, 3b, 3c and 3f suggest that an increase in the metal content causes a shift of the HP center toward shorter periods. In fact, as shown more in detail in Fig. 4, passing from $`Z=0.01`$ $`Y=0.26`$ (upper panel) to $`Z=0.04`$ $`Y=0.25`$ (bottom panel), for $`M/M=7`$, the period corresponding to the HP center moves from $`10.5`$ d to $`8.2`$ d, attaining $`9.5`$ d for $`Z=0.02`$ $`Y=0.25`$,0.26 (middle panels). We remind that this trend is in agreement with the observations. In fact, empirical data for Galactic Cepheids suggest that the HP center corresponds to a period PHP$`10.0`$ d (Moskalik, Buchler & Marom 1992; Moskalik et al. 2000), whereas in the LMC ($`Z=0.008`$) PHP$`10.5`$ d and in the SMC ($`Z=0.004`$) PHP$`11.0`$ d (Beaulieu 1998). Unfortunately, no empirical evidence is available for the HP phenomenon in sovrasolar Cepheid samples. The bolometric light curves of the new models were transformed into the observational bands ($`UBVRIJK`$) by means of the model atmospheres by Castelli, Gratton & Kurucz (1997a,b) and mean magnitudes and colors were then derived for each chemical composition and stellar mass. In Table 2 we report the model periods and intensity weighted mean magnitudes, but magnitude averaged values are also available upon request to the authors. The static magnitude values<sup>2</sup><sup>2</sup>2For static magnitude we mean the magnitude the star would have were it not pulsating have also been derived and used to obtain the boundaries of the instability strip, at each chemical composition, in the various period-magnitude planes and, in turn, to construct synthetic multiband PL relations (see below).
## 3 Predicted Cepheid relations
The results presented in the previous section allow us to derive all the relevant relations connecting the pulsation period to mean magnitudes and colors, as well as synthetic PL relations, following the same procedure as in our previous papers (see Caputo et al. 2000a, hereinafter C00,F02).
### 3.1 PLC and Wesenheit relations
Linear regression through the period and magnitude values reported in Table 2 provides the multiband PLC relations given in Table 3. In the same table we also report the PLC coefficients of our previous model sets. Similarly, the coefficients of the reddening free Wesenheit relations (see C00 and references therein) are reported in Table 4. These are defined by using the ratios between total extinction and the various color excesses given by Cardelli et al. (1989). In agreement with the recent empirical evidence by Ngeow & Kanbur (2005), we find that the predicted Wesenheit relations are well represented by linear functions, at variance with PL (see below) and Period-Color (see C00) relations. We notice that, even if PLC and Wesenheit relations have the advantage of being independent on the distribution of pulsators within the instability strip, holding for each individual star as a result of its period-density relation and black body behavior, they heavily rely on the assumption of an evolutionary ML relation. Without this assumption we would obtain tight mass-dependent PLC and Wesenheit relations (see C05) which allow to provide sound constraints on the pulsation mass of each individual Cepheid, once known the absolute magnitudes and the intrinsic colors, and, by comparison with evolutionary masses, to give an estimate of mass-loss during or before the Cepheid phase. The interested reader is referred to C05 for a detailed and updated investigation of this problem. Here we only notice that if we used a non-canonical mass-luminosity relation<sup>3</sup><sup>3</sup>3By non-canonical we mean based on evolutionary models including a mild core overshooting during the hydrogen burning phase (see Chiosi, Wood & Capitanio, 1993). In this scenario the luminosity level is brighter than the canonical case by 0.25 dex. the PLC relations would provide absolute magnitudes brighter than the ones obtained with the relations reported in Table 4 by $``$ 0.2 mag.
### 3.2 Synthetic PL relations
PL relations are well known to depend on the topology of the instability region and on the distribution of pulsators within the strip. For this reason we did not use the individual models but we populated the predicted instability strip by adopting the procedure suggested by Kennicutt et al. (1998) and already used by C00 and F02. In particular, 1000 pulsators were uniformly distributed from the blue to the red boundary of the instability strip, with a mass law as given by $`dn/dm=m^3`$ over the mass range 5-11$`M_{}`$ (see C00 for further details). The resulting synthetic distributions for the new model sequences are shown in Fig. 5. Inspection of these plots confirms the evidence shown in previous papers (Bono et al. 1999b; C00) that moving toward the shorter wavelengths the $`M_\lambda \mathrm{log}P`$ distribution of fundamental pulsators with periods longer than $``$ 3 days is much better represented by a quadratic relation with a clear dependence on both metallicity and the intrinsic width of the instability strip. On the other hand $`J`$ and $`K`$ band period-magnitude distributions are remarkably narrow, linear and only slightly dependent on chemical composition. The resulting synthetic multiband ($`BVRIJK`$) PL relations are given in Table 5 and Table 6 (quadratic and linear solutions, respectively) and overplotted as solid and dashed lines in each panel of the quoted figure. However, we wish to remind that, as remarked in C00, the present solutions refer to a specific pulsator distribution and that different populations may modify the results. In particular, if the longer periods ($`\mathrm{log}P`$1.5) are rejected in the final fit, then the predicted linear $`PL`$ relations become steeper and the intrinsic dispersion in the $`BVR`$ bands is reduced (see Table 7). Such a selection was also adopted by F02 because the slope of the predicted linear PL<sub>V</sub> and PL<sub>I</sub> relations for Cepheids with period $`\mathrm{log}P`$1.5 and the metallicity of the LMC ($`Z=0.008`$) was found to be $`2.75\pm 0.02`$ and $`2.98\pm 0.01`$, respectively, in very good agreement with the values ($`2.77\pm `$0.03 and $`2.98\pm `$0.02) inferred from the huge sample of LMC Cepheids in the OGLE-II catalog (Udalski 2000).
## 4 Theory versus observations
In order to test the predictive capabilities of current pulsation models, in this section we compare the theoretical multiband PL relations with recent observations for Cepheids belonging to the Milky Way and the LMC.
In the recent papers by Tammann, Sandage & Reindl (2003, hereinafter S03) and Sandage, Tammann & Reindl (2004, hereinafter S04), accurate $`BVI`$ Period-Color (PC), PL and PLC relations are derived on the basis of large databases for Galactic and LMC Cepheids, respectively. These authors show that the PL relations for Cepheids in the Galaxy, LMC and SMC have significantly different slopes (S03), and in particular S04 report about the experimental evidence of a change of the slope (a break) of the PL relation for the LMC Cepheids near 10 days. This last result is supported by the broken PC relation for the LMC Cepheids showed in S04. Indeed, any nonlinearity of the PC relations must be reflected in the PL relation. These results represent an important tool to test the accuracy of current model predictions concerning both the nonlinearity of optical PL relations and the dependence of Cepheid properties on metal abundance.
As for the first point, the theoretical evidence for the nonlinearity of $`BVI`$ PL and PC relations was already reported in our previous papers (see e.g. Bono et al. 1999b, C00) and for this reason we usually adopt linear PL relations for $`\mathrm{log}P1.5`$ (see previous section). In particular, Fig. 4 in C00 shows that the nonlinear behavior of the PC relation is more evident for $`Z=0.004`$ and $`Z=0.008`$ (representative of the LMC and SMC metallicity respectively), whereas it is significantly reduced for $`Z=0.02`$ (representative of the metallicity of Galactic Cepheids). This result is in agreement with S03 and S04 that found the break at 10 days only for the LMC Cepheids. On this basis, in order to compare our results with the data presented by S04, we also derived the theoretical linear PL relations for $`\mathrm{log}P1.0`$ and $`\mathrm{log}P>1.0`$, for the metal abundances $`Z=0.004`$ and $`Z=0.008`$. These theoretical relations are compared with the S04 LMC sample in Fig. 6. In this plot dots represent the OGLE sample selected as in S04 and open circles the longer period Cepheids collected by S04 from several sources (see S04 for details). Solid lines show the fits for $`\mathrm{log}P<1.0`$ and $`\mathrm{log}P>1.0`$ obtained by S04 (by using both samples), whereas short and long dashed lines represent the theoretical PL relations for $`Z=0.004`$ and $`Z=0.008`$ respectively, with the same period selection. We also plotted the theoretical PL relations for $`Z=0.004`$ in order to take into account the significant metallicity dispersion of LMC Cepheids (Luck et al. 1998) and the results by Bono et al. (1999b) that at longer periods the observed distribution of LMC Cepheids in the period-magnitude diagram is better represented by the theoretical one for $`Z=0.004`$ (see Bono et al. 1999b for details). Inspection of Fig. 6 and of Tables 7 and 8 suggests that the slopes obtained for $`Z=0.004`$ and $`Z=0.008`$ in the two period ranges are similar. However, the $`B`$ and $`V`$ PL relations for $`Z=0.008`$ are systematically fainter than the observational one, which shows a better agreement with the model predictions for $`Z=0.004`$. On the other hand, in the $`I`$ band we have a very good agreement between the empirical fits and the relations for $`Z=0.008`$, whereas the relations for $`Z=0.004`$ seem to be systematically brighter.
In order to investigate the dependence of Cepheid properties on metallicity, in Fig. 7 we plot the S04 Galactic sample with our PL relations for $`Z=0.02`$ and $`Z=0.01`$. Open circles represent the S04 sample with distances from Baade-Becker-Wesselink expansion parallaxes and filled circles indicate the S04 sample with distances from cluster/associations Cepheids (see S04 for details). The solid line represents the S04 fit (obtained by using both the samples, see S04 for details), the long dashed line is our PL relation (for $`\mathrm{log}P1.5`$) for $`Z=0.01`$, $`Y=0.26`$ and the other lines are the PL relations (for $`\mathrm{log}P1.5`$) for $`Z=0.02`$ and the labeled values of the Helium content. Our theoretical relations are flatter than the S04 one, but a better agreement with the data is found when we consider the model predictions for $`Z=0.01`$, $`Y=0.26`$ and $`Z=0.02`$ and $`Y=0.25`$,0.26, at least for $`\mathrm{log}P1.5`$. For the other chemical compositions, the discrepancy is particularly evident at the longer periods and for the $`B`$ and $`V`$ bands. In this context, we remind that current comparisons rely on the assumption of static model atmospheres and that theoretical colors can be affected by systematic uncertainties. Moreover, the $`B`$ and $`V`$ PL relations are very sensitive to the topology of the instability strip and, in turn, to the adopted input physics and to the treatment of convection in the pulsation models (see below). The effects of these uncertainties are likely more important for long period expanded structures. The better agreement obtained for models with $`Z=0.01`$ supports the recent suggestions by Asplund, Grevesse & Sauval (2005) that the solar metallicity is lower ($`Z`$ 0.01) than usually adopted ($`Z=0.02`$).
Finally, we consider the interferometric results for seven Galactic Cepheids by Kervella et al. (2004a,b,c), which are a subsample of the S04 dataset, but only with $`V`$ and $`K`$ band observations. These authors have presented accurate radius and distance determinations based on interferometric measurements of the angular diameter. On this basis they have derived new period-radius, as well as $`V`$ and K band PL relations (Kervella et al. 2004b, hereinafter K04b) by assuming the slopes by Gieren et al. (1998). All these objects have metal abundance close to the solar one. For $`l`$Car the metal abundance reported by K04b is about twice the solar value, but the recent spectroscopic measurement by R05 suggests a solar value also for this object. In Fig. 8 we show the comparison between our predicted $`K`$ band (upper panel) and $`V`$ band (lower panel) PL relations (for $`\mathrm{log}P`$1.5, see Table 7) at solar metallicity with the labeled helium abundances and the interferometric results by K04b. The intrinsic dispersion of the theoretical relations is represented by the vertical error bar in the labels, whereas the solid line represents the empirical PL relation obtained by fitting the data. For $`\eta `$Aql ($`\mathrm{log}P=0.8559`$) we have reported both the determinations used by K04b. Moreover, for variable $`l`$Car, the revised magnitude values by K04c, based on a more accurate interferometric determination of radius and distance, are also reported (empty circle)<sup>4</sup><sup>4</sup>4In C05 we discussed the possibility that $`l`$Car is a peculiar variable star. Indeed, on the basis of the comparison between the pulsational and evolutionary masses, it seems to be an object on the first crossing of the instability strip.. An inspection of this figure shows that in the $`K`$ band our theoretical relations are able to reproduce the data within the errors, whereas in the $`V`$ band the predicted PL relations are fainter than the empirical one and fail to match the location of the two pulsators with the smallest error on the absolute magnitude. In Fig. 9, we also compare K04b data with the theoretical relation at $`Z=0.01`$. We notice that, as already found for the comparison with S04 data, model predictions at this lower metal abundance better reproduce the interferometric results for Galactic Cepheids.
The discrepancies found in the comparison with the observational data by S04 and K04abc can be due, at least in part, to the uncertainties still affecting the adopted theoretical scenario. In particular we remind that current models are based on specific assumptions concerning both the evolutionary ML relations (see Bono et al. 1999a, 2000c, C05 for details) and the value of the mixing length ($`\alpha `$) parameter adopted in the treatment of convection to close the system of nonlinear dynamical and convective equations (see Bono & Stellingwerf 1994, Bono et al. 1999a). For the former point we refer the interested reader to the detailed discussion by C05. As for the mixing length parameter, even if recent results based on the theoretical fitting of observed Cepheid light curves suggest that the value of the $`\alpha `$ parameter should increase when moving from the blue to the red boundary of the instability strip (see Bono, Castellani & Marconi 2002a), in agreement with recent results obtained from the modeling of RR Lyrae stars (see e.g. Di Criscienzo, Marconi & Caputo 2004), all the models presented and adopted in this paper have been computed with $`\alpha =1.5`$. Specific model sets at $`Z=0.02`$ $`Y=0.28`$ and $`Z=0.01`$ $`Y=0.26`$ computed, by increasing $`\alpha `$ from 1.5 to 1.8, show that the instability strip gets significantly narrower with the red boundary getting bluer by at least 300-400 $`K`$ and a smaller redward shift of the blue boundary. This occurrence is due to the higher sensitivity of the red part of the instability strip to the efficiency of the convective transfer. On the basis of these results and taking into account the possibility that the $`\alpha `$ value is different at the blue and red edge of the strip, we expect that the PL relation may become brighter and steeper when $`\alpha `$ increases.
## 5 Metallicity and helium effects on the predicted distance scale
In order to test the results presented by F02 concerning the combined metallicity and helium effects on the Cepheid distance scale and provide a refined theoretical correction, we applied the same procedure adopted by the quoted authors to our extended model set. In particular, we considered our models with the various chemical composition as real Cepheids at the fixed distance modulus $`\mu _0=0`$ mag. By applying the predicted linear $`V`$ and $`I`$ band PL relations with $`Z=0.008`$ and $`Y=0.25`$ for $`\mathrm{log}P1.5`$, we determined the value $`\mu _{0,0.008}`$ for all the pulsators. This method simulates the HST Key project procedure (e.g. F01) which uses observations in the two bands $`V`$ and $`I`$ and adopts the LMC PL relations as universal (F01, Udalski et al. 1999). In this context, we adopted $`\mu _V\mu _I=E(VI)`$ and $`A_I/E(VI)=1.54`$ from Cardelli et al. (1989). The derived $`\mu _{0,0.008}`$ values confirm the results by F02: a) for periods shorter than 10 days, the discrepancy between $`\mu _{0,0.008}`$ and the real value ($`\mu _0=0`$ mag) is small enough ($`<0.1`$ mag) to support the adoption of universal LMC-referenced PL linear relations; b) for periods longer than 10 days, the discrepancy is larger than $`0.1`$ mag (up to 0.3 mag for $`Z=0.02`$ and $`Y=0.28`$ and period longer than 20 days) over the range $`Z0.010.04`$, so that a correction is required. In particular a dependence on chemical composition of the form suggested by F02 for longer period Cepheids is also found on the basis of the extended model set presented in this paper. The mean correction for $`\mathrm{log}P1.0`$ is
$$c=3.642+11.511Y1.697\mathrm{log}Z+5.334Y\mathrm{log}Z$$
whereas for Cepheid samples with $`\mathrm{log}P1.3`$ the mean correction is better reproduced by:
$$c=5.894+18.141Y2.792\mathrm{log}Z+8.576Y\mathrm{log}Z$$
both with an intrinsic uncertainty of $`\pm 0.02`$mag and by assuming $`\mathrm{\Delta }Y/\mathrm{\Delta }Z>1.5`$ (see Fig. 10). For $`\mathrm{\Delta }Y/\mathrm{\Delta }Z1.5`$ the dependence is more complicated but the correction is always lower than 0.1 mag and can be neglected as in the case of shorter periods. We also remind that on the basis of current estimates in the literature, we do not expect such low values for the $`\mathrm{\Delta }Y/\mathrm{\Delta }Z`$ parameter (Izotov & Thuan, 2004 and references therein). As shown in Fig. 10, the predicted mean metallicity correction implies that pulsators get fainter as their metallicity increases until a turnover point is reached close to the solar metal abundance. Such a behavior was already mentioned to be in agreement with the recent spectroscopic results by R05 and to reproduce the empirical metallicity correction by Kennicutt et al. (1998) for $`\mathrm{\Delta }Y/\mathrm{\Delta }Z3.5`$. We also notice that the mean correction gets smaller when the lowest period of the investigated sample decreases from 20 (lower panel) to 10 days (upper panel). In particular, the mean correction for $`\mathrm{log}P1`$ is higher than 0.1 mag only at the highest metallicities ($`Z0.03`$) and $`\mathrm{\Delta }Y/\mathrm{\Delta }Z`$ values ($`3`$). On the other hand, if the Cepheid periods are longer than or equal to 20 days, the mean correction is larger than 0.1 mag on a wide range of metallicities and helium contents.
## 6 Conclusions
We have presented an extended set of nonlinear convective pulsation models at varying the metallicity and $`\mathrm{\Delta }Y/\mathrm{\Delta }Z`$ ratio. On this basis, we obtain the following main results:
1. we have confirmed our previous results concerning the shift of the instability strip toward lower effective temperatures as the metal abundance increases at fixed $`\mathrm{\Delta }Y/\mathrm{\Delta }Z`$, at least up to $`Z=0.03`$. At the same time, when passing from $`Z=0.03`$ to $`Z=0.04`$, the strip narrows due to the reduced efficiency of pulsation. The effect of variation of the helium abundance at fixed metallicity is lower than the one obtained by varying the metallicity at fixed $`\mathrm{\Delta }Y/\mathrm{\Delta }Z`$. In particular, the fundamental red edge slightly moves toward higher effective temperatures as the helium content increases, whereas the fundamental blue edge does not show a clear trend;
2. inspection of the bolometric light curves, in the period range affected by the HP phenomenon, shows that passing from $`Z=0.01`$ to $`Z=0.04`$ the period corresponding to the HP center moves from $``$10.5 d to $``$8.2 d, in agreement with the empirical evidence of a decrease of PHP as the metallicity increases and with our previous theoretical results for $`Z0.02`$;
3. a comparison with the large database of Galactic and LMC Cepheids by Sandage et al. and Tamman et al. shows that, in agreement with the conclusions of these authors, the $`BVI`$ PL relations for LMC pulsators are well reproduced by linear theoretical relations with a break at $`\mathrm{log}P=1`$. As for the dependence on metallicity, we find that our theoretical PL relations for $`Z=0.02`$ are generally flatter than the empirical ones for Galactic Cepheids, with the discrepancy increasing toward the longer periods. A good agreement is obtained when, on the basis of suggestions in the recent literature, $`Z=0.01`$ is assumed as solar metal abundance in the models;
4. a comparison with recent accurate interferometric results by Kervella et al. shows that in the $`K`$ band our theoretical period-luminosity relations are able to reproduce the data within the errors, whereas in the $`V`$ band the predicted PL relations are fainter than the empirical one. Among the possible reasons for such a discrepancy, as well as for the one quoted in the previous point, we have identified the uncertainty on the mixing length parameter adopted in the treatment of convection and on the value of the solar metallicity;
5. we have derived the theoretical correction to the distance moduli inferred with the HST Key project procedure, when the effect of metallicity and helium abundance are taken into account. We find that this effect is smaller than 10% and can be neglected for Cepheid samples with $`\mathrm{log}P<1.0`$ and for $`\mathrm{\Delta }Y/\mathrm{\Delta }Z1.5`$. For longer periods and higher $`\mathrm{\Delta }Y/\mathrm{\Delta }Z`$ values, the dependence of the mean theoretical correction on chemical composition has the same analytical form of the one found by F02, but it is shown to become important only for periods longer than 20 days.
It is a pleasure to thank Filippina Caputo for a critical reading of the manuscript, useful discussions and suggestions. We also thank G. Tammann and B. Reindl for sending us the Galactic and LMC Cepheid samples adopted in the paper by Sandage, Tammann & Reindl (2004) and an anonymous referee for his useful comments. Financial support for this study was provided by MIUR, under the scientific projects “Stellar Populations in the Local Group” (P.I.: Monica Tosi), “Continuity and Discontinuity in the Milky Way Formation” (P.I.: Raffaele Gratton) and “On the evolution of stellar systems: fundamental step toward the scientific exploitation of VST” (P.I. Massimo Capaccioli). This project made use of computational resources granted by the “Consorzio di Ricerca del Gran Sasso” according to the “Progetto 6: Calcolo Evoluto e sue applicazioni (RSV6) - Cluster C11/B”.
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# Liquid-gas Phase Transition in Strange Hadronic Matter with Weak Y-Y Interaction
## I Introduction
Strangeness opens a new dimension to nuclear physics. Nuclear system with strangeness has many implications in astrophysics and cosmology and thus arouses much research interest. Nuclear systems with strangeness can be categorized into two species: strange quark matter(or strangelet) and strange hadronic matter(SHM). Quantum chromodynamics predicts that the strangenss-rich system could be strange quark matter or strangelet, consisting of *u*, *d* and *s* quarks. Strangelet has been a hot research topic for both theoretical physics and experimental physics, but there is still no convincing experimental result to confirm its existence. The other possible existence of strangeness is SHM, in which quarks are localized within nucleons and hyperons. The research of SHM is attractive because of the complexity of its structure and ingredients, especially, the phase transition including the quark deconfinement phase transition and liquid-gas(L-G) phase transition.
This paper evolves from an attempt to study the L-G phase transition in SHM at finite temperature. In a previous paperYang:2003up , we have studied the possibility of an L-G phase transition in SHM utilizing the extended Furnstahl-Serot-Tang(FST) model. After investigating the mechanical and chemical stabilities of SHM, we found a critical pressure on the binodal surface for $`T=13MeV`$ and a limit pressure on the binodal surface for lower temperature. After examining the continuities of entropy and specific heat capacity, we found the L-G phase transition in SHM is of second order. Employing the chiral $`SU(3)`$ quark mean field model, the same conclusion has also been drawn in ref.Wang:2004gy .
Recently, Takahashi *et al.*Takahashi:2001nm observed the $`\mathrm{\Lambda }\mathrm{\Lambda }`$ energy in a $`{}_{\mathrm{\Lambda }\mathrm{\Lambda }}{}^{6}He`$ double hypernucleus. Their observation deduced the $`\mathrm{\Lambda }\mathrm{\Lambda }`$ energy to be $`\mathrm{\Delta }B_{\mathrm{\Lambda }\mathrm{\Lambda }}=1.01\pm 0.20_{0.11}^{+0.18}MeV`$, which is much smaller than the previous estimation $`\mathrm{\Delta }B_{\mathrm{\Lambda }\mathrm{\Lambda }}45MeV`$ from earlier experimentsDanysz:1963 ; Dalitz:1989 ; Prowse:1966 ; Aoki:1991ip ; Dover:1991kf . This leads to a significant decrease of potential $`V_\mathrm{\Lambda }^{(\mathrm{\Lambda })}`$ from about $`20MeV`$ to $`5MeV`$Song:2003xe ; Qian:2004td . The new data imply that the hyperon-hyperon interaction(Y-Y interaction) can be weaker than known beforeSong:2003xe . We employed two of many candidates of phenomenological SHM modelsZhang:2000ug ; Qian:2003ct ; Qian:2004td ; Ikeda:1985um ; Schaffner:1993qj ; Schaffner-Bielich:2000wj ; Schulze:1998jf ; Zhang:2003fn , namely, extended FST modelZhang:2000ug ; Qian:2003ct ; Qian:2004td and modified quark meson coupling(MQMC) modelSong:2003xe , and the weak Y-Y interaction to reexamine the stability of SHM. We came to a conclusion that, while the system with the strong Y-Y interaction and in a quite large strangeness fraction region is more deeply bound than the ordinary nuclear matter due to the opening of new degrees of freedom, the system with weak Y-Y interaction is rather loosely bound compared to the latterSong:2003xe . This conclusion is not model-dependent on the extended FST model or MQMC model and true for both zero and finite temperatureQian:2004td .
Noting that the stability affects on the L-G phase transition directly, the stability of SHM with weak Y-Y interaction is quite different from that with strong Y-Y interaction. It is of interest to investigate the influence of this difference on the L-G phase transition in SHM. To employ the extended FST model and reexamine the L-G phase transition in SHM with weak Y-Y interaction is the main purpose of this paper. We will show that there is still a second order L-G phase transition in SHM with weak Y-Y interaction, but the binodal surface, limit pressure and the equation of state are very different from those of strong Y-Y interaction.
The paper is organized as follows. In Sec. II, the extended FST model is briefly introduced and the parameters for strong and weak Y-Y interactions are given. In Sec. III we reexamine the L-G phase transition in SHM with weak Y-Y interaction. The last section is the summary of the main results.
## II The extended FST model
In refs.Zhang:2000ug ; Qian:2003ct ; Qian:2004td , we extended the original FST model to include not only nucleons and $`\sigma `$, $`\omega `$ mesons, but also $`\mathrm{\Lambda }`$, $`\mathrm{\Xi }`$ hyperons and strange mesons $`\sigma ^{}`$ and $`\varphi `$. In ref.Yang:2003up , we used this model to study the L-G phase transition in SHM. The details of this model can be found in refs.Yang:2003up ; Zhang:2000ug ; Qian:2003ct ; Qian:2004td . Hereafter, we only give brief formulae, which are necessary for discussing the L-G phase transition.
In mean-field approximation, the Lagrangian of the extended FST model is presented as follows:
$`_{MFT}`$ $`=`$ $`\overline{\psi }_N(i\gamma ^\mu _\mu g_{\omega N}\gamma ^0V_0M_N+g_{sN}\sigma _0)\psi _N`$ (1)
$`+\overline{\psi }_\mathrm{\Lambda }(i\gamma ^\mu _\mu g_{\omega \mathrm{\Lambda }}\gamma ^0V_0g_{\varphi \mathrm{\Lambda }}\gamma ^0\varphi _0M_\mathrm{\Lambda }+g_{s\mathrm{\Lambda }}\sigma _0+g_{\sigma ^{}\mathrm{\Lambda }}\sigma _0^{})\psi _\mathrm{\Lambda }`$
$`+\overline{\psi }_\mathrm{\Xi }(i\gamma ^\mu _\mu g_{\omega \mathrm{\Xi }}\gamma ^0V_0g_{\varphi \mathrm{\Xi }}\gamma ^0\varphi _0M_\mathrm{\Xi }+g_{s\mathrm{\Xi }}\sigma _0+g_{\sigma ^{}\mathrm{\Xi }}\sigma _0^{})\psi _\mathrm{\Xi }`$
$`+{\displaystyle \frac{1}{2}}\left(1+\eta {\displaystyle \frac{\sigma _0}{S_0}}\right)m_\omega ^2V_0^2+{\displaystyle \frac{1}{4!}}\zeta \left(g_{\omega N}V_0\right)^4+{\displaystyle \frac{1}{2}}m_\varphi ^2\varphi _0^2{\displaystyle \frac{1}{2}}m_\sigma ^{}^2\sigma ^^2`$
$`H_q\left(1{\displaystyle \frac{\sigma _0}{S_0}}\right)^{4/d}\left[{\displaystyle \frac{1}{d}}ln\left(1{\displaystyle \frac{\sigma _0}{S_0}}\right){\displaystyle \frac{1}{4}}\right]`$
$`g_{ij}`$ are the coupling constants of baryon $`j`$ to meson $`i`$ field.
By using the standard technique of statistical mechanics, the thermodynamic potential $`\mathrm{\Omega }`$ is obtained
$`\mathrm{\Omega }`$ $`=`$ $`V\{H_g[(1{\displaystyle \frac{\sigma _0}{S_0}})^{\frac{4}{d}}({\displaystyle \frac{1}{d}}\mathrm{ln}(1{\displaystyle \frac{\sigma _0}{S_0}}){\displaystyle \frac{1}{4}})+{\displaystyle \frac{1}{4}}]`$ (2)
$`{\displaystyle \frac{1}{2}}(1+\eta {\displaystyle \frac{\sigma _0}{S_0}})m_\omega ^2V_0^2{\displaystyle \frac{1}{4!}}\zeta (g_{\omega N}V_0)^4{\displaystyle \frac{1}{2}}m_\varphi ^2\varphi _0^2+{\displaystyle \frac{1}{2}}m_\sigma ^{}^2\sigma _0^^2\}`$
$`2k_BT\{{\displaystyle \underset{i,𝐤}{}}\mathrm{ln}[1+e^{\beta (E_i^{}(k)\nu _i}]+{\displaystyle \underset{i,𝐤}{}}\mathrm{ln}[1+e^{\beta (E_i^{}(k)+\nu _i)}]\}`$
where $`\beta `$is the inverse temperature and $`V`$ is the volume of the system
$$E_i^{}(k)=\sqrt{M_i^2+k^2}$$
(3)
The effective masses of the hyperons and nucleons are
$$M_i^{}=M_ig_{si}\sigma _0g_{\sigma ^{}i}\sigma _0^{}(i=\mathrm{\Lambda },\mathrm{\Xi }),$$
(4)
$$M_i^{}=M_ig_{si}\sigma _0(i=N)$$
(5)
The mean-field values $`\varphi _0`$, $`V_0`$ , $`\sigma _0`$ and $`\sigma _0^{}`$ are determined by the corresponding extreme conditions of the thermodynamic potential.
The baryon densities $`\rho _{Bi}`$ is given by
$$\rho _{Bi}=\psi _i^+\psi _i=\frac{g_i}{\pi ^2}𝑑kk^2\left[n_i\left(k\right)\overline{n}_i\left(k\right)\right]$$
(6)
where $`g_i=2`$ for $`i=N`$ or $`\mathrm{\Xi }`$, $`g_i=1`$ for $`i=\mathrm{\Lambda }`$. The baryon and anti-baryon distributions are, respectively, expressed as
$$n_i(k)=\{exp[\beta (E_i^{}(k)\nu _i)]+1\}^1$$
(7)
and
$$\overline{n}_i(k)=\{exp[\beta (E_i^{}(k)+\nu _i)]+1\}^1$$
(8)
$`\nu _i`$ are related to chemical potential $`\mu _i`$ by
$`\mu _N`$ $`=`$ $`\nu _N+g_{\omega N}V_0,`$
$`\mu _\mathrm{\Lambda }`$ $`=`$ $`\nu _\mathrm{\Lambda }+g_{\omega \mathrm{\Lambda }}V_0+g_{\varphi \mathrm{\Lambda }}\varphi _0,`$
$`\mu _\mathrm{\Xi }`$ $`=`$ $`\nu _\mathrm{\Xi }+g_{\omega \mathrm{\Xi }}V_0+g_{\varphi \mathrm{\Xi }}\varphi _0.`$ (9)
Since the system has equal number of protons and neutrons, and equal number of $`\mathrm{\Xi }^0`$ and $`\mathrm{\Xi }^{}`$, the chemical equilibrium conditions for the reactions $`\mathrm{\Lambda }+\mathrm{\Lambda }n+\mathrm{\Xi }^0`$ and $`\mathrm{\Lambda }+\mathrm{\Lambda }p+\mathrm{\Xi }^{}`$ are unified as
$$2\mu _\mathrm{\Lambda }=\mu _N+\mu _\mathrm{\Xi }.$$
(10)
Eq.(10) implies that only two components of $`N`$, $`\mathrm{\Lambda }`$ and $`\mathrm{\Xi }`$ are independent. The SHM is regarded as a two-component system, for instance, the nucleon component and the $`\mathrm{\Xi }`$ component. The strangeness fraction$`f_S`$ is introduced as
$$f_S\frac{\rho _{B\mathrm{\Lambda }}+2\rho _{B\mathrm{\Xi }}}{\rho _B}$$
(11)
which plays the similar role as that of the asymmetric parameter $`\alpha =(\rho _n\rho _p)/(\rho _n+\rho _p)`$ in the asymmetric nuclear matter. We can use the same method as that of refs.Yang:2003up ; Muller:1995ji ; Qian:2000fq ; Qian:2001dw ; Qian:2002rp ; Qian:2002kj to address the L-G phase transition.
Following the usual procedure of statistical physics, we can calculate the other thermodynamic quantities from thermodynamic potential $`\mathrm{\Omega }`$. For example, the pressure and entropy density are calculated by formulas $`p=\mathrm{\Omega }/V`$ and $`S/V=(\mathrm{\Omega }/(1/\beta ))_{V,\mu _i}/V=(p/(1/\beta ))_{V,\mu _i,\sigma _0,V_0,\varphi _0,\sigma _0^{}}`$. The results are expressed as below
$`p`$ $`=`$ $`{\displaystyle \underset{i}{}}{\displaystyle \frac{g_i}{6\pi ^2}}{\displaystyle 𝑑k\frac{k^4}{E_i^{}(k)}[n_i(k)+\overline{n}_i(k)]}H_q\left\{\left(1{\displaystyle \frac{\sigma _0}{S_0}}\right)^{\frac{4}{d}}\left[{\displaystyle \frac{1}{d}}ln\left(1{\displaystyle \frac{\sigma _0}{S_0}}\right){\displaystyle \frac{1}{4}}\right]+{\displaystyle \frac{1}{4}}\right\}`$ (12)
$`+{\displaystyle \frac{1}{2}}\left(1+\eta {\displaystyle \frac{\sigma _0}{S_0}}\right)m_\omega ^2V_0^2+{\displaystyle \frac{1}{4!}}\zeta g_{\omega N}^4V_0^4+{\displaystyle \frac{1}{2}}m_\varphi ^2\varphi _0^2{\displaystyle \frac{1}{2}}m_\sigma ^{}^2\sigma _0^^2,`$
$`s`$ $``$ $`S/V={\displaystyle \underset{i}{}}{\displaystyle \frac{g_i}{6\pi ^2}}{\displaystyle }dk{\displaystyle \frac{k^4}{E_i^{}(k)}}\{{\displaystyle \frac{\beta ^2(E_i^{}\left(k\right)\nu _i)\mathrm{exp}\left[\beta \left(E_i^{}(k)\nu _i\right)\right]}{\left[\mathrm{exp}\left[\beta \left(E_i^{}(k)\nu _i\right)\right]+1\right]^2}}`$ (13)
$`+{\displaystyle \frac{\beta ^2(E_i^{}\left(k\right)+\nu _i)\mathrm{exp}\left[\beta \left(E_i^{}(k)+\nu _i\right)\right]}{\left[\mathrm{exp}\left[\beta \left(E_i^{}(k)+\nu _i\right)\right]+1\right]^2}}.\}`$
In the previous calculationYang:2003up , we employed the parameter set T1 given by ref.Zhang:2000ug , and the earlier data of strong Y-Y interaction. In refs.Song:2003xe ; Qian:2004td , we have the new data of weak Y-Y interaction and adjusted the parameters from $`g_{\sigma ^{}\mathrm{\Lambda }}^2=48.31`$ and $`g_{\sigma ^{}\mathrm{\Xi }}^2=154.62`$ for strong Y-Y interaction, to $`g_{\sigma ^{}\mathrm{\Lambda }}^2=28.73`$ and $`g_{\sigma ^{}\mathrm{\Xi }}^2=129.06`$ for weak Y-Y interaction. The details of the parameters are summarized in Table 1. We will use the parameters for weak Y-Y interaction to address the L-G phase transition in the next section.
## III L-G phase transition in SHM with weak Y-Y interaction
In this section, we employ the extended FST model to investigate the L-G phase transition in the SHM with weak Y-Y interaction.
Using the formulae in Sec.II and the parameters of Table 1, we calculate the pressure and study the equation of state first. The results for pressure vs. baryon density $`\rho _B`$ curves at $`T=10MeV`$ are shown in Fig.1, with weak and strong Y-Y interactions, respectively. We can see that the equation of state for two cases are quite different in two aspects. Firstly, with weak Y-Y interaction, when $`f_S>0.7`$, the pressure increases with density monotonically and there are no mechanical unstable regions. While with strong Y-Y interaction, there are always mechanical unstable regions for each $`f_S`$. Secondly, with weak Y-Y interaction, the minimum of the pressure-density curve rises monotonically with $`f_S`$, till $`f_S>0.7`$, where there is no minimum. While with strong Y-Y interaction, the minimum drops with $`f_S`$ and reaches the lowest value on the curve of $`f_S=1.3`$. When $`f_S>1.3`$, the pressure-density curves ascend with the increase of $`f_S`$. These features are consistent with the behavior of energy-density curves and the stability of SHM in ref.Qian:2003ct and Qian:2004td .
To discuss the chemical instability of SHM with weak Y-Y interaction, we show the chemical potential isobars for nucleons and $`\mathrm{\Xi }`$ against $`f_S`$ at temperature $`T=10MeV`$ in Fig.2 for $`p=0.05,0.10,0.20,0.40,0.50`$ and $`0.60MeVfm^3`$ respectively. An inflection point could be spotted on the curve with the pressure $`p^i=0.50MeVfm^3`$. There are chemical unstable regions on curves with $`p<p^i`$. This behavior is similar to that of strong Y-Y interaction but the inflection points are different for these two cases.
The chemical, thermal and mechanical Gibbs equilibrium conditions for two separate phases require
$$T^L=T^G=T$$
(14)
$$\mu _q^L(T,\rho ^L,f_s^L)=\mu _q^G(T,\rho ^G,f_s^G),(q=N,\mathrm{\Xi }),$$
(15)
$$p^L(T,\rho ^L,f_s^L)=p^G(T,\rho ^G,f_s^G),$$
(16)
where the superscripts $`L`$, $`G`$ denote the liquid and gas phases, respectively. We follow the procedures used in refs.Muller:1995ji ; Qian:2000fq ; Qian:2001dw ; Qian:2002rp ; Yang:2003up to solve out Eqs.(14),(15),(16). Collecting up all the solutions and we get the phase boundary, i.e., binodal surface. The binodal surface at $`T=10MeV`$ are shown in Fig.3, where the dotted line refers to weak Y-Y interaction and the solid line to strong Y-Y interaction. We see from Fig.3 that the differences between two binodal surfaces are very remarkable. Though the limit pressures, above which the L-G phase transition cannot take place, exist for both two cases, the values are very different. The limit pressure for weak Y-Y interaction $`p_{lim}^W=0.50MeVfm^3`$ is much higher than that of the strong Y-Y interaction as $`p_{lim}^S=0.095MeVfm^3`$. The enclosed region of binodal surface with weak Y-Y interaction is much larger than that of the strong Y-Y interaction. A higher limit pressure and a larger region of phase separation imply that the system is likely to be less stable and the L-G phase transition can take place more easily. This result is of course very reasonable if we notice the conclusion given by refsSong:2003xe ; Qian:2004td that the SHM with weak Y-Y interaction is less stable.
To determine the order of the L-G phase transition, we examine the entropy density and the specific heat capacity of SHM with weak Y-Y interaction. By Ehrenfest’s definition, the first order phase transition is characterized by the discontinuities of the first order derivatives of the chemical potential, such as the discontinuities of entropy and volume, while the second order phase transition unfolds the discontinuous behavior for the second order derivatives of the chemical potential, such as the specific heat capacity. The entropy per baryon is
$$s(T,p,f_s)=\frac{S(T,p,f_s)}{\rho _B}.$$
(17)
where $`S(T,p,f_s)`$ can be calculated from Eq.(13).
The specific heat capacity is
$`C_p`$ $`=`$ $`T({\displaystyle \frac{S}{T}})_{p,f_s}`$ (18)
$`=T{\displaystyle \underset{i}{}}{\displaystyle \frac{g_i}{6\pi ^2}}{\displaystyle 𝑑k\frac{k^4}{E_i^{}(k)}\left[\frac{d^2n_i}{dT^2}+\frac{d^2\overline{n}_i}{dT^2}\right]},`$
where
$`{\displaystyle \frac{d^2n_i}{dT^2}}={\displaystyle \frac{e^{\beta \left(E_i^{}(k)\nu _i\right)}\left[(E_i^{}(k)\nu _i)^2(e^{\beta \left(E_i^{}(k)\nu _i\right)}1)2T(E_i^{}(k)\nu _i)(e^{\beta \left(E_i^{}(k)\nu _i\right)}+1)\right]}{T^4(e^{\beta \left(E_i^{}(k)\nu _i\right)}+1)^3}},`$ (19)
$`{\displaystyle \frac{d^2\overline{n}_i}{dT^2}}={\displaystyle \frac{e^{\beta \left(E_i^{}(k)+\nu _i\right)}\left[(E_i^{}(k)+\nu _i)^2(e^{\beta \left(E_i^{}(k)+\nu _i\right)}1)2T(E_i^{}(k)+\nu _i)(e^{\beta \left(E_i^{}(k)+\nu _i\right)}+1)\right]}{T^4(e^{\beta \left(E_i^{}(k)+\nu _i\right)}+1)^3}}.`$ (20)
The entropy vs. temperature curves for $`f_S=0.0,0.1`$ and $`0.4`$ at $`p=0.11MeVfm^3`$ are shown in Fig.4. We see from Fig.4 that the entropy evolves continuously through the phase transition when $`f_S>0.0`$. But when $`f_S=0.0`$, our model reduces to the one-component symmetric nuclear matter. The entropy becomes discontinuous. The L-G phase transition becomes the first order. This result is consistent with theoretical and experimental results about L-G phase transition in symmetric nuclear matter.
The specific heat capacity vs. temperature curves for $`f_S=0.1`$ and $`0.4`$ at $`p=0.11MeVfm^3`$ are shown in Fig.5. The specific heat capacity curves are discontinuous for $`f_S=0.1`$ and $`0.4`$. The continuity of entropy together with the discontinuity of specific heat capacity demonstrates that the L-G phase transition in SHM with weak Y-Y interaction is still of second order, as expected.
One of the important signals for L-G phase transition is the investigation of Caloric curvesKolomietz:2001gd ; Ma:2003dc ; Ma:2004ey ; Sil:2004bu . The Caloric curve reflects the relationship between excited energy of hadronic matter and the temperature. The energy of SHM is
$`\epsilon `$ $`=`$ $`{\displaystyle \underset{i}{}}{\displaystyle \frac{g_i}{\pi ^2}}{\displaystyle 𝑑kk^2E_i^{}(k)[n_i(k)+\overline{n}_i(k)]}+H_q\left\{\left(1{\displaystyle \frac{\sigma _0}{S_0}}\right)^{\frac{4}{d}}\left[{\displaystyle \frac{1}{d}}ln\left(1{\displaystyle \frac{\sigma _0}{S_0}}\right){\displaystyle \frac{1}{4}}\right]+{\displaystyle \frac{1}{4}}\right\}`$ (21)
$`{\displaystyle \frac{1}{2}}\left(1+\eta {\displaystyle \frac{\sigma _0}{S_0}}\right)m_\omega ^2V_0^2{\displaystyle \frac{1}{4!}}\zeta g_{\omega N}^4V_0^4{\displaystyle \frac{1}{2}}m_\varphi ^2\varphi _0^2+{\displaystyle \frac{1}{2}}m_\sigma ^{}^2\sigma _0^^2`$
$`+g_{\omega N}V_0\rho _{BN}+(g_{\omega \mathrm{\Lambda }}V_0+g_{\varphi \mathrm{\Lambda }}\varphi _0)\rho _{B\mathrm{\Lambda }}+(g_{\omega \mathrm{\Xi }}V_0+g_{\varphi \mathrm{\Xi }}\varphi _0)\rho _{B\mathrm{\Xi }}`$
and the excited energy $`E^{}/\rho _B`$ isKolomietz:2001gd
$`E^{}/\rho _B={\displaystyle \frac{\lambda \epsilon ^L(\rho ^L,fs^L,T)+(1\lambda )\epsilon ^G(\rho ^G,fs^G,T)}{\lambda \rho ^L+(1\lambda )\rho ^G}}\left({\displaystyle \frac{\epsilon ^L(\rho ^L,fs^L,T)}{\rho ^L}}\right)_{T=0}`$ (22)
We present the Caloric curves at pressure $`p=0.11MeVfm^3`$ for $`f_S=0.0,0.1`$ and $`0.4`$ in Fig.6. There are slopes during the L-G phase transition for $`f_S=0.1`$ and $`0.4`$, while for $`f_S=0.0`$ there is a platform. This platform-shaped Caloric curve has been got in experiments of the L-G phase transition in normal symmetric nuclear matterMa:2003dc ; Ma:2004ey . For $`f_S>0.0`$, the slopes mean that temperature does not remain constant during phase transition. This is the very feature of second order phase transition as was first pointed out by Müller and SerotMuller:1995ji .
To exhibit the behavior of Caloric curves of SHM with strong and weak Y-Y interactions clearly, we plot the temperature $`T`$ vs. $`E^{}/\rho _B`$ curves at $`p=0.11MeVfm^3`$, $`f_S=0.4`$ for strong and weak Y-Y interactions in Fig.7. We find the temperature region of phase transition for weak Y-Y interaction is lower than that of strong Y-Y interaction and the slope for weak Y-Y interaction is more acclivitous. Perhaps these distinctions can provide us a possibility to check the Y-Y interaction be strong or weak, and then the $`\mathrm{\Lambda }\mathrm{\Lambda }`$ energy.
## IV Summary
We have used the new parameters about the weak Y-Y interaction derived from recent experiment data to reexamine the L-G phase transition in SHM. We have employed the extended FST model with nucleons and $`\mathrm{\Lambda }`$, $`\mathrm{\Xi }`$ hyperons and simplified the SHM into a two-component system to discuss its L-G phase transition. We come to the following conclusions:
1. No matter the Y-Y interaction is strong or weak, the L-G phase transition can take place in SHM. The L-G phase transition is of second order for $`f_S0`$, the entropy is continuous and the heate capacity is discontinuous at the phase transition point. But for $`f_S=0`$ the L-G phase transition becomes the first order. The entropy is discontinuous.
2. Due to the difference between the stability of SHM with strong and weak Y-Y interactions, the physical features of L-G phase transition, including the limit pressure, the regions enclosed by the binodal surface, the Caloric curves, the slope of phase transition curves are very different for weak and strong Y-Y interactions.
3. After reexaming the features of SHM, we come to a conclusion that the L-G phase transition can occur in SHM with weak Y-Y interaction more easily than that with strong Y-Y interaction.
## V Acknowledgements
This work is supported in part by National Natural Science Foundation of China under Nos. 10235030, 10247001, 10375013, 10347107 10047005, 10075071, the National Basic Research Programme 2003CB716300, the Foundation of Education Ministry of China under contract 2003246005 and CAS Knowledge Innovation Project No. KJCX2-N11.
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# Untitled Document
On one inverse problem relatively domain for the plate
Gasimov Y.S.
Institute of Applied mathematics Baku State University, Z. Khalilov 23,
AZ1148 Baku, Azerbaijan, e-mail: .ysfgasimov@yahoo.com
Abstract. Inverse problem relatively domain for the plate under across vibrations is considered. The definition of s-functions is interoduced. The construction for defining of the domain of the plate by given s-functions is offered.
Plates are elements of the constructions, which are widely used in various technical solutions. In this connection investigation of the different characteristics of the plates is one of the actual problems of the optimal projecting theory . We consider the problem of finding of the form of the plate under across vibrations by given characteristics of the system- inverse problem relatively domain.
In traditional inverse spectral problems by given some experimental data (scattering data, normalizing numbers etc.) the potential is reconstructed or the necessary and sufficient conditions are proved providing unequivocal determination of the seek functions .
In differ from this, inverse spectral problem relatively domain has another specification. Firstly, these problems require to find rather a function but domain. Secondly, the choice of data (results of the observation), sufficient for reconstruction of the domain is also enough difficult problem \[3-4\].
Let $`D`$ be the domain of the plate with boundary $`S_D`$.
It is known that the function $`\omega (x_1x_2,t)`$ describing across vibrations of the plate satisfies equation
$$\omega _{x_1x_1x_1x_1}+2\omega _{x_1x_1x_2x_2}+\omega _{x_2x_2x_2x_2}+\omega _{tt}=0.$$
(1)
Assuming the process stabilized the solution - eigen-vibration is seeking as
$$\omega (x_1,x_2,t)=u(x_1,x_2)\mathrm{cos}\lambda t,$$
where $`\lambda `$-is an eigen-frequency.
Substituting to (1) we get
$$\mathrm{\Delta }^2=\lambda u,xD,$$
(2)
where $`\mathrm{\Delta }^2=\mathrm{\Delta }\mathrm{\Delta }`$, $`\mathrm{\Delta }`$\- Laplace operator.
For different cases different boundary conditions may be put. The object under investigation is the freezed plate with boundary conditions
$$u=0,\frac{u}{n}=0,xS_D,$$
(3)
here $`DR^2`$ -bounded convex domain with boundary $`S_DC^2`$. It is known , that eigen-frequency $`\lambda _j`$ \- is positive and may be numbered as $`0\lambda _1\lambda _2\mathrm{}.`$ The set of all convex bounded domains $`DE^2`$ we denote by $`M`$. Let
$$K=\{DM:S_D\dot{C}^2\},$$
where $`\dot{C}^2`$ is a class of the piece-wise twice continuous differentiable functions.
$`K`$ may be defined by various ways, for example, by fixing of area or length of boundary, or by condition $`D_1DD_2`$, where $`D_1,D_2M`$ are given domains. The problem is: To find a domain $`DK`$ , such that
$$\frac{\left|\mathrm{\Delta }u_j(x)\right|^2}{\lambda _j}=s_j(x),xS_D,j=1,2,\mathrm{},$$
(4)
where $`u_j\left(x\right)`$ and $`\lambda _j`$ are eigen-vibration and eigen-frequency of the problem (2)-(3) in the domain D correspondingly, $`s_j\left(x\right)`$-given functions. Note that $`s_j\left(x\right),j=1,2,\mathrm{}`$ we call s-functions of the problem (2),(3).
In the following formula is obtained for the eigenfrequency of the freezed plate under across vibrations
$$\lambda _j=\frac{1}{4}\underset{u_j}{max}\underset{S_D}{}\left|\mathrm{\Delta }u_j(\xi )\right|^2P_D(n(\xi ))𝑑s,$$
(5)
where $`P_D\left(x\right)=\underset{lD}{max}(l,x),xE^n`$is a support function of $`D`$, and $`max`$ is taken over all eigen-vibrations $`u_j`$ corresponding to eigen-frequency $`\lambda _j`$ in the case of its multiplicity. As we see, the boundary values of the function $`\left|\mathrm{\Delta }u_j\left(x\right)\right|^2`$ unequivocally define $`\lambda _j`$. From (5) considering (4) we get
$$\underset{S_D}{}S_j\left(\xi \right)P_D\left(n\left(x\right)\right)𝑑s=4,j=1,2,\mathrm{}$$
(6)
This relation is basic for solving of the considering problem by given $`s`$-functions.
Now we prove the lemma that will be used in further.
Lemma1. Let $`f(x)`$ be continuous function on $`S_B`$. Then for any
$$D_1,D_2K$$
$$\underset{S_{D_1+D_2}}{}f(n(\xi ))𝑑\xi =\underset{S_{D_1}}{}f(n(\xi ))𝑑\xi +\underset{S_{D_2}}{}f(n(\xi ))𝑑\xi ,$$
(7)
where $`D_1+D_2`$ is taken in the sence of Minkovsky i.e. $`D_1+D{}_{2}{}^{}=\{x:x=x_1+x_2,x{}_{1}{}^{}D_1,x_2D_2\}`$, B-unit sphere.
Proof. It is known, that $`f\left(x\right)`$ may be continuously, positively defined extended over all the space and presented as a limit of the difference of two convex functions
$$f\left(x\right)=\underset{n\mathrm{}}{lim}\left[g_n\left(x\right)h_n\left(x\right)\right].$$
(8)
Not disrupting integrity we can write
$$f\left(x\right)=g\left(x\right)h\left(x\right),$$
(9)
where $`g\left(x\right),h\left(x\right)`$ are convex, positively defined functions.
It is known that for any continuous, convex, positively-defined function $`P(x)`$ there exists the only convex bounded domain $`D`$ such, that $`P\left(x\right)`$ is a support function of $`D`$, i.e. $`P\left(x\right)=P_D\left(x\right)`$. The opposed statement is also true. It is also known that $`D`$ is a subdifferential of its support function at the point $`x=0`$
$$D=P\left(0\right)=\{lE^n:P\left(x\right)(l,x),xE^n\}.$$
So there exist the domains $`G`$and $`H`$ such that
$$g(x)=P_G(x),h(x)=P_H(x),xB$$
(10)
Considering (9), (10) we get
$$\begin{array}{c}\underset{S_{D_1+D_2}}{}f(n(x))ds=\underset{S_{D_1+D_2}}{}[g(n(x))h(n(x)))ds]=\hfill \\ =\underset{S_{D_1+D_2}}{}P_G(n(x))𝑑s\underset{S_{D_1+D_2}}{}P_H(n(x))𝑑s.\hfill \end{array}$$
(11)
As for any $`D_1`$,$`D_2`$
$$\underset{S_{D_1}}{}P_{D_2}(n(x))ds=\underset{S_{D_2}}{}P_{D_1}(n(x)))ds$$
(12)
from (11) one may obtain $`\underset{S_{D_1+D_2}}{}f(n(x))𝑑s=\underset{S_G}{}P_{D_1+D_2}(n(x))𝑑s\underset{S_H}{}P_{D_1+D_2}(n(x))𝑑s`$.
As $`P_{D_1+D_2}\left(x\right)=P_{D_1}\left(x\right)+P_{D_2}\left(x\right)`$ applying (12) again we get (7).
Lemma is proved.
Now we investigate the main problem of the work- reconstraction of $`D`$ by given $`s`$\- functions.
Let $`BE^2`$ and $`S_B`$\- its boundary. By $`\phi _k\left(x\right),k=1,2,..`$ we denote some basis in $`C\left(S_B\right)`$-space of continuous functions on $`S_B`$. These functions may be continuously, positive homogeneously extended to . It may be done as:
$$\stackrel{~}{\phi }_k(x)=\{\begin{array}{c}\phi _k\left(\frac{x}{x}\right)x,xB,x0,\hfill \\ 0,x=0.\hfill \end{array}$$
(13)
One may test, that these functions are continuous and satisfy the positive homogeneity condition
$$\stackrel{~}{\phi }_j(\alpha x)=\alpha \stackrel{~}{\phi }_k(x),\alpha >0$$
Not disrupting integrity we can denote $`\stackrel{~}{\phi }_j\left(x\right)`$ by $`\phi _k\left(x\right)`$.
Thus we obtain the system of continuous, positive homogeneous functions defined on $`B`$.
As we noted above each positive homogeneous, continuous function $`\phi _j\left(x\right)`$may be presented in the form
$$\phi _k(x)=\underset{n\mathrm{}}{lim}\left[g_n^k(x)h_n^k(x)\right]$$
(14)
and there exist satisfying above mentioned properties domains $`G_n^k`$ and $`H_n^k`$ such, that
$$g_n^k\left(x\right)=P_{G_n^k}\left(x\right),h_n^k(x)=P_{H_n^k}(x).$$
(15)
These domains we call basic domains. Substituting these into (14) we get
$$\phi _k\left(x\right)=\underset{n\mathrm{}}{lim}\left[P_{G_n^k}^k\left(x\right)P_{H_n^k}^k\left(x\right)\right].$$
(16)
For similarity let’s suppose that
$$\phi _k\left(x\right)=P_{G^k}\left(x\right)P_{H^k}\left(x\right),$$
(17)
where $`G^k`$and $`H^k`$ are closed, bounded convex domains.
As $`n(x)S_B`$, for any $`xS_D`$, we can decomposite $`P_D\left(x\right)`$, $`xS_B`$ by basic functions $`\phi _k\left(x\right)`$
$$P_D(x)=\underset{k=1}{\overset{\mathrm{}}{}}\alpha _k\phi _k(x),xS_B,\alpha R.$$
(18)
Considering (17) from this one may get
$$P_D\left(x\right)=\underset{k=1}{\overset{\mathrm{}}{}}\alpha _k\left(P_{G^k}\left(x\right)P_{H^k}\left(x\right)\right),xS_B.$$
(19)
The set of all indexes for which $`\alpha _k0\left(\alpha _k<0\right)`$ denote by $`I^+\left(I^{}\right)`$.
Then the relation (19) may be written as
$$P_D(x)=\underset{k=1}{\overset{\mathrm{}}{}}\alpha _k\phi _k(x),xS_B,\alpha R$$
(20)
From last taking into account the properties of support functions we obtain
$$D\underset{kI^{}}{}\alpha _kG^k+\underset{kI^+}{}\alpha _kH^k=\underset{kI^+}{}\alpha _kG^k\underset{kI^{}}{}\alpha _kH^k.$$
The use (20) and the lemma give
$$\begin{array}{c}\underset{S_D}{}s_j(\xi )P_D(n(\xi )d\xi +\underset{\underset{kI^{}}{}(\alpha _k)S_{G^k}}{}s_j(\xi )P_D(n(\xi ))d\xi +\hfill \\ +\underset{\underset{kI^+}{}\alpha _kS_{H^k}}{}s_j(\xi )P_D(n(\xi )d\xi =\underset{\underset{kI^+}{}\alpha _kS_{G^k}}{}s_j(\xi )P_D(n(\xi ))d\xi +\hfill \\ +\underset{\underset{kI^{}}{}(\alpha _k)S_{H^k}}{}s_j(\xi ))P_D(n(\xi ))d\xi .\hfill \end{array}$$
From this considering (7) we have
$$\begin{array}{c}\underset{S_D}{}s_j(\xi )P_D(n(\xi )d\xi =\underset{k=1}{\overset{\mathrm{}}{}}\alpha _k[\underset{S_{G^k}}{}s_j(\xi )P_D(n(\xi ))d\xi \hfill \\ \underset{S_{H^k}}{}s_j(\xi )P_D(n(\xi )d\xi ]=4.\hfill \end{array}$$
(21)
Substituting here (18) one may get
$$\underset{k,m=1}{\overset{\mathrm{}}{}}A_{k,m}\left(j\right)\alpha _k\alpha _m=4,j=1,2,\mathrm{},$$
(22)
where
$$\begin{array}{c}A_{k,m}(j)=\underset{S_{G^k}}{}s_j(x)\left[P_{G^m}(n(x))P_{H^m}(n(x))\right]𝑑s\hfill \\ \underset{S_{H^k}}{}s_j(x)\left[P_{G^m}(n(x))P_{H^m}(n(x))\right]𝑑s.\hfill \end{array}$$
The equation (22) has, generally speaking, no only solution. The function $`P_D\left(x\right)`$ is reconstructed by the help of the solution of this equation using (18).
As we noted above domain $`D`$ is unequivocally defined by its support function $`P_D\left(x\right)`$. Suppose that (22) has the only solution providing convexity of the support function of $`D`$.
Let’s show that the expressions $`\frac{\left|\mathrm{\Delta }u_j(\xi )\right|^2}{\lambda _j},j=1,2,\mathrm{},`$ for the problem (2), (3) in the reconstructed by the help of this solution using (18) domain $`D`$ indeed are $`s`$-functions. Really, if $`D`$ is a domain in which the problem (2), (3) has given by formulae (18) $`s`$-functions then decomposition $`\overline{D}`$ by formulae (18) and making above done transformations we get the equation (22) with the same coefficients.
From the assumption that this equation has the only solution, it follows $`\overline{D}=D`$.
If (22) has more than one solution then the searching domain is in among the ones, constracted by (18) using these solutions, providing convexity of $`P\left(x\right)`$.
References
1. S.H. Gould. Variational Methods for Eigenvalue Problems. Oxford Univerrsty Press, 1966, 328 p.
2. S.A. Avdonin, M.I. Belishev. Boundary control and inverse problem for non-selfadjoint Sturm-Liouville operator (BC-method), Control and Cybernetics, 25, 1996, p.429-440.
3. Pesaint Fabrisio, Zolesio Jean-Paul. Derivees par rapport au domaine des valeurs propres du Laplacien. C.r.Acad. sci. Ser,1.-1995-321, 10.,p.1337-1310.
4. J.Elschner, G.Schimdt, M.Yamamoto. An inverse problem in periodic diffractive optics: global uniqueness with a single wave number. Inverse problems, 2003, 19, p.779-787.
5. Vladimirov V.S. Equations of mathematical physics. M.: Nauka, 1988, 512 p.
6. A.A. Niftieyev, Y.S. Gasimov. Control by boundaries and eigenvalue problems with variable domain. Publishing house “BSU”, 2004, 185 p.
7. D.M. Burago, V.A.Zalgarmer. Geometrical inequalities. M.: Nauka, 1981, 400 p.
8. V.F. Demyanov, A.M. Rubinov. Basises of non-smooth analyses and quasidifferential calculas. M.: Nauka, 1990.
9. Gasimov Y.S. Niftiev A.A. On a minimization of the eigenvalues of Shrodinger operator over domains. Doclady RAS, 2001, v. 380, 3, p. 305-307.
10. Y.S. Gasimov. On some properties of the eigenvalues by the variation of the domain. Mathematical physics, analyses, geometry, 2003, v.10, 2, p.249-255.
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# Untitled Document
June 9, 2005
Gaussian Quadrature without Orthogonal Polynomials
Ilan Degani
Department of Chemical Physics
Weizmann Institute of Science, Rehovot 76100, Israel
ilan.degani@weizmann.ac.il
and
Jeremy Schiff
Department of Mathematics
Bar-Ilan University, Ramat Gan 52900, Israel
schiff@math.biu.ac.il
> Abstract
>
> A novel development is given of the theory of Gaussian quadrature, not relying on the theory of orthogonal polynomials. A method is given for computing the nodes and weights that is manifestly independent of choice of basis in the space of polynomials. This method can be extended to compute nodes and weights for Gaussian quadrature on the unit circle and Gauss type quadrature rules with some fixed nodes.
The aim of this letter is to show how the theory of Gaussian quadrature can be developed using elementary linear algebra, without relying on the theory of orthogonal polynomials. This is certainly useful from a pedagogical point of view, showing how basic results in linear algebra can have powerful applications and allowing Gaussian quadrature to be taught in introductory courses before the rather heavier topic of orthogonal polynomials. But it is also important from a more fundamental point of view. Orthogonal polynomials are an invaluable technical tool, but they are, after all, just a choice of basis in the vector space of polynomials. In this paper we give a method for computing nodes and weights that is independent of the choice of basis.
A number of the formulae in this letter can also be found in the papers of Sack and Donovan and Gautschi on the “modified moment method” for computing nodes and weights. The emphasis, however, is somewhat different: these papers focus on efficient, stable computation, while the focus here is more on the conceptual issue of basis indepedence. It also turns out that occasional benefit can be obtained by a choice of basis not covered in and . Our basis independent method for computing nodes and weights is described in theorem 3. It allows immediate generalization to the cases of Gaussian quadrature on the unit circle and Gauss type quadrature with some fixed nodes, as we shall show. Our view of quadrature formulae as being related to blinear forms allows a generalization to multiple dimensions, though this is discussed in a separate paper .
There is, of course, already a completely elementary approach to Gaussian quadrature that avoids mention of orthogonal polynomials. We seek $`x_i,w_i`$, $`i=1,\mathrm{},n+1`$, such that the quadrature formula
$$_a^bw(x)f(x)𝑑x\underset{i=1}{\overset{n+1}{}}w_if(x_i),w_i0,$$
(1)
is exact when $`f(x)`$ is any polynomial of degree not more than $`2n+1`$. Here $`w(x)`$ is a known non-negative weight function on the interval $`[a,b]`$. Substituting in turn $`f(x)=1,x,x^2,\mathrm{},x^{2n+1}`$ we see that the nodes and weights must satisfy the equations
$$\underset{i=1}{\overset{n+1}{}}w_ix_i^\alpha =m_\alpha ,\alpha =0,1,\mathrm{},2n+1,$$
(2)
where $`m_\alpha `$ is the $`\alpha `$th moment
$$m_\alpha =_a^bw(x)x^\alpha 𝑑x.$$
(3)
In principle we should be able to obtain the nodes and weights by solving the system (2) of nonlinear equations. The drawback of this approach is that there is no evident way to solve (2) except in very simple cases; it is far from clear that there even exists a solution. Also, dependence on choice of basis is hardly resolved in this approach.
A linear algebra approach to Gaussian quadrature. Let $`𝒫_n`$ denote the vector space of polynomials of degree up to $`n`$ (so $`\mathrm{dim}𝒫_n=n+1`$). Let $`w(x)`$ be a given non-negative weight function on the interval $`[a,b]`$, and write
$$<f|g>=_a^bw(x)f(x)g(x)𝑑x$$
(4)
for any polynomials $`f,g`$. For any $`n`$, $`<|>`$ determines an inner product on the space $`𝒫_n`$. Define the symmetric bilinear form $`X(,)`$ on $`𝒫_n`$ by
$$X(f,g)=_a^bw(x)xf(x)g(x)𝑑x.$$
(5)
For any symmetric bilinear form on an inner product space there exists an orthonormal basis in which the form is diagonal. In other words, we can find eigenvalues $`x_1,\mathrm{},x_{n+1}𝐑`$ and eigenvectors $`u_1(x),\mathrm{},u_{n+1}(x)𝒫_n`$ such that
$`<u_i|u_j>`$ $`=`$ $`\delta _{ij},`$ (6)
$`X(u_i,u_j)`$ $`=`$ $`x_i\delta _{ij}.`$ (7)
Theorem 1. Let $`x_1,\mathrm{},x_{n+1}`$ be the eigenvalues of the bilinear form $`X(,)`$ with corresponding orthonormal eigenvectors $`u_1(x),\mathrm{},u_{n+1}(x)`$. Then for any polynomial $`p(x)𝒫_{2n+1}`$
$$_a^bw(x)p(x)𝑑x=\underset{i=1}{\overset{n+1}{}}<1|u_i>^2p(x_i),$$
(8)
i.e. the quadrature rule
$$_a^bw(x)f(x)𝑑x\underset{i=1}{\overset{n+1}{}}w_if(x_i),$$
(9)
where $`w_i=<1|u_i>^2`$, is exact whenever $`f(x)𝒫_{2n+1}`$.
The proof proceeds via a lemma.
The $`\delta `$ Lemma. For any polynomial $`p𝒫_{n+1}`$, we have
$$<p|u_i>=<1|u_i>p(x_i).$$
(10)
In other words, up to a normalization factor, the projection of $`p(x)`$ on $`u_i(x)`$ is found by simply evaluating $`p(x)`$ at $`x=x_i`$.
Proof of The $`\delta `$ Lemma. We prove the result for $`p(x)=1,x,\mathrm{},x^{n+1}`$, the general result follows by linearity. For $`p(x)=1`$ the result is trivial. For $`j=1,\mathrm{},n+1`$ we have
$$<x^j|u_i>=_a^bw(x)x^ju_i(x)𝑑x=X(x^{j1},u_i)=x_i<x^{j1}|u_i>.$$
(11)
Iterating this procedure $`j`$ times we obtain $`<x^j|u_i>=x_i^j<1|u_i>`$, as required. $``$
Proof of Theorem 1. We prove the result for $`p(x)=1,x,\mathrm{},x^{2n+1}`$, the general result follows by linearity. For $`p(x)=1`$ we have
$$_a^bw(x)𝑑x=<1|1>=\underset{i=1}{\overset{n+1}{}}<1|u_i><u_i|1>=\underset{i=1}{\overset{n+1}{}}<1|u_i>^2$$
(12)
as required. For $`p(x)=x^j`$, $`j=1,\mathrm{},2n+1`$, we can write $`p(x)=x^{j_1+j_2+1}`$, where $`j_1,j_2\{1,2,\mathrm{},n\}`$, and thus
$`{\displaystyle _a^b}w(x)x^j𝑑x`$ $`=`$ $`X(x^{j_1},x^{j_2})`$ (13)
$`=`$ $`{\displaystyle \underset{i_1=1}{\overset{n+1}{}}}{\displaystyle \underset{i_2=1}{\overset{n+1}{}}}<x^{j_1}|u_{i_1}><x^{j_2}|u_{i_2}>X(u_{i_1},u_{i_2})`$
$`=`$ $`{\displaystyle \underset{i=1}{\overset{n+1}{}}}<1|u_i>^2x_i^j,`$
where in the last step we have used the $`\delta `$ Lemma. $``$
Theorem 1 proves the existence of a Gaussian quadrature formula for any interval $`[a,b]`$ and non-negative weight function $`w(x)`$. The proof is based just on diagonalization of the bilinear form $`X(,)`$ with no prior knowledge needed of orthogonal polynomials. The next theorem shows that some of the main results on Gaussian quadrature are easily obtained within our approach.
Theorem 2.
1) The Gaussian quadrature formula of theorem 1 is unique.
2) The nodes of the Gaussian quadrature formula lie in the interval $`[a,b]`$.
3) The nodes of the Gaussian quadrature formula are roots of any nontrivial degree $`n+1`$ polynomial that is orthogonal to all polynomials of degree at most $`n`$.
To motivate the first part of the proof we make the observation that by setting $`p(x)=u_j(x)`$ in the $`\delta `$ Lemma we obtain
$$u_j(x_i)=\frac{\delta _{ij}}{<1|u_i>},$$
(14)
implying the interesting result that the functions $`u_i(x)`$ are actually proportional to the Lagrange cardinal functions for the points $`x_1,\mathrm{},x_{n+1}`$. Note that by setting $`p(x)=u_i(x)`$ in (10) we see that $`<1|u_i>`$ cannot be zero. Note also that the last equation implies the $`x_i`$ are distinct. We mention in passing that the fact that $`u_i(x_j)`$ is proportional to $`\delta _{ij}`$ is one reason we gave the $`\delta `$ Lemma its name.
Proof. 1) Suppose $`X_i,W_i`$, $`i=1,\mathrm{},n+1`$ are nodes and weights for a Gaussian quadrature formula, i.e. for any polynomial $`p(x)𝒫_{2n+1}`$
$$_a^bw(x)p(x)𝑑x=\underset{i=1}{\overset{n+1}{}}W_ip(X_i),$$
(15)
with the $`X_i`$ distinct and $`W_i>0`$. Define the $`n`$th degree polynomials $`U_1(x),\mathrm{},U_{n+1}(x)`$ by requiring
$$U_j(X_i)=\frac{\delta _{ij}}{\sqrt{W_i}},i,j=1,\mathrm{},n+1.$$
(16)
We then have
$$<U_i|U_j>=_a^bw(x)U_i(x)U_j(x)𝑑x=\underset{k=1}{\overset{n+1}{}}W_kU_i(X_k)U_j(X_k)=\delta _{ij}$$
(17)
and
$$X\left(U_i|U_j\right)=_a^bw(x)xU_i(x)U_j(x)𝑑x=\underset{k=1}{\overset{n+1}{}}W_kX_kU_i(X_k)U_j(X_k)=X_i\delta _{ij}.$$
(18)
Thus the $`X_i`$ are necesarilly the eigenvalues of the bilinear form $`X`$, and the $`U_i`$ are the orthonormal eigenvectors. Furthermore we have
$$<1|U_i>=_a^bw(x)U_i(x)𝑑x=\underset{k=1}{\overset{n+1}{}}W_kU_i(X_k)=\sqrt{W_i},$$
(19)
so $`W_i=<1|U_i>^2`$, as before.
2) We have
$$_a^bw(x)au_i^2(x)𝑑x_a^bw(x)xu_i^2(x)𝑑x_a^bw(x)bu_i^2(x)𝑑x,$$
so $`aX(u_i,u_i)b`$ and therefore $`ax_ib`$.
3) Note that the $`\delta `$ Lemma holds for all $`p(x)𝒫_{n+1}`$. If $`p(x)`$ is a polynomial of degree $`n+1`$ that is orthogonal to all polynomials of degree at most $`n`$, then it is orthogonal to all the $`u_i(x)`$, and thus, by the $`\delta `$ Lemma, $`p(x_i)=0`$ for all $`i`$. $``$
Calculation of the nodes and weights. Having presented our nonstandard development of the theory of Gaussian quadrature, we turn our attention to the question of computing nodes and weights using an algorithm that allows for an arbitrary choice of basis in $`𝒫_n`$. To find the nodes and weights we clearly need information on the inner product $`<|>`$ and the blinear form $`X(,)`$ on $`𝒫_n`$. So assume the functions $`q_1(x),q_2(x),\mathrm{},q_{n+1}(x)`$ are a basis of $`𝒫_n`$ and that we are given the $`(n+1)\times (n+1)`$ matrices $`B,A`$ defined by
$`B_{ij}`$ $`=`$ $`<q_i|q_j>={\displaystyle _a^b}w(x)q_i(x)q_j(x)𝑑x,i,j=1,\mathrm{},n+1,`$ (20)
$`A_{ij}`$ $`=`$ $`X(q_i,q_j)={\displaystyle _a^b}w(x)xq_i(x)q_j(x)𝑑x,i,j=1,\mathrm{},n+1.`$ (21)
We are about to see that knowledge of the matrices $`B,A`$ and just one of the functions $`q_i(x)`$ is sufficient to determine the nodes $`x_i`$ and weights $`w_i`$. As a prelude, note that there exists a diagonal matrix $`D`$ and a matrix $`V`$ satisfying
$$V^TBV=I,AV=BVD.$$
(22)
This is nothing but the statement that $`X(,)`$ has orthonormal eigenvectors, written in the basis $`\{q_i\}`$. In Matlab Release 14 the matrices $`D,V`$ are computed from $`A,B`$ by the command \[V D\]=eig(A,B,’chol’).
Theorem 3. If $`D,V`$ are matrices satisfying (22), with $`D`$ diagonal, then the nodes and weights of the Gaussian quadrature formula (that is exact for polynomials in $`𝒫_{2n+1}`$) are determined by
$$x_i=D_{ii},w_i=\left(\frac{(V^1)_{ij}}{q_j(x_i)}\right)^2.$$
(23)
(The second relation holds for arbitrary $`j`$).
Proof. Since all the integrals needed to compute $`B`$ and $`A`$ are integrals of polynomials of degree up to $`2n+1`$, we have
$`B_{ij}`$ $`=`$ $`{\displaystyle \underset{k=1}{\overset{n+1}{}}}w_kq_i(x_k)q_j(x_k)={\displaystyle \underset{k=1}{\overset{n+1}{}}}Q_{ik}W_{kk}Q_{jk},`$ (24)
$`A_{ij}`$ $`=`$ $`{\displaystyle \underset{k=1}{\overset{n+1}{}}}w_kx_kq_i(x_k)q_j(x_k)={\displaystyle \underset{k=1}{\overset{n+1}{}}}Q_{ik}W_{kk}X_{kk}Q_{jk},`$ (25)
where we have introduced diagonal matrices $`X,W`$ with the nodes $`x_i`$ and the weights $`w_i`$ along their diagonals, and the matrix $`Q`$ with entries $`Q_{ik}=q_i(x_k)`$. More concisely, we have
$$B=QWQ^T,A=QWXQ^T.$$
(26)
Inserting these into the second relation in (22) we obtain
$$QWXQ^TV=QWQ^TVD,$$
(27)
and therefore
$$X=(Q^TV)D(Q^TV)^1.$$
(28)
Since diagonal matrices are conjugate if and only if they are equal, up to reordering of the diagonal entries, we deduce the first statement in the theorem. Furthermore, after reordering the nodes if necessary, we see that the matrix $`Q^TV`$ commutes with $`X`$. But $`X`$ is a diagonal matrix with distinct entries on the diagonal, so $`Q^TV`$ must also be diagonal. Write $`\sqrt{W}Q^TV=Y`$. Then $`Y^TY=V^TQWQ^TV=I`$ (by (26) and the first relation in (22)). So the diagonal entries of $`Y`$ are all plus or minus $`1`$. Writing $`\sqrt{W}Q^T=YV^1`$ we obtain $`w_iQ_{ji}^2=((V^1)_{ij})^2`$ for all $`i,j`$, giving the second statement in the theorem. $``$
Notes: 1. Under a change of the basis $`\{q_i\}`$ we will have $`QMQ`$, $`AMAM^T`$, $`BMBM^T`$, $`VM^TV`$, where $`M`$ is a suitable nonsingular matrix. The nodes and the weights, however, remain unchanged.
2. Once the nodes $`x_i`$ are known, then the weights can be obtained from the first equation in (26). This, however, requires full knowlegde of the basis $`q_i`$, whereas the formula for the weights in the theorem requires knowledge of just a single basis element.
3. The simplest cases of theorem 3 are well-known. If we choose $`q_i(x)`$ to be the orthormal basis of $`𝒫_n`$ formed by Gram-Schmidt orthonormalization of the basis $`1,x,\mathrm{},x^n`$, then $`B=I`$ and $`A`$ is tridiagonal. The method for finding nodes and weights reduces to the classic and widely used method of (see for an exposition). If we choose $`q_i(x)=x^{i1}`$ then the matrices $`A`$ and $`B`$ are Hankel matrices ($`A_{ij}`$ and $`B_{ij}`$ both depend only on $`i+j`$), and the computation of $`D`$ and $`V`$ is numerically unstable (at least for all but the smallest values of $`n`$). For more general bases the method is essentially equivalent to that of and , but the formulation here makes the issue of basis independence much clearer. The papers and focus on bases in which the $`q_i(x)`$ satisfy a recursion of the form
$$xq_i(x)=a_iq_{i+1}(x)+b_iq_i(x)+c_iq_{i1}(x).$$
(29)
This makes most of the entries of $`A`$ expressible in terms of the entries of $`B`$ and the coefficients $`a_i,b_i,c_i`$.
Example: For the weight function $`w(x)=(1+x)^1`$ on $`[0,1]`$, the construction of orthogonal polynomials by Gram-Schmidt orthonormalization of the standard basis $`1,x,\mathrm{},x^n`$ is awkward. It is much more convenient to work in a basis in which most of the elements have a factor $`(1+x)`$. As a first attempt, take
$$q_i(x)=\{\begin{array}{cc}(1+x)x^{i1}& i=1,\mathrm{},n\\ 1& i=n+1\end{array}.$$
(30)
It is straightforward to compute the necessary integrals to obtain
$$B_{ij}=\{\begin{array}{cc}\frac{1}{i+j1}+\frac{1}{i+j}& i,j=1,\mathrm{},n\\ \frac{1}{i}& i=1,\mathrm{},n,j=n+1\\ \frac{1}{j}& j=1,\mathrm{},n,i=n+1\\ \mathrm{ln}2& i=j=n+1\end{array},$$
(31)
$$A_{ij}=\{\begin{array}{cc}\frac{1}{i+j}+\frac{1}{i+j+1}& i,j=1,\mathrm{},n\\ \frac{1}{i+1}& i=1,\mathrm{},n,j=n+1\\ \frac{1}{j+1}& j=1,\mathrm{},n,i=n+1\\ 1\mathrm{ln}2& i=j=n+1\end{array}.$$
(32)
For fairly small $`n`$, it is straightforward to find $`D,V`$ and the nodes and weights using Matlab. As $`n`$ gets larger and the conditioning of the Hankel matrices $`B`$ and $`A`$ becomes a concern, it is preferable to use the basis
$$q_i(x)=\{\begin{array}{cc}(1+x)p_i(x)& i=1,\mathrm{},n\\ 1& i=n+1\end{array}$$
(33)
where the $`\{p_i(x)\}`$ are the standard orthonormal polynomials on $`[0,1]`$ (orthonormal with respect to the weight function $`1`$). Then $`B`$ and $`A`$ are tridiagonal and pentadiagonal matrices respectively.
The point in this example is that we can exploit the freedom of basis to simplify evaluation of the necesary integrals. The cost is a little more linear algebra, but this is hardly a challenge by modern standards.
Generalization 1. Gaussian Quadrature on the Unit Circle . Let $`w(\theta )`$ be a $`2\pi `$periodic non-negative weight function. Assume there exists a quadrature formula
$$\frac{1}{2\pi }_0^{2\pi }f\left(e^{i\theta }\right)w(\theta )𝑑\theta \underset{r=1}{\overset{n+1}{}}w_rf(z_r),w_r0,$$
(34)
with distinct nodes $`z_r`$, which is exact for the $`2n+2`$ functions $`f(z)=z^n,\mathrm{},z^n,z^{n+1}`$ and all their linear combinations. Let $`q_r(z)`$, $`r=1,\mathrm{},n+1`$ be a basis for the vector space of (complex) polynomials of degree at most $`n`$ in $`z`$. Define the $`(n+1)\times (n+1)`$ matrices $`B,A`$ by
$`B_{rs}`$ $`=`$ $`{\displaystyle \frac{1}{2\pi }}{\displaystyle _0^{2\pi }}w(\theta )\overline{q_r(e^{i\theta })}q_s(e^{i\theta })𝑑\theta ,r,s=1,\mathrm{},n+1,`$ (35)
$`A_{rs}`$ $`=`$ $`{\displaystyle \frac{1}{2\pi }}{\displaystyle _0^{2\pi }}w(\theta )\overline{q_r(e^{i\theta })}e^{i\theta }q_s(e^{i\theta })𝑑\theta ,r,s=1,\mathrm{},n+1.`$ (36)
Note that $`B`$ is Hermitian and positive definite, but $`A`$ need not be Hermitian. We assume the pair $`A,B`$ is diagonalizable, i.e. that there exists a diagonal matrix $`D`$ and a matrix $`V`$ such that
$$AV=BVD.$$
(37)
The matrices $`D,V`$ can be obtained in Matlab by the command \[V D\]=eig(A,B). We then have the following analog of theorem 3:
Theorem 3. The nodes and weights of the quadrature formula (34) (that is exact for $`f(z)=z^n,\mathrm{},z^n,z^{n+1}`$) are determined by
$$z_r=D_{rr},w_r=\frac{(V^1)_{rs}(BV)_{sr}}{q_s(z_r)\overline{q_s\left(\frac{1}{\overline{z_r}}\right)}}.$$
(38)
(The second relation holds for arbitrary $`s`$).
Proof: Proceeding as in the real case we obtain
$$B=QW\stackrel{~}{Q},A=QWX\stackrel{~}{Q},$$
(39)
where $`Q,\stackrel{~}{Q}`$ are matrices with entries $`Q_{rs}=q_r(z_s)`$ and $`\stackrel{~}{Q}_{rs}=\overline{q_s\left(\frac{1}{\overline{z_r}}\right)}`$, respectively, and $`X,W`$ are diagonal matrices with diagonal entries equal to the nodes and the weights, respectively. Using these relations in (37) we rapidly obtain
$$X=(\stackrel{~}{Q}V)D(\stackrel{~}{Q}V)^1,$$
(40)
giving us the first result, that $`X=D`$, and also that the matrix $`Y=\stackrel{~}{Q}V`$ is diagonal.
From $`\stackrel{~}{Q}=YV^1`$ we deduce that
$$\stackrel{~}{Q}_{rs}=Y_{rr}(V^1)_{rs},$$
(41)
for each $`r,s`$. Multiplying the first equation in (39) on the right by $`V`$, gives $`BV=QWY`$, implying
$$(BV)_{sr}=Y_{rr}w_rQ_{sr},$$
(42)
for each $`r,s`$. Finally we divide (42) by (41) to obtain the second result in the theorem. $``$
Notes: 1. The result in the theorem shows how to compute the weights given $`A,B`$ and one of the $`q_r`$. It may be easier, if all of the $`q_r`$ are known, to directly use the first formula in (39) to find the weights.
2. For each $`k`$ we have $`A𝐯_k=z_kB𝐯_k`$, where $`𝐯_k`$ denotes the vector that is the $`k`$th column of $`V`$. Thus $`𝐯_k^{}(Az_kB)𝐯_k=0`$, or
$$\frac{1}{2\pi }_0^{2\pi }(e^{i\theta }z_k)|a(\theta )|^2w(\theta )𝑑\theta =0,$$
(43)
where $`a(\theta )=_r(𝐯_k)_rq_r(e^{i\theta })`$. Since
$$\left|_0^{2\pi }e^{\sqrt{1}\theta }|a(\theta )|^2w(\theta )𝑑\theta \right|_0^{2\pi }|a(\theta )|^2w(\theta )𝑑\theta ,$$
(44)
it follows that $`|z_k|1`$, i.e. the nodes are inside the unit circle.
Example: Taking $`w(\theta )=\mathrm{sin}^2\theta `$ and using the standard basis $`\{q_r(z)\}=\{1,z,\mathrm{},z^n\}`$ we have
$`A_{rs}`$ $`=`$ $`{\displaystyle \frac{1}{2\pi }}{\displaystyle _0^{2\pi }}e^{i(rs)\theta }\mathrm{sin}^2\theta d\theta ={\displaystyle \frac{1}{4}}\left(2\delta _{rs}\delta _{rs+2}\delta _{rs2}\right),`$ (45)
$`B_{rs}`$ $`=`$ $`{\displaystyle \frac{1}{2\pi }}{\displaystyle _0^{2\pi }}e^{i(rs+1)\theta }\mathrm{sin}^2\theta d\theta ={\displaystyle \frac{1}{4}}\left(2\delta _{rs+1}\delta _{rs+3}\delta _{rs1}\right),`$ (46)
where $`\delta _r`$ is $`1`$ if $`r=0`$ and $`0`$ otherwise, and we have used $`\frac{1}{2\pi }_0^{2\pi }e^{ir\theta }𝑑\theta =\delta _r`$. With $`n=7`$ and these simple choices of $`A`$ and $`B`$, and with a little help from Matlab, theorem 3 rapidly reproduces the nodes and weights given in the third part of table 1 in .
For a further reference on the use of Gaussian quadrature for complex integrals see .
Generalization 2. Gauss-type rules with some fixed nodes. We seek nodes and weights $`x_i,w_i`$, $`i=1,\mathrm{},n+1`$, and weights $`v_\alpha `$, $`\alpha =1,\mathrm{},m`$, such that the quadrature formula
$$_a^bf(x)w(x)𝑑x\underset{\alpha =1}{\overset{m}{}}v_\alpha f(y_\alpha )+\underset{i=1}{\overset{n+1}{}}w_if(x_i),w_i,v_\alpha 0,$$
(47)
is exact when $`f(x)`$ is any polynomial of degree not more than $`2n+m+1`$. Here the $`y_\alpha `$, $`\alpha =1,\mathrm{},m`$, are given nodes in $`[a,b]`$, and we assume the $`x_i`$ and $`y_\alpha `$ are all distinct. There is no particular reason to assume $`w(x)`$ is non-negative. Let $`q_i(x)`$, $`i=1,\mathrm{},n+1`$ be an arbitrary basis for the vector space of polynomials of degree at most $`n`$. Let the $`(n+1)\times (n+1)`$ matrices $`B,A`$ be defined by
$`B_{ij}`$ $`=`$ $`{\displaystyle _a^b}w(x){\displaystyle \underset{\alpha =1}{\overset{m}{}}}(xy_\alpha )q_i(x)q_j(x)dx,i,j=1,\mathrm{},n+1,`$ (48)
$`A_{ij}`$ $`=`$ $`{\displaystyle _a^b}w(x){\displaystyle \underset{\alpha =1}{\overset{m}{}}}(xy_\alpha )xq_i(x)q_j(x)dx,i,j=1,\mathrm{},n+1.`$ (49)
Both $`B`$ and $`A`$ are symmetric, but in this generalization $`B`$ need not be positive definite (though in the important cases $`m=1`$, $`y_1=a`$ or $`b`$, and $`m=2`$, $`y_1=a`$, $`y_2=b`$, it is positive or negative definite). Assume the pair $`A,B`$ is diagonalizable, i.e. that we can find a diagonal matrix $`D`$ and a matrix $`V`$ satisfying
$$AV=BVD.$$
(50)
We then have the following:
Theorem 3<sup>′′</sup>. If the quadrature formula (47) is exact when $`f(x)`$ is any polynomial of degree not more than $`2n+m+1`$, and the pair $`A,B`$ is diagonalizable, then the nodes and weights are determined by
$$x_i=D_{ii},w_i\underset{\alpha =1}{\overset{m}{}}(x_iy_\alpha )=\frac{(V^1)_{ij}(BV)_{ji}}{q_j(x_i)^2}.$$
(51)
(The second relation holds for arbitrary $`j`$).
The proof of this is sufficiently similar to that of theorem 3 that we omit it.
Once the nodes $`x_i`$ and weights $`w_i`$ have been determined using this theorem, the weights $`v_\alpha `$ can be determined by solving the linear system arising from exactness of the quadrature formula for $`f(x)=1,x,\mathrm{},x^{m1}`$.
Acknowledgement: We thank David Tannor for repeatedly pushing us to understand the results in , which led to this work. We thank Gene Golub and Walter Gautschi for useful comments. This letter is a revised version, with additions, of a manuscript with the same title by the second author written in April 2003 but not submitted for publication.
References.
P.J.Davis and P.Rabinowitz, Methods of Numerical Integration, Academic Press (1975).
L.Darius, P.González-Vera and F.Marcellan, Gaussian Quadrature Formulae on the Unit Circle, J.Comp.Appl.Math. 140 (2002) 159-183.
I.Degani, J.Schiff and D.J.Tannor, Commuting Extensions and Cubature Formulae, Numer.Math., to appear. arXiv:math.NA/0408076.
A.S.Dickinson and P.R.Certain, Calculation of Matrix Elements for One-Dimensional Quantum-Mechanical Problems, J.Chem.Phys. 49 (1968) 4209-4211.
W.Gautschi, On the Construction of Gaussian Quadrature Rules from Modified Moments, Math.Comp. 24 (1970) 245-260.
G.H.Golub and J.H.Welsch, Calculation of Gauss Quadrature Rules, Math.Comp. 23 (1969) 221-230.
R.A.Sack and A.F.Donovan, An Algorithm for Gaussian Quadrature given Modified Moments, Numer.Math. 18 (1972) 465-478.
P.E.Saylor and D.C.Smolarski, Why Gaussian Quadrature in the Complex Plane?, Numerical Algorithms 26 (2001) 251-280.
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# Glueball spectrum based on a rigorous three-dimensional relativistic equation for two-gluon bound states I: Derivation of the relativistic equation
## I Introduction
It is a prominent feature of Quantum Chromodynamics (QCD) that besides the quark-gluon interaction, there are also the interactions between gluons. The self-interaction of gluons suggests that the gluons may form glueballs (the bound states of gluons) through their interactions . This is an essential prediction which is of decisive significance for testing QCD and, therefore, has been arising great interest in searching for the glueballs in experiment \[3-7\]. But, there have been no faithful evidence to be established so far for their existence \[8-15\]. On the other hand, the property, the mass spectrum, the production and decay of the glueballs have extensively been investigated theoretically. Many approaches were proposed to serve such investigations such as the potential model \[16-18\], the bag model \[19-20\], the sum rule , the Bethe-Salpeter (B-S) equation \[22-24\] and the lattice simulation \[25-30\]. However, the theoretical results given by different approaches are different and even contradictory with each other \[31-32\]. This situation is attributed mainly to the fact that the quark-gluon confinement has not clearly been understood so far. Commonly, it is believed that the lattice gauge approach would give a more reliable prediction because the approach is grounded on the first principle of QCD and essentially nonperterbative although certain approximations are inevitably made in practical calculations. In addition, it is widely recognized that the B-S equation which is set up on the basis of quantum field theory is a rigorous formalism for the bound state problem and suitable to study the glueballs \[23-24\]. Nevertheless, there are two difficulties in the previous application of this equation. One difficulty arises from the interaction kernel in the equation. This kernel was not given a closed expression in the past although the expression can be derived by the procedure as demonstrated in a recent publication of one author of this paper . Ordinarily, the kernel is defined by a sum of all two-particle irreducible Feynman graphs and can only be calculated by the perturbation method. Another difficulty is ascribed to the four-dimensional nature of the equation in which the relative time (or the relative energy) is unphysical and would lead to unphysical solutions . So, many efforts in the past are paid to recast the equation into a three-dimensional one in the instantaneous approximation or the quasipotential approach \[37-38\].
In this paper, we are devoted to deriving a rigorous three-dimensional relativistic equation satisfied by two-gluon bound states so as to provide a firm basis for further study. A practical application of this equation to calculate the glueball spectrum will be presented in the next paper. The distinctive features of the equation derived are as follows. (1) The equation is exactly relativistic, containing all the retardation effect in it, unlike the B-S equation given in the instantaneous approximation in which the retardation effect is completely neglected. The equation is derived in the equal-time formalism by the consideration that a bound state is space-like and can exists in the equal-time Lorentz frame. In this frame, the relativistic equation naturally becomes a three-dimensional one without loss of any rigorism. Moreover, different from the B-S equation, the three-dimensional equation derived in this paper is a standard eigen-equation of Schrödinger-type. In the position space, it appears to be first-order differential equations. In particular, the interaction kernel in the equation is given a closed expression which is derived by the procedure proposed first in Ref. for a two-fermion system and subsequently demonstrated in Ref. for quark-antiquark bound states. The kernel derived contains exactly all the interactions taking place in the bound states and is represented in terms of only a few types of Green’s functions. Such a kernel can not only be easily calculated by the perturbation method, but also is suitable for nonperturbative investigations. (2) The three-dimensional equation is established in the angular momentum representation. Since a glueball is the state of definite spin and parity, obviously, to investigate the glueballs, it is much more convenient to work in the angular momentum representation. In order to express the relativistic equation in such a representation, it is necessary to express the QCD in the same representation. This can be done by expanding the quantized gluon and quark fields in terms of the gluon multipole fields and the spherical Dirac spinors, respectively \[41-42\]. With these expansions, the vertices in the interaction Hamiltonian can all be given explicit and analytical expressions, as will be shown in the next paper. This achievement is due to that the integrals containing three and four spherical Bessel functions in the vertices are all calculated analytically and expressed explicitly. We would like to note that in comparison with the momentum representation in which every glueball state must be separately constructed according to a certain requirement for the Lorentz and CTP transformation properties and the vertices would involve the gluon polarization vectors which are not convenient to deal with, the angular momentum representation has an advantage that the glueball state can easily be written out in a consistent manner and the QCD vertices exhibit the spin structures much clearly.(3) The relativistic equation is set up by starting from the QCD with massive gluons. According to the conventional concept of QCD, in order to keep the Lagrangian to be gauge-invariant, the gluons must be massless. On the contrary, in most of the previous investigations of the glueballs, an effective gluon mass was phenomenologically introduced so as to get reasonable theoretical results . The gluon mass was supposed to be generated dynamically from the interaction with the physical vacuum of the Yang-Mills theory or through strong gluon-binding force . Apparently, these arguments would not be considered to be stringent and logically consistent with the concept of the ordinary QCD. One of the authors of this paper in his recent article gave a different reasoning that the QCD with massive gluons can, actually, be set up on the principle of gauge-invariance without the need of introducing the Higgs mechanism or the Stückelberg fields. The essential points to achieve this conclusion are: (a) The gluon fields must be viewed as a constrained system in the whole space of vector potentials and the Lorentz condition, as a necessary constraint, must be introduced from the beginning and imposed on the Lagrangian; (b) The gauge-invariance of a gauge field should be generally examined from the action of the field other than from the Lagrangian because the action is of more fundamental dynamical meaning than the Lagrangian. Particularly, for a constrained system such as the gluon field, the gauge-invariance should be seen from its action given in the physical space defined by the Lorentz condition. This concept is well-known in Mechanics; (c) In the physical space, only infinitesimal gauge transformations are possibly allowed and necessary to be considered. This fact was clarified originally in Ref.. Based on these points of view, it is easy to prove that the QCD with massive gluons is gauge-invariant. Moreover, the renormalizability and unitarity of the theory have been proved to be no problems .
The remainder of this paper is arranged as follows. In Section II, the massive QCD and its Lagrangian are briefly described. Section III is used to formulate the angular momentum representation and give the expansions for vector, spinor and ghost fields in this representation. In Section IV, the expression of QCD Hamiltonian in the angular momentum space will be described and discussed. Section V serves to derive the three-dimensional relativistic equation satisfied by the glueballs. In Section VI, we are devoted to derive a closed expression of the interaction kernel included in the relativistic equation. In the last section, some remarks will be made. In Appendix, we present a breif derivation of the spherical Dirac spinors given in the angular momentum representation.
## II QCD Lagrangian with massive gluons
In the previous attempt of building up the massive non-Abelian gauge field theory, the following massive Yang-Mills Lagrangian density was chosen to be the starting point :
$$=\frac{1}{4}F^{a\mu \nu }F_{\mu \nu }^a+\frac{1}{2}\mu ^2A^{a\mu }A_\mu ^a,$$
(1)
where $`A_\mu ^a`$ is the vector potential for a gluon field,
$$F_{\mu \nu }^a=_\mu A_\nu ^a_\nu A_\mu ^a+gf^{abc}A_\mu ^bA_\nu ^c$$
(2)
is the field strength tensor in which $`g`$ is the QCD coupling constant and $`f^{abc}`$ are the structure constants of color SU(3) group and $`\mu `$ is the gluon mass. The first term in the Lagrangian is the ordinary Yang-Mills Lagrangian which is gauge-invariant under the whole Lie group and used to determine the form of interactions among the gluon fields themselves. The second term in the Lagrangian is the mass term which is not gauge-invariant and only affects the kinematic property of the fields. The above Lagrangian itself was ever considered to give a complete description of the massive gauge field dynamics. This consideration is not correct because the Lagrangian is not only not gauge-invariant, but also contains redundant unphysical degrees of freedom which must be eliminated by introducing a suitable constraint condition. As we know, a massive gauge field has three polarization states which need only three spatial components of the four-dimensional vector potential $`A_\mu ^a`$ to describe them. In Lorentz-covariant formulation, a full vector potential $`A^\mu `$ can be split into two Lorentz-covariant parts: the transverse vector potential $`A_T^\mu `$ and the longitudinal vector potential $`A_L^\mu `$
$$A^\mu =A_T^\mu +A_L^\mu .$$
(3)
Since the Lorentz-covariant transverse vector potential $`A_T^{a\mu }`$ contains three-independent spatial components, it is sufficient to represent the polarization states of a massive gluon. Whereas, the Lorentz-covariant longitudinal vector potential $`A_L^{a\mu }`$ appears to be a redundant unphysical variable which must be constrained by introducing the Lorentz condition
$$^\mu A_\mu ^a=0,$$
(4)
whose solution is
$$A_L^{a\mu }=0.$$
(5)
With this solution, the massive Yang-Mills Lagrangian may be expressed in terms of the independent dynamical variables $`A_T^{a\mu }`$
$$=\frac{1}{4}F_T^{a\mu \nu }F_{T\mu \nu }^a+\frac{1}{2}\mu ^2A_T^{a\mu }A_{T\mu }^a,$$
(6)
which gives a complete description of the massive gluon field dynamics. If we want to represent the dynamics in the whole space of the full vector potential as described by the massive Yang-Mills Lagrangian in Eq.(2.1), the massive gluon field must be treated as a constrained system. In this case, the Lorentz condition in Eq.(2.4), as a constraint, is necessarily introduced from the onset and imposed on the Lagrangian in Eq.(2.1) so as to guarantee the redundant degrees of freedom to be eliminated from the Lagrangian. Since the action is of more dynamical significance than the Lagrangian, the gauge-invariance of QCD should generally be seen from the action given by the Lagrangian in Eq.(2.6) or the Lagrangian in Eq.(2.1) constrained by the Lorentz condition in Eq.(2.4). Under the gauge transformations
$$\delta A_\mu ^a=D_\mu ^{ab}\theta ^b,$$
(7)
where
$$D_\mu ^{ab}=\delta ^{ab}_\mu gf^{abc}A_\mu ^c.$$
(8)
Noticing the identity $`f^{abc}A^{a\mu }A_\mu ^b=0`$, it is easy to prove that the action given by the Lagrangian in Eq.(2.1) and constrained by he Lorentz condition in Eq.(2.4) is gauge-invariant,
$$\delta S=d^4x\delta =\mu ^2d^4x\theta ^a^\mu A_\mu ^a=0.$$
(9)
This suggests that the QCD with massive gluons may also be set up on the basis of gauge-invariance principle.
Now, let us briefly describe quantization of the QCD with massive gluons, This quantization was carried out by different approaches in Ref.. A simpler quantization is performed in the Lagrangian path-integral formalism by means of the Lagrange undetermined multiplier method which was shown to be equivalent to the Faddeev-Popov approach of quantization . For this quantization, it is convenient to generalize the QCD Lagrangian and the Lorentz condition to the following forms:
$`_\lambda `$ $`=`$ $`\overline{\psi }\{i\gamma ^\mu (_\mu igT^aA_\mu ^a)m\}\psi {\displaystyle \frac{1}{4}}F^{a\mu \nu }F_{\mu \nu }^a`$ ()
$`+{\displaystyle \frac{1}{2}}\mu ^2A^{a\mu }A_\mu ^a{\displaystyle \frac{1}{2}}\alpha (\lambda ^a)^2`$
and
$$^\mu A_\mu ^a+\alpha \lambda ^a=0,$$
(11)
where for completeness, the quark fields have been included in the Lagrangian in which $`\overline{\psi }`$ and $`\psi `$ stand for the quark fields, $`T^a`$ are the color matrices and $`m`$ is the quark mass, $`\lambda ^a(x)`$ are the extra functions which will be identified with the Lagrange multipliers and $`\alpha `$ is an arbitrary constant playing the role of gauge parameter. According to the general procedure for constrained systems, the constraint in Eq.(2.11) may be incorporated into the Lagrangian in Eq.(2.10) by the Lagrange multiplier method, giving a generalized Lagrangian such that
$`_\lambda `$ $`=`$ $`\overline{\psi }\{i\gamma ^\mu (_\mu igT^aA_\mu ^a)m\}\psi {\displaystyle \frac{1}{4}}F^{a\mu \nu }F_{\mu \nu }^a+{\displaystyle \frac{1}{2}}\mu ^2A^{a\mu }A_\mu ^a`$ (13)
$`{\displaystyle \frac{1}{2}}\alpha (\lambda ^a)^2+\lambda ^a(^\mu A_\mu ^a+\alpha \lambda ^a)`$
$`=`$ $`\overline{\psi }\{i\gamma ^\mu (_\mu igT^aA_\mu ^a)m\}\psi {\displaystyle \frac{1}{4}}F^{a\mu \nu }F_{\mu \nu }^a+{\displaystyle \frac{1}{2}}\mu ^2A^{a\mu }A_\mu ^a`$ (14)
$`+\lambda ^a^\mu A_\mu ^a+{\displaystyle \frac{1}{2}}\alpha (\lambda ^a)^2.`$
This Lagrangian is obviously not gauge-invariant. However, for building up a correct gauge field theory, it is necessary to require the dynamics of the system, i.e. the action given by the Lagrangian (2.12) to be invariant under the gauge transformations denoted in Eqs.(2.7) and (2.8). By this requirement, noticing the identity $`f^{abc}A^{a\mu }A_\mu ^b=0`$ and applying the constraint condition in Eq.(2.11), we find
$$\delta S_\lambda =\frac{1}{\alpha }d^4x^\nu A_\nu ^a(x)^\mu (𝒟_\mu ^{ab}(x)\theta ^b(x))=0,$$
(15)
where
$$𝒟_\mu ^{ab}(x)=\delta ^{ab}\frac{\sigma ^2}{\mathrm{}_x}_\mu ^x+D_\mu ^{ab}(x),$$
(15)
in which $`\sigma ^2=\alpha \mu ^2`$ and $`D_\mu ^{ab}(x)`$ was defined in Eq.(2.8). From equation (2.11) we see $`\frac{1}{\alpha }^\nu A_\nu ^a=\lambda ^a0`$. Therefore, to ensure the action to be gauge-invariant, the following constraint condition on the gauge group is necessary to be required
$$_x^\mu (𝒟_\mu ^{ab}(x)\theta ^b(x))=0.$$
(16)
These are the coupled equations satisfied by the parametric functions $`\theta ^a(x)`$ of the gauge group. Since the Jacobian is not singular
$$detM0,$$
(17)
where
$`M^{ab}(x,y)`$ $`=`$ $`{\displaystyle \frac{\delta (_x^\mu 𝒟_\mu ^{ac}(x)\theta ^c(x))}{\delta \theta ^b(y)}}_{\theta =0}`$ (18)
$`=`$ $`\delta ^{ab}(\mathrm{}_x+\sigma ^2)\delta ^4(xy)gf^{abc}_x^\mu (A_\mu ^c(x)\delta ^4(xy)),`$ ()
the above equations are solvable and would give a set of solutions which express the parametric functions $`\theta ^a(x)`$ as functionals of the vector potentials $`A_\mu ^a(x)`$. The constraint conditions in equation (2.15) may also be incorporated into the Lagrangian (2.12) by the Lagrange undetermined multiplier method. In doing this, it is convenient, as usually done, to introduce ghost field variables $`C^a(x)`$ in such a fashion
$$\theta ^a(x)=\xi C^a(x),$$
(19)
where $`\xi `$ is an infinitesimal Grassmann’s number. In accordance with Eq.(2.18), the constraint condition in Eq.(2.15) can be rewritten as
$$^\mu (𝒟_\mu ^{ab}C^b)=0,$$
(20)
where the number $`\xi `$ has been dropped. This constraint condition usually is called ghost equation. When the condition in Eq.(2.19) is incorporated into the Lagrangian in Eq.(2.12) by the Lagrange multiplier method, we obtain a more generalized Lagrangian as follows
$`_\lambda `$ $`=`$ $`\overline{\psi }\{i\gamma ^\mu (_\mu igT^aA_\mu ^a)m\}\psi {\displaystyle \frac{1}{4}}F^{a\mu \nu }F_{\mu \nu }^a+{\displaystyle \frac{1}{2}}\mu ^2A^{a\mu }A_\mu ^a`$ ()
$`+\lambda ^a^\mu A_\mu ^a+{\displaystyle \frac{1}{2}}\alpha (\lambda ^a)^2+\overline{C}^a^\mu (𝒟_\mu ^{ab}C^b),`$
where $`\overline{C}^a(x)`$, acting as Lagrange undetermined multipliers, are the new scalar variables conjugate to the ghost variables $`C^a(x).`$
At present, we are ready to formulate the quantization of the QCD with massive gluons. As we learn from the Lagrange undetermined multiplier method, the dynamical and constrained variables as well as the Lagrange multipliers in the Lagrangian (2.20) can all be treated as free ones, varying arbitrarily. Therefore, we are allowed to use this kind of Lagrangian to construct the generating functional of Green’s functions
$`Z[J^{a\mu },\overline{\eta },\eta ,\overline{\xi }^a,\xi ^a]`$ (22)
$`=`$ $`{\displaystyle \frac{1}{N}}{\displaystyle }D(A_\mu ^a,\overline{\psi },\psi ,\overline{C}^a,C^a,\lambda ^a)exp\{i{\displaystyle }d^4x[_\lambda (x)+J^{a\mu }(x)A_\mu ^a(x)`$ ()
$`+\overline{\psi }\eta +\overline{\eta }\psi +\overline{\xi }^a(x)C^a(x)+\overline{C}^a(x)\xi ^a(x)]\},`$
where $`D(A_\mu ^a,\mathrm{},\lambda ^a)`$ denotes the functional integration measure, $`J_\mu ^a,\overline{\eta },\eta ,\overline{\xi }^a`$ and $`\xi ^a`$ are the external sources coupled to the gluon, quark and ghost fields and $`N`$ is a normalization constant. Looking at the expression of the Lagrangian (2.20), we see, the integral over $`\lambda ^a(x)`$ is of Gaussian-type. Upon completing the calculation of this integral, we finally arrive at
$`Z[J^{a\mu },\overline{\eta },\eta ,\overline{\xi }^a,\xi ^a]`$ (24)
$`=`$ $`{\displaystyle \frac{1}{N}}{\displaystyle }D(A_\mu ^a,\overline{\psi },\psi ,\overline{C}^a,C^a,)exp\{i{\displaystyle }d^4x[_{eff}(x)`$ ()
$`+J^{a\mu }(x)A_\mu ^a(x)+\overline{\psi }\eta +\overline{\eta }\psi +\overline{\xi }^a(x)C^a(x)+\overline{C}^a(x)\xi ^a(x)]\},`$
where
$`_{eff}`$ $`=`$ $`\overline{\psi }\{i\gamma ^\mu (_\mu igT^aA_\mu ^a)m\}\psi {\displaystyle \frac{1}{4}}F^{a\mu \nu }F_{\mu \nu }^a+{\displaystyle \frac{1}{2}}\mu ^2A^{a\mu }A_\mu ^a`$ ()
$`{\displaystyle \frac{1}{2\alpha }}(^\mu A_\mu ^a)^2^\mu \overline{C}^a𝒟_\mu ^{ab}C^b`$
is the effective Lagrangian given in the general gauges.
From the generating functional shown in Eqs.(2.22) and ( 2.23), one may derive the free gluon propagator as follows
$$iD_{\mu \nu }^{cd}(k)=\frac{i\delta ^{cd}}{k^2\mu ^2+i\epsilon }[g_{\mu \nu }(1\alpha )\frac{k_\mu k_\nu }{k^2\sigma ^2+i\epsilon }].$$
(27)
It is emphasized that when the gluon mass tends to zero, this propagator together with the effective Lagrangian in Eq.(2.23) and the generating functional in Eq.(2.22) all immediately go over to the results given in the QCD with massless gluons.
We would like to point out that the above propagator may also be derived from the Lagrangian in Eq.(2.23) by the method of canonical quantization. In doing this, we need to use the Fourier representation of the free gluon field operator
$`𝐀_\mu ^c(x)`$ $`=`$ $`{\displaystyle }{\displaystyle \frac{d^3k}{(2\pi )^{3/2}}}\{{\displaystyle \frac{1}{2\omega (\stackrel{}{k})}}{\displaystyle \underset{\lambda =1}{\overset{3}{}}}[ϵ_\mu ^\lambda (\stackrel{}{k})𝐚_\lambda ^c(\stackrel{}{k})e^{ikx}+ϵ_\mu ^\lambda (\stackrel{}{k})𝐚_\lambda ^{c+}(\stackrel{}{k})e^{ikx}]`$ ()
$`{\displaystyle \frac{1}{2\omega _0(\stackrel{}{k})}}{\displaystyle \frac{\stackrel{~}{k}}{\mu }}[𝐚_0^c(\stackrel{}{k})e^{i\stackrel{~}{k}x}+𝐚_0^{c+}(\stackrel{}{k})e^{i\stackrel{~}{k}x}],`$
where $`k=(k_0,\stackrel{}{k})`$ and $`\stackrel{~}{k}=(\stackrel{~}{k}_0,\stackrel{}{k})`$ with $`k_0\omega (\stackrel{}{k})=\sqrt{\stackrel{}{k}^2+\mu ^2}`$ and $`\stackrel{~}{k}_0\omega _0(\stackrel{}{k})=\sqrt{\stackrel{}{k}^2+\sigma ^2}`$, $`ϵ_\mu ^\lambda (\stackrel{}{k})`$ are the polarization vectors satisfying the transversity condition
$$k^\mu ϵ_\mu ^\lambda (\stackrel{}{k})=0,$$
(29)
which means that only three spatial components of $`ϵ_\mu ^\lambda (\stackrel{}{k})`$ are independent and the creation and annihilation operators $`𝐚_\lambda ^c(\stackrel{}{k})`$ and $`𝐚_\lambda ^{c+}(\stackrel{}{k})`$ ($`\lambda =0,1,2,3)`$ are subject to the following commutation relations
$`[𝐚_\lambda ^c(\stackrel{}{k}),𝐚_\lambda ^{}^{d+}(\stackrel{}{k^{}})]`$ $`=`$ $`\delta ^{cd}g^{\lambda \lambda ^{}}\delta ^3(\stackrel{}{k}\stackrel{}{k^{}}),`$ (30)
$`[𝐚_\lambda ^c(\stackrel{}{k}),𝐚_\lambda ^{}^d(\stackrel{}{k^{}})]`$ $`=`$ $`[𝐚_\lambda ^{c+}(\stackrel{}{k}),𝐚_\lambda ^{}^{d+}(\stackrel{}{k^{}})]=0.`$ ()
By making use of the expression in Eq.(2.25) and the above commutation relations, it is not difficult to derive
$$iD_{\mu \nu }^{cd}(xy)=0\left|T\{𝐀_\mu ^c(x)𝐀_\nu ^d(y)\}\right|0=\frac{d^4k}{(2\pi )^4}iD_{\mu \nu }^{cd}(k)e^{ik(xy)},$$
(31)
where $`iD_{\mu \nu }^{cd}(k)`$ is just as that written in Eq.(2.24).
The propagator in Eq.(2.24) is written in arbitrary gauges. It has been proved that the S-matrices given by the QCD with massive gluons is independent of the gauge parameter $`\alpha `$ . For example, in the tree diagram approximation, noticing the transversity condition denoted in Eq.(2.26) and the on-shell property of the gluon states, it is easy to verify that the $`\alpha `$-dependent term proportional to $`k_\mu k_\nu `$ in Eq.(2.24) gives no contribution to the S-matrix elements. This fact was early pointed out in Ref.. In view of this fact, we may simply take the Feynman gauge ($`\alpha =1`$) in practical calculations. In this gauge, the effective Lagrangian in Eq.(2.23) becomes
$$=_0+_I,$$
(32)
where
$$_0=\frac{1}{2}A_\nu ^a(\mathrm{}+\mu ^2)A^{a\nu }+\overline{\psi }(i\gamma ^\mu _\mu m)\psi +\overline{C}^a(\mathrm{}+\mu ^2)C^a$$
(33)
and
$`_I`$ $`=`$ $`{\displaystyle \frac{1}{2}}gf^{abc}(_\mu A_\nu ^a_\nu A_\mu ^a)A^{b\mu }A^{c\nu }{\displaystyle \frac{1}{4}}g^2f^{abc}f^{ade}A^{b\mu }A^{c\nu }A_\mu ^dA_\nu ^e`$ ()
$`+\overline{\psi }\gamma ^\mu T^a\psi A_\mu ^a+gf^{abc}^\mu \overline{C}^aC^bA_\mu ^c`$
are the free and interaction parts of the Lagrangian respectively. Correspondingly, the gluon propagator in Eq.(2.24) is reduced to
$$iD_{\mu \nu }^{cd}(k)=\frac{i\delta ^{cd}g_{\mu \nu }}{k^2\mu ^2+i\epsilon }.$$
(35)
## III Expressions of the QCD fields in angular momentum representation
The fields in the Lagrangian in Eqs.(2.30) and (2.31) may be expressed in the angular momentum representation in terms of the eigenfunctions of total angular momenta for the vector, spinor and scalar fields. These eigenfunctions are described below. For the vector field, the complete set of the eigenfunctions were already found in the literature. They include the scalar multipole field $`A_{JM}^S(k\stackrel{}{x})`$ and the vectorial multipole fields $`\stackrel{}{A}_{JM}^M(k\stackrel{}{x}),\stackrel{}{A}_{JM}^E(k\stackrel{}{x})`$ and $`\stackrel{}{A}_{JM}^L(k\stackrel{}{x}).`$ They are displayed in the following.
$$A_{JM}^S(k\stackrel{}{x})=\sqrt{\frac{2}{\pi }}kj_J(kr)Y_{JM}(\widehat{x}),$$
(36)
$`\stackrel{}{A}_{JM}^M(k\stackrel{}{x})`$ $`=`$ $`{\displaystyle \frac{i}{\sqrt{J(J+1)}}}\widehat{L}A_{JM}^S(k\stackrel{}{x})`$ (37)
$`=`$ $`\sqrt{{\displaystyle \frac{2}{\pi }}}kj_J(kr)\stackrel{}{Y}_{JJM}(\widehat{x}),`$ ()
$`\stackrel{}{A}_{JM}^E(k\stackrel{}{x})`$ $`=`$ $`{\displaystyle \frac{1}{k\sqrt{J(J+1)}}}\times \widehat{L}A_{JM}^S(k\stackrel{}{x})`$ (38)
$`=`$ $`i\sqrt{{\displaystyle \frac{2}{\pi }}}{\displaystyle \frac{k}{\sqrt{2J+1}}}[\sqrt{J}j_{J+1}(kr)\stackrel{}{Y}_{JJ+1M}(\widehat{x})`$ (39)
$`\sqrt{J+1}j_{J1}(kr)\stackrel{}{Y}_{JJ1M}(\widehat{x})],`$
$`\stackrel{}{A}_{JM}^L(k\stackrel{}{x})`$ $`=`$ $`{\displaystyle \frac{i}{k}}A_{JM}^S(k\stackrel{}{x})`$ (40)
$`=`$ $`i\sqrt{{\displaystyle \frac{2}{\pi }}}{\displaystyle \frac{k}{\sqrt{2J+1}}}[\sqrt{J+1}j_{J+1}(kr)\stackrel{}{Y}_{JJ+1M}(\widehat{x})`$ (41)
$`+\sqrt{J}j_{J1}(kr)\stackrel{}{Y}_{JJ1M}(\widehat{x})],`$
where $`J,M`$ mark the total angular momentum and its third component of the vector field eigenfunction, $`\stackrel{}{k}`$ is the momentum of a single particle, $`k=|\stackrel{}{k}|,`$ $`j_l(kr)`$ is the $`l`$-th spherical Bessel function with $`r=\left|\stackrel{}{x}\right|`$, $`\widehat{L}=i\stackrel{}{x}\times \mathrm{}`$ is the orbital angular momentum operator and $`\stackrel{}{Y}_{JlM}(\widehat{x})`$ denotes the vectorial spherical harmonic function with total angular momentum $`JM`$ and orbital angular momentum $`l=J1,J,J+1.`$ This function is defined as
$$\stackrel{}{Y}_{JlM}(\widehat{x})=\underset{mq}{}C_{lm1q}^{JM}Y_{lm}(\widehat{x})\stackrel{}{e}_q,$$
(42)
where $`C_{lm1q}^{JM}`$ are the Clebsch-Gordan (C-G) coefficients, $`Y_{lm}(\widehat{x})`$ with $`\widehat{x}=(\theta ,\varphi )`$ are the eigenfunction of orbital angular momentum operator $`\widehat{L}`$ and $`\stackrel{}{e}_q`$ is the eigenfunction of the spin operator for a vector particle . Since the function defined by
$$\stackrel{}{A}_{JM}(\stackrel{}{x})=j_J(kr)\stackrel{}{Y}_{JlM}(\widehat{x})$$
(42)
and
$$A_{JM}^0(\stackrel{}{x})=j_J(kr)Y_{JM}(\widehat{x})$$
(43)
respectively satisfy the following equations of motion derived from the first term of the Lagrangian in Eq.(2.30)
$$(\mathrm{}+\mu ^2)\stackrel{}{A}=0$$
(44)
and
$$(\mathrm{}+\mu ^2)A^0=0$$
(45)
as we see, the functions in Eqs.(3.1)-(3.4) completely describe the eigenstates of the total angular momentum for a free massive gluon field (or for a massless gluon field when the mass $`\mu =0)`$. Furthermore, as the functions $`\stackrel{}{A}_{JM}^M(k\stackrel{}{x})`$ and $`\stackrel{}{A}_{JM}^E(k\stackrel{}{x})`$ satisfy the transversity condition
$$\mathrm{}\stackrel{}{A}=0$$
(46)
and are related to each other by
$$\times \stackrel{}{A}_{JM}^M(k\stackrel{}{x})=k\stackrel{}{A}_{JM}^E(k\stackrel{}{x}),$$
(47)
the $`\stackrel{}{A}_{JM}^M(k\stackrel{}{x})`$ is usually called transverse magnetic multipole field and $`\stackrel{}{A}_{JM}^E(k\stackrel{}{x})`$ transverse electric multipole field. While the field $`\stackrel{}{A}_{JM}^L(k\stackrel{}{x})`$ is referred to as longitudinal multipole field because it obeys the condition of longitudinal fields $`\mathrm{}\times \stackrel{}{A}=0.`$ It is easy to verify that these vectorial eigenfunctions meet the orthogonality relation
$$d^3x\stackrel{}{A}_{JM}^\lambda (k\stackrel{}{x})\stackrel{}{A}_{J^{^{}}M^{^{}}}^\lambda ^{^{}}(k^{^{}}\stackrel{}{x})=\delta _{\lambda \lambda ^{^{}}}\delta (kk^{^{}})\delta _{JJ^{^{}}}\delta _{MM^{^{}}},$$
(48)
where $`\lambda ,\lambda ^{}=M,E,L`$ label the different modes of the multipole fields. Similarly, for the scalar multipole field, we have
$$d^3xA_{JM}^s(k\stackrel{}{x})A_{J^{^{}}M^{^{}}}^s(k\stackrel{}{x})=\delta (kk^{^{}})\delta _{JJ^{^{}}}\delta _{MM^{^{}}}.$$
(49)
Therefore, the multipole fields defined in Eqs.(3.1)-(3.4) may suitably be chosen as the basis functions to establish the angular momentum representation for the gluon fields. A gluon field operator $`𝐀_\mu ^c(x)`$ may be expanded as
$$\stackrel{}{𝐀}^c(x)=\underset{\lambda JM}{}_0^{\mathrm{}}\frac{dk}{\sqrt{2\omega }}[𝐚_{JM}^{c\lambda }(k)\stackrel{}{A}_{JM}^\lambda (k\stackrel{}{x})e^{i\omega t}+𝐚_{JM}^{c\lambda ^+}(k)\stackrel{}{A}_{JM}^\lambda ^{}(k\stackrel{}{x})e^{i\omega t}],$$
(50)
where $`\lambda =M,E,L`$,
$$𝐀_0^c(x)=\underset{JM}{}_0^{\mathrm{}}\frac{dk}{\sqrt{2\omega }}[𝐚_{JM}^{cs}(k)A_{JM}^s(k\stackrel{}{x})e^{i\omega t}+𝐚_{JM}^{cs^+}(k)A_{JM}^s^{}(k\stackrel{}{x})e^{i\omega t}].$$
(51)
The creation and annihilation operators in the above expansions satisfy the following commutation relations
$`[𝐚_{JM}^{c\lambda }(k),𝐚_{J^{}M^{}}^{c^{}\lambda ^+}(k^{})]`$ $`=`$ $`\delta (kk^{})\delta _{cc^{}}\delta _{\lambda \lambda ^{}}\delta _{JJ^{}}\delta _{MM^{}},`$ (52)
$`[𝐚_{JM}^{c\lambda }(k),𝐚_{J^{}M^{}}^{c^{}\lambda ^{}}(k^{})]`$ $`=`$ $`[𝐚_{JM}^{c\lambda +}(k),𝐚_{J^{}M^{}}^{c^{}\lambda ^+}(k^{})]=0,`$ ()
$`[𝐚_{JM}^{cs}(k),𝐚_{J^{}M^{}}^{c^{}s^+}(k^{})]`$ $`=`$ $`\delta (kk^{})\delta _{cc^{}}\delta _{JJ^{}}\delta _{MM^{\prime \prime }},`$ (53)
$`[𝐚_{JM}^{cs}(k),𝐚_{J^{}M^{}}^{c^{}s^{}}(k^{})]`$ $`=`$ $`[𝐚_{JM}^{cs+}(k),𝐚_{J^{}M^{}}^{c^{}s^+}(k^{})]=0.`$ ()
From the transformation of a gluon field under the space inversion
$$P\stackrel{}{𝐀}^c(t,\stackrel{}{x})P^1=\stackrel{}{𝐀}^c(t,\stackrel{}{x}),P𝐀_o^c(t,\stackrel{}{x})P^1=𝐀_o^c(t,\stackrel{}{x}),$$
(54)
where $`P`$ is the space inversion operator and the parity of the multipole fields
$$A_{JM}^S(k\stackrel{}{x})=(1)^JA_{JM}^S(k\stackrel{}{x}),\stackrel{}{A}_{JM}^\lambda (k\stackrel{}{x})=(1)^{J+\pi _\lambda }\stackrel{}{A}_{JM}^\lambda (k\stackrel{}{x}),$$
(55)
where $`\pi _\lambda =0`$ if $`\lambda =M`$ and $`\pi _\lambda =1`$ if $`\lambda =E,L,`$ one may find from Eqs.(3.14) and (3.15) the parities of the annihilation operators
$$P𝐚_{JM}^{cs}(k)P^1=(1)^J𝐚_{JM}^{cs}(k),P𝐚_{JM}^{c\lambda }(k)P^1=(1)^{J+1+\pi _\lambda }𝐚_{JM}^{c\lambda }(k)$$
(56)
and the same parity for the creation operators. In addition, from the charge conjugation of the gluon field $`CA_\mu ^c(t,\stackrel{}{x})C^1=A_\mu ^c(t,\stackrel{}{x})`$, we obtain a minus C-parity for the operators $`𝐚_{JM}^{c\lambda }(k)`$ and $`𝐚_{JM}^{c\lambda +}(k)`$
$$C𝐚_{JM}^{c\lambda }(k)C^1=𝐚_{JM}^{c\lambda }(k),C𝐚_{JM}^{c\lambda +}(k)C^1=𝐚_{JM}^{c\lambda +}(k),$$
(57)
where $`\lambda =S,M,E,L.`$
Now let us turn to the angular momentum representation of spinor fields. As shown in Appendix, this representation may be set up by means of the spherical Dirac spinors. These spinors are shown in the following.
$$u_{JM}^\sigma (p\stackrel{}{x})=\left(\begin{array}{c}\sqrt{\frac{\epsilon +m}{2\epsilon }}u_J^\sigma (pr)\mathrm{\Omega }_{JM}^\sigma (\widehat{x})\\ \sqrt{\frac{\epsilon m}{2\epsilon }}u_J^\sigma (pr)\mathrm{\Omega }_{JM}^\sigma (\widehat{x})\end{array}\right),$$
(58)
$$v_{JM}^\sigma (p\stackrel{}{x})=(1)^{J+M+\sigma }\left(\begin{array}{c}\sqrt{\frac{\epsilon m}{2E}}u_J^\sigma (pr)\mathrm{\Omega }_{JM}^\sigma (\widehat{x})\\ \sqrt{\frac{\epsilon +m}{2\epsilon }}u_J^\sigma (pr)\mathrm{\Omega }_{JM}^\sigma (\widehat{x})\end{array}\right),$$
(59)
where $`JM`$ denote the total angular momentum and its third component of a free fermion, $`\sigma =\pm 1`$, $`\epsilon =\sqrt{p^2+m^2}`$ is the energy of the fermion in which $`m`$ is the mass and $`\stackrel{}{p}`$ is the momentum, $`p=|\stackrel{}{p}|`$, $`\mathrm{\Omega }_{JM}^\sigma (\widehat{x})`$ is the spherical harmonic spinor defined as
$$\mathrm{\Omega }_{JM}^\sigma (\widehat{x})=\left(\begin{array}{c}\sigma \sqrt{\frac{J+\sigma (M\frac{1}{2})+\frac{1}{2}}{2Jl+1}}Y_{J\frac{\sigma }{2},M\frac{1}{2}}(\widehat{x})\\ \sqrt{\frac{J\sigma (M+\frac{1}{2})+\frac{1}{2}}{2J\sigma +1}}Y_{J\frac{\sigma }{2},M+\frac{1}{2}}(\widehat{x})\end{array}\right)$$
(60)
and $`u_J^\sigma (pr)`$ is defined by
$$u_J^\sigma (pr)=i^{J\frac{\sigma }{2}}\sqrt{\frac{2}{\pi }}pj_{J\frac{\sigma }{2}}(pr).$$
(61)
The spinors $`u_{JM}^\sigma (p\stackrel{}{x})`$ and $`v_{JM}^\sigma (p\stackrel{}{x})`$ are the eigenstates of total angular momentum for a fermion. They respectively obey the free Dirac equations of positive energy state and negative energy state. Based on the above expressions of the spherical Dirac spinors, it is easy to prove the following orthonormality and completeness relations
$$\begin{array}{c}d^3xu_{JM}^{\sigma ^+}(p\stackrel{}{x})u_{J^{}M^{}}^\sigma (p^{}\stackrel{}{x})=\delta (pp^{})\delta _{JJ^{^{}}}\delta _{\sigma \sigma ^{}}\delta _{MM^{}},\hfill \\ d^3xv_{JM}^{\sigma ^+}(p\stackrel{}{x})v_{J^{}M^{}}^\sigma ^{}(p^{}\stackrel{}{x})=\delta (pp^{})\delta _{JJ^{^{}}}\delta _{\sigma \sigma ^{}}\delta _{MM^{}},\hfill \\ d^3xu_{JM}^{\sigma ^+}(p\stackrel{}{x})v_{J^{}M^{}}^\sigma ^{}(p^{}\stackrel{}{x})=d^3xv_{JM}^{\sigma ^+}(p\stackrel{}{x})u_{J^{}M^{}}^\sigma ^{}(p^{}\stackrel{}{x})=0,\hfill \\ \underset{JM\sigma }{}_0^{\mathrm{}}𝑑p[u_{JM}^\sigma (p\stackrel{}{x})u_{JM}^{\sigma ^+}(p\stackrel{}{x}^{})+v_{JM}^\sigma (p\stackrel{}{x})v_{JM}^{\sigma ^+}(p\stackrel{}{x}^{})]=\delta ^3(\stackrel{}{x}\stackrel{}{x}^{}).\hfill \end{array}$$
(62)
Clearly, the free quark field operators may be expanded in terms of the spherical Dirac spinors
$$\psi (x)=\underset{s\sigma JM}{}_0^{\mathrm{}}𝑑p[𝐛_{JM}^{s\sigma }(p)u_{JM}^\sigma (p\stackrel{}{x})e^{i\epsilon t}+𝐝_{JM}^{s\sigma ^+}(p)v_{JM}^\sigma (p\stackrel{}{x})e^{i\epsilon t}],$$
(63)
$$\overline{\psi }(x)=\underset{s\sigma JM}{}_0^{\mathrm{}}𝑑p[𝐛_{JM}^{s\sigma ^+}(p)\overline{u}_{JM}^\sigma (p\stackrel{}{x})e^{i\epsilon t}+𝐝_{JM}^{s\sigma }(p)\overline{v}_{JM}^\sigma (p\stackrel{}{x})e^{i\epsilon t}],$$
(64)
where $`s=(c,f),`$ $`c`$ and $`f`$ are the color and flavor indices for a quark,$`\overline{u}_{JM}^\sigma (p\stackrel{}{x})=u_{JM}^\sigma (p\stackrel{}{x})^+\gamma _0,\overline{v}_{JM}^\sigma (p\stackrel{}{x})=v_{JM}^\sigma (p\stackrel{}{x})^+\gamma _0`$ and $`𝐛_{JM}^{s\sigma ^+}(p),𝐝_{JM}^{s\sigma ^+}(p)`$ and $`𝐛_{JM}^{s\sigma }(p),𝐝_{JM}^{s\sigma }(p)`$ are the creation and annihilation operators. These operators satisfy the following anticommutation relations
$$\begin{array}{c}\{𝐛_{JM}^{s\sigma }(p),𝐛_{J^{}M^{}}^{s^{}\sigma ^{}+}(p^{})\}=\delta (pp^{})\delta _{ss^{^{}}}\delta _{\sigma \sigma ^{}}\delta _{JJ^{}}\delta _{MM^{}},\hfill \\ \{𝐝_{JM}^{s\sigma }(p),𝐝_{J^{}M^{}}^{s^{}\sigma ^{}+}(p^{})\}=\delta (pp^{})\delta _{ss^{^{}}}\delta _{\sigma \sigma ^{}}\delta _{JJ^{}}\delta _{MM^{}},\hfill \end{array}$$
(65)
with the other anticommutators being zero.
From the space inversion transformation of the quark field
$$P\psi (\stackrel{}{x},t)P^1=\eta _P\gamma ^0\psi (\stackrel{}{x},t)$$
(66)
and the relations
$`\gamma ^0u_{JM}^\sigma (p\stackrel{}{x})`$ $`=`$ $`(1)^{J\frac{\sigma }{2}}u_{JM}^\sigma (p\stackrel{}{x}),\gamma ^0v_{JM}^\sigma (p\stackrel{}{x})`$ (67)
$`=`$ $`(1)^{J\frac{\sigma }{2}+1}v_{JM}^\sigma (p\stackrel{}{x}),`$ ()
it is easy to find from Eqs.(3.27) and (3.28) the parity of the quark operators such that
$$\begin{array}{c}P𝐛_{lm}^{s\lambda }(p)P^1=\eta _P(1)^{J\frac{\sigma }{2}}𝐛_{lm}^{s\lambda }(p),\hfill \\ P𝐝_{lm}^{s\lambda }(p)P^1=\eta _P(1)^{J+\frac{\sigma }{2}+1}𝐝_{lm}^{s\lambda }(p),\hfill \\ P𝐛_{lm}^{s\lambda ^+}(p)P^1=\eta _P(1)^{J\frac{\sigma }{2}}𝐛_{lm}^{s\lambda ^+}(p),\hfill \\ P𝐝_{lm}^{s\lambda ^+}(p)P^1=\eta _P(1)^{J+\frac{\sigma }{2}+1}𝐝_{lm}^{s\lambda ^+}(p),\hfill \end{array}$$
(68)
where the phase factor $`\eta _P`$ usually is chosen to be $`\eta _P`$ $`=1`$. By noticing the charge conjugation transformations
$$\psi (x)^1=\eta _CC\overline{\psi }(x),$$
(69)
where $`C=i\gamma ^2\gamma ^0`$ and
$$v_{JM}^\sigma (p\stackrel{}{x})=C\overline{u}_{JM}^\sigma (p\stackrel{}{x})^T,u_{JM}^\sigma (p\stackrel{}{x})=C\overline{v}_{JM}^\sigma (p\stackrel{}{x})^T,$$
(70)
the $`C`$-parities of the creation and annihilation operators are easily found from Eqs.(3.27) and (3.28)
$$\begin{array}{c}b_{lm}^{s\lambda }(p)^1=\eta _Cd_{lm}^{s\lambda }(p),d_{lm}^{s\lambda }(p)^1=\eta _Cb_{lm}^{s\lambda }(p),\hfill \\ b_{lm}^{s\lambda ^+}(p)^1=\eta _Cd_{lm}^{s\lambda ^+}(p),d_{lm}^{s\lambda ^+}(p)^1=\eta _Cb_{lm}^{s\lambda ^+}(p).\hfill \end{array}$$
(71)
For the ghost fields, considering their scalar character, we can write their expressions in the angular momentum representation as follows
$$𝐂^a(x)=\underset{JM}{}_0^{\mathrm{}}\frac{dk}{\sqrt{2\omega }}[𝐜_{JM}^a(k)A_{JM}^s(k\stackrel{}{x})e^{i\omega t}+𝐝_{JM}^{a^+}(k)A_{JM}^s^{}(k\stackrel{}{x})e^{i\omega t}],$$
(72)
$$\overline{𝐂}^a(x)=\underset{JM}{}_0^{\mathrm{}}\frac{dk}{\sqrt{2\omega }}[𝐝_{JM}^a(k)A_{JM}^s(k\stackrel{}{x})e^{i\omega t}+𝐜_{JM}^{a^+}(k)A_{JM}^s^{}(k\stackrel{}{x})e^{i\omega t}].$$
(73)
The anticommutation relations for the ghost particle operators in the above are
$`\{𝐜_{JM}^a(k),𝐜_{J^{}M^{}}^{a^{}+}(k^{})\}`$ $`=`$ $`\delta _{aa^{}}\delta _{JJ^{}}\delta _{MM^{}}\delta (kk^{}),`$ (74)
$`\{𝐝_{JM}^a(k),𝐝_{J^{}M^{}}^{a^{}+}(k^{})\}`$ $`=`$ $`\delta _{aa^{}}\delta _{JJ^{}}\delta _{MM^{}}\delta (kk^{}),`$ ()
with the other anticommutators being vanishing.
## IV QCD Hamiltonian in angular momentum representation
In this section, we are devoted to discussing the QCD Hamiltonian in the angular momentum representation. By virtue of the expansions given in Eqs.(3.14), (3.15), (3.27), (3.28), (3.36) and (3.37), it is not difficult to formulate the Hamiltonian derived from the Lorentz-covariant Lagrangian written in Eqs.(2.30) and (2.31). However, due to presence of the field $`A^0(x)`$, the expression of the Hamiltonian is rather complicated. Considering that the field $`A^0(x)`$ is unphysical and will eventually be eliminated from the S-matrix and the B-S equation by the ghost fields as ensured by the unitarity of the theory , for the sake of simplifying the representation of Hamiltonian, we would rather to start with the Lagrangian obtained from the Lagrangian in Eqs.(2.30) and (2.31) by setting the unphysical fields $`A^0(x),\overline{C}^a(x)`$ and $`C^a(x)`$ to vanish, as was commonly done in the lattice gauge calculations \[25-30\] and similarly done in Ref. where a massive QCD Lagrangian taken in the temporal gauge was used for calculating the glueball spectrum in the framework of B-S equation. The Lagrangian we start with is
$$=_0+_I,$$
(75)
where
$`_0`$ $`=`$ $`{\displaystyle \frac{1}{2}}A_i^a(\mathrm{}+\mu ^2)A_i^a+\overline{\psi }(i\gamma ^\mu _\mu m)\psi `$ (76)
$`=`$ $`{\displaystyle \frac{1}{2}}(\stackrel{}{A_k^a})^2+{\displaystyle \frac{1}{2}}A_i^a(^2\mu ^2)A_i^a+\overline{\psi }(i\gamma ^\mu _\mu m)\psi `$ ()
and
$`_I`$ $`=`$ $`{\displaystyle \frac{1}{2}}gf^{abc}(_iA_j^a_jA_i^a)A_i^bA_j^c`$ ()
$`{\displaystyle \frac{1}{4}}g^2f^{abc}f^{ade}A_i^bA_j^cA_i^dA_j^e\overline{\psi }\gamma _iT^a\psi A_i^a.`$
By making use of the canonical variables conjugate to $`A_k^a`$, $`\mathrm{\Psi }`$ and $`\overline{\mathrm{\Psi }}`$ which are defined by
$`\mathrm{\Pi }_k^a`$ $`=`$ $`{\displaystyle \frac{}{\stackrel{}{A_k^a}}}=A_k^a,`$ (78)
$`\mathrm{\Pi }_\psi `$ $`=`$ $`{\displaystyle \frac{}{\stackrel{}{\psi }}}=i\overline{\psi }\gamma ^0,`$ (79)
$`\mathrm{\Pi }_{\overline{\psi }}`$ $`=`$ $`{\displaystyle \frac{}{\stackrel{}{\overline{\psi }}}}=0,`$ ()
we can write the Hamiltonian density as
$``$ $`=`$ $`\mathrm{\Pi }_k^a\stackrel{}{A_k^a}+\mathrm{\Pi }_\psi \stackrel{}{\psi }`$ (80)
$`=`$ $`_0+_I,`$ ()
where
$$_0=\frac{1}{2}(\stackrel{}{A_k^a})^2\frac{1}{2}A_i^a(^2\mu ^2)A_i^a+\overline{\psi }(\stackrel{}{\alpha }\mathrm{}+\beta m)\psi $$
(81)
and
$`_I`$ $`=`$ $`{\displaystyle \frac{g}{2}}f^{abc}(\stackrel{}{A^a}\times \stackrel{}{A^b})(\times \stackrel{}{A^c})`$ (83)
$`+{\displaystyle \frac{1}{8}}g^2f^{abe}f^{cde}(\stackrel{}{A^a}\times \stackrel{}{A^b})(\stackrel{}{A^c}\times \stackrel{}{A^d})`$
$`+g\overline{\psi }T^a\stackrel{}{\gamma }\stackrel{}{A^a}\psi .`$
The first two terms in Eq.(4.6) denote the energy of free gluon fields. The second term gives the energy of free quark fields. In the interaction Hamiltonian density (4.7), the first, second and third terms represent the gluon three-line vertex, the gluon four-line vertex and quark-gluon vertex respectively. It should be noted that the above Hamiltonian includes not only the transverse field $`\stackrel{}{A}_T`$ , but also the longitudinal field $`\stackrel{}{A}_L`$ whose presence can be seen from the constraint condition in Eq.(2.11) because in the case of $`A^0=0`$ and $`\alpha =1`$, the condition becomes $`\stackrel{}{A_L^c}=\lambda 0.`$This is different from the theory proposed in Ref. where only the transverse field exists.
Now let us derive the expression of the Hamiltonian given by the Hamiltonian density shown above in the angular momentum representation. First, we derive the expression of the free Hamiltonian. On substituting Eqs.(3.14),(3.15),(3.27),(3.28) into Eq.(4.6), one may get
$`H_0`$ $`=`$ $`{\displaystyle d^3x_0}`$ (84)
$`=`$ $`{\displaystyle \underset{c\lambda JM}{}}{\displaystyle _0^{\mathrm{}}}𝑑k\omega (\stackrel{}{k})𝐚_{JM}^{c\lambda ^+}(k)𝐚_{JM}^{c\lambda }(k)`$ (85)
$`+{\displaystyle \underset{s\sigma lm}{}}{\displaystyle _0^{\mathrm{}}}𝑑p\epsilon (\stackrel{}{p})[𝐛_{JM}^{s\sigma ^+}(p)𝐛_{JM}^{s\sigma }(p)+𝐝_{JM}^{s\sigma ^+}(p)𝐝_{JM}^{s\sigma }(p)].`$
In later derivations, it is convenient to introduce compact notations for various indices so as to simplify the expressions of the operators and formulas. For the gluon field, we define
$$𝐚_{JM}^{c\lambda \xi }(k)=\{\begin{array}{c}𝐚_{JM}^{c\lambda }(k)\text{ if }\xi =1\\ \text{ }𝐚_{JM}^{c\lambda ^+}(k)\text{ if }\xi =1\end{array},$$
(86)
$$\stackrel{}{A}_{JM}^{\lambda \xi }(k\stackrel{}{x})=\{\begin{array}{c}\stackrel{}{A}_{JM}^\lambda (k\stackrel{}{x})\text{ if }\xi =1\\ \text{ }\stackrel{}{A}_{JM}^\lambda ^{}(k\stackrel{}{x})\text{ if }\xi =1\end{array},$$
(86)
then, the commutators in Eqs.(3.16) and (3.17) can be unified to write as
$$[𝐚_{JM}^{c\lambda \xi }(k),𝐚_{J^{}M^{}}^{c^{}\lambda ^{}\xi ^{}}(k^{})]=\delta (kk^{})\delta _{cc^{}}\delta _{\lambda \lambda ^{}}\delta _{JJ^{}}\delta _{MM^{}}\mathrm{sin}[(\xi \xi ^{})\frac{\pi }{4}].$$
(87)
Furthermore, we define $`\alpha (c,\lambda ,J,M,k,\xi )`$ in which the $`\xi `$ will be written as $`\xi _\alpha `$ later on. In this notation, the commutator in Eq.(4.11) will be represented as
$$[𝐚_\alpha ,𝐚_\beta ]=\mathrm{}_{\alpha \beta },$$
(88)
where
$$\mathrm{}_{\alpha \beta }=\delta (kk^{})\delta _{cc^{}}\delta _{\lambda \lambda ^{}}\delta _{JJ^{}}\delta _{MM^{}}\mathrm{sin}[(\xi \xi ^{})\frac{\pi }{4}]$$
(89)
and the expansion in Eq.(3.14) becomes
$$\stackrel{}{A^c}(x)=\underset{\alpha }{}𝐚_\alpha \stackrel{}{A}_\alpha (k\stackrel{}{x})e^{i\xi _\alpha \omega _\alpha t},$$
(90)
where
$$\underset{\alpha }{}\underset{c\lambda lm}{}_0^{\mathrm{}}𝑑k,\omega _\alpha =\sqrt{k^2+\mu ^2},$$
(91)
$$\stackrel{}{A}_\alpha (k\stackrel{}{x})=\frac{1}{\sqrt{2\omega _\alpha }}\stackrel{}{A}_{JM}^{\lambda \xi }(k\stackrel{}{x}).$$
(92)
For the quark field, with the definition
$`𝐜_{lm}^{s\sigma \eta }(p)`$ $`=`$ $`\{\begin{array}{cc}𝐛_{lm}^{s\sigma }(p)& if\text{ }\eta =1\\ 𝐝_{lm}^{s\sigma ^+}(p)& if\text{ }\eta =1\end{array},`$ ()
$`w_{lm}^{\sigma \eta }(p\stackrel{}{x})`$ $`=`$ $`\{\begin{array}{cc}u_{lm}^\sigma (p\stackrel{}{x})& if\text{ }\eta =1\\ v_{lm}^\sigma (p\stackrel{}{x})& if\text{ }\eta =1\end{array},`$ ()
the anticommutators in Eq.(3.29) becomes
$$\{𝐜_{JM}^{s\sigma \eta }(p),𝐜_{J^{^{}}M^{}}^{s^{}\sigma ^{}\eta ^{}+}(p^{})\}=\delta (pp^{})\delta _{ss^{^{}}}\delta _{\sigma \sigma ^{^{}}}\delta _{JJ^{^{}}}\delta _{MM^{}}\delta _{\eta \eta ^{}}.$$
(97)
When we set $`\alpha =(s,\sigma ,l,m,p,\eta )`$ in which the $`\eta `$ will be written as $`\eta _\alpha `$ later on, the above anticommutator can be simply written as
$$\{𝐜_\alpha ,𝐜_\beta ^+\}=\delta _{\alpha \beta },$$
(98)
where
$$\delta _{\alpha \beta }\delta (pp^{})\delta _{ss^{^{}}}\delta _{\sigma \sigma ^{}}\delta _{JJ^{^{}}}\delta _{MM^{}}\delta _{\eta \eta ^{}},$$
(99)
and the expansions in Eqs.(3.27) and (3.28) will be represented by
$$\begin{array}{c}\psi (x)=\underset{\alpha }{}𝐜_\alpha w_\alpha (p\stackrel{}{x})e^{i\xi _\alpha \epsilon _\alpha t},\hfill \\ \overline{\psi }(x)=\underset{\alpha }{}𝐜_\alpha ^+\overline{w}_\alpha (p\stackrel{}{x})e^{i\xi _\alpha \epsilon _\alpha t},\hfill \end{array}$$
(100)
where
$$\underset{\alpha }{}\underset{s\lambda lm}{}_0^{\mathrm{}}𝑑p,\epsilon _\alpha =\sqrt{p^2+m^2}.$$
(101)
In view of the notations defined above, the free Hamiltonian will compactly be rewritten as follows
$$H^0=H_g^0+H_q^0,$$
(102)
$$H_g^0=\frac{1}{2}\underset{\alpha \beta }{}\omega _{\alpha \beta }:𝐚_\alpha 𝐚_\beta :,$$
(103)
$$H_q^0=\underset{\alpha }{}\xi _\alpha \epsilon _\alpha :𝐜_\alpha ^+𝐜_\alpha :,$$
(104)
where the symbol $`::`$ denotes the normal product and
$$\omega _{\alpha \beta }\omega (k)\delta (kk^{})\delta _{cc^{}}\delta _{\lambda \lambda ^{}}\delta _{JJ^{}}\delta _{MM^{}}(1\delta _{\xi \xi ^{}}).$$
(105)
Next, we derive the expression of the interaction Hamiltonian in the angular momentum representation. In doing this, It is helpful to use the transformations for an operator ( a Hamiltonian, a field operator, a creation or annihilation operator) among the Heisenberg picture, the Schrödinger picture and the interaction picture
$$𝐅_H(t)=e^{iHt}𝐅_se^{iHt}=e^{iHt}e^{H_0t}𝐅_I(t)e^{iH_0t}e^{iHt},$$
(106)
where $`𝐅_H(t),𝐅_s`$ and $`𝐅_I(t)`$ stand for the operators given in the the Heisenberg, Schrödinger and interaction pictures. From Eq.(4.28), it is clearly seen that $`𝐅_s=`$ $`𝐅_H(0)=𝐅_I(0)`$ which is independent of time. According to the relation in Eq.(4.28), we only need to write out the interaction Hamiltonian in the Schrödinger picture (or in the interaction picture). In this picture, on inserting the expressions in Eqs.(4.14) and (4.22) into Eq.(4.7), we obtain
$`H_I(0)`$ $`=`$ $`{\displaystyle d^3x_I(\stackrel{}{x},0)}`$ (107)
$`=`$ $`{\displaystyle \underset{\alpha \beta \gamma }{}}A(\alpha \beta \gamma ):𝐚_\alpha 𝐚_\beta 𝐚_\gamma :`$ (108)
$`+{\displaystyle \underset{\alpha \beta \gamma \delta }{}}B(\alpha \beta \gamma \delta )`$ $`:`$ $`𝐚_\alpha 𝐚_\beta 𝐚_\gamma 𝐚_\delta :`$ ()
$`+{\displaystyle \underset{\alpha \beta \gamma }{}}C(\alpha \beta \gamma )`$ $`:`$ $`𝐜_\alpha ^+𝐜_\beta 𝐚_\gamma :,`$ (109)
where
$$A(\alpha \beta \gamma )=\frac{g}{2}f^{abc}d^3x(\stackrel{}{A_\alpha }\times \stackrel{}{A_\beta })(\times \stackrel{}{A_\gamma }),$$
(110)
$$B(\alpha \beta \gamma \delta )=\frac{1}{8}g^2f^{abe}f^{cde}d^3x(\stackrel{}{A_\alpha }\times \stackrel{}{A_\beta })(\stackrel{}{A_\gamma }\times \stackrel{}{A_\delta })$$
(110)
and
$$C(\alpha \beta \gamma )=gd^3x\overline{w}_\alpha T^c\stackrel{}{\gamma }w_\beta \stackrel{}{A_\gamma }.$$
(111)
By making use of the expressions of $`\stackrel{}{A_\alpha },\overline{w}_\alpha `$ and $`w_\beta `$, Completing the integrals over $`\stackrel{}{x}`$ and applying the formulas of angular momentum couplings, one may derive explicit expressions of the coefficients $`A(\alpha \beta \gamma ),B(\alpha \beta \gamma \delta )`$ and $`C(\alpha \beta \gamma )`$. The expressions of $`A(\alpha \beta \gamma )`$ and $`B(\alpha \beta \gamma \delta )`$ will be given in the next paper. As for the coefficient $`C(\alpha \beta \gamma ),`$its expression will be presented later in other publication.
## V Three-dimensional relativistic equation for two gluon bound states
The aim of this section is to derive a rigorous three-dimensional relativistic equation satisfied by two-gluon bound states. This equation may be derived from the equation satisfied by the following gluon two-time Green’s functions defined in the Heisenberg picture. These Green’s functions include
$`G(\alpha ^+\beta ^+;\gamma ^{}\delta ^{};t_1t_2)`$ $`=`$ $`0^+|T\{𝐚_\alpha (t_1)𝐚_\beta (t_1)𝐚_\gamma ^+(t_2)𝐚_\delta ^+(t_2)\}|0^{},`$ (112)
$`G(\alpha ^{}\beta ^{};\gamma ^+\delta ^+;t_1t_2)`$ $`=`$ $`0^+|T\{𝐚_\alpha ^+(t_1)𝐚_\beta ^+(t_1)𝐚_\gamma (t_2)𝐚_\delta (t_2)\}|0^{},`$ (113)
$`G(\alpha ^{}\beta ^+;\gamma ^{}\delta ^+;t_1t_2)`$ $`=`$ $`0^+|T\{𝐚_\alpha ^+(t_1)𝐚_\beta (t_1)𝐚_\gamma ^+(t_2)𝐚_\delta (t_2)\}|0^{},`$ ()
$`G(\alpha ^+\beta ^{};\gamma ^+\delta ^{};t_1t_2)`$ $`=`$ $`0^+|T\{𝐚_\alpha (t_1)𝐚_\beta ^+(t_1)𝐚_\gamma (t_2)𝐚_\delta ^+(t_2)\}|0^{},`$ (114)
where $`T`$ symbolizes the time-ordered product and $`|0^\pm `$ denote the physical vacuum states defined in Heisenberg picture. It should be noted that the creation and annihilation operators in Eq.(5.1) are those defined in Eq.(3.15). For later convenience, we will use the notation defined in Eq.(4.9) for the gluon operators. With this notation, instead of the above Green’s functions, we may start, in a consistent way, from the following Green’s function defined by
$$G(\alpha \beta ;\gamma \delta ;t_1t_2)=0^+|T\{𝐚_\alpha (t_1)𝐚_\beta (t_1)𝐚_\gamma (t_2)𝐚_\delta (t_2)\}|0^{},$$
(115)
where $`𝐚_\alpha `$ can be a creation operator or an annihilation one.
On differentiating Eq.(5.2) with respect to time $`t_{1,}`$ it is found that
$$\begin{array}{c}i\frac{}{t_1}G(\alpha \beta ;\gamma \delta ;t_1t_2)=i\delta (t_1t_2)S(\alpha \beta ,\gamma \delta )\hfill \\ +0^+|T\{[(i\frac{}{t_1}𝐚_\alpha (t_1))𝐚_\beta (t_1)\hfill \\ +𝐚_\alpha (t_1)(i\frac{}{t_1}𝐚_\beta (t_1))]𝐚_\gamma (t_2)𝐚_\delta (t_2)\}|0^{},\hfill \end{array}$$
(115)
where
$$\begin{array}{c}S(\alpha \beta ,\gamma \delta )=0^+|[𝐚_\alpha (t_1)𝐚_\beta (t_1),𝐚_\gamma (t_2)𝐚_\delta (t_2)]|0^{}\hfill \\ =0^+|[𝐚_\alpha 𝐚_\beta ,𝐚_\gamma 𝐚_\delta ]|0^{}\hfill \end{array}$$
(116)
is time-independent owing to the time displacement shown in Eq.(4.28) and the restriction of $`\delta (t_1t_2)`$. According to the relation in Eq.(4.28), we have
$$i\frac{}{t}𝐚_\alpha (t)=e^{iHt}[𝐚_\alpha ,H]e^{iHt}.$$
(117)
Since $`𝐚_\alpha `$ in the commutator is independent of time, we may use the Hamiltonian given in Eqs.(4.24)-(4.26) and (4.29) and the commutation relation in Eq.(4.12) to compute the commutator $`[𝐚_\alpha ,H]`$. When the result is substituted in Eq.(5.5), we obtain
$$\begin{array}{c}i\frac{}{t_1}𝐚_\alpha (t_1)=\underset{\rho \sigma }{}\omega _{\rho \sigma }\mathrm{\Delta }_{\alpha \sigma }𝐚_\rho (t_1)+\underset{\rho \sigma \tau }{}f_1(\rho \sigma \tau )\mathrm{\Delta }_{\alpha \tau }:𝐚_\rho (t_1)𝐚_\sigma (t_1):\hfill \\ +\underset{\rho \sigma \tau \lambda }{}f_2(\rho \sigma \tau \lambda )\mathrm{\Delta }_{\alpha \lambda }:𝐚_\rho (t_1)𝐚_\sigma (t_1)𝐚_\tau (t_1):\hfill \\ +\underset{\rho \sigma \tau }{}f_3(\rho \sigma \tau )\mathrm{\Delta }_{\alpha \tau }:𝐜_\rho ^+(t_1)𝐜_\sigma (t_1):,\hfill \end{array}$$
(118)
where $`\mathrm{\Delta }_{\alpha \sigma }`$ was defined in Eq.(4.13) and the coefficients $`f_i`$ are defined by
$$\begin{array}{c}f_1(\rho \sigma \tau )=A(\rho \sigma \tau )+A(\rho \tau \sigma )+A(\tau \rho \sigma ),\hfill \\ f_2(\rho \sigma \tau \lambda )=B(\rho \sigma \tau \lambda )+B(\rho \sigma \lambda \tau )+B(\rho \lambda \sigma \tau )+B(\lambda \rho \sigma \tau ),\hfill \\ f_3(\rho \sigma \tau )=C(\rho \sigma \tau ).\hfill \end{array}$$
(119)
Clearly, the expression of $`i\frac{}{t_1}𝐚_\beta (t_1)`$ may be written out from Eq.(5.6) by replacing $`\alpha `$ with $`\beta `$. Inserting this expression and that given in Eq.(5.6) into Eq.(5.3), one may find an equation of motion obeyed by the Green’s function denoted in Eq.(5.2) such that
$$\begin{array}{c}i\frac{}{t_1}G(\alpha \beta ;\gamma \delta ;t_1t_2)=i\delta (t_1t_2)S(\alpha \beta ,\gamma \delta )\hfill \\ +\underset{\rho \sigma \lambda }{}\omega _{\rho \sigma }\mathrm{\Delta }_{\alpha \beta ,\sigma \lambda }G(\rho \lambda ;\gamma \delta ;t_1t_2)\hfill \\ +\underset{\rho \sigma \lambda }{}g_1(\alpha \beta ;\rho \sigma \lambda )G_1(\rho \sigma \lambda ;\gamma \delta ;t_1t_2)\hfill \\ +\underset{\rho \sigma \tau \lambda }{}g_2(\alpha \beta ;\rho \sigma \tau \lambda )G_2(\rho \sigma \tau \lambda ;\gamma \delta ;t_1t_2)\hfill \\ +\underset{\rho \sigma \lambda }{}g_3(\alpha \beta ;\rho \sigma \lambda )G_3(\rho \sigma \lambda ;\gamma \delta ;t_1t_2),\hfill \end{array}$$
(120)
where
$$\mathrm{\Delta }_{\alpha \beta ,\sigma \lambda }=\mathrm{\Delta }_{\alpha \sigma }\delta _{\beta \lambda }+\mathrm{\Delta }_{\beta \sigma }\delta _{\alpha \lambda },$$
(121)
$$\begin{array}{c}g_1(\alpha \beta ;\rho \sigma \lambda )=\underset{\tau }{}f_1(\rho \sigma \tau )\mathrm{\Delta }_{\alpha \beta ,\tau \lambda },\hfill \\ g_2(\alpha \beta ;\rho \sigma \tau \lambda )=\underset{\theta }{}f_2(\rho \sigma \tau \theta )\mathrm{\Delta }_{\alpha \beta ,\theta \lambda },\hfill \\ g_3(\alpha \beta ;\rho \sigma \lambda )=\underset{\tau }{}f_3(\rho \sigma \tau )\mathrm{\Delta }_{\alpha \beta ,\tau \lambda },\hfill \end{array}$$
(122)
and the Green functions $`G_i`$ are defined by
$`G_1(\rho \sigma \lambda ;\gamma \delta ;t_1t_2)`$ (123)
$`=`$ $`0^+|T\{:𝐚_\rho (t_1)𝐚_\sigma (t_1):𝐚_\lambda (t_1)𝐚_\gamma (t_2)𝐚_\delta (t_2)\}|0^{},`$ ()
$`G_2(\rho \sigma \tau \theta ;\gamma \delta ;t_1t_2)`$ (124)
$`=`$ $`0^+|T\{:𝐚_\rho (t_1)𝐚_\sigma (t_1)𝐚_\tau (t_1):𝐚_\theta (t_1)𝐚_\gamma (t_2)𝐚_\delta (t_2)\}|0^{},`$ ()
$`G_3(\rho \sigma \lambda ;\gamma \delta ;t_1t_2)`$ (125)
$`=`$ $`0^+|T\{:𝐜_\rho ^+(t_1)𝐜_\sigma (t_1):𝐚_\lambda (t_1)𝐚_\gamma (t_2)𝐚_\delta (t_2)\}|0^{}.`$ ()
By the Fourier transformation
$$G_i(\alpha ,\mathrm{},\delta ;\omega )=\frac{1}{i}_{\mathrm{}}^+\mathrm{}𝑑te^{i\omega t}G_i(\alpha ,\mathrm{},\delta ;t),$$
(126)
the equation in Eq.(5.8) will be of the form in the energy representation
$$\begin{array}{c}\omega G(\alpha \beta ;\gamma \delta ;\omega )=S(\alpha \beta ,\gamma \delta )\hfill \\ +\underset{\rho \sigma \lambda }{}\omega _{\rho \sigma }\mathrm{\Delta }_{\alpha \beta ,\sigma \lambda }G(\rho \lambda ;\gamma \delta ;\omega )\hfill \\ +\underset{\rho \sigma \lambda }{}g_1(\alpha \beta ;\rho \sigma \lambda )G_1(\rho \sigma \lambda ;\gamma \delta ;\omega )\hfill \\ +\underset{\rho \sigma \tau \lambda }{}g_2(\alpha \beta ;\rho \sigma \tau \lambda )G_2(\rho \sigma \tau \lambda ;\gamma \delta ;\omega )\hfill \\ +\underset{\rho \sigma \lambda }{}g_3(\alpha \beta ;\rho \sigma \lambda )G_3(\rho \sigma \lambda ;\gamma \delta ;\omega ).\hfill \end{array}$$
(127)
Let us look at the second term on the right hand side (RHS) of the above equation. From the definitions of $`\omega _{\rho \sigma }`$ and $`\mathrm{\Delta }_{\alpha \beta ,\sigma \lambda }`$ given in Eqs.(4.27), (5.9) and (4.13), we see, only when $`\xi _\rho \xi _\sigma ,`$ $`\omega _{\rho \sigma }`$ and $`\mathrm{\Delta }_{\rho \sigma }`$ are nonvanishing. When we define
$$\omega _\alpha =\{\begin{array}{c}\omega (k)\text{ }if\text{ }\xi _\alpha =+1\\ \omega (k)\text{ }if\text{ }\xi _\alpha =1\end{array},$$
(128)
the second term mentioned above gives
$$\underset{\rho \sigma \lambda }{}\omega _{\rho \sigma }\mathrm{\Delta }_{\alpha \beta ,\sigma \lambda }G(\rho \lambda ;\gamma \delta ;\omega )=(\omega _\alpha +\omega _\beta )G(\alpha \beta ;\gamma \delta ;\omega ).$$
(129)
Thus, Eq.(5.15) can be written as
$$\begin{array}{c}(\omega \omega _\alpha \omega _\beta )G(\alpha \beta ;\gamma \delta ;\omega )\hfill \\ =S(\alpha \beta ,\gamma \delta )+\underset{\rho \sigma \lambda }{}g_1(\alpha \beta ;\rho \sigma \lambda )G_1(\rho \sigma \lambda ;\gamma \delta ;\omega )\hfill \\ +\underset{\rho \sigma \tau \lambda }{}g_2(\alpha \beta ;\rho \sigma \tau \lambda )G_2(\rho \sigma \tau \lambda ;\gamma \delta ;\omega )\hfill \\ +\underset{\rho \sigma \lambda }{}g_3(\alpha \beta ;\rho \sigma \lambda )G_3(\rho \sigma \lambda ;\gamma \delta ;\omega ).\hfill \end{array}$$
(130)
In order to obtain a closed equation, it is necessary (actually possible) to introduce effective interaction kernels $`\mathrm{\Lambda }_i`$ in such a fashion
$`{\displaystyle \underset{\rho \sigma \lambda }{}}g_1(\alpha \beta ;\rho \sigma \lambda )G_1(\rho \sigma \lambda ;\gamma \delta ;\omega )`$ (131)
$`=`$ $`{\displaystyle \underset{\rho \sigma }{}}\mathrm{\Lambda }_1(\alpha \beta ;\rho \sigma ;\omega )G(\rho \sigma ;\gamma \delta ;\omega ),`$ ()
$`{\displaystyle \underset{\rho \sigma \tau \lambda }{}}g_2(\alpha \beta ;\rho \sigma \tau \lambda )G_2(\rho \sigma \tau \lambda ;\gamma \delta ;\omega )`$ (132)
$`=`$ $`{\displaystyle \underset{\rho \sigma }{}}\mathrm{\Lambda }_2(\alpha \beta ;\rho \sigma ;\omega )G(\rho \sigma ;\gamma \delta ;\omega ),`$ ()
$`{\displaystyle \underset{\rho \sigma \lambda }{}}g_3(\alpha \beta ;\rho \sigma \lambda )G_3(\rho \sigma \lambda ;\gamma \delta ;\omega )`$ (133)
$`=`$ $`{\displaystyle \underset{\rho \sigma }{}}\mathrm{\Lambda }_3(\alpha \beta ;\rho \sigma ;\omega )G(\rho \sigma ;\gamma \delta ;\omega ).`$ ()
With the kernels defined above, Eq.(5.18) becomes
$`(\omega \omega _\alpha \omega _\beta )G(\alpha \beta ;\gamma \delta ;\omega )`$ (134)
$`=`$ $`S(\alpha \beta ,\gamma \delta )+{\displaystyle \underset{\rho \sigma }{}}K(\alpha \beta ;\rho \sigma ;\omega )G(\rho \sigma ;\gamma \delta ;\omega ),`$ ()
this is just the equation satisfied by the Green’s function $`G(\alpha \beta ;\gamma \delta ;\omega )`$ in which
$$K(\alpha \beta ;\rho \sigma ;\omega )=\underset{i=1}{\overset{3}{}}\mathrm{\Lambda }_i(\alpha \beta ;\rho \sigma ;\omega ).$$
(135)
is the total interaction kernel.
The equation obeyed by glueball states may be derived from the above equation with the aid of the Lehmann representation of the Green’s function. Suppose $`\{|n\}`$ form a complete set of glueball states, noticing the transformation in Eq.(4.28) and $`H|n=E_n|n`$, we can write
$$\begin{array}{c}G(\alpha \beta ;\gamma \delta ;t_1t_2)\hfill \\ =\underset{n}{}\{\theta (t_1t_2)0^+|𝐚_\alpha 𝐚_\beta |nn|𝐚_\gamma 𝐚_\delta |0^{}e^{iE_n(t_1t_2)}\hfill \\ +\theta (t_2t_1)0^+|𝐚_\gamma 𝐚_\delta |nn|𝐚_\alpha 𝐚_\beta |0^{}e^{iE_n(t_2t_1)}\}.\hfill \end{array}$$
(136)
Substituting in Eq.(5.24) the representation of the step function
$$\theta (t)=\frac{i}{2\pi }𝑑q\frac{e^{iqt}}{q+iϵ}$$
(137)
and then performing a Fourier transformation with respect to time, we can get from Eq.(5.24) the Lehmann representation such that
$$G(\alpha \beta ;\gamma \delta ;\omega )=\underset{n}{}[\frac{\chi _{\alpha \beta }(n)\overline{\chi }_{\gamma \delta }(n)}{\omega E_n+iϵ}\frac{\chi _{\gamma \delta }(n)\overline{\chi }_{\alpha \beta }(n)}{\omega +E_niϵ}],$$
(138)
where $`E_n`$ is the total energy of a glueball state and
$$\chi _{\alpha \beta }(n)=0^+|𝐚_\alpha 𝐚_\beta |n,\overline{\chi }_{\alpha \beta }(n)=n|𝐚_\alpha 𝐚_\beta |0^{}$$
(139)
are the B-S amplitudes describing the two-gluon bound states. On replacing the $`G(\alpha \beta ;\gamma \delta ;\omega )`$ in Eq.(5.22) with its Lehmann representation, then multiplying the both sides of Eq.(5.22) with $`\omega E_n`$ and finally taking the limit $`\omega E_n`$, we eventually arrive at
$$(E_n\omega _\alpha \omega _\beta )\chi _{\alpha \beta }(n)=\underset{\rho \sigma }{}K(\alpha \beta ;\rho \sigma ;E_n)\chi _{\rho \sigma }(n).$$
(140)
This is just the wanted three-dimensional relativistic equation for the two-gluon bound states. It should be pointed out that the above equation actually represents a set of coupled equations which may be written out by setting $`\xi _\alpha `$ in the index $`\alpha `$ and $`\xi _\beta `$ in the index $`\beta `$ to be $`\pm 1.`$Each of the equations is manifestly of Schrödinger-type, i.e., a standard eigen-equation. In the next paper, these coupled equations will be reduced to an equivalent equation satisfied only by the B-S amplitude for which the two gluons are in the positive energy states.
## VI Closed expression of the interaction kernel
In this section, we are devoted to derive the effective interaction kernel appearing in Eq.(5.28) and defined in Eq.(5.23), giving a closed expression of it. For this derivation, we need equations of motion which describe the evolution of the Green’s function $`G(\alpha \beta ;\gamma \delta ;t_1t_2)`$ and those shown in Eqs.(5.11)-(5.13) with time $`t_2`$. Taking the derivative of the Green’s function in Eq.(5.2) with respect to $`t_2`$, we have
$$\begin{array}{c}i\frac{}{t_2}G(\alpha \beta ;\gamma \delta ;t_1t_2)=i\delta (t_1t_2)S(\alpha \beta ,\gamma \delta )\hfill \\ +0^+|T\{𝐚_\beta (t_1)𝐚_\beta (t_1)[(i\frac{}{t_2}𝐚_\gamma (t_2))𝐚_\delta (t_2)+𝐚_\gamma (t_2)i\frac{}{t_2}𝐚_\delta (t_2)]\}|0^{},\hfill \end{array}$$
(141)
where
$$\begin{array}{c}i\frac{}{t_2}𝐚_\gamma (t_2)=\underset{\rho \sigma }{}\omega _{\rho \sigma }\mathrm{\Delta }_{\lambda \sigma }𝐚_\rho (t_2)\hfill \\ +\underset{\rho \sigma \tau }{}f_1(\rho \sigma \tau )\mathrm{\Delta }_{\gamma \tau }:𝐚_\rho (t_2)𝐚_\sigma (t_2):\hfill \\ +\underset{\rho \sigma \tau \lambda }{}f_2(\rho \sigma \tau \lambda )\mathrm{\Delta }_{\gamma \lambda }:𝐚_\rho (t_2)𝐚_\sigma (t_2)𝐚_\tau (t_2):\hfill \\ +\underset{\rho \sigma \tau }{}f_3(\rho \sigma \tau )\mathrm{\Delta }_{\gamma \tau }:𝐜_\rho ^+(t_2)𝐜_\sigma (t_2):,\hfill \end{array}$$
(142)
here the coefficients $`f_i`$ are the same as represented in Eqs.(5.7). The expression of $`i\frac{}{t_2}𝐚_\delta (t_2)`$ can directly be written out from Eq.(6.2) by replacing $`\gamma `$ by $`\delta .`$ When this expression and the one given above are inserted into Eq.(6.1), we are led to
$$\begin{array}{c}i\frac{}{t_2}G(\alpha \beta ;\gamma \delta ;t_1t_2)=i\delta (t_1t_2)S(\alpha \beta ,\gamma \delta )\hfill \\ +\underset{\rho \sigma \lambda }{}\omega _{\rho \sigma }\mathrm{\Delta }_{\gamma \delta ,\sigma \lambda }G(\alpha \beta ;\rho \lambda ;t_1t_2)\hfill \\ +\underset{\rho \sigma \lambda }{}G_1(\alpha \beta ;\rho \sigma \lambda ;t_1t_2)g_1(\rho \sigma \lambda ;\gamma \delta )\hfill \\ +\underset{\rho \sigma \tau \lambda }{}G_2(\alpha \beta ;\rho \sigma \tau \lambda ;t_1t_2)g_2(\rho \sigma \tau \lambda ;\gamma \delta )\hfill \\ +\underset{\rho \sigma \lambda }{}G_3(\alpha \beta ;\rho \sigma \lambda ;t_1t_2)g_3(\rho \sigma \lambda ;\gamma \delta ),\hfill \end{array}$$
(143)
where $`S(\alpha \beta ,\gamma \delta )`$ was defined in Eq.(5.4), the coefficients $`g_i`$ are the same as those represented in Eq.(5.10) since the following equality holds
$$g_i(\rho ,\mathrm{},\lambda ;\gamma \delta )=g_i(\gamma \delta ;\rho ,\mathrm{},\lambda )$$
(144)
and the Green functions $`G_i`$ are defined as follows
$$\begin{array}{c}G_1(\alpha \beta ;\rho \sigma \lambda ;t_1t_2)\hfill \\ =0^+|T\{𝐚_\alpha (t_1)𝐚_\beta (t_1):𝐚_\rho (t_2)𝐚_\sigma (t_2):𝐚_\lambda (t_2)\}|0^{},\hfill \\ G_2(\alpha \beta ;\rho \sigma \tau \lambda ;t_1t_2)\hfill \\ =0^+|T\{𝐚_\alpha (t_1)𝐚_\beta (t_1):𝐚_\rho (t_2)𝐚_\sigma (t_2)𝐚_\tau (t_2):𝐚_\lambda (t_2)\}|0^{},\hfill \\ G_3(\alpha \beta ;\rho \sigma \lambda ;t_1t_2)\hfill \\ =0^+|T\{𝐚_\alpha (t_1)𝐚_\beta (t_1):𝐜_\rho ^+(t_2)𝐜_\sigma (t_2):𝐚_\lambda (t_2)\}|0^{}.\hfill \end{array}$$
(145)
By the Fourier transformation denoted in Eq.(5.14) and noticing Eq.(5.17), we obtain
$$\begin{array}{c}(\omega +\omega _\gamma +\omega _\delta )G(\alpha \beta ;\gamma \delta ;\omega )=S(\alpha \beta ,\gamma \delta )\hfill \\ \underset{\rho \sigma \lambda }{}G_1(\alpha \beta ;\rho \sigma \lambda ;\omega )g_1(\rho \sigma \lambda ;\gamma \delta )\hfill \\ \underset{\rho \sigma \tau \lambda }{}G_2(\alpha \beta ;\rho \sigma \tau \lambda ;\omega )g_2(\rho \sigma \tau \lambda ;\gamma \delta )\hfill \\ \underset{\rho \sigma \lambda }{}G_3(\alpha \beta ;\rho \sigma \lambda ;\omega )g_3(\rho \sigma \lambda ;\gamma \delta ).\hfill \end{array}$$
(146)
According to the same procedure as formulated in Eqs.(6.1)-(6.6), one can derive the equations of motion obeyed by the Green’s functions denoted in Eqs.(5.11)-(5.13). Because the same expressions of the differentials $`i\frac{}{t_2}𝐚_\gamma (t_2)`$ and $`i\frac{}{t_2}𝐚_\delta (t_2)`$ are employed in all the derivations, the equations of motion for those Green’s functions may immediately be written down by referencing the equation in Eq.(6.6). The equation for the Green’s function in Eq.(5.11) is
$$\begin{array}{c}(\omega +\omega _\gamma +\omega _\delta )G_1(\rho \sigma \lambda ;\gamma \delta ;\omega )=S_1(\rho \sigma \lambda ,\gamma \delta )\hfill \\ \underset{\mu \nu \tau }{}G_{11}(\rho \sigma \lambda ;\mu \nu \tau ;\omega )g_1(\mu \nu \tau ;\gamma \delta )\hfill \\ \underset{\mu \nu \tau \kappa }{}G_{12}(\rho \sigma \lambda ;\mu \nu \tau \kappa ;\omega )g_2(\mu \nu \tau \kappa ;\gamma \delta )\hfill \\ \underset{\mu \nu \tau }{}G_{13}(\rho \sigma \lambda ;\mu \nu \tau ;\omega )g_3(\mu \nu \tau ;\gamma \delta ),\hfill \end{array}$$
(147)
where
$$S_1(\rho \sigma \lambda ,\gamma \delta )=0^+|[:𝐚_\rho 𝐚_\sigma :𝐚_\lambda ,𝐚_\gamma 𝐚_\delta ]|0^{},$$
(148)
$`G_{11}(\rho \sigma \lambda ;\mu \nu \tau ;\omega ),G_{12}(\rho \sigma \lambda ;\mu \nu \tau \kappa ;\omega )`$ and $`G_{13}(\rho \sigma \lambda ;\mu \nu \tau ;\omega )`$ are the Fourier transforms of the following Green’s functions
$$\begin{array}{c}G_{11}(\rho \sigma \lambda ;\mu \nu \tau ;t_1t_2)\hfill \\ =0^+|T\{:𝐚_\rho (t_1)𝐚_\sigma (t_1):𝐚_\lambda (t_1):𝐚_\mu (t_2)𝐚_\nu (t_2):𝐚_\tau (t_2)\}|0^{},\hfill \\ G_{12}(\rho \sigma \lambda ;\mu \nu \tau \kappa ;t_1t_2)\hfill \\ =0^+|T\{:𝐚_\rho (t_1)𝐚_\sigma (t_1):𝐚_\lambda (t_1):𝐚_\mu (t_2)𝐚_\nu (t_2)𝐚_\tau (t_2):𝐚_\kappa (t_2)\}|0^{},\hfill \\ G_{13}(\rho \sigma \lambda ;\mu \nu \tau ;t_1t_2)\hfill \\ =0^+|T\{:𝐚_\rho (t_1)𝐚_\sigma (t_1):𝐚_\lambda (t_1):𝐜_\mu ^+(t_2)𝐜_\nu (t_2):𝐚_\tau (t_2)\}|0^{}.\hfill \end{array}$$
(149)
For the the Green’s function in Eq.(5.12), we can write
$$\begin{array}{c}(\omega +\omega _\gamma +\omega _\delta )G_2(\rho \sigma \tau \lambda ;\gamma \delta ;\omega )=S_2(\rho \sigma \tau \lambda ,\gamma \delta )\hfill \\ \underset{\mu \nu \theta }{}G_{21}(\rho \sigma \tau \lambda ;\mu \nu \theta ;\omega )g_1(\mu \nu \theta ;\gamma \delta )\hfill \\ \underset{\mu \nu \theta \kappa }{}G_{22}(\rho \sigma \tau \lambda ;\mu \nu \kappa \theta ;\omega )g_2(\mu \nu \kappa \theta ;\gamma \delta )\hfill \\ \underset{\mu \nu \theta }{}G_{23}(\rho \sigma \tau \lambda ;\mu \nu \theta ;\omega )g_3(\mu \nu \theta ;\gamma \delta ),\hfill \end{array}$$
(150)
where
$$S_2(\rho \sigma \tau \lambda ,\gamma \delta )=0^+|[:𝐚_\rho 𝐚_\sigma 𝐚_\tau :𝐚_\lambda ,𝐚_\gamma 𝐚_\delta ]|0^{}$$
(151)
and the Green’s functions on the RHS of Eq.(6.10) are defined as
$$\begin{array}{c}G_{21}(\rho \sigma \tau \lambda ;\mu \nu \theta ;t_1t_2)\hfill \\ =0^+|T\{:𝐚_\rho (t_1)𝐚_\sigma (t_1)𝐚_\tau (t_1):𝐚_\lambda (t_1):𝐚_\mu (t_2)𝐚_\nu (t_2):𝐚_\theta (t_2)\}|0^{},\hfill \\ G_{22}(\rho \sigma \tau \lambda ;\mu \nu \kappa \theta ;t_1t_2)\hfill \\ =0^+|T\{:𝐚_\rho (t_1)𝐚_\sigma (t_1)𝐚_\tau (t_1):𝐚_\lambda (t_1):𝐚_\mu (t_2)𝐚_\nu (t_2)𝐚_\kappa (t_2):𝐚_\theta (t_2)\}|0^{},\hfill \\ G_{23}(\rho \sigma \tau \lambda ;\mu \nu \theta ;t_1t_2)\hfill \\ =0^+|T\{:𝐚_\rho (t_1)𝐚_\sigma (t_1)𝐚_\tau (t_1):𝐚_\lambda (t_1):𝐜_\mu ^+(t_2)𝐜_\nu (t_2):𝐚_\theta (t_2)\}|0^{}.\hfill \end{array}$$
(152)
For the Green’s function in Eq.(5.13), we have the equation
$$\begin{array}{c}(\omega +\omega _\gamma +\omega _\delta )G_3(\rho \sigma \lambda ;\gamma \delta ;\omega )=S_3(\rho \sigma \lambda ,\gamma \delta )\hfill \\ \underset{\mu \nu \tau }{}G_{31}(\rho \sigma \lambda ;\mu \nu \tau ;\omega )g_1(\mu \nu \tau ;\gamma \delta )\hfill \\ \underset{\mu \nu \tau \kappa }{}G_{32}(\rho \sigma \lambda ;\mu \nu \tau \kappa ;\omega )g_2(\mu \nu \tau \kappa ;\gamma \delta )\hfill \\ \underset{\mu \nu \tau }{}G_{33}(\rho \sigma \lambda ;\mu \nu \tau ;\omega )g_3(\mu \nu \tau ;\gamma \delta ),\hfill \end{array}$$
(153)
where
$$S_3(\rho \sigma \lambda ,\gamma \delta )=0^+|[:𝐜_\rho ^+𝐜_\sigma :𝐚_\lambda ,𝐚_\gamma 𝐚_\delta ]|0^{}$$
(154)
and the Green’s functions on the RHS of Eq.(6.13) are defined by
$$\begin{array}{c}G_{31}(\rho \sigma \lambda ;\mu \nu \tau ;t_1t_2)\hfill \\ =0^+|T\{:𝐜_\rho ^+(t_1)𝐜_\sigma (t_1):a_\lambda (t_1):a_\mu (t_2)a_\nu (t_2):a_\tau (t_2)\}|0^{},\hfill \\ G_{32}(\rho \sigma \lambda ;\mu \nu \tau \kappa ;t_1t_2)\hfill \\ =0^+|T\{:𝐜_\rho ^+(t_1)𝐜_\sigma (t_1):𝐚_\lambda (t_1):𝐚_\mu (t_2)𝐚_\nu (t_2)𝐚_\tau (t_2):𝐚_\kappa (t_2)\}|0^{},\hfill \\ G_{33}(\rho \sigma \lambda ;\mu \nu \tau ;t_1t_2)\hfill \\ =0^+|T\{:𝐜_\rho ^+(t_1)𝐜_\sigma (t_1):𝐚_\lambda (t_1):𝐜_\mu ^+(t_2)𝐜_\nu (t_2):𝐚_\tau (t_2)\}|0^{}.\hfill \end{array}$$
(155)
Now we are ready to derive the interaction kernels. Multiplying the both sides of Eq.(5.19) with $`(\omega +\omega _\gamma +\omega _\delta )`$ and then applying Eqs.(6.6) and (6.7), we have
$`{\displaystyle \underset{\rho \sigma }{}}\mathrm{\Lambda }_1(\alpha \beta ;\rho \sigma ;\omega )[S(\rho \sigma ,\gamma \delta ){\displaystyle \underset{\mu \nu \lambda }{}}G_1(\rho \sigma ;\mu \nu \lambda ;\omega )g_1(\mu \nu \lambda ;\gamma \delta )`$ (158)
$`{\displaystyle \underset{\mu \nu \tau \lambda }{}}G_2(\rho \sigma ;\mu \nu \tau \lambda ;\omega )g_2(\mu \nu \tau \lambda ;\gamma \delta )`$
$`{\displaystyle \underset{\mu \nu \lambda }{}}G_3(\rho \sigma ;\mu \nu \lambda ;\omega )g_3(\mu \nu \lambda ;\gamma \delta )]`$
$`=`$ $`{\displaystyle \underset{\rho \sigma \lambda }{}}g_1(\alpha \beta ;\rho \sigma \lambda )[S_1(\rho \sigma \lambda ,\gamma \delta ){\displaystyle \underset{\mu \nu \tau }{}}G_{11}(\rho \sigma \lambda ;\mu \nu \tau ;\omega )g_1(\mu \nu \tau ;\gamma \delta )`$ (160)
$`{\displaystyle \underset{\mu \nu \tau \kappa }{}}G_{12}(\rho \sigma \lambda ;\mu \nu \tau \kappa ;\omega )g_2(\mu \nu \tau \kappa ;\gamma \delta )`$
$`{\displaystyle \underset{\mu \nu \tau }{}}G_{13}(\rho \sigma \lambda ;\mu \nu \tau ;\omega )g_3(\mu \nu \tau ;\gamma \delta )].`$
In order to obtain the expression of the kernel $`\mathrm{\Lambda }_1(\alpha \beta ;\rho \sigma ;\omega )`$ from the above equation, it is necessary to eliminate the kernel in the second, third and fourth terms on the left hand side (LHS) of the above equation. Considering that the Green’s function $`G(\alpha \beta ;\gamma \delta ;\omega ),`$as a matrix $`G`$, has an inverse $`G^1(\alpha \beta ;\gamma \delta ;\omega )`$, then from Eq.(5.19) we can get
$$\mathrm{\Lambda }_1(\alpha \beta ;\gamma \delta ;\omega )=\underset{\rho \sigma \lambda }{}\underset{\mu \nu }{}g_1(\alpha \beta ;\rho \sigma \lambda )G_1(\rho \sigma \lambda ;\mu \nu ;\omega )G^1(\mu \nu ;\gamma \delta ;\omega ).$$
(161)
Substituting Eq.(6.17) into the second, third and fourth terms on the LHS of Eq.(6.16), then moving these terms to the RHS of Eq.(6.16) and finally acting on the both sides of that equation with the inverse of the matrix $`S(\rho \sigma ,\gamma \delta )`$ which is assumed to exist, we eventually arrive at
$`\mathrm{\Lambda }_1(\alpha \beta ;\gamma \delta ;\omega )`$ $`=`$ $`\mathrm{\Lambda }_1^{(1)}(\alpha \beta ;\gamma \delta ;\omega )+\mathrm{\Lambda }_1^{(2)}(\alpha \beta ;\gamma \delta ;\omega )`$ ()
$`+\mathrm{\Lambda }_1^{(3)}(\alpha \beta ;\gamma \delta ;\omega ),`$
where
$$\mathrm{\Lambda }_1^{(1)}(\alpha \beta ;\gamma \delta ;\omega )=\underset{\rho \sigma \lambda }{}\underset{\mu \nu }{}g_1(\alpha \beta ;\rho \sigma \lambda )S_1(\rho \sigma \lambda ,\mu \nu )S^1(\mu \nu ,\gamma \delta ),$$
(162)
$`\mathrm{\Lambda }_1^{(2)}(\alpha \beta ;\gamma \delta ;\omega )`$ (163)
$`=`$ $`{\displaystyle \underset{\rho \sigma \lambda }{}}{\displaystyle \underset{\theta \pi }{}}g_1(\alpha \beta ;\rho \sigma \lambda )\{{\displaystyle \underset{\mu \nu \tau }{}}G_{11}(\rho \sigma \lambda ;\mu \nu \tau ;\omega )g_1(\mu \nu \tau ;\theta \pi )`$ (165)
$`+{\displaystyle \underset{\mu \nu \tau \kappa }{}}G_{12}(\rho \sigma \lambda ;\mu \nu \tau \kappa ;\omega )g_2(\mu \nu \tau \kappa ;\theta \pi )`$
$`+{\displaystyle \underset{\mu \nu \tau }{}}G_{13}(\rho \sigma \lambda ;\mu \nu \tau ;\omega )g_3(\mu \nu \tau ;\theta \pi )\}S^1(\theta \pi ,\gamma \delta ),`$
and
$`\mathrm{\Lambda }_1^{(3)}(\alpha \beta ;\gamma \delta ;\omega )`$ (166)
$`=`$ $`{\displaystyle \underset{\rho \sigma \lambda }{}}{\displaystyle \underset{\mu \nu }{}}{\displaystyle \underset{ϵ\iota }{}}{\displaystyle \underset{\theta \pi }{}}g_1(\alpha \beta ;\rho \sigma \lambda )G_1(\rho \sigma \lambda ;\mu \nu ;\omega )G^1(\mu \nu ;ϵ\iota ;\omega )`$ (169)
$`\times \{{\displaystyle \underset{\xi \eta \zeta }{}}G_1(ϵ\iota ;\xi \eta \zeta ;\omega )g_1(\xi \eta \zeta ;\theta \pi )`$
$`+{\displaystyle \underset{\xi \eta \zeta \chi }{}}G_2(ϵ\iota ;\xi \eta \zeta \chi ;\omega )g_2(\xi \eta \zeta \chi ;\theta \pi )`$
$`+{\displaystyle \underset{\xi \eta \zeta }{}}G_3(ϵ\iota ;\xi \eta \zeta ;\omega )g_3(\xi \eta \zeta ;\theta \pi )\}S^1(\theta \pi ,\gamma \delta ).`$
By the same method as described above, the kernel $`\mathrm{\Lambda }_2(\alpha \beta ;\gamma \delta ;\omega )`$ can be derived from the relation in Eq.(5.20) by using the equations in Eq.(6.6) and (6.10). The result is as follows
$$\mathrm{\Lambda }_2(\alpha \beta ;\gamma \delta ;\omega )=\mathrm{\Lambda }_2^{(1)}(\alpha \beta ;\gamma \delta ;\omega )+\mathrm{\Lambda }_2^{(2)}(\alpha \beta ;\gamma \delta ;\omega )+\mathrm{\Lambda }_2^{(3)}(\alpha \beta ;\gamma \delta ;\omega ),$$
(170)
where
$$\mathrm{\Lambda }_2^{(1)}(\alpha \beta ;\gamma \delta ;\omega )=\underset{\rho \sigma \tau \lambda }{}\underset{\mu \nu }{}g_2(\alpha \beta ;\rho \sigma \tau \lambda )S_2(\rho \sigma \tau \lambda ;\mu \nu ;\omega )S^1(\mu \nu ,\gamma \delta ),$$
(170)
$`\mathrm{\Lambda }_2^{(2)}(\alpha \beta ;\gamma \delta ;\omega )`$ (171)
$`=`$ $`{\displaystyle \underset{\rho \sigma \tau \lambda }{}}{\displaystyle \underset{\theta \pi }{}}g_2(\alpha \beta ;\rho \sigma \tau \lambda )\{{\displaystyle \underset{\mu \nu \epsilon }{}}G_{21}(\rho \sigma \tau \lambda ;\mu \nu \epsilon ;\omega )g_1(\mu \nu \epsilon ;\theta \pi )`$ (173)
$`+{\displaystyle \underset{\mu \nu \epsilon \kappa }{}}G_{22}(\rho \sigma \tau \lambda ;\mu \nu \epsilon \varkappa ;\omega )g_2(\mu \nu \epsilon \varkappa ;\theta \pi )`$
$`{\displaystyle \underset{\mu \nu \epsilon }{}}G_{23}(\rho \sigma \tau \lambda ;\mu \nu \epsilon ;\omega )g_3(\mu \nu \epsilon ;\theta \pi )\}S^1(\theta \pi ,\gamma \delta ),`$
and
$`\mathrm{\Lambda }_2^{(3)}(\alpha \beta ;\gamma \delta ;\omega )`$ (174)
$`=`$ $`{\displaystyle \underset{\rho \sigma \tau \lambda }{}}{\displaystyle \underset{\mu \nu }{}}{\displaystyle \underset{ϵ\iota }{}}{\displaystyle \underset{\theta \pi }{}}g_2(\alpha \beta ;\rho \sigma \tau \lambda )G_2(\rho \sigma \tau \lambda ;\mu \nu ;\omega )G^1(\mu \nu ;ϵ\iota ;\omega )`$ (177)
$`\times \{{\displaystyle \underset{\xi \eta \zeta }{}}G_1(ϵ\iota ;\xi \eta \zeta ;\omega )g_1(\xi \eta \zeta ;\theta \pi )`$
$`+{\displaystyle \underset{\xi \eta \zeta \kappa }{}}G_2(ϵ\iota ;\xi \eta \zeta \chi ;\omega )g_2(\xi \eta \zeta \chi ;\theta \pi )`$
$`+{\displaystyle \underset{\xi \eta \zeta }{}}G_3(ϵ\iota ;\xi \eta \zeta ;\omega )g_3(\xi \eta \zeta ;\theta \pi )\}S^1(\theta \pi ,\gamma \delta ).`$
Analogously, the kernel $`\mathrm{\Lambda }_3(\alpha \beta ;\gamma \delta ;\omega )`$ may be derived from Eqs.(5.21), (6.6) and (6.13). The result is shown below
$`\mathrm{\Lambda }_3(\alpha \beta ;\gamma \delta ;\omega )`$ (178)
$`=`$ $`\mathrm{\Lambda }_3^{(1)}(\alpha \beta ;\gamma \delta ;\omega )+\mathrm{\Lambda }_3^{(2)}(\alpha \beta ;\gamma \delta ;\omega )+\mathrm{\Lambda }_3^{(3)}(\alpha \beta ;\gamma \delta ;\omega ),`$ ()
where
$$\mathrm{\Lambda }_3^{(1)}(\alpha \beta ;\gamma \delta ;\omega )=\underset{\rho \sigma \lambda }{}\underset{\mu \nu }{}g_3(\alpha \beta ;\rho \sigma \lambda )S_3(\rho \sigma \lambda ,\mu \nu )S^1(\mu \nu ,\gamma \delta ),$$
(179)
$`\mathrm{\Lambda }_3^{(2)}(\alpha \beta ;\gamma \delta ;\omega )`$ (180)
$`=`$ $`{\displaystyle \underset{\rho \sigma \tau \lambda }{}}{\displaystyle \underset{\theta \pi }{}}g_3(\alpha \beta ;\rho \sigma \lambda )\{{\displaystyle \underset{\mu \nu \tau }{}}G_{31}(\rho \sigma \lambda ;\mu \nu \tau ;\omega )g_1(\mu \nu \tau ;\theta \pi )`$ (182)
$`+{\displaystyle \underset{\mu \nu \tau \kappa }{}}G_{32}(\rho \sigma \lambda ;\mu \nu \tau \kappa ;\omega )g_2(\mu \nu \tau \kappa ;\theta \pi )`$
$`+{\displaystyle \underset{\mu \nu \tau }{}}G_{33}(\rho \sigma \lambda ;\mu \nu \tau ;\omega )g_3(\mu \nu \tau ;\theta \pi )\}S^1(\theta \pi ,\gamma \delta ),`$
$`\mathrm{\Lambda }_3^{(3)}(\alpha \beta ;\gamma \delta ;\omega )`$ (183)
$`=`$ $`{\displaystyle \underset{\rho \sigma \lambda }{}}{\displaystyle \underset{\mu \nu }{}}{\displaystyle \underset{ϵ\iota }{}}{\displaystyle \underset{\theta \pi }{}}g_3(\alpha \beta ;\rho \sigma \lambda )G_3(\rho \sigma \lambda ;\mu \nu ;\omega )G^1(\mu \nu ;ϵ\iota ;\omega )`$ (186)
$`\times \{{\displaystyle \underset{\xi \eta \zeta }{}}G_1(ϵ\iota ;\xi \eta \zeta ;\omega )g_1(\xi \eta \zeta ;\theta \pi )`$
$`+{\displaystyle \underset{\xi \eta \zeta \kappa }{}}G_2(ϵ\iota ;\xi \eta \zeta \chi ;\omega )g_2(\xi \eta \zeta \chi ;\theta \pi )`$
$`+{\displaystyle \underset{\xi \eta \zeta }{}}G_3(ϵ\iota ;\xi \eta \zeta ;\omega )g_3(\xi \eta \zeta ;\theta \pi )\}S^1(\theta \pi ,\gamma \delta ).`$
Based on the expressions described in Eqs.(6.18)-(6.29), the total kernel can easily be written out. Using the matrix notation for the Green’s functions and the other functions, the kernel is represented as
$`K`$ $`=`$ $`{\displaystyle \underset{i,j=1}{\overset{3}{}}}\mathrm{\Lambda }_i^{(i)}`$ (187)
$`=`$ $`\{{\displaystyle \underset{i=1}{\overset{3}{}}}g_iS_i{\displaystyle \underset{i,j=1}{\overset{3}{}}}g_iG_{ij}g_j+{\displaystyle \underset{i,j=1}{\overset{3}{}}}g_iG_iG^1G_jg_j\}S^1.`$ ()
This is just the closed expression of the interaction kernel. According to the general argument as presented in Refs., the third term in the above expression plays the role of eliminating all the B-S reducible (two-particle reducible) diagrams contained in the first two terms. In fact, the relations in Eqs.(5.19)-(5.21) and the following ones
$$G_jg_j=G\mathrm{\Lambda }_j$$
(188)
are inserted into the last term in Eq.(6.30), it is seen that
$$\underset{i,j=1}{\overset{3}{}}g_iG_iG^1G_jg_j=\underset{i,j=1}{\overset{3}{}}\mathrm{\Lambda }_iG\mathrm{\Lambda }_j=KGK,$$
(189)
which exhibits the typical structure of the B-S reducible part of the interaction kernel. Therefore, the kernel shown in Eq.(6.30) is truly B-S irreducible, consistent with the conventional concept. The equation in Eq.(5.28) and the kernel in Eq.(6.30) will be employed, in the next paper, to calculate the glueball spectrum in the ladder approximation.
## VII Concluding remarks
In this paper, the exact three-dimensional relativistic equation for two gluon glueball states and its interaction kernel have been derived from the QCD with massive gluons. When the gluon mass tends to zero, the equation and its kernel will naturally go over to the ones for the QCD with massless gluons. As shown in Eq.(5.28), The equation derived is a standard eigenvalue equation of Schrödinger-type. In the position space it appears to be a set of first-order differential equations. This kind of equation was given for fermion systems , but never formulated for boson systems in the past. It should be noted that since the equation is derived in a special equal-time Lorentz frame, it is certainly not Lorentz-covariant even though the equation is rigorous and includes all the retardation effect in it. In order to derive a Lorentz-covariant equation, one may start from the four-time Green’s function
$$G_{\alpha \beta \gamma \delta }(T,t;T^{},t^{})=0^+|T\{𝐚_\alpha (t_1)𝐚_\beta (t_2)𝐚_\gamma (t_3)𝐚_\delta (t_4)\}|0^{},$$
(190)
where
$`T`$ $`=`$ $`{\displaystyle \frac{1}{2}}(t_1+t_2),t=t_1t_2,`$ (191)
$`T^{}`$ $`=`$ $`{\displaystyle \frac{1}{2}}(t_3+t_4),t^{}=t_3t_4.`$ ()
Differentiating the above Green’s function with respect to $`t_1`$ and $`t_2`$ and employing the expression of the differential $`i\frac{}{t_1}𝐚_\alpha (t_1)`$ as shown in Eq.(5.6) and the similar expression for $`i\frac{}{t_2}𝐚_\beta (t_2)`$, one may obtain two equations for the Green’s function $`G_{\alpha \beta \gamma \delta }(T,t;T^{},t^{}).`$ Adding the both equations and subtracting one equation from another, following the same procedure as formulated in Sec.V, it is easy to derive two equations satisfied by the B-S amplitude $`\chi _{\alpha \beta }(T,t)`$ which respectively describe the evolutions of the state with respect to the center of mass time $`T`$ and the relative time $`t`$. In the equal-time frame, owing to the relative time being absent, we are left only with the equation with respect to the center of mass time as given in Sec.V. It is interesting to note that either the equation with respect to $`T`$ or the equation with respect to $`t`$, appears to be a first order differential equation of Schrödinger-type in the position space whose solutions are determined merely by the initial condition of the B-S amplitudes at the time origin. This is an essential feature of the equations mentioned above which is different from the B-S equation. The latter equation is a higher order differential equation and hence, like the Klein-Gordon equation, has unphysical solutions with negative norm as pointed out in the previous literature . This is because the solutions of B-S equation are determined not only by the initial amplitudes at the time origin, but also by the time-differentials of the amplitudes at the time origin.
Another point we would like to note is that unlike the Dyson-Schwinger equation, the relativistic equation derived in this paper is of a closed form. In particular, the interaction kernel in the equation is given a closed expression. The expression contains only a few types of Green’s functions and vacuum expectation values of the operator commutators. They are unambiguously defined in the Heisenberg picture and each of them can independently be calculated by the perturbation method without concerning other Green’s functions. Especially, the kernel represents all the interactions taking place in the bound states and, therefore, are suitable for nonperturbative investigations because the Green’s functions and vacuum expectation values of the commutators, in principle, are able to be evaluated by a certain nonperturbative method as suggested by the lattice gauge approach.
At last, it would be pointed out that although the equation in Sec.V and the kernel in Sec.VI are derived in the angular momentum representation, they suit to formulate the equation and the kernel in the momentum representation as long as the angular momentum quantum numbers in the indices $`\alpha ,\beta ,\mathrm{}`$ are replaced by the momentum ones. That is to say, the equation and the kernel formally remain unchanged in the both representations.
## VIII Acknowledgement
The authors are grateful to professor Shi-Shu Wu for useful discussions. This work was supported in part by National Natural Science Foundation of China.
## IX Appendix: Spherical spinors in the angular momentum representation
In this appendix, we intend to give a derivation of the spherical Dirac spinors which are used, as basis functions, to establish the angular momentum representation for fermion fields. It is well-known that in the relativistic case, unlike the helicity, the spin of a free fermion is not a good quantum number. However, the total angular momentum operator of the fermion commutes with the Hamiltonian. Therefore, it is meaningful to discuss eigenfunctions of the total angular momentum which satisfy Dirac equation. Let us start from the positive energy spinor $`u_s(\stackrel{}{p})`$ which is the solution to the Dirac equation $`(i^\mu p_\mu m)u(\stackrel{}{p})=0`$. This spinor is taken to be
$$u_s(\stackrel{}{p})=\sqrt{\frac{\epsilon +m}{2\epsilon }}\left(\begin{array}{c}1\hfill \\ \frac{\stackrel{}{\sigma }\stackrel{}{p}}{\epsilon +m}\hfill \end{array}\right)\phi _s$$
(192)
which is normalized in such a fashion: $`u_s^+(\stackrel{}{p})u_s(\stackrel{}{p})=1`$. The negative energy spinor can be given by the charge conjugation $`v_s(\stackrel{}{p})=C\overline{u}_s(\stackrel{}{p})^T`$. Suppose $`Y_{lm}(\widehat{p})`$ with $`\widehat{p}=\stackrel{}{p}/\left|\stackrel{}{p}\right|=(\theta ,\phi )`$ and $`\phi _s`$ are the orbital angular momentum and the spin eigenfunctions respectively; the total angular momentum eigenfunctions may be constructed in the momentum space by the C-G coupling
$$\mathrm{\Omega }_{JM}^l(\widehat{p})=\underset{ls}{}C_{lm\frac{1}{2}s}^{JM}Y_{lm}(\widehat{p})\phi _s,$$
(193)
where $`s=\frac{\sigma }{2},\sigma =\pm 1,l=J\pm \frac{\sigma }{2}.`$Noticing the representation
$$\phi _{\frac{1}{2}}=\left(\begin{array}{c}\hfill \text{0}\\ \hfill \text{1}\end{array}\right),\phi _{\frac{1}{2}}=\left(\begin{array}{c}\hfill \text{1}\\ \hfill \text{0}\end{array}\right)$$
(194)
and employing the explicit expressions of the C-G coupling coefficients for different values of $`\sigma `$ which may be found in the textbook, one may derive from (A.2) the expression as follows
$$\mathrm{\Omega }_{JM}^\sigma (\widehat{p})=\left(\begin{array}{c}\sigma \sqrt{\frac{J+\sigma (M\frac{1}{2})+\frac{1}{2}}{2Jl+1}}Y_{J\frac{\sigma }{2},M\frac{1}{2}}(\widehat{p})\\ \sqrt{\frac{J\sigma (M+\frac{1}{2})+\frac{1}{2}}{2J\sigma +1}}Y_{J\frac{\sigma }{2},M+\frac{1}{2}}(\widehat{p})\end{array}\right).$$
(195)
By making use of the Dirac spinor in (A.1) and the eigenfunctions represented in (A.2) or in (A.4) with respect to the momentum $`\stackrel{}{p}`$, we may construct the spherical Dirac spinor in the position space through the following Fourier transformation
$$u_{JM}^\sigma (p\stackrel{}{x})=𝑑\widehat{p}\frac{e^{i\stackrel{}{p}\stackrel{}{x}}}{(2\pi )^{3/2}}pu(\stackrel{}{p})\mathrm{\Omega }_{JM}^\sigma (\widehat{p}).$$
(196)
Substituting the expansion
$$e^{i\stackrel{}{p}\stackrel{}{x}}=4\pi \underset{lm}{}i^lj_l(pr)Y_{lm}^{}(\widehat{p})Y_{lm}(\widehat{x})$$
(197)
and the expressions written in (A.1) and (A.4) into (A.5), considering
$$\stackrel{}{\sigma }\widehat{p}=\left(\begin{array}{cc}\mathrm{cos}\theta \hfill & \mathrm{sin}\theta e^{i\phi }\hfill \\ \mathrm{sin}\theta e^{i\phi }\hfill & \mathrm{cos}\theta \hfill \end{array}\right)$$
(198)
and
$$\stackrel{}{\sigma }\widehat{p}\mathrm{\Omega }_{JM}^\sigma (\widehat{p})=\mathrm{\Omega }_{JM}^\sigma (\widehat{p}),$$
(199)
which is easily proved by utilizing the familiar recursion formulas for the spherical harmornic functions, it is not difficult to derive the expression shown in Eq.(3.22). The expression of the function $`v_{JM}^\sigma (p\stackrel{}{x})`$ may be derived by the charge conjugation denoted in Eq.(3.34). The result was written in Eq.(3.23). It would be noted that the eigenfunction $`\mathrm{\Omega }_{JM}^\sigma (\widehat{x})`$ in Eq.(3.24) which is defined in the position space and appears in Eqs.(3.22) and (3.23) is of the same form as the function $`\mathrm{\Omega }_{JM}^\sigma (\widehat{p})`$ in (A.4) which is defined in the momentum space.
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# On restricted Leibniz algebras
## Introduction
In characteristic $`p`$ the notion of Lie algebra is fruitfully replaced by the notion of *restricted Lie algebra*, that is a Lie algebra equipped with a formal $`p`$-th power called the Frobenius map. A non-antisymmetric version of Lie algebras has been introduced by the second author in under the name *Leibniz algebras*. In Dzhumadil’daev and S.A. Abdykassymova introduced the notion of *restricted Leibniz algebra* in characteristic $`p`$.
It is well-known that an associative algebra gives rise to a restricted Lie algebra by taking the $`p`$-th power as Frobenius map. In the Leibniz case, we prove that the Leibniz algebra associated to a diassociative algebra is in fact a restricted Leibniz algebra.
In the Lie case, the functor from associative algebras to restricted Lie algebras admits a left adjoint $`U_p`$ and it is known that, for a restricted Lie algebra $`𝔤`$, the associative algebra $`U_p(𝔤)`$ classifies the restricted Lie modules over $`𝔤`$. In the Leibniz case, these two roles of $`U_p`$ are played by two different objects. First, we prove that the restricted enveloping diassociative algebra $`Ud_p(𝔤)`$ provides a functor which is left adjoint to the forgetful functor. Second, we show that the restricted Leibniz $`𝔤`$-modules are classified by an associative algebra $`UL_p(𝔤)`$ that we construct explicitly.
In the last part we show that the tensor product of a Leibniz algebra and a Zinbiel algebra (dual notion for Koszul duality), which is known to be a Lie algebra, has a finer structure: it is a *pre-Lie algebra*. This result is valid in any characteristic. In characteristic $`p`$ it turns out that this pre-Lie algebra is in fact a restricted pre-Lie algebra in the sense of Dzhumadil’daev, cf. .
We thank Bruno Vallette for a careful reading of a first version of this paper.
In this paper we denote by $`k`$ a field of prime characteristic $`p`$ except at the beginning of section 4 where no characteristic hypothesis is made.
## 1. Diassociative algebras and restricted Leibniz algebras.
A Leibniz algebra is a non-commutative version of Lie algebra. Thus, before we proceed with Leibniz algebras let us recall some facts in the Lie framework which led N. Jacobson to introduce the notion of restricted Lie algebra. In many cases in which Lie algebras arise naturally in prime characteristic one finds structures richer than the ordinary Lie algebras. For example, any associative algebra $`A`$ over $`k`$ gives rise to a pair $`(A_{Lie},[p])`$ where $`A_{Lie}`$ is the Lie algebra with bracket given by $`[a,b]:=abba`$ where $`a,bA`$ and $`[p]:A_{Lie}A_{Lie}`$ is the ‘Frobenius mapping’ $`aa^p`$. Another example is obtained by considering the Lie algebra $`Der(𝔘)`$ of derivations of a not necessarily associative algebra $`𝔘`$ over $`k`$. It can be proved (see ) that $`Der(𝔘)`$ is closed under the map $`[p]:DD^p`$ where $`DDer(𝔘)`$ and the map $`[p]`$ verifies specific relations. These facts lead to the following definition.
### 1.1. Definition
A *restricted Lie algebra* is a pair $`(L,[p])`$ where $`L`$ is a Lie algebra over $`k`$ and $`[p]:LL,xx^{[p]}`$ is map which satisfies the following relations:
$`(\alpha x)^{[p]}`$ $`=`$ $`\alpha ^px^{[p]},\alpha k,xL`$
$`[x,y^{[p]}]`$ $`=`$ $`[\mathrm{}[x,\underset{p}{\underset{}{y],y]\mathrm{},y}}],x,yL`$
$`(x+y)^{[p]}`$ $`=`$ $`x^{[p]}+y^{[p]}+{\displaystyle \underset{i=1}{\overset{i=p1}{}}}s_i(x,y),`$
where $`is_i(x,y)`$ is the coefficient of $`\lambda ^{i1}`$ in $`[\mathrm{}[x,\underset{p1}{\underset{}{(\lambda x+y)],(\lambda x+y)]\mathrm{},(\lambda x+y)}}]`$.
A *restricted* morphism $`f:L_1L_2`$ is a Lie morphism such that $`f(x^{[p]})=f(x)^{[p]}`$. We denote by $`𝐩\text{-}\mathrm{𝐋𝐢𝐞}`$ the category of restricted Lie algebras.
### 1.2. Example
Any associative algebra $`A`$ over $`k`$ has the structure of restricted algebra with $`p`$-map given by $`aa^p`$. In particular, from the Leibniz rule we obtain that the derivation algebra $`Der(𝔘)gl(𝔘)`$ of a not necessarily associative algebra $`𝔘`$ is a restricted Lie algebra.
Let $`L`$ be a restricted Lie algebra, and let $`U(L)`$ be the universal enveloping algebra of $`L`$. The *universal restricted enveloping algebra* is defined by $`U_p(L):=U(L)/\{x^px^{[p]}|xL\}`$. A module over $`U_p(L)`$ is called *restricted* Lie module.
From example 1.2 above we see that there is a functor $`\mathrm{𝐀𝐬𝐬}𝐩\text{-}\mathrm{𝐋𝐢𝐞}`$. The left adjoint of this functor is the restricted enveloping functor $`U_p(L)`$. A natural question which arises is wether there is an analogous statement in the Leibniz framework in prime characteristic. Let us now recall the notion of Leibniz algebra, for more details about Leibniz algebras the reader could consult ,.
### 1.3. Definition
A Leibniz algebra over $`k`$ is a $`k`$-module $`𝔤`$ equipped with a bilinear map, called bracket,
$$[,]:𝔤\times 𝔤𝔤$$
satisfying the $`\mathrm{𝐿𝑒𝑖𝑏𝑛𝑖𝑧}\mathrm{𝑖𝑑𝑒𝑛𝑡𝑖𝑡𝑦}:`$
$$[x,[y,z]]=[[x,y],z][[x,z],y],$$
for all $`x,y,z𝔤`$
We denote by $`\mathrm{𝐋𝐞𝐢𝐛}`$ the category of Leibniz algebras over $`k`$.
### 1.4. Example
Any Lie algebra is a Leibniz algebra.
When we replace Lie algebras by Leibniz algebras the role of associative algebras is played by the diassociative algebras (see ). We recall the definition:
### 1.5. Definitions
A diassociative algebra is a $`k`$-module $`D`$ equipped with two $`k`$-linear maps $`,:DDD`$ called respectively the left product and the right product such that the products $``$ and $``$ are associative and satisfy the following relations:
(1) $`x(yz)`$ $`=`$ $`x(yz),`$
(2) $`(xy)z`$ $`=`$ $`x(yz),`$
(3) $`(xy)z`$ $`=`$ $`(xy)z.`$
Let $`D`$ and $`D^{}`$ be diassociative algebras. A *morphism* of diassociative algebras $`f:DD^{}`$ is a $`k`$-linear map such that:
$$f(xy)=f(x)f(y)andf(xy)=f(x)f(y)$$
for all $`x,yD`$.
The category of diassociative algebras over $`k`$ is denoted by $`\mathrm{𝐃𝐢𝐚𝐬}`$.
###### Proposition 1.6.
Let $`D`$ be a diassociative algebra. Then $`D`$ endowed with the following bracket
$$[x,y]:=xyyx$$
is a Leibniz algebra, denoted by $`D_{Leib}`$.
*Proof.* See proposition 4.2 in .
Here we would like to note that in the Lie context the solution of many problems in prime characteristic requires only the knowledge of whether a $`p`$-map can be introduced in a given Lie algebra $`L`$ over $`k`$. This led to the concept of *restrictable* Lie algebras (see ). In particular a Lie algebra $`L`$ over $`k`$ is called *restrictable* if $`(\text{Ad }x)^p`$ is an inner derivation for all $`xL`$ (here $`\text{Ad }x(y)=[y,x]`$). In other words for all $`xL`$ there exist $`x^{[p]}L`$ such that $`(\text{Ad }x)^p=(\text{Ad }x^{[p]})`$. From Theorem $`11`$ in we obtain that a Lie algebra $`L`$ is restrictable if and only if there is a $`p`$-map $`[p]:LL`$ which makes $`L`$ a restricted Lie algebra.
In A.S Dzhumadil’daev and S.A. Abdykassymova introduce an analogue of the notion of *restrictable* Lie algebra in the Leibniz framework. They call it restricted Leibniz algebra, and we keep their terminology.
### 1.7. Notation
Let $`𝔤`$ be a Leibniz algebra, we denote by $`r_x:𝔤𝔤`$ the right multiplication operator given by $`r_x(y):=[y,x]`$ for all $`x,y𝔤`$.
### 1.8. Definition
A *restricted Leibniz algebra* is a pair $`(𝔤,[p])`$ where $`𝔤`$ is a Leibniz algebra $`𝔤`$ over $`k`$, and $`[p]:𝔤𝔤,xx^{[p]}`$ is a map (called the $`p`$-map) such that $`r_x^p=r_{x^{[p]}}`$ for all $`x𝔤`$.
Let $`(𝔤,[p])`$ and $`(𝔥,[p])`$ be restricted Leibniz algebras. A *restricted morphism* $`f:𝔤𝔥`$ is a Leibniz morphism such that $`f(x^{[p]})=f(x)^{[p]}`$ for all $`x𝔤`$. We denote the category of restricted Leibniz algebras over $`k`$ by $`𝐩\text{-}\mathrm{𝐋𝐞𝐢𝐛}`$.
### 1.9. Example
Any restricted Lie algebra is restricted as a Leibniz algebra.
The next theorem shows that the associated Leibniz algebra $`D_{Leib}`$ of a diassociative algebra $`D`$ over $`k`$ has the structure of restricted Leibniz algebra. We first prove a technical lemma.
###### Lemma 1.10.
If $`D`$ is a diassociative algebra then we have
$$x(\mathrm{}(\underset{n}{\underset{}{yy)y)\mathrm{})y}})=x(\mathrm{}(\underset{n}{\underset{}{yy)y)\mathrm{})y}})$$
for all $`x,yD`$ and $`n1`$. We adopt the notation $`y^n:=(\mathrm{}(\underset{n}{\underset{}{yy)y)\mathrm{})y}}`$ and $`y^n:=(\mathrm{}(\underset{n}{\underset{}{yy)y)\mathrm{})y}}`$.
*Proof.* In order to prove the identity we proceed by induction on $`n`$. For $`n=2`$ this identity follows from axiom $`(1)`$ of the definition 1.5. We suppose that the relation is true for $`n`$. Then $`xy^{(n+1)}=x(y^ny)=(xy^n)y`$. By induction hypothesis we get $`(xy^n)y=(xy^n)y=x(y^ny)`$. Finally, by axiom $`(1)`$ of the definition 1.5, we obtain that $`x(y^ny)=xy^{(n+1)}`$. $`\mathrm{}`$
###### Theorem 1.11.
Let $`D`$ be a diassociative algebra over $`k`$. If $`[p]:DD`$ is the map given by $`xx^p`$, then $`(D_{Leib},[p])`$ is a restricted Leibniz algebra.
*Proof.* For all $`yD`$ we will denote by $`R_y^{},L_y^{}:DD`$ the maps defined respectively by $`R_y^{}(x):=xy`$ and $`L_y^{}(x):=yx`$. The Leibniz bracket, with this notation, is given by $`[x,y]=(R_y^{}L_y^{})(x)`$ for all $`x,yD`$. Moreover, from relation $`(2)`$ of the definition 1.5 we have that $`R_y^{}L_y^{}=L_y^{}R_y^{}`$. Therefore by the binomial formula we obtain that $`(R_y^{}L_y^{})^n=_{i=0}^{i=n}\left(\genfrac{}{}{0pt}{}{n}{i}\right)(R_y^{})^i(L_y^{})^{ni}`$.
Besides in prime characteristic $`p`$ we have $`\left(\genfrac{}{}{0pt}{}{p}{i}\right)=0(modp)`$ if $`0<i<p`$. Thus, for $`p=n`$ we get that $`(R_y^{}L_y^{})^p=(R_y^{})^p(L_y^{})^p`$. Therefore we have:
$`[\mathrm{}[x,\underset{p}{\underset{}{y],y]\mathrm{}y}}]`$ $`=`$ $`(R_y^{})^p(x)(L_y^{})^p(x)`$
$`=`$ $`(\mathrm{}(x\underset{p}{\underset{}{y)y)\mathrm{})y}})(\underset{p}{\underset{}{y\mathrm{}(y(y}}x)\mathrm{})`$
$`=`$ $`x(\mathrm{}(\underset{p}{\underset{}{(yy)y)\mathrm{})y}})(\underset{p}{\underset{}{y(\mathrm{}(y(yy}})\mathrm{})x`$
$`=`$ $`xy^py^px.`$
Finally, by lemma 1.10 above we obtain:
$$[\mathrm{}[x,\underset{p}{\underset{}{y],y]\mathrm{}y}}]=xy^py^px=[x,y^p].$$
$`\mathrm{}`$
### 1.12. Remark
Here we made of a choice to define a $`p`$-map. In Goichot introduced the quotient $`D_{as}`$ of $`D`$ by the relation $`=`$. In fact all the elements $`x\mathrm{}x\mathrm{}x`$ have the same image in $`D_{as}`$ and each of them can be taken as $`p`$-map.
###### Proposition 1.13.
Let $`D`$ be a diassociative algebra over $`k`$. The following formula holds in $`D_{p\text{-}Leib}`$
$$[z,(x+y)^{[p]}]=[z,x^{[p]}]+[z,y^{[p]}]+[z,\underset{i=1}{\overset{i=p1}{}}s_i(x,y)],$$
where the bracket involved in $`s_i(x,y)`$ is the Leibniz bracket $`xyyx`$.
*Proof.* Since $``$ is an associative product we have the Jacobson formula
$$(x+y)^{[p]}=x^{[p]}+y^{[p]}+\underset{i=1}{\overset{i=p1}{}}s_i(x,y),$$
where the bracket used in $`s_i`$ is the Lie bracket $`xyyx`$. Under left bracketing with $`z`$ we get the same element by replacing this bracket by the Leibniz bracket because of the relations $`z(xy)=z(xy)`$ and $`(xy)z=(xy)z`$. $`\mathrm{}`$
From Theorem 1.11 above we see that, in prime characteristic, the structure of restricted Leibniz algebra arises in a natural way from the structure of diassociative algebra. Moreover we have the following Proposition:
###### Proposition 1.14.
The following diagram of categories of algebras is commutative.
### 1.15. Example
Let $`D`$ be an diassociative algebra over $`k`$, we denote by $`(D)`$ the diassociative algebra of $`n\times n`$-matrices with entries in $`D`$. Then $`𝔤l_n(D):=(D)_{Leib}`$ is a restricted Leibniz algebra. Here we would like to mention that in zero characteristic the homology of $`𝔤l_n(D)`$ when $`D`$ has a bar-unit, has been computed by Frabetti (see ).
### 1.16. Example
Let $`A`$ be an associative $`k`$-algebra equipped with a $`k`$-module map $`D:AA`$ satisfying the condition:
$$D(a(Db))=DaDb=D((Da)b)$$
for all $`a,bA`$. We define a bilinear map on $`A`$ by
$$[a,b]:=a(Db)(Db)a.$$
If we set $`ab:=aDb`$ and $`ba:=(Db)a`$ for all $`a,bA`$, then we can easily check that $`(A,,)`$ is a diassociative algebra. Therefore by theorem 1.11 above $`A`$ has the structure of restricted Leibniz algebra with map $`[p]:AA`$ given by $`aa^p=(Da)^{p1}a`$ for all $`aA`$.
In case that $`D=id`$, we notice that $`(A,[,])`$ is a restricted Lie algebra with map given by $`aa^p`$ for all $`aA`$.
## 2. Restricted universal enveloping diassociative algebra.
Let us recall that the free diassociative algebra over the vector space $`V`$ is of the form $`T(V)VT(V)`$, cf. .
The functor $`()_{Leib}:\mathrm{𝐃𝐢𝐚𝐬}\mathrm{𝐋𝐞𝐢𝐛}`$ has a left adjoint functor $`Ud:\mathrm{𝐋𝐞𝐢𝐛}\mathrm{𝐃𝐢𝐚𝐬}`$ given by
$$Ud(𝔤):=T(𝔤)𝔤T(𝔤)/\{[x,y]xy+yx|x,y𝔤\}.$$
By theorem 1.11 we obtain a functor $`()_{p\text{-}Leib}:\mathrm{𝐃𝐢𝐚𝐬}𝐩\text{-}\mathrm{𝐋𝐞𝐢𝐛}`$. In the following we construct a left adjoint to the functor $`()_{p\text{-}Leib}.`$
### 2.1. Definition
Let $`(𝔤,[p])`$ be a restricted Leibniz algebra. Then the *restricted universal diassociative algebra* $`Ud_p(𝔤)`$ is defined by:
$$T(𝔤)𝔤T(𝔤)/\{[x,y]xy+yx,x^{[p]}x^p|x,y𝔤\},$$
where, in $`T(𝔤)𝔤T(𝔤)`$, we have made the following identification:
$`[x,y]`$ $`=`$ $`1[x,y]1,`$
$`xy`$ $`=`$ $`1xy,`$
$`yx`$ $`=`$ $`yx1,`$
$`x^{[p]}`$ $`=`$ $`1x^{[p]}1,`$
$`x^p`$ $`=`$ $`\underset{p}{\underset{}{xx\mathrm{}x}}1.`$
###### Proposition 2.2.
Let $`(𝔤,[p])`$ be a restricted Leibniz algebra and $`D`$ a diassociative algebra over $`k`$. Then we have the following natural bijection of sets.
$$Hom_{\mathrm{𝐃𝐢𝐚𝐬}}(Ud_p(𝔤),D)Hom_{𝐩\text{-}\mathrm{𝐋𝐞𝐢𝐛}}(𝔤,D_{p\text{-}Leib}).$$
In other words the functor $`Ud_p`$ is left adjoint to the functor $`()_{p\text{-}Leib}`$.
*Proof.* Let $`f:𝔤D_{𝐩\text{-}\mathrm{𝐋𝐞𝐢𝐛}}`$ be a morphism of restricted Leibniz algebras. Since $`T(𝔤)𝔤T(𝔤)`$ is the free diassociative algebra on $`𝔤`$ there is a unique extension of $`f`$ to a morphism $`\widehat{f}:T(𝔤)𝔤T(𝔤)D`$ of diassociative algebras. Moreover, we have $`\widehat{f}([x,y]xy+yx)=0`$ and $`\widehat{f}(x^{[p]}x^p)`$ because $`f`$ is a restricted Leibniz morphism. Therefore, $`\widehat{f}`$ induces a diassociative algebra morphism $`\varphi _f:Ud_p(𝔤)D`$.
Conversely, if $`\varphi :Ud_p(𝔤)D`$ is a diassociative algebra morphism then the restriction $`f_\varphi :k𝔤kD`$ of $`\varphi `$ to $`k𝔤k`$ is a restricted Leibniz morphism.
It is easy to check that $`f_{(\varphi _f)}=f`$ and $`\varphi _{(f_\varphi )}=\varphi `$ and these two constructions give rise to a bijection
$$Hom_{\mathrm{𝐃𝐢𝐚𝐬}}(Ud_p(𝔤),D)Hom_{𝐩\mathrm{𝐋𝐞𝐢𝐛}}(𝔤,D_{p\text{-}Leib}).$$
$`\mathrm{}`$
## 3. Restricted modules and extensions.
In order to compare with the Lie context let us recall that a strongly abelian extension of restricted Lie algebra $`𝔤`$ by a restricted Lie module $`M`$ is an exact sequence of restricted Lie algebras:
$$0M𝔟𝔤0$$
such that $`[M,M]=0`$ and $`M^{[p]}=0`$.
### 3.1. Remark
There is a bijection between the second Hochschild cohomology group $`H^2(𝔤,M)`$ (see ) and the set of equivalence classes of strong abelian extensions. We note also that we can characterize extensions for which the $`p`$-map is not $`0`$ using the general scheme of Quillen-Barr-Beck’s cohomology theory for details the reader may consult .
By analogy we give the following definition:
### 3.2. Definition
An *abelian extension of restricted Leibniz algebras* (extension of $`𝔤`$ by $`M`$) is an exact sequence of restricted Leibniz algebras:
$$():\mathrm{\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}0}M𝔟𝔤0$$
such that $`[M,M]=0`$ and $`M^{[p]}=0`$
Here we recall the definition of *restricted* Leibniz module given by A. S. Dzhumadil’daev and S. A. Abdykassymova in .
### 3.3. Definition
A *restricted* Leibniz $`𝔤`$-module is a Leibniz $`𝔤`$-module such that
$$[m,x^{[p]}]=[\mathrm{}[m,\underset{p}{\underset{}{x],x]\mathrm{},x}}].$$
###### Lemma 3.4.
If $`()`$ is an extension of the restricted Leibniz algebra $`𝔤`$ by $`M`$
$$():\mathrm{\hspace{0.33em}\hspace{0.33em}0}M𝔟𝔤0$$
with $`[M,M]=0`$ and $`M^{[p]}=0`$, then the vector space $`M`$ inherits the structure of restricted Leibniz module.
*Proof.* Given such an extension we see that $`M`$ is a $`k`$-module equipped with two actions (left and right) of $`𝔤`$
$$[,]:𝔤\times MMand[,]:M\times 𝔤M$$
such that the following relations hold:
(4) $`[m,[x,y]]`$ $`=`$ $`[[m,x],y][[m,y],x],`$
(5) $`[x,[m,y]]`$ $`=`$ $`[[x,m],y][[x,y],m],`$
(6) $`[x,[y,m]]`$ $`=`$ $`[[x,y],m][[x,m],y],`$
(7) $`[m,x^{[p]}]`$ $`=`$ $`[\mathrm{}[m,\underset{p}{\underset{}{x],x]\mathrm{},x}}],`$
for all $`x,y𝔤`$ and $`mM`$. Therefore $`M`$ is a restricted Leibniz module. $`\mathrm{}`$
### 3.5. Restricted universal enveloping algebra of a restricted Leibniz algebra
Let $`𝔤^l`$ and $`𝔤^r`$ be two copies of the Leibniz algebra $`𝔤`$. We denote by $`l_x`$ and $`r_x`$ the elements of $`𝔤^l`$ and $`𝔤^r`$ corresponding to $`x𝔤`$. Let $`T(𝔤^l𝔤^r)`$ be the tensor k-algebra, which is associative and unital.
### 3.6. Definition
Let $`𝔤`$ be a restricted Leibniz algebra. The *restricted universal enveloping algebra* $`UL_p(𝔤)`$ of $`𝔤`$ is defined by $`UL_p(𝔤):=T(𝔤^l𝔤^r)/I_p`$ where $`I_p`$ is the two-sided ideal corresponding to the relations:
(8) $`r_{[x,y]}r_xr_yr_yr_x`$ $`=`$ $`0,`$
(9) $`l_{[x,y]}l_xr_yr_yl_x`$ $`=`$ $`0,`$
(10) $`(r_y+l_y)l_x`$ $`=`$ $`0,`$
(11) $`r_{x^{[p]}}r_x^p`$ $`=`$ $`0,`$
for all $`x,y𝔤`$.
###### Theorem 3.7.
Let $`𝔤`$ be a restricted Leibniz algebra. The category of restricted Leibniz $`𝔤`$-modules is equivalent to the category of right modules over $`UL_p(𝔤)`$.
*Proof.* Let $`M`$ be a restricted $`𝔤`$-module. Then we define actions of $`𝔤^l`$ and $`𝔤^r`$ respectively by:
$$ml_x:=[x,m],andmr_x:=[m,x].$$
for all $`mM`$ and $`x,y𝔤`$. These actions are extended by composition and linearity to an action of $`T(𝔤^l𝔤^r)`$. Since $`M`$ is a restricted $`𝔤`$-module it is in particular a Leibniz $`𝔤`$-module, therefore the elements of type $`(8),(9),(10)`$ act trivially. Besides, from axiom $`(7)`$ elements of type $`(11)`$ act trivially as well. Therefore $`M`$ is endowed with the structure of right $`UL_p(𝔤)`$-module.
Conversely, let $`M`$ be a right $`UL_p(𝔤)`$-module. Then the restriction of the action on $`𝔤^l`$ and $`𝔤^r`$ makes $`M`$ into a Leibniz $`𝔤`$-module. Moreover, from relation $`(11)`$ of the definition 3.6 we obtain that $`M`$ is actually a restricted Leibniz $`𝔤`$-module. $`\mathrm{}`$
## 4. Lie property of Koszul duality for Leibniz algebras
Let us recall two notions of algebras which are closely related to Lie and Leibniz algebras.
By definition a *pre-Lie algebra* (also called Vinberg algebra, and right-symmetric algebra) is a vector space $`A`$ equipped with a binary operation denoted $`\{,\}`$ whose associator is right-symmetric, that is, satisfies the relation
$$\{\{x,y\},z\}\{x,\{y,z\}\}=\{\{x,z\},y\}\{x,\{z,y\}\}.$$
It is well-known that the antisymmetrized operation $`[x,y]:=\{x,y\}\{y,x\}`$ is a Lie bracket. In other words, pre-Lie algebras are Lie admissible algebras.
### 4.1. Example
Any associative algebra is a pre-Lie algebra.
### 4.2. Notation
Let $`A`$ be a pre-Lie algebra. We denote by $`R_a:AA`$ the right-multiplication operator given by $`R_a(b):=\{b,a\}`$ for all $`a,bA`$. Moreover, we denote the $`n`$-th power, by $`a^{\{n\}}:=\{\mathrm{}\{\{\underset{n}{\underset{}{a,a\},a\}\mathrm{}\},a}}\}`$ where $`aA`$ and $`n1`$.
Next, let us recall the notion of Zinbiel algebra which is the dual of Leibniz algebra in Koszul sense.
By definition a *Zinbiel algebra* is a vector space $``$ equipped with a binary operation denoted $``$ which satisfies the relation
$$(ab)c=a(bc)+a(cb).$$
The following lemma is useful for computations.
###### Lemma 4.3.
Let $``$ be a Zinbiel algebra over $`k`$. In any characteristic we have, for any $`a,b`$:
$$(\mathrm{}((a\underset{n}{\underset{}{b)b)\mathrm{})b}}=n!a\underset{n}{\underset{}{(b(b(\mathrm{}b)))}}.$$
In particular, in characteristic $`p`$, we get:
$$(\mathrm{}((a\underset{p}{\underset{}{b)b)\mathrm{})b}}=0$$
for all $`a,b`$.
###### Proof.
Let us denote $`b^n:=bb^{n1}`$ where $`b^1=b`$. We prove by induction that $`b^nb=nb^{n+1}`$. It is true for $`n=1`$. Then we compute:
(12) $`b^nb`$ $`=`$ $`(bb^{n1})b,`$
(13) $`=`$ $`b(b^{n1}b+bb^{n1}),`$
(14) $`=`$ $`b((n1)b^n+b^n),`$
(15) $`=`$ $`nbb^n=nb^{n+1}).`$
Similarly, we compute by induction:
(16) $`(\mathrm{}((a\underset{n}{\underset{}{b)b)\mathrm{})b}}`$ $`=`$ $`(n1)!(a(b^{n1}))b,`$
(17) $`=`$ $`(n1)!(a(b^{n1}b+bb^{n1})),`$
(18) $`=`$ $`n!ab^n.`$
$`\mathrm{}`$
It is known (cf. ) that the tensor product of a Leibniz algebra with a Zinbiel algebra is a Lie algebra. Here we prove a finer result which is valid without characteristic hypothesis.
###### Proposition 4.4.
Over any field $`k`$ the tensor product $`𝔤`$ of the Leibniz algebra $`𝔤`$ with the Zinbiel algebra $``$ is a pre-Lie algebra.
*Proof.* For $`x,y𝔤`$ and $`a,b`$ we define
$$\{xa,yb\}:=[x,y]ab.$$
We compute
$$\begin{array}{c}\{\{xa,yb\},zc\}\{xa,\{yb,zc\}\}\hfill \\ =[[x,y],z](ab)c[x,[y,z]]a(bc)\hfill \\ =[[x,y],z]\left(a(bc)+a(cb)\right)+\left([[x,y],z]+[[x,z],y]\right)a(bc)\hfill \\ =[[x,y],z]a(cb)+[[x,z],y]a(bc).\hfill \end{array}$$
As a consequence the associator of the binary operation $`\{,\}`$ is right-symmetric. We have proved that the operation $`\{,\}`$ is pre-Lie. $`\mathrm{}`$
The category of pre-Lie algebras in prime characteristic have been studied by A. Dzhumadil’daev in . In particular A. Dzumadil’daev, by theorem 1.1 in , proves that there is a Jacobson formula for the $`p`$-th power of a sum of two elements of a pre-Lie algebra. Moreover, he introduces the notion of restricted pre-Lie algebra:
### 4.5. Definition
A *restricted pre-Lie* algebra is a pre-Lie algebra $`A`$ over $`k`$ such that $`R_a^p=R_{a^{\{p\}}}`$ for all $`aA`$.
The next theorem give us the Lie property of the Koszul duality in prime characteristic for the case of Leibniz and Zinbiel algebras.
###### Theorem 4.6.
Let $`𝔤`$ be a Leibniz algebra and $``$ a Zinbiel algebra. Then the tensor product $`𝔤`$ is a restricted pre-Lie algebra with pre-Lie bracket given by:
$$\{xa,yb\}:=[x,y](ab)$$
and $`p`$-map given by
$$(yb)^{\{p\}}=[\mathrm{}[[\underset{p}{\underset{}{y,y],y]\mathrm{},y}}](\mathrm{}((\underset{p}{\underset{}{bb)b)\mathrm{}b)}}.$$
###### Proof.
From Proposition 4.4 it follows that the bracket $`\{,\}`$ defined on $`𝔤`$ is indeed a pre-Lie bracket.
Moreover, we claim that $`𝔤`$ is a restricted pre-Lie algebra. Indeed we will prove that $`R_{yb}^p=R_{(yb)^{\{p\}}}=0`$. Since $`𝔤`$ is a Leibniz algebra, we have the relations
$$[x,[y,z]]=[x,[z,y]],[x,[y,y]]=0$$
for all $`x,y𝔤`$. Using the relations above we get $`[x,[\mathrm{}[[\underset{p}{\underset{}{y,y],y]\mathrm{},y]}}]=0`$. Thus,
$$\{xa,(yb)^{\{p\}}\}=[x,[\mathrm{}[[\underset{p}{\underset{}{y,y],y]\mathrm{},y]}}]a((\mathrm{}((\underset{p}{\underset{}{bb)b)\mathrm{})b}})=0.$$
Besides, by Lemma 4.3 above we obtain that $`R_{(yb)}^p=0.`$ Therefore $`𝔤`$ is a restricted pre-Lie algebra as claimed. $`\mathrm{}`$
###### Corollary 4.7.
Let $`𝔤`$ be a Leibniz algebra and $``$ a Zinbiel algebra. Then the tensor product $`𝔤`$ is a restricted Lie algebra with Lie bracket given by:
$$[xa,yb]:=[x,y](ab)[y,x](ba)$$
and $`p`$-map given by
$$(yb)^{[p]}=[\mathrm{}[[\underset{p}{\underset{}{y,y],y]\mathrm{},y}}](\mathrm{}((\underset{p}{\underset{}{bb)b)\mathrm{}b)}}.$$
###### Proof.
By Corollary 2.4 in and Theorem 4.6 the corresponding Lie algebra $`(𝔤)_{Lie}`$ is a restricted Lie algebra with $`p`$-map given by $`(yb)(yb)^{\{p\}}`$ and the corollary is proved. $`\mathrm{}`$
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# Special solutions for Ricci flow equation in 2D using the linearization approach
##
In the last time the Ricci flow equations become a major tool for addressing a variety of problems in physics and mathematics . The Ricci flows are second order non-linear parabolic differential equations for the components of the metric $`g_{\mu \nu }`$ of an $`n`$-dimensional Riemannian manifold which are driven by the Ricci curvature tensor $`R_{\mu \nu }`$:
$$\frac{}{t}g_{\mu \nu }=R_{\mu \nu }.$$
(1)
This equation describes geometric deformations of the metric $`g_{\mu \nu }`$ with parameter $`t`$. In particular the Ricci flow equations on two dimensional manifolds have attracted considerable attention on the physical literature in connection with two-dimensional black hole geometry, exact solutions of the renormalization group equations that describe the decay of singularities in non-compact spaces, etc. .
As an attempt to quantize gravity, it is interesting to investigate quantum field theory in a curved space-time background. Solving relativistic field equations in $`(3+1)`$-dimensional curved space-time is generally a difficult process. An alternative approach is to consider lower-dimensional space-times models where exact solutions may be obtained. It is now long time since lower-dimensional gravity proved to exhibit many of the qualitative features of $`(3+1)`$-dimensional general relativity, high dimensional black holes, cosmological models and branes .
The purpose of this paper is to present new explicit solutions for $`2D`$ Ricci flow equation using a linearization approach.
In two dimensions it is useful to consider a local system of conformally flat coordinates in which the metric has the form
$$ds_t^2=\frac{1}{2}e^{\mathrm{\Phi }(x,y;t)}(dx^2+dy^2)=2e^{\mathrm{\Phi }(z_+,z_{};t)}dz_+dz_{}$$
using Cartesian coordinates $`x,y`$ or the complex conjugate variables $`2z_\pm =y\pm ix`$.
Having in mind that the only non-vanishing component of the Ricci tensor is
$$R_+=_+_{}\mathrm{\Phi }(z_+,z_{};t)$$
the Ricci flow equation (1) becomes
$$\frac{}{t}e^{\mathrm{\Phi }(z_+,z_{};t)}=_+_{}\mathrm{\Phi }(z_+,z_{};t).$$
(2)
In what follows we shall use the substitution
$$v(z_+,z_{};t)=e^{\mathrm{\Phi }(z_+,z_{};t)}$$
writing eq. (2) in the form
$$(v)_t=(\mathrm{ln}v)_{z_+z_{}}.$$
(3)
This equation has been studied in detail from the algebraic point of view in . It is considered as a ”continual” version of the general Toda-type equation for a given Lie algebra. Also in it is presented a formal power series solution by expanding the path-ordered exponentials. Although the proposed general solution provides a formal complete solution for (3), its form is quite intricate and difficult to handle.
Our approach to the equation is somewhat different. We use a direct nonlinear substitution, then split the resulting nonlinear equation. This results in a class of special solutions. Moreover we consider that the equation (3) is rather in the class of linearizable systems than Lax-pair solvable ones.
Our supposition comes from the following observations:
1. Starting with (3) and using the substitution $`v=\phi _{z_+}`$, after integrating once with respect to $`z_+`$ we get
$$\phi _t\phi _{z_+}=\phi _{z_+z_{}}+C\phi _{z_+}$$
(4)
where $`C`$ should be a function of $`z_{}`$ and $`t`$, but for the moment it is considered constant. Now making the substitution $`\phi =\mathrm{ln}F`$ we will end up with the following quadratic equation:
$$F_tF_{z_+}F_{z_+z_{}}F+F_{z_+}F_z_{}CFF_{z_+}=0.$$
(5)
This equation has the following multi-shock-like solution:
$$F=1+e^{\eta _1}+e^{\eta _2}+\mathrm{}+e^{\eta _N}$$
(6)
where $`\eta _i=k_iz_+C(z_{}t)`$ , for any $`k_i`$ and positive integer $`N`$. Of course the solution (6) is not a general solution because (5) is not a bilinear Hirota form (it is not gauge-invariant ) and $`C`$ is just a constant. But the lacking of interaction between exponentials is characteristic to linearizable systems (like Burgers, Liouville, etc.) .
2. Equation (3) does not pass the Painlevé test. Usually Painlevé test is considered an integrability detector. Of course it is not infallible. There are many equations which do not pass the Painlevé test but still are completely integrable. These are either in the class of Hamiltonian systems with separable Hamilton-Jacobi equation, or are again some linearizable systems .
Accordingly we are going to seek an underlying linear (or solvable) system for (3). Let us remark that eq. (3) is symmetric in variables $`z_+`$ and $`z_{}`$ and, consequently, in the rest of the paper, the role of these variables can be interchanged.
Let us assume that $`C`$ is no longer a constant, but a free function of $`z_{}`$ and $`t`$. In this case, defining
$$\phi =\psi +_{\mathrm{}}^tC(z_{},t)𝑑t$$
then (4) will have the form
$$\psi _{z_+}\psi _t=\psi _{z_+z_{}}.$$
(7)
Using the same nonlinear substitution $`\psi =\mathrm{ln}F`$ we get
$$F_{z_+}(F_t+F_z_{})=FF_{z_+z_{}}.$$
(8)
Equation (8) can be split in some linear or nonlinear solvable equations. Of course, all the possibilities we are going to analyze will give only special solutions and not general ones.
First of all, we shall split eq. (8) into a system of linear equations. Here we list the possibilities:
* Linearization Ia
$`F_{z_+z_{}}`$ $`=`$ $`0`$
$`F_t+F_z_{}`$ $`=`$ $`0`$
with the general solution
$$F(z_+,z_{},;t)=f(z_+)+g(tz_{})$$
where $`f,g`$ are arbitrary functions. The solution of (3) is
$$v(z_+,z_{};t)=\frac{f^{}(z_+)}{f(z_+)+g(tz_{})}$$
(9)
which is in fact the generalization of the multi-shock solution (6).
* Linearization Ib
$`F_{z_+}`$ $`=`$ $`F_{z_+z_{}}`$
$`F_t+F_z_{}`$ $`=`$ $`F.`$
This system is equivalent with the previous one by means of the transformation:
$$F(z_+,z_{};t)e^z_{}F(z_+,z_{};t)$$
(10)
and $`v(z_+,z_{};t)`$ is invariant.
* Linearization Ic
$`F_t+F_{z_+}`$ $`=`$ $`F_{z_+z_{}}`$
$`F_{z_+}`$ $`=`$ $`F`$
with the general solution
$$F(z_+,z_{},;t)=h(tz_{})e^{z_++z_{}}$$
for any arbitrary function $`h`$. Unfortunately this gives a trivial solution for (3), namely $`v=1`$.
The next attempt is to split eq. (8) in a solvable system of nonlinear equations. Like in the linearization of the type I, we have the possibilities:
* Linearization IIa
$`F_{z_+}`$ $`=`$ $`F^\alpha ,\alpha ,`$
$`F^{\alpha 1}(F_t+F_z_{})`$ $`=`$ $`F.`$ (11)
The advantage of this splitting is that the first equation of the system (Special solutions for Ricci flow equation in 2D using the linearization approach) is a Bernoulli one with the solution
$$F(z_+,z_{};t)=\{(1\alpha )[z_++h(z_{},t)]\}^{\frac{1}{1\alpha }}$$
where $`h(z_{},t)`$ is an arbitrary function. Introducing this expression in the second equation of the system one finds a linear equation for $`h(z_{},t)`$:
$$h_t+(1\alpha )h_z_{}=0.$$
In this way, the general solution for the system (Special solutions for Ricci flow equation in 2D using the linearization approach) is
$$F(z_+,z_{};t)=\{(1\alpha )[z_++Cz_+(t\frac{z_{}}{1\alpha })]\}^{\frac{1}{1\alpha }}$$
which gives
$$v(z_+,z_{},t)=\frac{1+C(t\frac{z_{}}{1\alpha })}{z_++Cz_+(t\frac{z_{}}{1\alpha })}=\frac{1}{z_+}.$$
(12)
Accordingly this nonlinear splitting gives a stationary solution, i.e. independent of the parameter of deformation $`t`$.
* Linearization IIb
$`F_{z_+z_{}}`$ $`=`$ $`F^\alpha F_{z_+}`$
$`F_t+F_z_{}`$ $`=`$ $`F^{\alpha +1}.`$ (13)
From the first equation of the system we get
$$F_z_{}=\frac{1}{\alpha +1}F^{\alpha +1}+\beta (z_{},t)$$
with $`\beta `$ an arbitrary function. Introducing in the second equation we get
$$F_t+F_z_{}=(\alpha +1)(F_z_{}\beta )$$
having the solution
$$F(z_+,z_{};t)=\frac{\alpha +1}{\alpha }^z_{}\beta (\xi ,z_{}+\alpha t\xi )𝑑\xi +h(z_+,z_{}+\alpha t).$$
Now solving the first equation for $`h`$ in the case of a general function $`\beta `$ is a difficult task. In any case, for $`\beta =0`$ the system can be solved having the solution
$$v(z_+,z_{};t)=\frac{f^{}(z_+)}{f(z_+)\alpha (1+\alpha )z_{}+\alpha t}.$$
(14)
We remark that the nonlinear splitting of the bilinear form gives solutions which are less general than (9).
A different approach to eq. (8) can be done using a special combination of the variables $`z_+`$ and $`z_{}`$ as a new independent variable. For example the well known solutions of cigar-type, or rational type, discussed in can be obtained by introducing the following substitution:
$$z_++z_{}=\xi .$$
(15)
Then (3) becomes
$$v_t=(\mathrm{ln}v)_{\xi \xi }.$$
This equation has been extensively studied by Rosenau in where, using the ”addition property”, he found cigar-type and rational type solutions.
Of course one can consider other symmetric combinations of $`z_+`$ and $`z_{}`$ to give new independent variables. The general method would imply the Lie-point symmetries (to find all the similarity reductions) but we are not going to do this here. Rather we will consider the following substitution:
$$z_+z_{}=\xi .$$
(16)
Using the same machinery we end up with the following bilinear equation:
$$F_\xi (F_t+\xi F_\xi )=\xi FF_{\xi \xi }$$
(17)
and, of course
$$v(z_+,z_{};t)=\frac{}{\xi }\mathrm{ln}F$$
.
As in the linearization I, the following alternatives are obvious:
* Linearization IIIa
Now, we can split (17) in the same way as before:
$`F_t+\xi F_\xi =0`$
$`F_{\xi \xi }=0`$
with the general solution
$$F(\xi ,t)=a\xi e^t+b$$
where $`a,b`$ are constants. In this case
$$v(z_+,z_{};t)=\frac{1}{z_+z_{}+\frac{b}{a}e^t}.$$
(18)
* Linearization IIIb
$`F_\xi =\xi F_{\xi \xi }`$
$`F_t+\xi F_\xi =F.`$
From the first equation of this system we have
$$F(\xi ,t)=\frac{1}{2}\xi ^2h(t)+\mu (t)$$
and from the second $`h(t)=ae^t,\mu (t)=ce^t`$ with $`a,c`$ constants. Accordingly
$$v(z_+,z_{};t)=\frac{z_+z_{}}{\frac{1}{2}(z_+z_{})^2+\frac{c}{a}e^{2t}}.$$
(19)
The next three splittings give trivial or stationary solutions:
* Linearization IIIc
$`F_\xi =F_{\xi \xi }`$
$`F_t+\xi F_\xi =\xi F.`$
We have $`F=ae^\xi `$ which leads to
$$v(z_+,z_{};t)=1.$$
(20)
* Linearization IIId
$`F_\xi =\xi F`$
$`F_t+\xi F_\xi =F_{\xi \xi }`$
with $`F(\xi ,t)=ae^{\frac{1}{2}\xi ^2+t}`$ and we get
$$v(z_+,z_{};t)=z_+z_{}.$$
(21)
Finally, let us consider the following splittings of eq. (17) in solvable systems of nonlinear equations:
* Linearization IVa
$`F_\xi =F^\alpha `$
$`F^{\alpha 1}(F_t+\xi F_\xi )=\xi F_{\xi \xi }.`$
From the first equation of the system we get
$$F=[(1\alpha )(\xi +g(t))]^{\frac{1}{1\alpha }},$$
but if we introduce it in the second equation we have $`g^{}+1=\alpha z`$ which gives $`g^{}=1`$ and $`\alpha =0`$ and consequently $`F`$ is trivial.
* Linearization IVb
Another incompatible splitting would be ($`\alpha `$ and $`\beta `$ are constants):
$`F_\xi =\xi ^\beta F^\alpha `$
$`\xi ^{\beta 1}F^{\alpha 1}(F_t+\xi F_\xi )=F_{\xi \xi }`$
which has no solution.
* Linearization IVc
Another choice could be
$`F_{\xi \xi }=FF_\xi `$
$`F_t+\xi F_\xi =\xi F^2.`$
From the first equation of the above system we have
$$F=h(t)\mathrm{tan}[\frac{1}{2}(\xi h(t)+g(t))]$$
with $`h(t)`$ and $`g(t)`$ arbitrary functions of $`t`$. From the second equation we obtain differential equations for $`h`$ and $`g`$ which are not compatible. Therefore we haven’t any solution in this case.
To conclude, the linearization approach represents an efficient procedure to generate solutions of the Ricci flow equation (3). We remark that the solution (9) is the most general and represents the largest linearizable sector.
The exhaustive study of the Lie symmetries and similarity reductions performed on the quadratic equation (8) will be the subject of forthcoming work .
### Acknowledgments
The authors have been partially supported by the MEC-AEROSPATIAL Program, Romania.
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# Knots, Braids and Hedgehogs from the Eikonal Equation
## 1 Introduction
It has been recently proved that the complex eikonal equation
$$(u)^2=0,$$
(1)
where $`u`$ is a complex scalar field plays a prominent role in high dimensional soliton theories. It had been recently shown that it is an integrability condition for a very wide class of $`CP^N`$ theories in (2+1) and (3+1) dimensional space-time , , (for further results in any dimension and for other models see , ). Indeed, when one assumes that the scalar field must fulfill (apart from the pertinent dynamical equations of motion) this constrain, then an integrable submodel can be defined. Here, integrability is understood as existence of an infinite family of local conserved currents. On the other hand, it is a conventional wisdom known form soliton theories in (1+1) dimensions that the appearance of an infinite family of the conserved currents is closely related with the existence of topological solitons. The reason is simple - solitons are observed in systems with highly non-trivial relations between degrees of freedom, what is equivalent to high degree of symmetries in the system. Usually, such symmetries result in appearance of the conserved currents.
In fact, it has been checked by analytical and numerical calculations that spectrum of the solutions of such integrable submodels can contain solitons . Let us only mention the Nicole model , where hopfion i.e. a soliton with a non-zero value of the Hopf index has been found.
This powerful approach has been recently applied also to the Faddeev-Niemi model . It is believed that this model, based on an unit, three component vector field $`\stackrel{}{n}`$ living in (3+1) Minkowski space-time can be a good candidate for the effective model for the low energy quantum gluodynamics , , , . Here, effective particles i.e. glueballs are described as knotted configurations of the gauge field <sup>1</sup><sup>1</sup>1Knots found application also in other branches of physics as well as chemistry and biology .. Indeed, such knotted solitons have been obtained in numerical work , , . In the case of the Faddeev-Niemi model, the construction of the integrable submodel can be easily done after taking into account the standard stereographic projection, which relates the original vector field with the complex scalar field
$$\stackrel{}{n}=\frac{1}{1+|u|^2}(u+u^{},i(uu^{}),|u|^21).$$
(2)
Unfortunately, even the submodel is too complicated and no exact, knotted solutions have been reported. In spite of that, the eikonal equation appears to be very helpful in the construction of an approximation to the Faddeev-Niemi hopfions . It provides a framework which enable us to analytically investigate qualitative (topology, shape) as well as quantitative (energy) features of the Faddeev-Niemi solitons. For example, using knotted solutions of the eikonal equation, so-called eikonal knots , approximated but analytical solutions of the Faddeev-Niemi model have been constructed. These solutions can represent the numerically obtained Faddeev-Niemi hopfions with approximately $`20\%`$ accuracy . Additionally, it was observed that energy of the eikonal knots (calculated in the Faddeev-Niemi model) with a fixed topological charge is bounded from below. All these facts might suggest that the topological content of the Faddeev-Niemi model and the eikonal equation is similar.
One can also notice that the eikonal knots , unlikely hopfions obtained in the dynamical toy-models , , , are really knotted objects. Because of that they can be applied in other physical or biological contexts as well.
Nonetheless, due to the complicatedness of the toroidal symmetry many properties of the eikonal knots and, in consequence, the Faddeev-Niemi hopfions have not been understood sufficiently.
First of all, the existence of non-torus knots, that is knots which cannot be plotted on a torus (for instance the figure-eight knot) is still an open question. Secondly, it is also unknown whether equation (1) admits knots which are located on arbitrary closed surfaces topologically equivalent to torus. These ’deformed knots’ might be very helpful in obtaining less energy approximation to the Faddeev-Niemi hopfions.
Moreover, since the vector field $`\stackrel{}{n}`$ can, in principle, form other topological objects (referred as monopoles and open strings), it would be important to know what kinds of topological structures are admitted by the eikonal equation.
In the present paper we would like to face some of the upper mentioned problems. It is done in the language of braided string solutions of the eikonal equations, which are found in the next section. Then, taking advantage of the relationship between braids (open strings) and knots (closed string), we can extend our results on the eikonal knots.
## 2 Topological Strings
In order to deal with strings it is natural to introduce the cylindrical coordinates. Then the eikonal equation is given as follows
$$(_\rho u)^2+\frac{1}{\rho ^2}(_\varphi u)^2+(_zu)^2=0,$$
(3)
where the solution is assumed to have the following form
$$u(\rho ,\varphi ,z)=R(\rho )\mathrm{\Phi }(\varphi )Z(z).$$
(4)
Using this Ansatz one can rewrite the eikonal equation and find
$$\frac{1}{R}(_\rho R)^2+\frac{1}{\rho ^2\mathrm{\Phi }}(_\varphi \mathrm{\Phi })^2+\frac{1}{Z}(_zZ)^2=0.$$
(5)
This equation can be expressed in terms of three first order ordinary differential equations. Let us start with the equation for the $`z`$ variable. Namely,
$$Z_z^2(z)+k^2Z^2(z)=0,$$
(6)
where $`k^2`$ is a positive constant. Thus
$$Z(z)=Ae^{\pm ikz},$$
(7)
where $`A`$ is a complex constant. The pertinent equation for the angular coordinate takes this form
$$\mathrm{\Phi }_\varphi ^2(\varphi )+n^2\mathrm{\Phi }(\varphi )=0$$
(8)
and possesses solutions
$$\mathrm{\Phi }(\varphi )=Be^{\pm in\varphi },$$
(9)
where $`B`$ is a complex constant. To guarantee the uniqueness of the solution, the parameter $`n`$ has to be an integer number. Finally, taking into account (6) and (8) we obtain the equation for $`\rho `$
$$R_\rho ^2\left(\frac{n^2}{\rho ^2}+k^2\right)R^2=0.$$
(10)
It can be solved and we find that
$$R(\rho )=C\left(\frac{\rho }{n+\sqrt{k^2\rho ^2+n^2}}\right)^{\pm n}e^{\pm \sqrt{k^2\rho ^2+n^2}}.$$
(11)
Here $`C`$ in a complex constant. Now, inserting found solutions (7), (9) and (11) into (4) one can derive a solution of the eikonal equation
$$u(\rho ,\varphi ,z)=C\left(\frac{\rho }{n+\sqrt{k^2\rho ^2+n^2}}\right)^{\pm n}e^{\pm \sqrt{k^2\rho ^2+n^2}}e^{\pm i(n\varphi +kz)},$$
(12)
with a new constant $`C`$. Moreover, one can observe that by means of (12) we are able to construct more general solution. Such a solution is just a sum of (12) solutions
$$u(\rho ,\varphi ,z)=\underset{j=1}{\overset{N}{}}C_j\left(\frac{\rho }{n_j+\sqrt{k_j^2\rho ^2+n_j^2}}\right)^{\pm n_j}e^{\pm \sqrt{k_j^2\rho ^2+n_j^2}}e^{\pm i(n_j\varphi +k_jz)}+c$$
(13)
with the parameters $`k_j`$ and $`n_j`$ obeying the condition
$$\frac{k_j}{n_j}=\text{const}.$$
(14)
Here, $`c`$ is a new complex constant. One can recognize in this formula an equivalent of the constant value of the ratio between ’winding’ numbers in $`\xi `$ and $`\varphi `$ directions for the knotted configurations . In both cases this condition guarantees that only regular, non-intersecting objects can be constructed. No singular configurations, neither knots nor strings, are allowed by the eikonal equation.
Of course, upper obtained configuration (13) is topological non-trivial only if the pertinent topological index does not vanish. This calculation cab be performed taking into account the well-known fact that the topological charge of any rational map of the form $`R(z)=p(z)/q(z)`$, where $`p`$ and $`q`$ are polynomials without a common divisor, is equal to maximum of the degree of this polynomial. In our case, for fixed value of the $`z`$ coordinate, the Ansatz is nothing else but a modified rational map. Thus, one can check that the standard winding number reads
$$Q=\text{max}\{n_j,j=1..N\}.$$
(15)
Once again the analogy to the knots is clearly visible. The global topology is fixed by the part of solution (13) with the biggest $`n_j`$ whereas local topological structure i.e. a particular distribution of the topological charge in the space depends on all $`n_j`$ (see ).
Let us now make a digression and notice that our solution (12) is also a solution of the massive eikonal equation
$$(u)^2=m^2u^2,$$
(16)
$`m^2>0`$, if we substitute $`k^2k^2+m^2`$. Then
$$u(\rho ,\varphi ,z)=C\left(\frac{\rho }{n+\sqrt{(k^2+m^2)\rho ^2+n^2}}\right)^{\pm n}e^{\pm \sqrt{(k^2+m^2)\rho ^2+n^2}}e^{\pm i(n\varphi +kz)}.$$
(17)
It is straightforward to see that (unlikely the simple massless eikonal equation) a sum of such solutions does not fulfill the massive equation.
This solution can be immediately adopted in the case of the massive eikonal equation in $`(2+1)`$ dimensions, where we find that
$$u(\rho ,\varphi )=\left(\frac{\rho }{n+\sqrt{m^2\rho ^2+n^2}}\right)^{\pm n}e^{\pm \sqrt{m^2\rho ^2+n^2}}e^{\pm in\varphi }.$$
(18)
This might be interesting in the context of the non-linear sigma model in $`(2+1)`$ dimensions and baby Skyrme models , where solutions of the (massless) eikonal equation are used to build approximated solutions. Thus all baby skyrmions are long-range objects with power-like asymptotic behavior. Topological defects which emerge from the massive eikonal equation are much more well-localized. The unit vector field fades exponentially as it is expected for massive fields. In fact, this feature allows us to think about equation (16) as a massive counterpart of the standard eikonal equation. Quite interesting, the massive solution can be achieved also from a dynamical equation. Namely,
$$^2u=m_{eff}^2u,$$
(19)
where $`m_{eff}`$ denotes an effective mass given as
$$m_{eff}^2=m^2\left(1+\frac{1}{\sqrt{m^2r^2+n^2}}\right).$$
(20)
We see that for all obtained solutions (massive and massless) the vector field tends to a constant value at the spatial infinity $`\stackrel{}{n}\stackrel{}{n}_{\mathrm{}}`$. Thus the position of the string is defined as a curve where the vector field points in the opposite direction i.e.
$$\stackrel{}{n}_0=\stackrel{}{n}^{\mathrm{}}.$$
(21)
Solution (13) gives
$$\stackrel{}{n}^{\mathrm{}}=\underset{\rho \mathrm{}}{lim}\stackrel{}{n}=(0,0,1),$$
(22)
where the solution with $`(+)`$ sing has been chosen. The position of the string is given by the formula
$$|u|^2=0.$$
(23)
Inserting solution (13) one can find that it is equivalent to the following algebraic equation
$$\underset{i,j=1}{\overset{N}{}}R_iR_j\mathrm{cos}[(k_ik_j)z+(n_in_j)\varphi ]+2c_0\underset{i=1}{\overset{N}{}}R_i\mathrm{cos}[k_iz+n_i\varphi ]+c_0^2=0.$$
(24)
For simplicity reasons we put $`C_i=1`$ and $`c=c_0R`$. In the next subsections this equation is solved in the simplest but generic cases.
### 2.1 N=1 case
We start with the simplest case and take into consideration only one-component solution (13). Then equation (24) takes the form
$$R^2+2c_0R\mathrm{cos}[kz+n\varphi ]+c_0^2=0$$
(25)
and can be easily solved. We obtain that the defects are located in
$$R(\rho _0)=c_0,kz+n\varphi =\pi +2\pi l,$$
(26)
where $`l=0,1..n1`$. In other words, as we expected, our solutions appear to be strings located on a cylinder with a constant radius $`\rho =\rho _0`$. For any fixed $`z=z_0`$ the vector field winds $`n`$ times around the strings in the $`z_0`$ plane. Simultaneously, these $`n`$ strings wrap infinitely many times around the infinite long cylinder. There are $`k`$ coils on the interval $`\mathrm{\Delta }z=2\pi n`$.
In Fig. 1 and Fig. 2 we present a part of two- and three-string solutions with $`n=2,k=1`$ and $`n=3,k=\frac{2}{3}`$ respectively. It is easy to observe that, after identification of the bases of the cylinder (i.e. upper and lower ends of the strings), both configurations can be associated with the trefoil knot. One can apply this procedure to other multi-string configurations and find various topologically inequivalent knots. This fact can be understood from the knot theory point of view. Indeed, it is a well-known result from knot and braid theory that every torus knot has a representation in terms of a braid diagram, obtained by the action of the braiding operators on a trivial braid.
Moreover, the observation that all string configurations are located on a cylinder with constant radius is equivalent to the fact that the corresponding knots are torus knots. For example, one can see that our string solutions do not allow us to represent the figure-eight knot (fig. 3).
### 2.2 N=2 case
The second simple but interesting case we will investigate, is the case with two-component solution (13) and with $`c_0=0`$. Now one gets
$$2R_1^2R_2^2\mathrm{cos}[(k_1k_2)z+(n_1n_2)\varphi ]+R_1^2+R_2^2=0.$$
(27)
Then there is a central string located at the origin
$$\rho =0$$
(28)
as well as $`n_1n_2`$ satellite strings
$$R_1(\rho _0)=R_2(\rho _0),(k_1k_2)z+(n_1n_2)\varphi =\pi ,+2\pi l$$
(29)
$`l=0,1..n_1n_2`$ which lie on the cylinder with a constant radius $`\rho _0`$. It can be checked that winding number of the central string is $`Q_c=\text{min}\{n_1,n_2\}`$ whereas each of the satellite strings carries the unit topological charge.
### 2.3 Elliptic configurations
Now, we demonstrate how one can construct a solution describing strings located on an arbitrary surface topologically equivalent to cylinder. For example, we adopt before obtained solutions in the case of the cylindrically-elliptic coordinates $`(\eta ,\varphi ,z)`$, where
$$x=\frac{1}{2}a\mathrm{cos}\varphi \mathrm{cosh}\eta ,$$
$$y=\frac{1}{2}a\mathrm{sin}\varphi \mathrm{sinh}\eta ,$$
$$z=z.$$
(30)
Then the eikonal equation possesses the following solutions
$$u(\eta ,\varphi ,z)=Ce^{\pm ikz}e^{\pm i\sqrt{\lambda ^2+\frac{k^2a^2}{4}}E(i\eta ,1\frac{4\lambda ^2}{4\lambda ^2+k^2a^2})}e^{\pm i\sqrt{\lambda ^2+\frac{k^2a^2}{4}}E(\varphi ,1\frac{4\lambda ^2}{4\lambda ^2+k^2a^2})}+c_0,$$
(31)
where $`E`$ is elliptic integral the second kind. In order to guarantee uniqueness of the solution, parameters $`\lambda `$ and $`k`$ must fulfill relation
$$\frac{2}{\pi }\sqrt{\lambda ^2+\frac{k^2a^2}{4}}E\left(1\frac{4\lambda ^2}{4\lambda ^2+k^2a^2}\right)=nN.$$
(32)
This new number $`n`$ fixes the topological contents of the configuration. We see that strings lie on a elliptic cylinder with a constant radius $`\eta _0`$
$$e^{\pm i\sqrt{\lambda ^2+\frac{k^2a^2}{4}}E(i\eta _0,1\frac{4\lambda ^2}{4\lambda ^2+k^2a^2})}=c_0$$
(33)
and are given by the formula
$$kz+\sqrt{\lambda ^2+\frac{k^2a^2}{4}}E(\varphi ,1\frac{4\lambda ^2}{4\lambda ^2+k^2a^2})=\pi +2\pi l,$$
(34)
where $`l=0,1..n1`$. The generalization to a configuration which describes strings located on any cylinder (or concentric cylinders) with an arbitrary shaped base is straightforward. One has to just introduce the pertinent coordinates and solve the eikonal equation.
It is also easy to notice that using the simplest topologically nontrivial solutions $`u_0`$ (12), (31) with $`n=1`$ we are able to obtain their multi-string generalizations (13). Let us for example consider cylindric string
$$u_0=C\frac{\rho e^{\sqrt{k^2\rho ^2+1}}e^{i(\varphi +kz)}}{1+\sqrt{k^2\rho ^2+1}}.$$
(35)
We see that every solution in the form
$$u=F(u_0),$$
(36)
where $`F`$ is a real and differentiable function, also satisfies the eikonal equation. In particular, multi-string configurations are derived by taking $`F`$ as a polynomial function.
## 3 Topological Hedgehogs
In this section we take under consideration another type of topological objects with the non-trivial $`\pi _2(S^2)`$ homotopy group i.e. hedgehogs or monopoles. The topology of these configurations emerges from the index of the map $`\stackrel{}{n}_{\mathrm{}}:S^2S^2`$.
Since the standard hedgehog possesses the rotational symmetry it is natural to apply the spherical coordinates $`(r,\theta ,\varphi )`$. The eikonal equation can be rewritten as
$$u_r^2+\frac{1}{r^2}u_\theta ^2+\frac{1}{r^2\mathrm{sin}^2\theta }u_\varphi ^2=0.$$
(37)
Once again we use the separable variables method and assume the solution in the form
$$u=f(r)g(\theta )h(\varphi ).$$
(38)
Then equation (1) reads
$$\frac{f_r^2}{f^2}+\frac{1}{r^2}\left[\frac{g_\theta ^2}{g^2}+\frac{1}{\mathrm{sin}^2\theta }\frac{h_\varphi ^2}{h^2}\right]=0.$$
(39)
In the standard manner we obtain that function $`h(\varphi )`$ satisfies
$$\frac{h_\varphi ^2}{h^2}=n^2$$
(40)
and have solutions
$$h(\varphi )=Ae^{\pm in\varphi },$$
(41)
where $`nN`$. Analogously the shape function is given by
$$\frac{f_r^2}{f^2}\frac{m^2}{r^2}=0.$$
(42)
The pertinent solutions are
$$f(r)=Br^{\pm m},$$
(43)
here $`m`$ is any real number. Now, inserting equations (40) and (42) we are able to write down equation for the remaining angular variable
$$\frac{g_\theta ^2}{g^2}\frac{n^2}{\mathrm{sin}^2\theta }=m^2.$$
(44)
One can integrate it and find the following solutions
$$g(\theta )=C\left(\frac{\sqrt{n^2+m^2\mathrm{sin}^2\theta }|m|\mathrm{cos}\theta }{\sqrt{n^2+m^2\mathrm{sin}^2\theta }+|m|\mathrm{cos}\theta }\right)^{\pm \frac{|n|}{2}}e^{|m|\mathrm{arctan}\sqrt{\frac{n^2}{m^2}+\mathrm{tan}^2\theta }},$$
(45)
parameterize by two, previously introduced numbers $`n,m`$. Thus, the general solution of the eikonal equation in the spherical coordinates reads
$$u(r,\theta ,\varphi )=$$
$$Ar^{\pm m}\left(\frac{\sqrt{n^2+m^2\mathrm{sin}^2\theta }|m|\mathrm{cos}\theta }{\sqrt{n^2+m^2\mathrm{sin}^2\theta }+|m|\mathrm{cos}\theta }\right)^{\pm \frac{|n|}{2}}e^{|m|\mathrm{arctan}\sqrt{\frac{n^2}{m^2}+\mathrm{tan}^2\theta }}e^{\pm in\varphi }.$$
(46)
Because of the fact that we are mainly interested in finding of solutions describing topological hedgehogs the following map must possess non-zero topological index
$$\stackrel{}{n}:S_{(r=\mathrm{},\theta ,\varphi )}^2S_{|\stackrel{}{n}|=1}^2.$$
(47)
In other words, behavior of the unit vector field at the spatial infinity cannot be trivial. For example all configurations with $`\stackrel{}{n}_{\mathrm{}}=\stackrel{}{\text{const}}.`$ are, from our point of view, uninteresting since they lead to vanishing topological charge. It results in the fact that any configuration with the non-zero topological index shoul have $`m=0`$ in (43). Then the function $`g`$ can be rewritten in simpler form
$$g(\theta )=C\left(\mathrm{tan}\frac{\theta }{2}\right)^{\pm |n|}.$$
(48)
Thus, the general hedgehog configuration takes the form
$$u(\varphi ,\theta )=Ce^{\pm in\varphi }\left(\mathrm{tan}\frac{\theta }{2}\right)^{\pm |n|}+c.$$
(49)
Identically as in the string case this solution can be generalized to a collection of the hedgehogs
$$u(\varphi ,\theta )=\underset{j=1}{\overset{N}{}}Ce^{\pm in_j\varphi }\left(\mathrm{tan}\frac{\theta }{2}\right)^{\pm |n_j|}+c,$$
(50)
where $`N=1,2..`$. One can check that the total topological charge, which tells us how many times the sphere $`S_{(r=\mathrm{},\theta ,\varphi )}^2`$ is winded on the sphere $`S_{|\stackrel{}{n}|=1}^2`$, is given by the expression
$$Q=\text{max}\{n_j,j=1..N\}.$$
(51)
Let us now notice that obtained eikonal hedgehog configurations (50) can be derived also from a dynamical system i.e. from $`O(3)`$ invariant action in the $`(3+1)`$ dimensional space-time
$$S=d^4x\frac{1}{2}(_\mu \stackrel{}{n})^2.$$
(52)
After expressing the Lagrangian in terms of the complex scalar field $`u`$ we find that
$$L=\frac{2}{(1+|u|^2)^2}_\mu u^\mu u^{}.$$
(53)
Then the pertinent equation of motion reads
$$\frac{_\mu ^\mu u}{(1+|u|^2)^2}2(_\nu u)^2\frac{u^{}}{(1+|u|^2)^3}=0.$$
(54)
One can check that our configurations (50) fulfill this equation. It is due to the fact that they solve not only the eikonal equation but also the Laplace equation
$$^2u=0.$$
(55)
It should be underline that it is unlikely the string solutions found in the previous section. In this case no action, from which strings might be obtained as solutions of the pertinent field equations, is known.
## 4 Conclusions
In the present work two kinds of the static topological defects (strings and hedgehogs) living in three dimensional space have been discussed. These structures have been investigated by means of the complex eikonal equation.
In particular, many braided multi-string configurations have been found. Such multi-string solutions possess well defined topological charge fixed by the biggest value of the $`n_i`$ parameter in Ansatz (13), which is nothing else but the winding number of the vector field around the strings. All solutions have been presented in the analytical form. The most important feature of our strings is that they lie on a cylinder with a constant radius. Since every knot (knotted closed string) possesses its representation in terms of a braid (knots appear when we cut off a finite cylinder and identify its bases) we see that only torus knots (i.e. knots which can be plotted on a torus) can be constructed. In fact, it was observed in our recent paper. At this stage it is hard to justify whether problem of deriving of non-cylindrical strings and non-toroidal knots in the framework of the eikonal equation is only an artefact of the applied method (and such solutions should be found), or due to some yet unknown reasons these solutions are forbidden here. Unfortunately, it is still an open problem.
On the other hand an important problem associated with a property of the eikonal knots has been successfully solved. Namely, it has been proved that the eikonal strings (and in consequence eikonal knots as well) can be located on any surface topologically equivalent to a cylinder i.e. a tube with the base given by any reasonable closed curve, for instance ellipse.
This fact seems to be very important in the context of the Faddeev-Niemi effective action for the low energy gluodynamics where particles (glueballs) are described as knotted flux-tubes of the gauge field. As we discussed it before the properties of the eikonal knots appear to be quite similar to the properties of the Faddeev-Niemi hopfions. In particular, the eikonal knots serve as approximated solutions of the Faddeev-Niemi model. Thus, the results of our paper might find a physical application if we use obtained here elliptic string solutions (or other strings located on a’deformed’ cylinder)) and try to construct some new approximated solutions of the Faddeev-Niemi model.
As we said it before, the problem of the existence of the non-torus knots is still unsolved. However, there are some hints indicating that only torus knots can appear in the Faddeev-Niemi model . Moreover, no non-torus knot has been observed using numerical methods , . Undoubtedly, solving the problem of the existence of non-cylindrical open strings (or equivalently, non-toroidal closed strings) in the eikonal equation, we could shed any light on the appearance (or not) of such structures in the Faddeev-Niemi model. This issue requires farther studies.
This work is partially supported by Foundation for Polish Science FNP and ESF ”COSLAB” program.
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# NON SUSY SEARCHES AT THE TEVATRON
## 1 Introduction
Although supersymmetry (SUSY) is one of the most popular models of physics beyond the Standard Model, alternative models have been proposed. Models of additional spatial dimensions are actively developed. Other models include technicolor, compositeness and models with additional Z and/or W like bosons.
The Tevatron is currently at the high energy frontier with 1.96 TeV center-of-mass energy and thus is one of the major facilities for the discovery of new particles. Searches for new particles at the Tevatron use one or several tool(s) to control the huge QCD background : presence of electrons and/or muons, missing transverse energy (MET), identified heavy flavor jet, identified tau decays, etc…
In the following, the results of searches for new non SUSY particles will be briefly described, more details can be found elsewhere $`^{\mathrm{?},\mathrm{?}}`$.
## 2 Leptonic final states
Unless explicitly mentioned, the leptonic final states only cover the final states with electrons and/or muons.
### 2.1 Dilepton final states
The di-electron and di-muon samples have low background contamination and a well understood Drell-Yan contribution. They form a good sample for searching for new particles in direct search (particles decaying to di-leptons like extra gauge bosons, technicolor particles, Kalusa-Klein excitations in some extra dimensions models, …) or in indirect effects in di-lepton production (compositeness, extra dimension models, …).
New physics would often be visible only in the high end of the invariant mass spectrum. Since electrons are measured in calorimeters, the di-electron sample benefit from a good mass and transverse momentum (pT) resolution, and easy triggering conditions. On the other hand, one has to control and understand the contamination from QCD events where jets are mis-identified as electrons. This has to be measured from real data as the simulation is not reliable enough at that level of precision. CDF and DØ have selected samples based on approximately $`200\mathrm{pb}^1`$, selecting events with two electrons with pT in excess of $`25\text{GeV/}c`$, one of which being required to be in the central part of the calorimeter. Similarly, di-photon samples are selected. The QCD background is less important in the isolated di-muon sample, but special care has to be taken to understand and reduce the background from cosmic rays. The CDF sample is based on $`\mathrm{\hspace{0.17em}200}\mathrm{pb}^1`$ and requires two muons with pT of excess of $`20\text{GeV/}c`$, one with $`|\eta |<1.0`$ and one with $`|\eta |<1.5`$. The DØ sample is based on $`\mathrm{\hspace{0.17em}250}\mathrm{pb}^1`$ and requires two muons with pT of excess of $`15\text{GeV/}c`$. A good agreement is found between the data and the prediction from Standard Model MC complemented by background estimates from real data (QCD and cosmics).
Since no significant deviation from the Standard Model has been observed, the results are interpreted as limits within various models of new physics. A few selected results are shown here, additional results can be found elsewhere $`^{\mathrm{?},\mathrm{?}}`$. Figure 1(a) show the excluded domain in the coupling versus mass of the first excitation of the graviton in the Randall-Sundrum model of extra dimensions from DØ. An interpretation of CDF results combining the electron and muons channels allow to set a limit on the mass of an additional sequential Z boson at $`815\text{GeV/}c^2`$. This search has been extended to the tau channel as can be seen in figure 1(b).
#### Search for a Z from long lived parents
In some models, Z bosons can be produced in the decays of long lived particles. CDF looks for such a striking signature in a di-muon sample from $`163\mathrm{pb}^1`$ of data. In addition to selecting the Z mass region, additional criteria are used to ensure a good reconstruction of the Z vertex, namely tight track quality criteria and the requirement of a minimum acoplanarity between the two muons. Two independent analyses are then applied : in the first, a tight cut of 0.1 cm is applied on the transverse decay length $`L_{xy}`$ of the di-muon system (with 2 events observed for $`0.72\pm 0.27`$ expected) and, in the second, a loose $`L_{xy}>0.03`$cm cut is applied while the di-muon system is required to have a pT$`>30\text{GeV/}c`$ (with 3 events observed and $`1.1\pm 0.8`$ expected). The results are interpreted in the model of a sequential b’ quark lighter than the top quark in figure 2.
### 2.2 Search for a W’
The di-lepton samples allowed to search for a new Z-like gauge boson. CDF has extended that search to that of a W-like boson (W’) in the electron channel. Events with one isolated electron with pT$`>25\text{GeV/}c`$, MET$`>25`$GeV and $`0.4<\mathrm{pT}/\mathrm{MET}<2.5`$ (to reject QCD background) are selected. Assuming Standard Model couplings strength, the analysis result in a limit on the W’ mass of $`842\text{GeV/}c^2`$ using $`205\mathrm{pb}^1`$ of data, with systematics dominated by the PDF uncertainties and the understanding of the electron energy scale.
## 3 Lepton plus jets final states : leptoquarks
Leptoquarks arise in various models beyond the Standard Model. They decay to a lepton and a quark or a neutrino and a quark. First generation leptoquarks are supposed to decay only to electron and/or electron neutrino while second generation leptoquark would decay only to muons and/or muon neutrinos. The branching fraction to the charged lepton is model dependent, so the searches are designed to cover all three final states : di-lepton and jets; lepton, MET and jets and jets plus MET.
Limits on first and second generation scalar leptoquarks from DØ and CDF are shown in figure 3.
## 4 Anomalous W plus heavy flavor production
Studying the associated production of a W with heavy flavor jet(s) is important as it constitutes a background to e.g. top and Higgs selections. Conversely, any deviation from Standard Model expectations would be a hint of new physics. DØ has looked for anomalous production of heavy flavor jet in association with a W in the electron and muons channels using $`\mathrm{\hspace{0.17em}150}\mathrm{pb}^1`$ of data. The events are required to have one electron (resp. one muon) with pT$`>20\text{GeV/}c`$ and $`|\eta |<1.1`$ (resp. $`|\eta |<1.6`$), with MET$`>20GeV`$ and not aligned to the lepton. The transverse mass of the lepton and MET is required to be between $`40\text{GeV/}c^2`$ and $`120\text{GeV/}c^2`$. Jets are potentially identified as b-jets using two independent algorithms, one based on the presence of a muon in or nearby the jet and the other on the reconstruction of secondary vertices. As can be seen in figure 4, a good agreement is found on the number of tagged jets between data and expectation for MC events. As no excess is seen, limits are set on the production cross-section of a $`Wb\overline{b}`$ type signal at 26.3 pb and on a top-like signal at 14.9 pb.
## References
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# Lepton Flavour Violating Leptonic/Semileptonic Decays of Charged Leptons in the Minimal Supersymmetric Standard Model
## 1 Introduction
The neutrino oscillation experiments gave a first experimental evidence beyond the standard model (SM) of the electroweak interactions. In the SM, neutrinos are massless, purely left-handed particles, so there is no leptonic analogy of the Cabibbo-Kobayashi-Maskawa (CKM) matrix in the SM. The neutrino oscillation experiments proved that the neutrinos do mix and that they do have mass. The mixing matrix in the lepton sector, the Maki-Nakagawa-Sakata (MNS) matrix has bi-large mixing structure, indicating that the source of the lepton-flavour mixing is different from the corresponding mixing in the quark sector. The lepton-flavour mixing observed in neutrino oscillation is the first confirmation that the lepton flavour is not conserved quantity. Therefore, experimental observation of the other lepton-flavour violating (LFV) processes is naturally expected. Theoretical study of such processes has a long history before the observation of neutrino oscillations. The model independent study of the operators using the SM fields shows that there are no LFV operators of dimension less or equal to four. There is one dimension five LFV operator that induces neutrino oscillations. The LFV decays can be induced only with the operators of dimension six or more. As the new physics is expected to appear at the scale much larger than the electroweak scale $`246`$ GeV, the LFV decay effects are expected to be much more suppressed than the neutrino oscillation effects. Model independent study of the LFV processes gives the limits on LFV which every model has to satisfy. A model dependent analysis depends on the structure of the model, but are much more predictive than the corresponding model independent analysis. Therefore, the both approaches are indispensable for a theoretical study of LFV. Although the leptonic LFV processes have been studied extensively both in model independent way and using various models , the semileptinic (SL) LFV processes have been studied only in few models .
The supersymmetric (SUSY) models have much nicer theoretical properties than their non-SUSY counterparts. For example, quadratically divergent contributions to the Higgs boson mass from heavy (e.g. GUT scale) particles cancel with its SUSY partners and as a result the gauge-hierarchy problem is much better resolved. The supersymmetrization of the SM cannot be done without additional assumptions. For instance, in the supersymmetric version of the SM there are dimension four operators violating both the lepton number ( / $`L`$ ) and the baryon number ( / $`B`$ ), leading to very fast proton decay. That led us to the introduction of a discrete $`_2`$ symmetry, so called, $`R`$-parity to forbid such undesirable terms. The SUSY breaking also has to be done in such a way not to induce too large flavour violation effects. There are few successful SUSY breaking mediation mechanisms, such as gravity mediation , gauge mediation , anomaly mediation , gaugino mediation , radion mediation , etc. The best established among them is the minimal supergravity model (mSUGRA) which assumes that the SUSY breaking occurs in the hidden sector at very high scale, which communicates with the visible sector (containing SM) with flavour-blind gravitational interactions. The induced soft SUSY breaking mass terms are requested to be universal at the SUSY breaking mediation scale (say, the (reduced) Planck scale), and are therefore flavour-diagonal. The magnitude of the soft SUSY breaking mass terms obtained is in the range to induce potentially observable consequences in the visible sector. The renormalization group (RG) flow from the (reduced) Planck scale to the mass scale of the right-handed neutrinos, induces the flavour non-diagonal terms in the SUSY soft breaking terms for the sleptons, through the flavor-non diagonal Dirac-neutrino Yukawa matrices they contain . They can lead to considerable LFV effects, which, depending on the model parameters can be in the range of the forthcoming LFV experiments .
In this paper, we assume the MSSM with three right-handed neutrinos as the low-energy effective theory below the GUT scale. In such a framework, neutrino oscillation data suggest the existence of very massive right-handed neutrinos which give rise to the small left-handed neutrino masses through the see-saw mechanism . In $`SO(10)`$ models, the required right-handed neutrinos can be naturally embedded into the common multiplet together with the SM particles for each generation. In this paper, the minimal renormazible SUSY $`SO(10)`$ model will be taken as a starting theoretical frame. One of the advantageous points of this model is the automatically conserved $`R`$-parity defined as $`R=(1)^{3(BL)+2S}`$, where $`S`$ represents the spin of a field. Namely, the $`SO(10)`$ model discussed here spontaneously breaks the gauged $`BL`$ symmetry by two units, leading to an automatic $`R`$-parity conservation. The breaking of $`SO(10)`$ group to the SM gauge group, $`SU(3)_C\times SU(2)_L\times U(1)_Y`$ and its phenomenological consequences has already been discussed in our previous publications.
The main goal of this paper is an analysis of neutrinoless SL LFV decays of charged leptons within the MSSM model where parameters are obtained from the underlying $`SO(10)`$ model. At the same time we intend to see how the previous phenomenological analyses constrain the LFV parameters. The paper consists basically of three parts which are given in three sections. In section two, we give the MSSM form factors comprised in the LFV amplitudes at the quark-lepton level. We rederive these form factors, because some of them were not derived completely in the previous literature. In section three, the charged-lepton SL LFV amplitudes at lepton-meson level are derived using a simple hadronization procedure for the quark currents. The branching ratios corresponding to these amplitudes are also given. In section four, the minimal renormalizable SUSY $`SO(10)`$ model is described, the parameters of the MSSM model are derived from the $`SO(10)`$ model, too. Using the predicted Yukawa couplings in the minimal $`SO(10)`$ model, we perform a numerical estimation for the SL LFV processes. The last section is devoted to a summary. In Appendix A, we give our notation for the mass matrices of the neutralino, chargino and sfermions. The MSSM Lagrangian for fermion-sfermion-(gaugino, Higgsino) interaction and the trilinear interactions with $`Z`$ boson are given in Appendix B and C, respectively. In appendix D we present the loop functions needed to evaluate the SL LFV processes. The quark content of meson states, essential for the hadronization of quark currents, is listed in Appendix E, together with constants that define the hadronized quark current in $`\gamma `$-penguin and $`Z`$-boson-penguin amplitude.
## 2 Effective Lagrangian for LFV Interactions
### 2.1 Sources for LFV Interactions
Even though the soft SUSY breaking parameters are flavour blind at the scale of the SUSY breaking mediation, the LFV interactions in the model can induce the LFV sources at low-energy through the renormalization effects. In the following analyses we assume the mSUGRA scenario as the SUSY breaking mediation mechanism. At the original scale of the SUSY breaking mediation we impose the boundary conditions on the soft SUSY breaking parameters, which are characterized by only five parameters: $`m_0`$, $`M_{1/2}`$, $`A_0`$, $`B`$ and $`\mu `$. Here, $`m_0`$ is the universal scalar mass, $`M_{1/2}`$ is the universal gaugino mass, and $`A_0`$ is the universal coefficient of the trilinear couplings. The parameters in the Higgs potential, $`B`$ and $`\mu `$, are determined at the electroweak scale so that the Higgs doublets obtain the correct electroweak symmetry breaking VEV’s through the radiative breaking scenario . The soft SUSY breaking parameters at low energies are obtained through the RGE evolutions with the boundary conditions at the GUT scale.
Although the SUSY breaking mediation scale is normally taken to be the (reduced) Planck scale or the string scale ($`10^{18}`$ GeV), in the following calculations we impose the boundary conditions at the GUT scale ($`10^{16}`$ GeV), and analyze the RGE evolutions from the GUT scale to the electroweak scale. This ansatz is the same as the one in the so-called constrained MSSM (CMSSM).
The effective theory which we analyze below the GUT scale is the MSSM with the right-handed neutrinos. The superpotential in the leptonic sector is given by
$`W_Y`$ $`=`$ $`Y_u^{ij}(u_R^c)_iq_jH_u+Y_d^{ij}(d_R^c)_iq_jH_d`$ (1)
$`+`$ $`Y_\nu ^{ij}(\nu _R^c)_i\mathrm{}_jH_u+Y_e^{ij}(e_R^c)_i\mathrm{}_jH_d+{\displaystyle \frac{1}{2}}M_{R_{ij}}(\nu _R^c)_i(\nu _R^c)_j+\mu H_dH_u,`$
where the indices $`i`$, $`j`$ run over three generations, $`H_u`$ and $`H_d`$ denote the up-type and down-type MSSM Higgs doublets, respectively, and $`M_{R_{ij}}`$ is the heavy right-handed Majorana neutrino mass matrix. We work in the basis where the charged-lepton Yukawa matrix $`Y_e`$ and the mass matrix $`M_{R_{ij}}`$ are real-positive and diagonal matrices: $`Y_e^{ij}=Y_{e_i}\delta _{ij}`$ and $`M_{R_{ij}}=\text{diag}(M_{R_1},M_{R_2},M_{R_3})`$. Thus, the LFV is originated from the off-diagonal components of the neutrino-Dirac-Yukawa coupling matrix $`Y_\nu `$. The soft SUSY breaking terms in the leptonic sector is described as
$`_{\text{soft}}`$ $`=`$ $`\stackrel{~}{q}_i^{}\left(m_{\stackrel{~}{q}}^2\right)_{ij}\stackrel{~}{q}_j+\stackrel{~}{u}_{Ri}^{}\left(m_{\stackrel{~}{u}}^2\right)_{ij}\stackrel{~}{u}_{Rj}+\stackrel{~}{d}_{Ri}^{}\left(m_{\stackrel{~}{d}}^2\right)_{ij}\stackrel{~}{d}_{Rj}`$ (2)
$`+`$ $`\stackrel{~}{\mathrm{}}_i^{}\left(m_\stackrel{~}{\mathrm{}}^2\right)_{ij}\stackrel{~}{\mathrm{}}_j+\stackrel{~}{\nu }_{Ri}^{}\left(m_{\stackrel{~}{\nu }}^2\right)_{ij}\stackrel{~}{\nu }_{Rj}+\stackrel{~}{e}_{Ri}^{}\left(m_{\stackrel{~}{e}}^2\right)_{ij}\stackrel{~}{e}_{Rj}`$
$`+`$ $`m_{H_u}^2H_u^{}H_u+m_{H_d}^2H_d^{}H_d+(B\mu H_dH_u+{\displaystyle \frac{1}{2}}B_\nu M_{R_{ij}}\stackrel{~}{\nu }_{Ri}^{}\stackrel{~}{\nu }_{Rj}+h.c.)`$
$`+`$ $`(A_u^{ij}\stackrel{~}{u}_{Ri}^{}\stackrel{~}{q}_jH_u+A_d^{ij}\stackrel{~}{d}_{Ri}^{}\stackrel{~}{q}_jH_d+h.c.)`$
$`+`$ $`(A_\nu ^{ij}\stackrel{~}{\nu }_{Ri}^{}\stackrel{~}{\mathrm{}}_jH_u+A_e^{ij}\stackrel{~}{e}_{Ri}^{}\stackrel{~}{\mathrm{}}_jH_d+h.c.)`$
$`+`$ $`({\displaystyle \frac{1}{2}}M_1\stackrel{~}{B}\stackrel{~}{B}+{\displaystyle \frac{1}{2}}M_2\stackrel{~}{W}^a\stackrel{~}{W}^a+{\displaystyle \frac{1}{2}}M_3\stackrel{~}{G}^a\stackrel{~}{G}^a+h.c.).`$
As discussed above, we impose the universal boundary conditions at the GUT scale such that
$`\left(m_{\stackrel{~}{q}}^2\right)_{ij}=\left(m_{\stackrel{~}{u}}^2\right)_{ij}=\left(m_{\stackrel{~}{d}}^2\right)_{ij}=m_0^2\delta _{ij},`$
$`\left(m_\stackrel{~}{\mathrm{}}^2\right)_{ij}=\left(m_{\stackrel{~}{\nu }}^2\right)_{ij}=\left(m_{\stackrel{~}{e}}^2\right)_{ij}=m_0^2\delta _{ij},`$
$`m_{H_u}^2=m_{H_d}^2=m_0^2,`$
$`A_u^{ij}=A_0Y_u^{ij},A_d^{ij}=A_0Y_d^{ij},`$
$`A_\nu ^{ij}=A_0Y_\nu ^{ij},A_e^{ij}=A_0Y_e^{ij},`$
$`M_1=M_2=M_3=M_{1/2},`$ (3)
and evolve the soft SUSY breaking parameters to the electroweak scale according to their RGE’s . The $`\mu `$ parameter and the $`B`$ parameter are determined at the electroweak scale so as to minimize the Higgs potential,
$`|\mu |^2`$ $`=`$ $`{\displaystyle \frac{m_{H_d}^2m_{H_u}^2\mathrm{tan}^2\beta }{\mathrm{tan}^2\beta 1}}{\displaystyle \frac{1}{2}}M_Z^2,`$
$`B\mu `$ $`=`$ $`{\displaystyle \frac{1}{2}}\left(m_{H_d}^2+m_{H_u}^2+2|\mu |^2\right)\mathrm{sin}2\beta \left(={\displaystyle \frac{1}{2}}m_A^2\mathrm{sin}2\beta \right).`$ (4)
The LFV sources in the soft SUSY breaking parameters such as the off-diagonal components of $`\left(m_\stackrel{~}{\mathrm{}}^2\right)_{ij}`$ and $`A_e^{ij}`$ are induced by the LFV interactions through the neutrino Dirac Yukawa couplings. For example, the LFV effect most directly emerges in the left-handed slepton mass matrix through the RGE’s such as
$`\mu {\displaystyle \frac{d}{d\mu }}\left(m_\stackrel{~}{\mathrm{}}^2\right)_{ij}`$ $`=`$ $`\mu {\displaystyle \frac{d}{d\mu }}\left(m_\stackrel{~}{\mathrm{}}^2\right)_{ij}|_{\text{MSSM}}`$ (5)
$`+`$ $`{\displaystyle \frac{1}{16\pi ^2}}\left(m_\stackrel{~}{\mathrm{}}^2Y_\nu ^{}Y_\nu +Y_\nu ^{}Y_\nu m_\stackrel{~}{\mathrm{}}^2+2Y_\nu ^{}m_{\stackrel{~}{\nu }}^2Y_\nu +2m_{H_u}^2Y_\nu ^{}Y_\nu +2A_\nu ^{}A_\nu \right)_{ij}`$
where the first term on the right-hand side denotes the MSSM term with no LFV. In the leading-logarithmic approximation, the off-diagonal components ($`ij`$) of the left-handed slepton mass matrix are estimated as
$`\left(\mathrm{\Delta }m_\stackrel{~}{\mathrm{}}^2\right)_{ij}{\displaystyle \frac{3m_0^2+A_0^2}{8\pi ^2}}\left(Y_\nu ^{}LY_\nu \right)_{ij},`$ (6)
where the distinct thresholds for the right-handed Majorana neutrinos are taken into account by the matrix $`L_{ij}=\mathrm{log}\left(\frac{M_G}{M_{R_i}}\right)\delta _{ij}`$. We can see that the neutrino-Dirac-Yukawa coupling matrix plays the crucial role in calculations of the LFV processes.
### 2.2 Effective Lagrangian in terms of quark fields and LFV form factors
In any model containing standard model as the low-energy effective theory, an effective Lagrangian for the SL LFV decays of a lepton contains only three terms: the photon-penguin, the Z-boson-penguin and the box term,
$$i_{\mathrm{eff}}\left(\mathrm{}_i\mathrm{}_j+\overline{q}+q^{}\right)=i_{\mathrm{eff}}^\gamma +i_{\mathrm{eff}}^Z+i_{\mathrm{eff}}^{\mathrm{box}}.$$
(7)
These terms have the following generic structure,
$`i_{\mathrm{eff}}^\gamma (x)`$ $`=`$ $`ie^2{\displaystyle }d^4y\overline{\mathrm{}}_j(x)[(_x^2\gamma _\mu +\begin{array}{c}\text{/}\\ \hfill \end{array}_x_{x\mu })D(xy)(𝒫_{1\gamma }^LP_L+𝒫_{1\gamma }^RP_R)`$
$`+`$ $`\sigma _{\mu \nu }_x^\nu D(xy)(𝒫_{2\gamma }^LP_L+𝒫_{2\gamma }^RP_R)]\mathrm{}_i(x)`$
$`\times `$ $`{\displaystyle \underset{q=u,d,s}{}}Q_q\overline{q}(y)\gamma ^\mu q(y),`$ (11)
$`i_{\mathrm{eff}}^Z(x)`$ $`=`$ $`i{\displaystyle \frac{g^2}{m_Z^2c_W^2}}\overline{\mathrm{}}_j(x)\gamma _\mu (𝒫_Z^LP_L+𝒫_Z^RP_R)\mathrm{}_i(x)`$ (12)
$`\times `$ $`{\displaystyle \underset{q=u,d,s}{}}\overline{q}(I_{3q}2Q_qs_W^2)\gamma ^\mu I_{3q}\gamma ^\mu \gamma _5)q(x),`$
$`i_{\mathrm{eff}}^{\mathrm{box}}`$ $`=`$ $`i{\displaystyle \underset{\overline{q}_aq_b=\overline{u}u,\overline{d}d,\overline{s}s,\overline{d}s,\overline{s}d}{}}[_{1\overline{q}_aq_b}^L(\overline{\mathrm{}}_j\gamma ^\mu P_L\mathrm{}_i)(\overline{q}_a\gamma _\mu P_Lq_b)+_{1\overline{q}_aq_b}^R(\overline{\mathrm{}}_j\gamma ^\mu P_R\mathrm{}_i)(\overline{q}_a\gamma _\mu P_Rq_b)`$ (13)
$`+`$ $`_{2\overline{q}_aq_b}^L(\overline{\mathrm{}}_j\gamma ^\mu P_L\mathrm{}_i)(\overline{q}_a\gamma _\mu P_Rq_b)+_{2\overline{q}_aq_b}^R(\overline{\mathrm{}}_j\gamma ^\mu P_R\mathrm{}_i)(\overline{q}_a\gamma _\mu P_Lq_b)`$
$`+`$ $`_{3\overline{q}_aq_b}^L(\overline{\mathrm{}}_jP_L\mathrm{}_i)(\overline{q}_aP_Lq_b)+_{3\overline{q}_aq_b}^R(\overline{\mathrm{}}_jP_R\mathrm{}_i)(\overline{q}_aP_Rq_b)`$
$`+`$ $`\overline{}_{3\overline{q}_aq_b}^L(\overline{\mathrm{}}_jP_L\mathrm{}_i)(\overline{q}_aP_Rq_b)+\overline{}_{3\overline{q}_aq_b}^R(\overline{\mathrm{}}_jP_R\mathrm{}_i)(\overline{q}_aP_Lq_b)`$
$`+`$ $`_{4\overline{q}_aq_b}^L(\overline{\mathrm{}}_j\sigma _{\mu \nu }P_L\mathrm{}_i)(\overline{q}_a\sigma ^{\mu \nu }P_Rq_b)+_{4\overline{q}_aq_b}^R(\overline{\mathrm{}}_j\sigma _{\mu \nu }P_R\mathrm{}_i)(\overline{q}_a\sigma ^{\mu \nu }P_Lq_b)],`$
where $`s_W=\mathrm{sin}\theta _W`$ and $`c_W=\mathrm{cos}\theta _W`$, $`Q_q`$ is the quark charge in units of $`e`$, and $`I_{3q}`$ is weak quark isospin. $`D(xy)`$ is the Green function for the massless scalar particle, contained in the photon propagator. The structure of the photon-penguin term in the effective Lagrangian is a consequence of the gauge invariance. Especially, the first term must contain $`\begin{array}{c}\text{/}\\ \hfill \end{array}_x_{x\mu }`$, which was neglected in reference . The information on the model under consideration is contained in the form factors $`𝒫_{a\gamma }^{L,R}`$, $`a=1,2`$, $`𝒫_Z^{L,R}`$, $`\overline{}_{3\overline{q}_aq_b}^L`$, and $`_{a\overline{q}_aq_b}^{L,R}`$, $`a=1,2,3,4`$. In the following three subsections these form factors are given for the MSSM.
#### 2.2.1 The photon-penguin form factors
The amplitude for $`\mathrm{}_i\mathrm{}_j\gamma ^{}`$ for an off-mass-shell photon process is obtained the from corresponding part of the effective Lagrangian neglecting the quark current and the photon propagator,
$$_{\mu }^{}{}_{}{}^{\gamma }=iT_\mu ^\gamma =e\overline{u}_\mathrm{}_j\left[(q^2\gamma _\mu q_\mu \begin{array}{c}\text{/}\\ q\hfill \end{array})(𝒫_{1\gamma }^LP_L+𝒫_{1\gamma }^RP_R)+i\sigma _{\mu \nu }q^\nu (𝒫_{2\gamma }^LP_L+𝒫_{2\gamma }^RP_R)\right]u_\mathrm{}_i.$$
(14)
The amplitude is written without photon polarization vector.
In the MSSM the photon-penguin amplitude has two contributions, a chargino and a neutralino contribution. That reflects in the structure of the form factors,
$`𝒫_{a\gamma }^{L,R}`$ $`=`$ $`𝒫_{a\gamma }^{(C)L,R}+𝒫_{a\gamma }^{(N)L,R},a=1,2`$ (15)
with $`C`$ and $`N`$ subscript denoting the chargino and neutralino part of a form factor.
Because of the gauge invariance, the zeroth-order and the first-order term in Taylor expansion in momenta and masses of incoming and outgoing particles are equal zero. Here, the second order term in Taylor expansion is presented, and higher order terms are neglected.
The neutralino contributions are
$`𝒫_{1\gamma }^{(N)L}`$ $`=`$ $`{\displaystyle \frac{i}{576\pi ^2}}N_{jAX}^{R(e)}N_{iAX}^{R(e)}{\displaystyle \frac{1}{m_{\stackrel{~}{e}_X}^2}}`$ (16)
$`\times `$ $`{\displaystyle \frac{11(x_{AX}^0)^318(x_{AX}^0)^2+9(x_{AX}^0)26(x_{AX}^0)^3\mathrm{ln}(x_{AX}^0)}{(1x_{AX}^0)^4}},`$
$`𝒫_{1\gamma }^{(N)R}`$ $`=`$ $`𝒫_{1\gamma }^{(N)L}|_{LR},`$ (17)
$`𝒫_{2\gamma }^{(N)L}`$ $`=`$ $`{\displaystyle \frac{i}{32\pi ^2}}[N_{jAX}^{R(e)}N_{iAX}^{R(e)}(1)m_j{\displaystyle \frac{1}{m_{\stackrel{~}{e}_X}^2}}`$ (18)
$`\times `$ $`{\displaystyle \frac{2(x_{AX}^0)^3+3(x_{AX}^0)^26(x_{AX}^0)+16(x_{AX}^0)^2\mathrm{ln}(x_{AX}^0)}{6(1x_{AX}^0)^4}}`$
$`+`$ $`N_{jAX}^{L(e)}N_{iAX}^{L(e)}(1)m_i{\displaystyle \frac{1}{m_{\stackrel{~}{e}_X}^2}}`$
$`\times `$ $`{\displaystyle \frac{2(x_{AX}^0)^3+3(x_{AX}^0)^26(x_{AX}^0)+16(x_{AX}^0)^2\mathrm{ln}(x_{AX}^0)}{6(1x_{AX}^0)^4}}`$
$`+`$ $`N_{jAX}^{L(e)}N_{iAX}^{R(e)}(1)m_{\stackrel{~}{\chi }_A^0}{\displaystyle \frac{1}{m_{\stackrel{~}{e}_X}^2}}{\displaystyle \frac{(x_{AX}^0)^2+1+2(x_{AX}^0)\mathrm{ln}(x_{AX}^0)}{(1x_{AX}^0)^3}}],`$
$`𝒫_{2\gamma }^{(N)R}`$ $`=`$ $`𝒫_{2\gamma }^{(N)L}|_{LR},`$ (19)
where $`x_{AX}^0=M_{\stackrel{~}{\chi }_A^0}^2/m_{\stackrel{~}{e}_X}^2`$. The chargino contributions are given by
$`𝒫_{1\gamma }^{(C)L}`$ $`=`$ $`{\displaystyle \frac{i}{576\pi ^2}}C_{jAX}^{R(e)}C_{iAX}^{R(e)}{\displaystyle \frac{1}{m_{\stackrel{~}{\nu }_X}^2}}`$ (20)
$`\times `$ $`{\displaystyle \frac{1645(x_{AX}^{})+36(x_{AX}^{})^27(x_{AX}^{})^3+6(23(x_{AX}^{}))\mathrm{ln}(x_{AX}^{})}{(1x_{AX}^{})^4}},`$
$`𝒫_{1\gamma }^{(C)R}`$ $`=`$ $`𝒫_{1\gamma }^{(C)L}|_{LR},`$ (21)
$`𝒫_{2\gamma }^{(C)L}`$ $`=`$ $`{\displaystyle \frac{i}{32\pi ^2}}[C_{jAX}^{R(e)}C_{iAX}^{R(e)}m_j{\displaystyle \frac{1}{m_{\stackrel{~}{\nu }_X}^2}}`$ (22)
$`\times `$ $`{\displaystyle \frac{2+3(x_{AX}^{})6(x_{AX}^{})^2+(x_{AX}^{})^3+6(x_{AX}^{})\mathrm{ln}(x_{AX}^{})}{6(1x_{AX}^{})^4}}`$
$`+`$ $`C_{jAX}^{L(e)}C_{iAX}^{L(e)}m_i{\displaystyle \frac{1}{m_{\stackrel{~}{\nu }_X}^2}}`$
$`\times `$ $`{\displaystyle \frac{2+3(x_{AX}^{})6(x_{AX}^{})^2+(x_{AX}^{})^3+6(x_{AX}^{})\mathrm{ln}(x_{AX}^{})}{6(1x_{AX}^{})^4}}`$
$`+`$ $`C_{jAX}^{L(e)}C_{iAX}^{R(e)}m_{\stackrel{~}{\chi }_A^{}}{\displaystyle \frac{1}{m_{\stackrel{~}{\nu }_X}^2}}{\displaystyle \frac{3+4(x_{AX}^{})(x_{AX}^{})^22\mathrm{ln}(x_{AX}^{})}{(1x_{AX}^{})^3}}],`$
$`𝒫_{2\gamma }^{(C)R}`$ $`=`$ $`𝒫_{2\gamma }^{(C)L}|_{LR},`$ (23)
where $`x_{AX}^{}=M_{\stackrel{~}{\chi }_A^{}}^2/m_{\stackrel{~}{\nu }_X}^2`$. The form factor contributions are written in the same way as in Ref. , including the loop functions into the expressions, to make the comparison with the results of Ref. easy. Both chargino and neutralino parts of the form factors agree with the corresponding form factors in the reference if the terms proportional to the mass of the lighter mesons $`m_j`$ are neglected. Nevertheless, these terms cannot be neglected, because the constants $`N_{jAX}^{L,R(e)}`$ and $`C_{jAX}^{L,R(e)}`$ (see Appendix B) also depend on the lepton masses is such a way that in some cases the term proportional to $`m_j`$ is larger than the term proportional to the mass of decaying lepton $`m_i`$.
#### 2.2.2 The Z-penguin form factors
The amplitude for $`\mathrm{}_i\mathrm{}_jZ^{}`$, for the off-mass-shell Z-boson process, (obtained analogously as the photon-penguin amplitude) is
$`_\mu ^Z`$ $`=`$ $`{\displaystyle \frac{i}{(4\pi )^2}}g\overline{u}_{l_j}\left[\gamma _\mu P_L(𝒫_Z^{(C)L}+𝒫_Z^{(N)L})+\gamma _\mu P_R(𝒫_Z^{(C)R}+𝒫_Z^{(N)R})\right]u_{l_i},`$ (24)
where $`𝒫_Z^{(C)L,R}`$ and $`𝒫_Z^{(N)L,R}`$ are the chargino and neutralino part of the total form factors, $`𝒫_Z^{L,R}`$. Here are the expressions for these form factors:
$`𝒫_Z^{(C)L}`$ $`=`$ $`C_{jBX}^{R(e)}C_{iAX}^{R(e)}[E_{BA}^{L(\stackrel{~}{\chi }^{})}m_{\stackrel{~}{\chi }_B^{}}m_{\stackrel{~}{\chi }_A^{}}F_1(m_{\stackrel{~}{\nu }_X}^2,m_{\stackrel{~}{\chi }_A^{}}^2,m_{\stackrel{~}{\chi }_B^{}}^2)`$ (25)
$``$ $`2E_{BA}^{R(\stackrel{~}{\chi }^{})}F_2(m_{\stackrel{~}{\nu }_X}^2,m_{\stackrel{~}{\chi }_A^{}}^2,m_{\stackrel{~}{\chi }_B^{}}^2)+\delta _{AB}G_{Ze}^Lf_1(m_{\stackrel{~}{\nu }_X}^2,m_{\stackrel{~}{\chi }_A^{}}^2)]`$
$`+`$ $`\left\{C_{jBX}^{R(e)}E_{BA}^{L(\stackrel{~}{\chi }^{})}C_{iAX}^{L(e)}[2F_2(m_{\stackrel{~}{\nu }_X}^2,m_{\stackrel{~}{\chi }_A^{}}^2,m_{\stackrel{~}{\chi }_B^{}}^2)]\right\},`$
$`𝒫_Z^{(N)L}`$ $`=`$ $`N_{jBY}^{R(e)}N_{iAX}^{R(e)}\left[2D_{YX}^{(\stackrel{~}{e})}F_2(m_{\stackrel{~}{\chi }_A^0}^2,m_{\stackrel{~}{e}_X}^2,m_{\stackrel{~}{e}_Y}^2)+\delta _{YX}G_{Ze}^Lf_2(m_{\stackrel{~}{\chi }_A^0}^2,m_{\stackrel{~}{e}_X}^2)\right]`$ (26)
$`+`$ $`\{N_{jBX}^{L(e)}E_{BA}^{L(\stackrel{~}{\chi }^0)}N_{iAX}^{R(e)}[m_{\stackrel{~}{\chi }_B^0}m_{\stackrel{~}{\chi }_A^0}F_1(m_{\stackrel{~}{e}_X}^2,m_{\stackrel{~}{\chi }_A^0}^2,m_{\stackrel{~}{\chi }_B^0}^2)]`$
$`+`$ $`N_{jBX}^{R(e)}E_{BA}^{L(\stackrel{~}{\chi }^0)}N_{iAX}^{L(e)}[2F_2(m_{\stackrel{~}{e}_X}^2,m_{\stackrel{~}{\chi }_A^0}^2,m_{\stackrel{~}{\chi }_B^0}^2)]\},`$
$`𝒫_Z^{(C)R}`$ $`=`$ $`𝒫_Z^{(C)L}(LR),`$ (27)
$`𝒫_Z^{(N)R}`$ $`=`$ $`𝒫_Z^{(N)L}(LR).\text{ }E_{BA}^{R(\stackrel{~}{\chi }^0)}=E_{BA}^{L(\stackrel{~}{\chi }^0)}.`$ (28)
$`G_{Ze}^L`$ and $`G_{Ze}^R`$ are constants appearing in the SM $`Ze_ie_j`$ vertices,
$`_{Ze_ie_j}`$ $`=`$ $`g\gamma _\mu \delta _{ij}\{G_{Ze}^LP_L+G_{Ze}^RP_R\}`$ (29)
$`g\gamma _\mu \delta _{ij}\left\{\left[{\displaystyle \frac{1}{2}}{\displaystyle \frac{1}{c_W}}+{\displaystyle \frac{s_W^2}{c_W}}\right]P_L+\left[{\displaystyle \frac{s_W^2}{c_W}}\right]P_R\right\},`$
while $`E_{BA}^{R(\stackrel{~}{\chi }^{})}`$ and $`E_{BA}^{L,R(\stackrel{~}{\chi }^0)}`$ are constants in the $`Z`$-boson—chargino and $`Z`$-boson—neutralino vertices, and $`D_{YX}^{(\stackrel{~}{e})}`$ is a constant in the $`Z`$-boson—selectron vertex. These constants are defined in Appendix C. $`F_1(a,b,c)`$ and $`F_2(a,b,c)`$ are loop functions contained in the triangle-diagram part of the amplitude, and $`f_1`$ and $`f_2`$ are the loop functions coming from the self-energy part of the amplitude. They are given in Appendix D.
The terms in Eqs. (25) and (26) which have a corresponding contribution in the photon amplitude (the leading order photon-penguin amplitude comes from six Feynman diagrams, while the $`Z`$-boson-penguin amplitude has eight Feynman-diagram contributions) have been compared by replacing $`Z`$-boson vertices with corresponding photon vertices, and an agreement was found. The remaining two Feynman diagram contributions which are embraced by curly brackets in Eqs. (25) and (26) have been checked carefully. The new terms in $`𝒫_Z^{(C)L}`$ in comparison with Ref. are third (self-energy–type term) and fourth term. Further, neither of our terms in $`𝒫_Z^{(C)R}`$ does agree with amplitude in Ref. , although the expression in the curly brackets is almost equal to it.
#### 2.2.3 The box form factors
The box contribution to the SL LFV $`\mathrm{}\mathrm{}_i\overline{q}_aq_b`$ amplitude comes from two box-diagrams in the leading order of perturbation theory. The box-amplitude reads
$`_{\mathrm{box}}`$ $`=`$ $`{\displaystyle \frac{i}{(4\pi )^2}}{\displaystyle \underset{\overline{q}_aq_b=\overline{u}u,\overline{d}d,\overline{s}s,\overline{d}s,\overline{s}d}{}}[_{1\overline{q}_aq_b}^L(\overline{u}_\mathrm{}_j\gamma ^\mu P_Lu_\mathrm{}_i)(\overline{u}_{q_a}\gamma _\mu P_Lv_{q_b})+_{1\overline{q}_aq_b}^R(\overline{u}_\mathrm{}_j\gamma ^\mu P_Ru_\mathrm{}_i)(\overline{u}_{q_a}\gamma _\mu P_Rv_{q_b})`$ (30)
$`+`$ $`_{2\overline{q}_aq_b}^L(\overline{u}_\mathrm{}_j\gamma ^\mu P_Lu_\mathrm{}_i)(\overline{u}_{q_a}\gamma _\mu P_Rv_{q_b})+_{2\overline{q}_aq_b}^R(\overline{u}_\mathrm{}_j\gamma ^\mu P_Ru_\mathrm{}_i)(\overline{u}_{q_a}\gamma _\mu P_Lv_{q_b})`$
$`+`$ $`_{3\overline{q}_aq_b}^L(\overline{u}_\mathrm{}_jP_Lu_\mathrm{}_i)(\overline{u}_{q_a}P_Lv_{q_b})+_{3\overline{q}_aq_b}^R(\overline{u}_\mathrm{}_jP_Ru_\mathrm{}_i)(\overline{u}_{q_a}P_Rv_{q_b})`$
$`+`$ $`\overline{}_{3\overline{q}_aq_b}^L(\overline{u}_\mathrm{}_jP_Lu_\mathrm{}_i)(\overline{u}_{q_a}P_Rv_{q_b})+\overline{}_{3\overline{q}_aq_b}^R(\overline{u}_\mathrm{}_jP_Ru_\mathrm{}_i)(\overline{u}_{q_a}P_Lv_{q_b})`$
$`+`$ $`_{4\overline{q}_aq_b}^L(\overline{u}_\mathrm{}_j\sigma _{\mu \nu }P_Lu_\mathrm{}_i)(\overline{u}_{q_a}\sigma ^{\mu \nu }P_Lv_{q_b})`$
$`+`$ $`_{4\overline{q}_aq_b}^R(\overline{u}_\mathrm{}_j\sigma _{\mu \nu }P_Ru_\mathrm{}_i)(\overline{u}_{q_a}\sigma ^{\mu \nu }P_Rv_{q_b})].`$
The very rich structure of the box-diagram amplitude is a consequence of the Fiertz transformation of the terms with product of lepton-quark and quark-lepton vector and axial-vector currents. All currents permitted by the Dirac algebra do appear. Each box-amplitude form factor has a chargino (C) and a neutralino (N) contribution.
$`_{i\overline{q}_aq_b}^{L,R}`$ $`=`$ $`_{i\overline{q}_aq_b}^{(N)L,R}+_{i\overline{q}_aq_b}^{(C)L,R},`$ (31)
$`\overline{}_{3\overline{q}_aq_b}^{L,R}`$ $`=`$ $`\overline{}_{3\overline{q}_aq_b}^{(N)L,R}+\overline{}_{3\overline{q}_aq_b}^{(C)L,R}.`$ (32)
Here and in the following equations $`\overline{q}_aq_b`$ assume the values appearing in the sum in Eq. (30). Neutralino contributions read
$`_{1\overline{q}_aq_b}^{(N)L}`$ $`=`$ $`{\displaystyle \frac{1}{4}}d_2(M_{\stackrel{~}{\chi }_A^0}^2,M_{\stackrel{~}{\chi }_B^0}^2,m_{\stackrel{~}{e}_X}^2,m_{\stackrel{~}{q}_Y}^2)N_{iAX}^{R(e)}N_{jBX}^{R(e)}N_{bBY}^{R(q)}N_{aAY}^{R(q)}`$ (33)
$`+`$ $`{\displaystyle \frac{1}{2}}d_0(M_{\stackrel{~}{\chi }_A^0}^2,M_{\stackrel{~}{\chi }_B^0}^2,m_{\stackrel{~}{e}_X}^2,m_{\stackrel{~}{q}_Y}^2)M_{\stackrel{~}{\chi }_A^0}M_{\stackrel{~}{\chi }_B^0}N_{iAX}^{R(e)}N_{jBX}^{R(e)}N_{bAY}^{R(q)}N_{aBY}^{R(q)},`$
$`_{2\overline{q}_aq_b}^{(N)L}`$ $`=`$ $`{\displaystyle \frac{1}{4}}d_2(M_{\stackrel{~}{\chi }_A^0}^2,M_{\stackrel{~}{\chi }_B^0}^2,m_{\stackrel{~}{e}_X}^2,m_{\stackrel{~}{q}_Y}^2)N_{iAX}^{R(e)}N_{jBX}^{R(e)}N_{bAY}^{L(q)}N_{aBY}^{L(q)}`$ (34)
$``$ $`{\displaystyle \frac{1}{2}}d_0(M_{\stackrel{~}{\chi }_A^0}^2,M_{\stackrel{~}{\chi }_B^0}^2,m_{\stackrel{~}{e}_X}^2,m_{\stackrel{~}{q}_Y}^2)M_{\stackrel{~}{\chi }_A^0}M_{\stackrel{~}{\chi }_B^0}N_{iAX}^{R(e)}N_{jBX}^{R(e)}N_{bAY}^{L(q)}N_{aBY}^{L(q)},`$
$`_{3\overline{q}_aq_b}^{(N)L}`$ $`=`$ $`d_0(M_{\stackrel{~}{\chi }_A^0}^2,M_{\stackrel{~}{\chi }_B^0}^2,m_{\stackrel{~}{e}_X}^2,m_{\stackrel{~}{q}_Y}^2)M_{\stackrel{~}{\chi }_A^0}M_{\stackrel{~}{\chi }_B^0}\{{\displaystyle \frac{1}{2}}N_{iAX}^{R(e)}N_{jBX}^{L(e)}N_{bBY}^{R(q)}N_{aAY}^{L(q)}`$ (35)
$``$ $`{\displaystyle \frac{1}{2}}N_{iAX}^{R(e)}N_{jBX}^{L(e)}N_{bAY}^{R(q)}N_{aBY}^{L(q)}\},`$
$`\overline{}_{3\overline{q}_aq_b}^{(N)L}`$ $`=`$ $`d_2(M_{\stackrel{~}{\chi }_A^0}^2,M_{\stackrel{~}{\chi }_B^0}^2,m_{\stackrel{~}{e}_X}^2,m_{\stackrel{~}{q}_Y}^2)\{{\displaystyle \frac{1}{2}}N_{iAX}^{R(e)}N_{jBX}^{L(e)}N_{bBY}^{L(q)}N_{aAY}^{R(q)}`$ (36)
$``$ $`{\displaystyle \frac{1}{2}}N_{iAX}^{R(e)}N_{jBX}^{L(e)}N_{bAY}^{L(q)}N_{aBY}^{R(q)}\},`$
$`_{4\overline{q}_aq_b}^{(N)L}`$ $`=`$ $`{\displaystyle \frac{1}{8}}d_0(M_{\stackrel{~}{\chi }_A^0}^2,M_{\stackrel{~}{\chi }_B^0}^2,m_{\stackrel{~}{e}_X}^2,m_{\stackrel{~}{q}_Y}^2)M_{\stackrel{~}{\chi }_A^0}M_{\stackrel{~}{\chi }_B^0}\{N_{iAX}^{R(e)}N_{jBX}^{L(e)}N_{bBY}^{R(q)}N_{aAY}^{L(q)}`$ (37)
$``$ $`N_{iAX}^{R(e)}N_{jBX}^{L(e)}N_{bAY}^{R(q)}N_{aBY}^{L(q)}\},`$
$`_{i\overline{q}_aq_b}^{(N)R}`$ $`=`$ $`_{i\overline{q}_aq_b}^{(N)L}|_{LR}(i=1,\mathrm{},4)`$ (38)
$`\overline{}_{3\overline{q}_aq_b}^{(N)R}`$ $`=`$ $`\overline{}_{3\overline{q}_aq_b}^{(N)L}|_{LR}.`$ (39)
The chargino contributions are
$`_{1\overline{q}_aq_b}^{(C)L}`$ $`=`$ $`{\displaystyle \frac{1}{4}}d_2(M_{\stackrel{~}{\chi }_A^{}}^2,M_{\stackrel{~}{\chi }_B^{}}^2,m_{\stackrel{~}{\nu }_X}^2,m_{\stackrel{~}{q^{}}_Y}^2)C_{iAX}^{R(e)}C_{jBX}^{R(e)}C_{bBY}^{R(q)}C_{aAY}^{R(q)}\delta _{qd}`$ (40)
$`+`$ $`{\displaystyle \frac{1}{2}}d_0(M_{\stackrel{~}{\chi }_A^{}}^2,M_{\stackrel{~}{\chi }_B^{}}^2,m_{\stackrel{~}{\nu }_X}^2,m_{\stackrel{~}{q^{}}_Y}^2)M_{\stackrel{~}{\chi }_A^{}}M_{\stackrel{~}{\chi }_B^{}}C_{iAX}^{R(e)}C_{jBX}^{R(e)}C_{bAY}^{R(q)}C_{aBY}^{R(q)}\delta _{qu},`$
$`_{2\overline{q}_aq_b}^{(C)L}`$ $`=`$ $`{\displaystyle \frac{1}{4}}d_2(M_{\stackrel{~}{\chi }_A^{}}^2,M_{\stackrel{~}{\chi }_B^{}}^2,m_{\stackrel{~}{\nu }_X}^2,m_{\stackrel{~}{q^{}}_Y}^2)C_{iAX}^{R(e)}C_{jBX}^{R(e)}C_{bAY}^{L(q)}C_{aBY}^{L(q)}\delta _{qd}`$ (41)
$``$ $`{\displaystyle \frac{1}{2}}d_0(M_{\stackrel{~}{\chi }_A^{}}^2,M_{\stackrel{~}{\chi }_B^{}}^2,m_{\stackrel{~}{\nu }_X}^2,m_{\stackrel{~}{q^{}}_Y}^2)M_{\stackrel{~}{\chi }_A^{}}M_{\stackrel{~}{\chi }_B^{}}C_{iAX}^{R(e)}C_{jBX}^{R(e)}C_{bAY}^{L(q)}C_{aBY}^{L(q)}\delta _{qu},`$
$`_{3\overline{q}_aq_b}^{(C)L}`$ $`=`$ $`d_0(M_{\stackrel{~}{\chi }_A^{}}^2,M_{\stackrel{~}{\chi }_B^{}}^2,m_{\stackrel{~}{\nu }_X}^2,m_{\stackrel{~}{q^{}}_Y}^2)M_{\stackrel{~}{\chi }_A^{}}M_{\stackrel{~}{\chi }_B^{}}\{{\displaystyle \frac{1}{2}}C_{iAX}^{R(e)}C_{jBX}^{L(e)}C_{bBY}^{R(q)}C_{aAY}^{L(q)}\delta _{qd}`$ (42)
$``$ $`{\displaystyle \frac{1}{2}}C_{iAX}^{R(e)}C_{jBX}^{L(e)}C_{bAY}^{R(q)}C_{aBY}^{L(q)}\delta _{qu}\},`$
$`\overline{}_{3\overline{q}_aq_b}^{(C)L}`$ $`=`$ $`d_2(M_{\stackrel{~}{\chi }_A^{}}^2,M_{\stackrel{~}{\chi }_B^{}}^2,m_{\stackrel{~}{\nu }_X}^2,m_{\stackrel{~}{q^{}}_Y}^2)\{{\displaystyle \frac{1}{2}}C_{iAX}^{R(e)}C_{jBX}^{L(e)}C_{bBY}^{L(q)}C_{aAY}^{R(q)}\delta _{qd}`$ (43)
$``$ $`{\displaystyle \frac{1}{2}}C_{iAX}^{R(e)}C_{jBX}^{L(e)}C_{bAY}^{L(q)}C_{aBY}^{R(q)}\delta _{qu}\},`$
$`_{4\overline{q}_aq_b}^{(C)L}`$ $`=`$ $`{\displaystyle \frac{1}{8}}d_0(M_{\stackrel{~}{\chi }_A^{}}^2,M_{\stackrel{~}{\chi }_B^{}}^2,m_{\stackrel{~}{\nu }_X}^2,m_{\stackrel{~}{q^{}}_Y}^2)M_{\stackrel{~}{\chi }_A^{}}M_{\stackrel{~}{\chi }_B^{}}\{C_{iAX}^{R(e)}C_{jBX}^{L(e)}C_{bBY}^{R(q)}C_{aAY}^{L(q)}\delta _{qd}`$ (44)
$``$ $`C_{iAX}^{R(e)}C_{jBX}^{L(e)}C_{bAY}^{R(q)}C_{aBY}^{L(q)}\delta _{qu}\},`$
$`_{i\overline{q}_aq_b}^{(C)R}`$ $`=`$ $`_{i\overline{q}_aq_b}^{(C)L}|_{LR}(i=1,\mathrm{},4)`$ (45)
where for $`q=u(d)`$, $`q^{}=d(u)`$. The sum over paired indices is assumed. The Kronecker function $`\delta _{qu}`$ \[$`\delta _{qu}`$\] denotes that $`q`$ quark is one of up ($`u`$, $`c`$ or $`t`$) quarks \[one of the down quarks\]. The loop functions $`d_0`$ and $`d_2`$ are evaluated by neglecting the momenta of incoming and outgoing particles. They are listed in Appendix D.
## 3 Amplitudes and branching ratios
### 3.1 Hadronization of currents
The effective Lagrangians (the matrix elements) for photon and $`Z`$-boson part of the amplitudes for SL LFV lepton decays comprise vector and axial-vector currents, but box amplitude contains all possible quark currents permitted by Dirac algebra, that is scalar-, pseudoscalar-, vector-, axial-vector- and tensor-quark currents. To perform calculation of charged-lepton SL LFV decays rates these currents have be converted into meson currents comprising mesons that appear in the possible final products of the charged-lepton SL LFV decays we study here. The hadronization procedure we use here is not exact in the sense that we do not include the sea-quark and gluon content of the meson fields, but it is precise enough to give much better than order of magnitude decay rates of the processes considered here. The quark content of the meson states is given in Appendix E. The hadronization of axial-vector current is achieved through PCAC (see e.g. (); for normalization of pseudoscalar coupling constants used here and for further details see ). The hadronization of the vector current is achieved using vector-meson-dominance assumption (see (); for normalization of the vector meson-decay constant and details see ). Hadronization of scalar currents is achieved by comparing the quark sector of the SM Lagrangian and the corresponding effective meson Lagrangian (for applications in content of LFV and details see ). Hadronization of the pseudoscalar current is obtained by the same procedure as for the scalar current. The results obtained by using this procedure is equal to the result obtained by using equation of motion for current-quark masses (see e.g. ) and results for hadronization of axial-vector current up to the difference of up and down quark masses or up to the difference of pseudoscalar decay constants. The hadronization of the tensor-quark currents is obtained by comparing the derivative of tensor-quark current with vector-quark current and using the equations of motion for current quark masses. The difference between terms, one containing the derivative of the incoming quark field and the other containing derivative of the outgoing quark field, have been neglected. The error expected from this approximation is proportional to the amount of breaking of $`SU(3)_{\mathrm{flavour}}`$ symmetry. The tensor currents are proportional to the current quark masses, and therefore give smaller contribution than the other quark currents. Therefore, the error introduced by this approximation in the total SL LFV amplitude is small.
Here we summarize the basic quantities needed to describe the hadronization of quark currents.
1. The pseudoscalar meson decay constants ($`f_P`$, $`P=\pi ^0,\eta ,\eta ^{},K^0,\overline{K}^0`$);
2. The constants $`\gamma _V`$ ($`V=\rho ^0,\varphi ,\omega ,K^0,\overline{K}^0`$) defining the vector meson decay constants ($`f_Vm_V^2/\gamma _V`$);
3. The mixing angles $`\theta _P`$ and $`\theta _V`$ defining the physical meson-nonet states in terms of the unphysical singlet and octet meson states;
4. The parameter $`r`$ ($`m_u`$, $`m_d`$ and $`m_s`$ are current quark masses),
$`r={\displaystyle \frac{2m_{\pi ^+}^2}{m_u+m_d}}={\displaystyle \frac{2m_{K^0}^2}{m_d+m_s}}={\displaystyle \frac{2m_{K^+}^2}{m_u+m_s}},`$ (46)
that appears in the hadronization procedure for scalar and pseudoscalar currents.
Having the identification of the quark currents with the corresponding meson currents, achieved by the above hadronization procedure, one can write down the effective Lagrangian as a sum of terms with an incoming lepton field $`\mathrm{}_i`$ and outgoing lepton field $`\mathrm{}_j`$ and pseudoscalar meson ($`P`$) or vector meson ($`V`$) field(s). This Lagrangian directly gives the amplitudes for the $`\mathrm{}_i\mathrm{}_jP(V)`$ processes. Amplitudes with a pseudoscalar meson in the final state contributions come from the pseudoscalar and axial-vector coupling part of the effective Lagrangian, while the amplitudes with vector mesons have vector and tensor coupling contributions. Only the scalar coupling gives no contribution to the one-meson processes in the final state, $`\mathrm{}_i\mathrm{}_jP(V)`$. They contribute only to the processes with two pseudoscalar mesons in the final state, $`\mathrm{}_i\mathrm{}_jP_1P_2`$.
### 3.2 Vector-meson–pseudoscalar-meson interactions
The processes with two pseudoscalar mesons in the final state are generated by the scalar-quark-current part of the effective Lagrangian, and vector-quark- and tensor-quark-current of the effective Lagrangian. The scalar-quark-current part of the effective Lagrangian produces two pseudoscalar fields directly. The vector-quark- and tensor-quark-current parts produce a resonant vector meson state (V), which decays into two pseudoscalar mesons (P). The $`VPP`$ interactions necessary for description of the $`VPP`$ interactions appearing in the charged-lepton SL LFV decays are described by the part of the meson Lagrangian containing these $`VPP`$ vertices
$`_{VPP}`$ $`=`$ $`{\displaystyle \frac{ig_{\rho \pi \pi }}{2}}\{\rho ^{0,\mu }(2\pi ^+\underset{\mu }{\overset{}{}}\pi ^{}+K^+\underset{\mu }{\overset{}{}}K^{}K^0\underset{\mu }{\overset{}{}}\overline{K}^0)`$ (47)
$`+`$ $`\sqrt{3}s_V\omega ^\mu \left(K^+\underset{\mu }{\overset{}{}}K^{}+K^0\underset{\mu }{\overset{}{}}\overline{K}^0\right)+\sqrt{3}c_V\varphi ^\mu \left(K^+\underset{\mu }{\overset{}{}}K^{}+K^0\underset{\mu }{\overset{}{}}\overline{K}^0\right)`$
$`+`$ $`K^{0,\mu }\left(\sqrt{2}\pi ^+\underset{\mu }{\overset{}{}}K^{}+\pi ^0\underset{\mu }{\overset{}{}}\overline{K}^0+\sqrt{3}c_P\overline{K}^0\underset{\mu }{\overset{}{}}\eta +\sqrt{3}s_P\overline{K}^0\underset{\mu }{\overset{}{}}\eta ^{}\right)`$
$`+`$ $`\overline{K}^{0,\mu }(\sqrt{2}\pi ^{}\underset{\mu }{\overset{}{}}K^+\pi ^0\underset{\mu }{\overset{}{}}K^0\sqrt{3}c_PK^0\underset{\mu }{\overset{}{}}\eta \sqrt{3}s_PK^0\underset{\mu }{\overset{}{}}\eta ^{})\}+\mathrm{}.`$
This Lagrangian is a part of the nonlinear $`(U(3)_L\times U(3)_R)/U(3)_V`$ symmetric sigma-model Lagrangian. $`U(3)_V`$ symmetry corresponds to the vector mesons in the linear realization of the gauge equivalent $`(U(3)_L\times U(3)_R)_{\mathrm{global}}\times U(3)_V`$ linear sigma-model . One can include the $`(U(3)_L\times U(3)_R)/U(3)_V`$ breaking terms, too . That was applied to SL LFV tau-lepton decays in Ref. , but for the estimates of the charged-lepton SL LFV decays it is unnecessary complication, and we will not consider it here.
From Eq. (47) one can read of the $`c_{VP_1P_2}`$ couplings in terms of $`g_{\rho \pi \pi }`$ coupling. For instance, $`c_{\rho ^0K^+K^{}}=\frac{1}{2}g_{\rho \pi \pi }`$.
When the amplitudes with vector-meson resonance(s) are formed, the square of the vector meson mass in the $`m_V^2/\gamma _V`$, appearing in every vector meson decay constant, has to be replaced with $`(m_V^2im_V\mathrm{\Gamma }_V)/\gamma _V`$, where $`\mathrm{\Gamma }_V`$ is the decay width of the vector meson .
### 3.3 Charged-lepton SL LFV with one meson in the final state
Now we can write down all amplitudes we are interested in. (Notice that the axial-vector mesons ($`A`$) are not included. They decay into three pseudoscalar mesons, and therefore the amplitudes are much more complicated. New vertices with $`AVP`$ couplings should be included, kinematics is much more involved etc. The scalar mesons are also not included. They also lead to complications.). They are given as follows.
$`i^{\mathrm{}_i\mathrm{}_jV}`$ $`=`$ $`i\overline{u}_\mathrm{}_j[(\gamma _\mu {\displaystyle \frac{q_\mu \gamma q}{q^2}})P_L𝒫_{1\gamma L}^{ijV}+(\gamma _\mu {\displaystyle \frac{q_\mu \gamma q}{q^2}})P_R𝒫_{1\gamma R}^{ijV}`$ (48)
$`+`$ $`{\displaystyle \frac{i\sigma _{\mu \nu }P_Lq^\nu }{q^2}}𝒫_{2\gamma L}^{ijV}+{\displaystyle \frac{i\sigma _{\mu \nu }P_Rq^\nu }{q^2}}𝒫_{2\gamma R}^{ijV}`$
$`+`$ $`\gamma _\mu P_L\left(𝒫_{ZL}^{ijV}+_{1L}^{ijV}\right)+\gamma _\mu P_R\left(𝒫_{ZR}^{ijV}+_{1R}^{ijV}\right)`$
$`+`$ $`i\sigma _{\mu \nu }P_Lq^\nu _{2L}^{ijV}+i\sigma _{\mu \nu }P_Rq^\nu _{2R}^{ijV}]u_\mathrm{}_i\epsilon ^\mu _V.`$
$`i^{\mathrm{}_i\mathrm{}_jP}`$ $`=`$ $`i\overline{u}_\mathrm{}_j\{[\gamma _\mu P_L(𝒫_{ZL}^{ijP}+_{1L}^{ijP})+\gamma _\mu P_R(𝒫_{ZR}^{ijP}+_{1R}^{ijP})]q^\mu `$ (49)
$`+`$ $`[P_L_{2L}^{ijP}+P_R_{2R}^{ijP}]\}u_\mathrm{}_i.`$
The form factors in Eqs. (48) and (49) are defined as follows:
$`𝒫_{a\gamma L,R}^{ijV}`$ $`=`$ $`e^2𝒫_{a\gamma }^{L,R}{\displaystyle \frac{m_V}{\sqrt{2}\gamma _V}}k_\gamma ^V,(a=1,2),(V=\rho ^0,\varphi ,\omega )`$ (50)
$`𝒫_{ZL,R}^{ijV}`$ $`=`$ $`{\displaystyle \frac{g_Z^2}{m_Z^2c_W^2}}𝒫_Z^{L,R}{\displaystyle \frac{m_V^2}{\sqrt{2}\gamma _V}}k_Z^V,(V=\rho ^0,\varphi ,\omega )`$ (51)
$`_{1L,R}^{ijV}`$ $`=`$ $`{\displaystyle \frac{m_V^2}{\sqrt{2}\gamma _V}}{\displaystyle \frac{1}{2}}[k_{\overline{u}u}^V(_{1\overline{u}u}^{L,R}+_{2\overline{u}u}^{L,R})+k_{\overline{d}d}^V(_{1\overline{d}d}^{L,R}+_{2\overline{d}d}^{L,R})`$ (52)
$`+`$ $`k_{\overline{s}s}^V\left(_{1\overline{s}s}^{L,R}+_{2\overline{s}s}^{L,R}\right)+k_{\overline{d}s}^V\left(_{1\overline{d}s}^{L,R}+_{2\overline{d}s}^{L,R}\right)`$
$`+k_{\overline{s}d}^V(_{1\overline{s}d}^{L,R}+_{2\overline{s}d}^{L,R})],(V=\rho ^0,\varphi ,\omega ,K^0,\overline{K^0}),`$
$`_{2L,R}^{ijV}`$ $`=`$ $`{\displaystyle \frac{2\sqrt{2}}{\gamma _V}}[k_{\overline{d}s}^V(m_dm_s)_{4\overline{d}s}^{L,R}`$ (53)
$`+k_{\overline{s}d}^V(m_sm_d)_{4\overline{s}d}^{L,R}],(V=K^0,\overline{K^0}),`$
$`𝒫_{ZL,R}^{ijP}`$ $`=`$ $`{\displaystyle \frac{g^2}{m_Z^2c_W^2}}𝒫_Z^{L,R}(\sqrt{2}f_P)k_Z^P,(P=\pi ^0,\eta ,\eta ^{})`$ (54)
$`_{1L,R}^{ijP}`$ $`=`$ $`s_{L,R}(\sqrt{2}f_P){\displaystyle \frac{1}{2}}[k_{\overline{u}u}^P(_{1\overline{u}u}^{L,R}+_{2\overline{u}u}^{L,R})+k_{\overline{d}d}^P(_{1\overline{d}d}^{L,R}+_{2\overline{d}d}^{L,R})`$ (55)
$`+k_{\overline{s}s}^P\left(_{1\overline{s}s}^{L,R}+_{2\overline{s}s}^{L,R}\right)+k_{\overline{d}s}^P\left(_{1\overline{d}s}^{L,R}+_{2\overline{d}s}^{L,R}\right)`$
$`+k_{\overline{s}d}^P(_{1\overline{s}d}^{L,R}+_{2\overline{s}d}^{L,R})],(P=\pi ^0,\eta ,\eta ^{},K^0,\overline{K^0}),`$
$`_{2L,R}^{ijP}`$ $`=`$ $`s_{L,R}({\displaystyle \frac{ir}{2}})(\sqrt{2}f_P){\displaystyle \frac{1}{2}}[k_{\overline{u}u}^P(_{3\overline{u}u}^{L,R}+\overline{}_{3\overline{u}u}^{L,R})+k_{\overline{d}d}^P(_{3\overline{d}d}^{L,R}+\overline{}_{3\overline{d}d}^{L,R})`$ (56)
$`+k_{\overline{s}s}^P\left(_{3\overline{s}s}^{L,R}+\overline{}_{3\overline{s}s}^{L,R}\right)+k_{\overline{d}s}^P\left(_{3\overline{d}s}^{L,R}+\overline{}_{3\overline{d}s}^{L,R}\right)`$
$`+k_{\overline{s}d}^P(_{3\overline{s}d}^{L,R}+\overline{}_{3\overline{s}d}^{L,R})],(P=\pi ^0,\eta ,\eta ^{},K^0,\overline{K^0}).`$
In Eqs. (55) and (56) $`s_L=1`$ and $`s_R=1`$. The constants $`k_\gamma ^V`$, $`k_Z^V`$, $`k_Z^P`$, and $`k_{\overline{q}_aq_b}^V`$ are defined in Appendix E.
A branching ratio for the processes $`\mathrm{}_i\mathrm{}_jP`$ with unpolarized initial and final particles reads
$`B(\mathrm{}_i\mathrm{}_jP)`$ $`=`$ $`{\displaystyle \frac{1}{8\pi }}{\displaystyle \frac{1}{m_i^2}}{\displaystyle \frac{1}{\mathrm{\Gamma }_{l_i}}}{\displaystyle \frac{\lambda ^{\frac{1}{2}}(m_i^2,m_j^2,m_P^2)}{2m_i}}`$ (57)
$`\times `$ $`[(|𝒫_{ZL}^{ijP}+_{1L}^{ijP}|^2+|𝒫_{ZR}^{ijP}+_{1R}^{ijP}|^2)i_{P1}`$
$`+`$ $`\left(|_{2L}^{ijP}|^2+|_{2R}^{ijP}|^2\right)i_{P2}`$
$`+`$ $`((𝒫_{ZL}^{ijP}+_{1L}^{ijP})(𝒫_{ZR}^{ijP}+_{1R}^{ijP})^{}+c.c.)m_jm_ii_{P3}`$
$`+`$ $`((𝒫_{ZL}^{ijP}+_{1L}^{ijP})(_{2L}^{ijP})^{}+(𝒫_{ZR}^{ijP}+_{1R}^{ijP})(_{2R}^{ijP})^{}+c.c.)m_ji_{P4}`$
$`+`$ $`((𝒫_{ZL}^{ijP}+_{1L}^{ijP})(_{2R}^{ijP})^{}+(𝒫_{ZR}^{ijP}+_{1R}^{ijP})(_{2L}^{ijP})^{}+c.c.)m_ii_{P5}`$
$`+`$ $`((_{2L}^{ijP})(_{2R}^{ijP})^{}+c.c.)m_im_j],`$
where $`\mathrm{\Gamma }_{l_i}`$ is a total decay rate of the lepton $`l_i`$ and
$`i_{P1}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left((m_i^2m_j^2)^2(m_i^2+m_j^2)m_P^2\right),`$
$`i_{P2}`$ $`=`$ $`{\displaystyle \frac{1}{2}}(m_i^2+m_j^2m_P^2),`$
$`i_{P3}`$ $`=`$ $`m_P^2,`$
$`i_{P4}`$ $`=`$ $`{\displaystyle \frac{1}{2}}(m_i^2+m_P^2m_j^2),`$
$`i_{P5}`$ $`=`$ $`{\displaystyle \frac{1}{2}}(m_i^2m_P^2m_j^2),`$ (58)
and
$$\lambda (x,y,z)=x^2+y^2+z^22xy2xz2yz.$$
(59)
A branching ratio for the processes $`\mathrm{}_i\mathrm{}_jV`$ with unpolarized initial and final particles reads.
$`B(\mathrm{}_i\mathrm{}_jV)`$ $`=`$ $`{\displaystyle \frac{1}{8\pi }}{\displaystyle \frac{1}{\mathrm{\Gamma }_\mathrm{}_i}}{\displaystyle \frac{1}{m_i^2}}{\displaystyle \frac{\lambda ^{\frac{1}{2}}(m_i^2,m_j^2,m_V^2)}{2m_i}}`$ (60)
$`\times `$ $`[(|𝒫_{1\gamma L}^{ijV}+𝒫_{ZL}^{ijV}+_{1L}^{ijV}|^2+|𝒫_{1\gamma R}^{ijV}+𝒫_{ZR}^{ijV}+_{1R}^{ijV}|^2)i_{V1}`$
$`+`$ $`\left(\left|{\displaystyle \frac{𝒫_{2\gamma L}^{ijV}}{m_V^2}}+_{2L}^{ijV}\right|^2+\left|{\displaystyle \frac{𝒫_{2\gamma R}^{ijV}}{m_V^2}}+_{2R}^{ijV}\right|^2\right)i_{V2}`$
$`+`$ $`((𝒫_{1\gamma L}^{ijV}+𝒫_{ZL}^{ijV}+_{1L}^{ijV})(𝒫_{1\gamma R}^{ijV}+𝒫_{ZR}^{ijV}+_{1R}^{ijV})^{}+c.c.)(m_im_j)`$
$`+`$ $`(({\displaystyle \frac{𝒫_{2\gamma L}^{ijV}}{m_V^2}}+_{2L}^{ijV})({\displaystyle \frac{𝒫_{2\gamma R}^{ijV}}{m_V^2}}+_{2R}^{ijV})^{}+c.c.)(m_V^2m_im_j)`$
$`+`$ $`((𝒫_{1\gamma L}^{ijV}+𝒫_{ZL}^{ijV}+_{1L}^{ijV})({\displaystyle \frac{𝒫_{2\gamma L}^{ijV}}{m_V^2}}+_{2L}^{ijV})^{}`$
$`+`$ $`(𝒫_{1\gamma R}^{ijV}+𝒫_{ZR}^{ijV}+_{1R}^{ijV})({\displaystyle \frac{𝒫_{2\gamma R}^{ijV}}{m_V^2}}+_{2R}^{ijV})^{}+c.c.)(m_ji_{V3})`$
$`+`$ $`((𝒫_{1\gamma L}^{ijV}+𝒫_{ZL}^{ijV}+_{1L}^{ijV})({\displaystyle \frac{𝒫_{2\gamma R}^{ijV}}{m_V^2}}+_{2R}^{ijV})^{}`$
$`+`$ $`(𝒫_{1\gamma R}^{ijV}+𝒫_{ZR}^{ijV}+_{1R}^{ijV})({\displaystyle \frac{𝒫_{2\gamma L}^{ijV}}{m_V^2}}+_{2L}^{ijV})^{}+c.c.)(m_ii_{V4})],`$
where
$`i_{V1}`$ $`=`$ $`{\displaystyle \frac{1}{2m_V^2}}\left[m_V^2(m_i^2+m_j^2)+(m_i^2m_j^2)^22m_V^4\right],`$
$`i_{V2}`$ $`=`$ $`(m_i^2m_j^2)^2{\displaystyle \frac{1}{2}}m_V^2(m_i^2+m_j^2){\displaystyle \frac{1}{2}}m_V^4,`$
$`i_{V3}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left(m_i^2m_j^2+m_V^2\right),`$
$`i_{V4}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left(m_i^2m_j^2m_V^2\right).`$ (61)
### 3.4 SL LFV decays of a lepton with two pseudoscalar mesons in the final state
Amplitude for a general $`\mathrm{}_i\mathrm{}_jP_1P_2`$ decay rate is a sum of scalar-coupling contribution and resonance contributions (coming from vector- and tensor-coupling contributions)
$$i^{\mathrm{}_i\mathrm{}_jP_1P_2}=i_{res}^{\mathrm{}_i\mathrm{}_jP_1P_2}+i_1^{\mathrm{}_i\mathrm{}_jP_1P_2},$$
(62)
where
$`_{res}^{\mathrm{}_i\mathrm{}_jP_1P_2}`$ $`=`$ $`i\overline{u}_\mathrm{}_j[D_{1L}^{P_1P_2}((\begin{array}{c}\text{/}\\ p\hfill \end{array}_2\begin{array}{c}\text{/}\\ p\hfill \end{array}_1){\displaystyle \frac{m_2^2m_1^2}{q^2}}\begin{array}{c}\text{/}\\ q\hfill \end{array})P_L`$ (69)
$`+`$ $`D_{1R}^{P_1P_2}\left((\begin{array}{c}\text{/}\\ p\hfill \end{array}_2\begin{array}{c}\text{/}\\ p\hfill \end{array}_1){\displaystyle \frac{m_2^2m_1^2}{q^2}}\begin{array}{c}\text{/}\\ q\hfill \end{array}\right)P_R`$ (76)
$`+`$ $`E_{1L}^{P_1P_2}i\sigma _{\mu \nu }P_L(p_1p_2)^\mu q^\nu +E_{1R}^{P_1P_2}i\sigma _{\mu \nu }P_R(p_1p_2)^\mu q^\nu ]u_\mathrm{}_i,`$ (77)
$`i_1^{\mathrm{}_i\mathrm{}_jP_1P_2}`$ $`=`$ $`i\overline{u}_\mathrm{}_j\left[P_LA_{1P_1P_2}^L+P_RA_{1P_1P_2}^R\right]u_\mathrm{}_i.`$ (78)
In Eq. (69) $`D_{1L,R}^{P_1P_2}`$ and $`E_{1L,R}^{P_1P_2}`$ are form factors built from the trilinear $`c_{VP_1P_2}`$ couplings (defined by the Lagrangian (47)), normalized vector-meson propagators,
$$\frac{m_V^2im_V\mathrm{\Gamma }_V}{q^2m_V^2+im_V\mathrm{\Gamma }_V}$$
(79)
and form factors for $`\mathrm{}_i\mathrm{}_jV`$ processes divided by the the mass of the resonant vector meson, e.g.
$`\stackrel{~}{𝒫}_{1\gamma L,R}^{ijV}`$ $`=`$ $`{\displaystyle \frac{𝒫_{1\gamma L,R}^{ijV}}{m_V^2}}.`$ (80)
The expressions for the $`D_{1L,R}^{P_1P_2}`$ and $`E_{1L,R}^{P_1P_2}`$ form factors are
$`D_{1L,R}^{P_1P_2}`$ $`=`$ $`{\displaystyle \underset{V}{}}\left(\stackrel{~}{𝒫}_{1\gamma L,R}^{ijV}+\stackrel{~}{𝒫}_{ZL,R}^{ijV}+\stackrel{~}{}_{1L,R}^{ijV}\right){\displaystyle \frac{m_V^2im_V\mathrm{\Gamma }_V}{q^2m_V^2+im_V\mathrm{\Gamma }_V}}c_{VP_1P_2},`$ (81)
$`E_{1L,R}^{P_1P_2}`$ $`=`$ $`{\displaystyle \underset{V}{}}\left({\displaystyle \frac{\stackrel{~}{𝒫}_{2\gamma L,R}^{ijV}}{q^2}}+\stackrel{~}{}_{2L,R}^{ijV}\right){\displaystyle \frac{m_V^2im_V\mathrm{\Gamma }_V}{q^2m_V^2+im_V\mathrm{\Gamma }_V}}c_{VP_1P_2}.`$ (82)
The sum goes over neutral vector mesons only $`(V=\rho ^0,\varphi ,\omega ,K^0,\overline{K^0})`$. The coefficients of the non-resonant part of the amplitude, the constants $`A_{1P_1P_2}`$, are defined as the coefficients of the $`P_1P_2`$ product of fields contained in the matrix-valued operator
$`{\displaystyle \frac{r}{4}}{\displaystyle \underset{q_a,q_b=u,d,s}{}}(\mathrm{\Pi }^2)_{\overline{q}_bq_a}(_{3\overline{q}_aq_b}^{L,R}+\overline{}_{3\overline{q}_aq_b}^{L,R})`$ (83)
where $`\mathrm{\Pi }`$ is a matrix of pseudoscalar fields
$`\mathrm{\Pi }`$ $`=`$ $`\left(\begin{array}{ccc}\pi ^0+\frac{1}{\sqrt{3}}\eta _8+\frac{\sqrt{2}}{\sqrt{3}}\eta _1& \sqrt{2}\pi ^+& \sqrt{2}K^+\\ \sqrt{2}\pi ^{}& \pi ^0+\frac{1}{\sqrt{3}}\eta _8+\frac{\sqrt{2}}{\sqrt{3}}\eta _1& \sqrt{2}K^0\\ \sqrt{2}K^{}& \sqrt{2}\overline{K}^0& \frac{2}{\sqrt{3}}\eta _8+\frac{\sqrt{2}}{\sqrt{3}}\eta _1\end{array}\right).`$ (87)
For example,
$`A_{1\pi ^0\pi ^0}^{L,R}`$ $`=`$ $`{\displaystyle \frac{r}{4}}\left[{\displaystyle \frac{1}{2}}\left(_{3\overline{u}u}^{L,R}+\overline{}_{3\overline{u}u}^{L,R}+_{3\overline{d}d}^{L,R}+\overline{}_{3\overline{d}d}^{L,R}\right)\right].`$ (88)
Having the amplitudes one can easily evaluate the branching fractions. We assume that incoming and outgoing particles are not polarised.
$`B(\mathrm{}_i\mathrm{}_jP_1P_2)`$ $`=`$ $`{\displaystyle \frac{1}{(2\pi )^2}}{\displaystyle \frac{1}{\mathrm{\Gamma }_\mathrm{}_i}}{\displaystyle \frac{1}{32m_i^3}}{\displaystyle _{(m_1+m_2)^2}^{(m_im_j)^2}}𝑑s_{12}{\displaystyle _{s_{j2}^{min}}^{s_{j2}^{max}}}𝑑s_{j2}|^{\mathrm{}_i\mathrm{}_jP_1P_2}|^2`$ (89)
$`=`$ $`{\displaystyle \frac{1}{(2\pi )^2}}{\displaystyle \frac{1}{32m_i^3}}{\displaystyle _{(m_1+m_2)^2}^{(m_im_j)^2}}𝑑s_{12}`$
$`\times `$ $`[(|D_{1L}^{P_1P_2}|^2+|D_{1R}^{P_1P_2}|^2)\overline{I}_1`$
$`+`$ $`\left(|E_{1L}^{P_1P_2}|^2+|E_{1R}^{P_1P_2}|^2\right)\overline{I}_2`$
$`+`$ $`\left(|A_{1P_1P_2}^L|^2+|A_{1P_1P_2}^R|^2\right)\overline{I}_3`$
$`+`$ $`((D_{1L}^{P_1P_2})(D_{1R}^{P_1P_2})^{}+c.c.)\overline{I}_4`$
$`+`$ $`((E_{1L}^{P_1P_2})(E_{1R}^{P_1P_2})^{}+c.c.)\overline{I}_5`$
$`+`$ $`((A_{1P_1P_2}^L)(A_{1P_1P_2}^R)^{}+c.c.)\overline{I}_6`$
$`+`$ $`((D_{1L}^{P_1P_2})(E_{1L}^{P_1P_2})^{}+(D_{1R}^{P_1P_2})(E_{1R}^{P_1P_2})^{}+c.c.)\overline{I}_7`$
$`+`$ $`((D_{1L}^{P_1P_2})(E_{1R}^{P_1P_2})^{}+(D_{1R}^{P_1P_2})(E_{1L}^{P_1P_2})^{}+c.c.)\overline{I}_8`$
$`+`$ $`((D_{1L}^{P_1P_2})(A_{1P_1P_2}^L)^{}+(D_{1R}^{P_1P_2})(A_{1P_1P_2}^R)^{}+c.c.)\overline{I}_9`$
$`+`$ $`((D_{1L}^{P_1P_2})(A_{1P_1P_2}^R)^{}+(D_{1R}^{P_1P_2})(A_{1P_1P_2}^L)^{}+c.c.)\overline{I}_{10}`$
$`+`$ $`((E_{1L}^{P_1P_2})(A_{1P_1P_2}^L)^{}+(E_{1R}^{P_1P_2})(A_{1P_1P_2}^R)^{}+c.c.)\overline{I}_{11}].`$
The Mandelstam variables are defined as $`s_{ab}=(p_ap_b)^2`$, e.g. $`s_{12}=(p_1p_2)^2`$. The kinematical bounds on the Mandelstam variables $`s_{j2}^{min}`$ and $`s_{j2}^{max}`$ are well known . The $`\overline{I}`$ integrals read,
$`\overline{I}_1`$ $`=`$ $`2\overline{s_{j2}^2}+2\overline{s_{j2}}(e_1+e_2)+(2e_1e_2e_3e_4)\overline{1},`$
$`\overline{I}_2`$ $`=`$ $`2\overline{s_{j2}^2}(e_{10})+2\overline{s_{j2}}(e_5e_{10}+e_9e_{10}e_6e_7e_8e_7)`$
$`+`$ $`2((e_5e_6e_7+e_8e_9e_7e_5e_9e_{10}e_8e_6e_{11})+e_3(e_{11}e_{10}e_7^2))\overline{1},`$
$`\overline{I}_3`$ $`=`$ $`e_3\overline{1},`$
$`\overline{I}_4`$ $`=`$ $`m_im_je_4\overline{1},`$
$`\overline{I}_5`$ $`=`$ $`m_im_j(e_{10}e_{11}e_7^2)\overline{1},`$
$`\overline{I}_6`$ $`=`$ $`m_im_j\overline{1},`$
$`\overline{I}_7`$ $`=`$ $`m_j(e_{12}e_6)\overline{1},`$
$`\overline{I}_8`$ $`=`$ $`m_i(e_{12}e_8)\overline{1},`$
$`\overline{I}_9`$ $`=`$ $`m_j\left(\overline{s_{j2}}+(e_2)\overline{1}\right),`$
$`\overline{I}_{10}`$ $`=`$ $`m_i\left(\overline{s_{j2}}+(e_1)\overline{1}\right),`$
$`\overline{I}_{11}`$ $`=`$ $`\overline{s_{j2}}(e_8e_6)+(e_8e_9e_5e_6)\overline{1},`$ (90)
where
$$\overline{s_{j2}^n}=_{s_{j2}^{min}}^{s_{j2}^{max}}𝑑s_{j2}s_{j2}^n.$$
(91)
The quantities $`e_i`$ read
$`e_1`$ $`=`$ $`e_5{\displaystyle \frac{m_2m_1}{s_{12}}}e_8,`$
$`e_2`$ $`=`$ $`e_9{\displaystyle \frac{m_2m_1}{s_{12}}}e_6,`$
$`e_3`$ $`=`$ $`{\displaystyle \frac{1}{2}}(m_i^2+m_j^2s_{12}),`$
$`e_4`$ $`=`$ $`e_{11}{\displaystyle \frac{(m_2^2m_1^2)^2}{s_{12}}},`$
$`e_5`$ $`=`$ $`m_2^2+{\displaystyle \frac{1}{2}}(m_i^2+m_j^2s_{12}),`$
$`e_6`$ $`=`$ $`{\displaystyle \frac{1}{2}}(m_i^2m_j^2+s_{12}),`$
$`e_7`$ $`=`$ $`m_1^2m_2^2,`$
$`e_8`$ $`=`$ $`{\displaystyle \frac{1}{2}}(m_i^2m_j^2s_{12}),`$
$`e_9`$ $`=`$ $`m_1^2+{\displaystyle \frac{1}{2}}(m_i^2+m_j^2s_{12}),`$
$`e_{10}`$ $`=`$ $`s_{12},`$
$`e_{11}`$ $`=`$ $`2m_1^2+2m_2^2s_{12},`$
$`e_{12}`$ $`=`$ $`e_4.`$ (92)
## 4 Minimal SO(10) model and its predictions
Now we are ready to estimate the SL LFV processes. In order to perform a concrete evaluation for the SL LFV processes, we need an information on the Yukawa couplings. In this paper, we make use of the minimal $`SO(10)`$ model, as an example which gives a precise information for the neutrino-Dirac-Yukawa couplings. We begin with an overreview of the minimal SUSY $`SO(10)`$ model proposed in and recently analysed in detail in Ref. . Even when we concentrate our discussion on the issue how to reproduce the realistic fermion mass matrices in the $`SO(10)`$ model, there are lots of possibilities for introduction of Higgs multiplets. The minimal supersymmetric $`SO(10)`$ model is the one where only one 10 and one $`\overline{\mathrm{𝟏𝟐𝟔}}`$ Higgs multiplet have Yukawa couplings (superpotential) with 16 matter multiplets. Therefore, the quark and lepton mass matrices can be described as
$`M_u`$ $`=`$ $`c_{10}M_{10}+c_{126}M_{126},`$
$`M_d`$ $`=`$ $`M_{10}+M_{126},`$
$`M_D`$ $`=`$ $`c_{10}M_{10}3c_{126}M_{126},`$
$`M_e`$ $`=`$ $`M_{10}3M_{126},`$
$`M_R`$ $`=`$ $`c_RM_{126},`$ (93)
where $`M_u`$, $`M_d`$, $`M_D`$, $`M_e`$ and $`M_R`$ denote up-type quark, down-type quark, neutrino Dirac, charged-lepton and right-handed neutrino Majorana mass matrices, respectively. Note that all the quark and lepton mass matrices are characterized by only two basic mass matrices, $`M_{10}`$ and $`M_{126}`$, and three complex coefficients $`c_{10}`$, $`c_{126}`$ and $`c_R`$.
The mass matrix formulas in Eq. (93) lead to the GUT relation among the quark and lepton mass matrices,
$`M_e=c_d\left(M_d+\kappa M_u\right),`$ (94)
where
$`c_d`$ $`=`$ $`{\displaystyle \frac{3c_{10}+c_{126}}{c_{10}c_{126}}},`$ (95)
$`\kappa `$ $`=`$ $`{\displaystyle \frac{4}{3c_{10}+c_{126}}}.`$ (96)
Without loss of generality, we can start with the basis where $`M_u`$ is real and diagonal, $`M_u=D_u`$. Since $`M_d`$ is the symmetric matrix, it can be described as $`M_d=V_{\mathrm{CKM}}^{}D_dV_{\mathrm{CKM}}^{}`$ by using the CKM matrix $`V_{\mathrm{CKM}}`$ and the real diagonal mass matrix $`D_d`$. <sup>1</sup><sup>1</sup>1 In general, $`M_d=U^{}D_dU^{}`$ by using a general unitary matrix $`U=e^{i\alpha }e^{i\beta T_3}e^{i\gamma T_8}V_{\mathrm{CKM}}e^{i\beta ^{}T_3}e^{i\gamma ^{}T_8}`$. We omit the diagonal phases to keep the number of the free parameters in the model as small as possible. Considering the basis-independent quantities, $`\mathrm{tr}[M_e^{}M_e]`$, $`\mathrm{tr}[(M_e^{}M_e)^2]`$ and $`\mathrm{det}[M_e^{}M_e]`$, and eliminating $`|c_d|`$, we obtain two independent equations,
$`\left({\displaystyle \frac{\mathrm{tr}[\stackrel{~}{M_e}^{}\stackrel{~}{M_e}]}{m_e^2+m_\mu ^2+m_\tau ^2}}\right)^2`$ $`=`$ $`{\displaystyle \frac{\mathrm{tr}[(\stackrel{~}{M_e}^{}\stackrel{~}{M_e})^2]}{m_e^4+m_\mu ^4+m_\tau ^4}},`$ (97)
$`\left({\displaystyle \frac{\mathrm{tr}[\stackrel{~}{M_e}^{}\stackrel{~}{M_e}]}{m_e^2+m_\mu ^2+m_\tau ^2}}\right)^3`$ $`=`$ $`{\displaystyle \frac{\mathrm{det}[\stackrel{~}{M_e}^{}\stackrel{~}{M_e}]}{m_e^2m_\mu ^2m_\tau ^2}},`$ (98)
where $`\stackrel{~}{M_e}V_{\mathrm{CKM}}^{}D_dV_{\mathrm{CKM}}^{}+\kappa D_u`$. With input data of six quark masses, three angles and one CP-phase in the CKM matrix and three charged-lepton masses, we can solve the above equations and determine $`\kappa `$ and $`|c_d|`$, but one parameter, the phase of $`c_d`$, is left undetermined . The original basic mass matrices, $`M_{10}`$ and $`M_{126}`$, are described by
$`M_{10}`$ $`=`$ $`{\displaystyle \frac{3+|c_d|e^{i\sigma }}{4}}V_{\mathrm{CKM}}^{}D_dV_{\mathrm{CKM}}^{}+{\displaystyle \frac{|c_d|e^{i\sigma }\kappa }{4}}D_u,`$ (99)
$`M_{126}`$ $`=`$ $`{\displaystyle \frac{1|c_d|e^{i\sigma }}{4}}V_{\mathrm{CKM}}^{}D_dV_{\mathrm{CKM}}^{}{\displaystyle \frac{|c_d|e^{i\sigma }\kappa }{4}}D_u,`$ (100)
as the functions of $`\sigma `$, the phase of $`c_d`$, with the solutions, $`|c_d|`$ and $`\kappa `$, determined by the GUT relation.
Now let us solve the GUT relation and determine $`|c_d|`$ and $`\kappa `$. Since the GUT relation of Eq. (94) is valid only at the GUT scale, we first evolve by renormalization group equations (RGE’s) the data from the weak scale to the corresponding quantities at the GUT scale for a given $`\mathrm{tan}\beta `$ Then we and use them as input data at the GUT scale. Note that it is non-trivial to find a solution of the GUT relation, since the number of free parameters (fourteen) is almost the same as the number of inputs (thirteen). The solution of the GUT relation exists only if we take appropriate input parameters. By using the experimental data at the $`M_Z`$ scale , we get the following absolute values for charged fermion masses (in units of GeV) and the CKM matrix at the GUT scale, $`M_G`$, with $`\mathrm{tan}\beta =45`$, $`m_s=0.072`$ and $`\delta =1.518`$:
$`m_u=0.00103,m_c=0.299,m_t=133`$
$`m_d=0.00170,m_s=0.0263,m_b=1.55`$
$`m_e=0.000411,m_\mu =0.0868,m_\tau =1.69`$
and
$`V_{\mathrm{CKM}}=\left(\begin{array}{ccc}0.975& 0.222& 0.0001460.00279i\\ 0.2220.000121i& 0.974+0.000129i& 0.0320\\ 0.006970.00272i& 0.03120.000626i& 0.999\end{array}\right),`$ (104)
in the standard parameterization. The phases of the fermion masses are not determined by the diadonalization procedure. Here the masses are chosen be real. The signs of the fermion masses have been chosen to be equal, $``$ for $`m_u`$, $`m_c`$, $`m_d`$ and $`m_s`$, and $`+`$ for $`m_t`$ and $`m_b`$. So determined masses at the GUT scale are used as input parameters in order to solve Eqs. (97) and (98). As an example, here we show one of two solutions,
$`\kappa `$ $`=`$ $`0.01340.000791i,`$
$`|c_d|`$ $`=`$ $`6.39.`$ (105)
Once the parameters, $`|c_d|`$ and $`\kappa `$, are determined, we can describe all the fermion mass matrices as a functions of $`\sigma `$ from the mass matrix formulas of Eqs. (93), (99) and (100). Interestingly, in the minimal $`SO(10)`$ model even light Majorana neutrino mass matrix, $`M_\nu `$, can be determined as a function of the phase $`\sigma `$ and $`c_R`$ through the seesaw mechanism $`M_\nu =M_D^TM_R^1M_D`$.
Now we give an example of the neutrino-Dirac-Yukawa coupling matrix which can fit the GUT relation. In the basis where both the charged-lepton and right-handed Majorana neutrino mass matrix are diagonal with real and positive eigenvalues, the neutrino Dirac Yukawa coupling matrix at the GUT scale for fixed $`\sigma =3.223[\mathrm{rad}]`$ is found to be
$$Y_\nu =\left(\begin{array}{ccc}0.000310+0.00348i& 0.0008940.000249i& 0.0447+0.0531i\\ 0.005900.0103i& 0.01640.0427i& 0.308+0.116i\\ 0.00215+0.00126i& 0.05580.0559i& 0.381+0.604i\end{array}\right).$$
(106)
Using these Yukawa coupling constant matrices, we proceed with numerical calculations. In evaluating the LFV branching ratios, we first solve the RGE’s for the soft SUSY breaking parameters and the Yukawa couplings in the MSSM to determine the masses and mixings for the SUSY particles. Then we input these data into the formulae presented in the previous sections.
In the following, we list the input parameters what we used for the hadronization processes. For the pseudoscalar meson decay constants, we take as an input the following values ,
$`f_{\pi ^0}`$ $`=`$ $`0.119[\mathrm{GeV}],f_\eta =0.131[\mathrm{GeV}],f_\eta ^{}=0.118[\mathrm{GeV}],`$ (107)
and for the vector-meson decay constants we use the values, extracted from the vector meson decay, $`Ve^+e^{}`$,
$`\gamma _{\rho ^0}`$ $`=`$ $`2.518,\gamma _\varphi =2.993,\gamma _\omega =3.116.`$ (108)
The mixing angles between singlet and octet states for the vector mesons and for the pseudoscalar mesons that we use are ,
$`\theta _V`$ $`=`$ $`35^{},\theta _P=17.3^{}.`$ (109)
In order to investigate the dependence of SL LFV branching ratios on model parameters we plot their dependence on $`m_{\stackrel{~}{\tau }_R}`$: $`\tau e\pi ^0`$, $`\tau e\eta `$, $`\tau e\eta ^{}`$, $`\tau \mu \pi ^0`$, $`\tau \mu \eta `$, and $`\tau \mu \eta ^{}`$ graphs in Fig. 1, $`\tau e\rho ^0`$, $`\tau \mu \rho ^0`$ $`\tau e\varphi `$, $`\tau \mu \varphi `$ $`\tau e\omega `$, and $`\tau \mu \omega `$ graphs in Fig 2, $`\tau e\pi ^+\pi ^{}`$, $`\tau \mu \pi ^{}\pi ^+`$, $`\tau eK^0\overline{K}^0`$, and $`\tau \mu K^0\overline{K}^0`$, $`\tau eK^+K^{}`$, and $`\tau \mu K^+K^{}`$, graphs in Fig 3, and $`\tau e\gamma `$ and $`\tau \mu \gamma `$ graphs in Fig 4. In Fig 5. we plot the $`m_{\stackrel{~}{\tau }_R}`$ dependence of the quantities used to determine the alowed region of $`m_{\stackrel{~}{\tau }_R}`$ mass, the MSSM contribution to the anomalous magnetic moment $`\mathrm{\Delta }a_\mu `$ and $`\mu e\gamma `$ branching ratio, The parameters are chosen so as to satisfy the WMAP constraint on the cold dark matter (CDM) relic density ,
$$\mathrm{\Omega }_{\mathrm{CDM}}h^2=0.1126.$$
(110)
We can transmute this value into the approximate relation between $`m_0`$ and $`M_{1/2}`$, such as
$$m_0[\mathrm{GeV}]=\frac{9}{28}M_{1/2}[\mathrm{GeV}]+150[\mathrm{GeV}],$$
(111)
for $`\mathrm{tan}\beta =45`$, and $`A_0=0`$. The maximal values for the branching ratios of the $`\tau e/\mu \pi ^0`$ processes are found to be
$`\mathrm{BR}(\tau e\pi ^0)`$ $``$ $`1.7\times 10^{14},`$ (112)
$`\mathrm{BR}(\tau \mu \pi ^0)`$ $``$ $`2.4\times 10^{12},`$ (113)
and for the $`\tau e/\mu \eta `$ processes $`\mathrm{BR}(\tau e/\mu \eta )0.15\times \mathrm{BR}(\tau e/\mu \pi ^0)`$, with suitably chosen CMSSM parameters which can realize the neutralino dark matter scenario by the WMAP data. It can be realized with the following set of parameters: $`\mathrm{tan}\beta =45`$, $`\mu >0`$ and $`A_0=0`$, with $`M_{1/2}=600[\mathrm{GeV}]`$, $`m_0=343[\mathrm{GeV}]`$. This parameter set can also predict the muon $`g2`$ within the range of the recent result of Brookhaven E821 experiment and also provides the $`\tau \mu \gamma `$ and $`\mu e\gamma `$ branching ratios close to the current experimental bound . The ratio between two processes $`\tau e/\mu \pi ^0`$ and $`\tau e/\mu \eta `$ is a result of the dominance of the $`Z`$-boson-penguin amplitude in these processes, and reflects the difference in the form factors and the mixings between the singlet state and the octet state of the $`\eta `$ mesons,
$$\frac{\mathrm{BR}(\tau e/\mu \eta )}{\mathrm{BR}(\tau e/\mu \pi )}\left(\frac{f_\eta }{f_\pi }\right)^2\times \left(\frac{c_P}{\sqrt{3}}+\frac{s_P}{\sqrt{6}}\right)^20.15.$$
(114)
We can also see the correlation between the branching ratios for the processes $`\tau e/\mu \rho ^0`$ and $`\tau e/\mu \gamma `$ as $`\mathrm{BR}(\tau e/\mu \rho ^0)5.5\times 10^3\times \mathrm{BR}(\tau e/\mu \gamma )`$. It can be estimated on the basis of the photon-penguin-amplitude dominance in these amplitudes, giving
$$\frac{\mathrm{BR}(\tau e/\mu \rho ^0)}{\mathrm{BR}(\tau e/\mu \gamma )}\frac{1}{2}\left(\frac{e}{\gamma _{\rho ^0}}\right)^27\times 10^3.$$
(115)
When we impose the constraints from discrepancy of $`\tau `$ and $`e^+e^{}`$ data in the muon $`g2`$ measurements and from upper limits on the $`\mu e\gamma `$ branching ratio, we obtain that the model permits the of $`m_{\stackrel{~}{\tau }_R}`$ values satisfying,
$$m_{\stackrel{~}{\tau }_R}>204[\mathrm{GeV}].$$
(116)
The lower bound comes from the muon $`g2`$ constraint. The $`\mu e\gamma `$ gives also the lower bound but it is below the lower bound from the muon $`g2`$ constraint. The $`g2`$ curve has uprising behaviour above $`m_{\stackrel{~}{\tau }_R}=340[\mathrm{GeV}]`$, but at $`m_{\stackrel{~}{\tau }_R}=1000[\mathrm{GeV}]`$ it is almost independent on $`m_{\stackrel{~}{\tau }_R}`$ and has a value $`5.4\times 10^{10}`$, slightly below the present experimental $`g2`$ uncertainty. Therefore, one can expect that the improvement of the $`g2`$ measurements will give the upper limit on $`m_{\stackrel{~}{\tau }_R}`$, too. Using the lower bound on $`m_{\stackrel{~}{\tau }_R}`$ values one can find the theoretical upper bounds for all leptonic and SL LFV branching ratios. Leptonic and dominant SL LFV deacays are given in Table I.
## 5 Summary
The evidence for the neutrino masses and flavour mixings implies the non-conservation of the lepton-flavour symmetry. Thus, the LFV processes in the charged-lepton sector are expected. In supersymmetric model based on the minimal $`SO(10)`$ model, the values for the rates of the LFV processes are generally still several orders of magnitudes below the accessible current experimental bounds. In this paper, we have presented the detailed theoretical description for the SL LFV decays of the charged leptons with one or two pseudoscalar mesons or one vector meson in the final state. Also, some previous formulae have been corrected. The $`\gamma `$-penguin amplitude is corrected to assure the gauge invariance, the Z-penguin amplitude is corrected, new box contributions to the box amplitude has been found and previously neglected terms are given.
To evaluate the decay rates of the LFV processes within the MSSM, the parameters and the LFV interactions of the MSSM have to be specified. It has been shown that the minimal SUSY $`SO(10)`$ model can simultaneously accommodate all the observed quark-lepton mass matrix data involving the neutrino oscillation data with appropriately fixed free parameters. In this model, the neutrino-Dirac-Yukawa coupling matrix are completely determined, and its off-diagonal components are the primary source of the lepton-flavour violation in the basis where the charged-lepton and the right-handed neutrino mass matrices are real and diagonal. Using this Yukawa coupling matrix, we have calculated the rate of the LFV processes assuming the mSUGRA scenario. The analytical formulae of various SL LFV processes, $`\mathrm{}_i\mathrm{}_jP,\mathrm{}_i\mathrm{}_jV,\mathrm{}_i\mathrm{}_jPP`$ are given. Using these formulae, we have numerically evaluated $`\mathrm{}_i\mathrm{}_jP,\mathrm{}_i\mathrm{}_jV`$ and $`\mathrm{}_i\mathrm{}_jP_1P_2`$ branching ratios. Among these, the branching ratios of $`\tau \mu \gamma `$ and $`\mu e\gamma `$ may be interesting for the near future experiments. The typical CMSSM parameters used in calculations can realize the neutralino dark matter scenario by the WMAP data.
###### Acknowledgments.
A.I. would like to thank to I. Picek and S. Fajfer for the discussions on current-quark masses and evaluation of tensor-quark current. A part of the work was presented by A.I. at the workshop, “The 8th International Workshop on Tau-Lepton Physics (Tau04)”, held in Nara-ken New Public Hall, Japan. We are grateful to all organizers of this workshop and particularly to Prof. Ohshima for his kind hospitality extended to A.I and T.K. during their stay at Nagoya University. T.F. would like to thank S.T. Petcov for his hospitality at SISSA. This work of T.F. and T.K. was supported by the Grant in Aid for Scientific Research from the Ministry of Education, Science and Culture and the work of T.K. was supported by the Research Fellowship of the Japan Society for the Promotion of Science (# 16540269 and # 7336). The work of A.I is supported by the Ministry of Science and Technology of Republic of Croatia under contract 0119261.
## Appendix A Notation for the MSSM Lagrangian
Here we summarize our notation necessary for defining the masses of the sparticles in the MSSM Lagrangian.
$$v\sqrt{H_u^0^2+H_d^0^2}=174.1[\mathrm{GeV}],$$
(117)
and
$$\mathrm{tan}\beta \frac{H_u^0}{H_d^0}.$$
(118)
Then the charged fermion mass matrices are given by
$`M_u^{ij}`$ $`=`$ $`Y_u^{ij}v\mathrm{sin}\beta ,`$ (119)
$`M_d^{ij}`$ $`=`$ $`Y_d^{ij}v\mathrm{cos}\beta ,`$ (120)
$`M_e^{ij}`$ $`=`$ $`Y_e^{ij}v\mathrm{cos}\beta .`$ (121)
The mass matrix of the charginos is written as:
$`=\left(\begin{array}{cc}\overline{\stackrel{~}{W}_R^{}},\hfill & \overline{\stackrel{~}{H}_{uR}^{}}\hfill \end{array}\right)M_{\stackrel{~}{\chi }^\pm }\left(\begin{array}{c}\stackrel{~}{W}_L^{}\hfill \\ \stackrel{~}{H}_{dL}^{}\hfill \end{array}\right)+h.c.,`$ (125)
$`M_{\stackrel{~}{\chi }^\pm }=\left(\begin{array}{cc}M_2& \sqrt{2}M_W\mathrm{cos}\beta \\ \sqrt{2}M_W\mathrm{sin}\beta & \mu \end{array}\right).`$ (128)
The mass matrix of the neutralinos is written as:
$`={\displaystyle \frac{1}{2}}\left(\begin{array}{cccc}\stackrel{~}{B}_L,\hfill & \stackrel{~}{W}_L^3,\hfill & \stackrel{~}{H}_{dL}^0,\hfill & \stackrel{~}{H}_{uL}^0\hfill \end{array}\right)M_{\stackrel{~}{\chi }^0}\left(\begin{array}{c}\stackrel{~}{B}_L\hfill \\ \stackrel{~}{W}_L^3\hfill \\ \stackrel{~}{H}_{dL}^0\hfill \\ \stackrel{~}{H}_{uL}^0\hfill \end{array}\right)+h.c.,`$ (134)
$`M_{\stackrel{~}{\chi }^0}=\left(\begin{array}{cccc}M_1& 0& M_Z\mathrm{sin}\theta _W\mathrm{cos}\beta & M_Z\mathrm{sin}\theta _W\mathrm{sin}\beta \\ 0& M_2& M_Z\mathrm{cos}\theta _W\mathrm{cos}\beta & M_Z\mathrm{cos}\theta _W\mathrm{sin}\beta \\ M_Z\mathrm{sin}\theta _W\mathrm{cos}\beta & M_Z\mathrm{cos}\theta _W\mathrm{cos}\beta & 0& \mu \\ M_Z\mathrm{sin}\theta _W\mathrm{sin}\beta & M_Z\mathrm{cos}\theta _W\mathrm{sin}\beta & \mu & 0\end{array}\right).`$ (139)
The mass matrices of the squarks are written as follows:
$``$ $`=`$ $`(m_{\stackrel{~}{u}}^2)^{ij}\stackrel{~}{u}_i^{}\stackrel{~}{u}_j(m_{\stackrel{~}{d}}^2)^{ij}\stackrel{~}{d}_i^{}\stackrel{~}{d}_j,`$
$`m_{\stackrel{~}{u}}^2`$ $`=`$ $`\left(\begin{array}{cc}m_{\stackrel{~}{q}}^2+M_u^{}M_u& A_u^{}v\mathrm{sin}\beta M_u^{}\mu ^{}\mathrm{cot}\beta \\ A_uv\mathrm{sin}\beta M_u\mu \mathrm{cot}\beta & m_{\stackrel{~}{u}}^2+M_uM_u^{}M_Z^2\mathrm{cos}2\beta \mathrm{sin}^2\theta _W\end{array}\right)`$ (143)
$`+`$ $`\left(\begin{array}{cc}M_Z^2\mathrm{cos}2\beta (\frac{1}{2}\frac{2}{3}\mathrm{sin}^2\theta _W)\mathrm{𝟏}_{3\times 3}& 0\\ 0& \frac{2}{3}M_Z^2\mathrm{cos}2\beta \mathrm{sin}^2\theta _W\mathrm{𝟏}_{3\times 3}\end{array}\right),`$ (146)
$`m_{\stackrel{~}{d}}^2`$ $`=`$ $`\left(\begin{array}{cc}m_{\stackrel{~}{q}}^2+M_d^{}M_d& A_d^{}v\mathrm{cos}\beta M_d^{}\mu ^{}\mathrm{tan}\beta \\ A_dv\mathrm{cos}\beta M_d\mu \mathrm{tan}\beta & m_{\stackrel{~}{d}}^2+M_dM_d^{}M_Z^2\mathrm{cos}2\beta \mathrm{sin}^2\theta _W\end{array}\right)`$ (149)
$`+`$ $`\left(\begin{array}{cc}M_Z^2\mathrm{cos}2\beta (\frac{1}{2}+\frac{1}{3}\mathrm{sin}^2\theta _W)\mathrm{𝟏}_{3\times 3}& 0\\ 0& \frac{1}{3}M_Z^2\mathrm{cos}2\beta \mathrm{sin}^2\theta _W\mathrm{𝟏}_{3\times 3}\end{array}\right).`$ (152)
The mass matrices of the sleptons are written as follows:
$``$ $`=`$ $`(m_{\stackrel{~}{\nu }}^2)^{ij}\stackrel{~}{\nu }_i^{}\stackrel{~}{\nu }_j(m_{\stackrel{~}{e}}^2)^{ij}\stackrel{~}{e}_i^{}\stackrel{~}{e}_j,`$
$`m_{\stackrel{~}{\nu }}^2`$ $`=`$ $`m_\stackrel{~}{\mathrm{}}^2+{\displaystyle \frac{1}{2}}M_Z^2\mathrm{cos}2\beta \mathrm{𝟏}_{3\times 3},`$
$`m_{\stackrel{~}{e}}^2`$ $`=`$ $`\left(\begin{array}{cc}m_\stackrel{~}{\mathrm{}}^2+M_e^{}M_e& A_e^{}v\mathrm{cos}\beta M_e^{}\mu ^{}\mathrm{tan}\beta \\ A_ev\mathrm{cos}\beta M_e\mu \mathrm{tan}\beta & m_{\stackrel{~}{e}}^2+M_eM_e^{}M_Z^2\mathrm{cos}2\beta \mathrm{sin}^2\theta _W\end{array}\right)`$ (155)
$`+`$ $`\left(\begin{array}{cc}M_Z^2\mathrm{cos}2\beta (\frac{1}{2}+\mathrm{sin}^2\theta _W)\mathrm{𝟏}_{3\times 3}& 0\\ 0& M_Z^2\mathrm{cos}2\beta \mathrm{sin}^2\theta _W\mathrm{𝟏}_{3\times 3}\end{array}\right).`$ (158)
They are diagonalized with unitary matrices as follows:
$`O_RM_{\stackrel{~}{\chi }^\pm }O_L^{}`$ $`=`$ $`\mathrm{diag}(M_{\stackrel{~}{\chi }_1^{}},M_{\stackrel{~}{\chi }_2^{}}),`$
$`O_N^{}M_{\stackrel{~}{\chi }^0}O_N^{}`$ $`=`$ $`\mathrm{diag}(M_{\stackrel{~}{\chi }_1^0},M_{\stackrel{~}{\chi }_2^0},M_{\stackrel{~}{\chi }_3^0},M_{\stackrel{~}{\chi }_4^0}),`$
$`U_{\stackrel{~}{f}}m_{\stackrel{~}{f}}^2U_{\stackrel{~}{f}}^{}`$ $`=`$ $`\mathrm{diag}(m_{\stackrel{~}{f}_1}^2,\mathrm{}.,m_{\stackrel{~}{f}_6}^2),(f=u,d,e),`$
$`U_{\stackrel{~}{\nu }}m_{\stackrel{~}{\nu }}^2U_{\stackrel{~}{\nu }}^{}`$ $`=`$ $`\mathrm{diag}(m_{\stackrel{~}{\nu }_1}^2,m_{\stackrel{~}{\nu }_2}^2,m_{\stackrel{~}{\nu }_3}^2).`$ (159)
## Appendix B Lagrangian for fermion-sfermion-gaugino/Higgsino interactions in MSSM
The LFV interactions in the MSSM include the fermion-sfermion-gaugino/Higgsino vertices. These vertices, and corresponding coupling constants ($`C_{iAX}^{L,R(f)}`$, $`N_{iAX}^{L,R(f)}`$, $`f=\nu ,e,u,d`$) are defined by the following Lagrangian
$``$ $`=`$ $`\overline{u_i}\left[C_{iAX}^{L(u)}P_L+C_{iAX}^{R(u)}P_R\right]\stackrel{~}{\chi }_A^+\stackrel{~}{d}_X+\overline{d_i}\left[C_{iAX}^{L(d)}P_L+C_{iAX}^{R(d)}P_R\right]\stackrel{~}{\chi }_A^{}\stackrel{~}{u}_X`$ (160)
$`+`$ $`\overline{\nu _i}C_{iAX}^{R(\nu )}P_R\stackrel{~}{\chi }_A^+\stackrel{~}{e}_X+\overline{e_i}\left[C_{iAX}^{L(e)}P_L+C_{iAX}^{R(e)}P_R\right]\stackrel{~}{\chi }_A^{}\stackrel{~}{\nu }_X`$
$`+`$ $`\overline{u_i}\left[N_{iAX}^{L(u)}P_L+N_{iAX}^{R(u)}P_R\right]\stackrel{~}{\chi }_A^0\stackrel{~}{u}_X+\overline{d_i}\left[N_{iAX}^{L(d)}P_L+N_{iAX}^{R(d)}P_R\right]\stackrel{~}{\chi }_A^0\stackrel{~}{d}_X`$
$`+`$ $`\overline{\nu _i}N_{iAX}^{R(\nu )}P_R\stackrel{~}{\chi }_A^0\stackrel{~}{\nu }_X+\overline{e_i}\left[N_{iAX}^{L(e)}P_L+N_{iAX}^{R(e)}P_R\right]\stackrel{~}{\chi }_A^0\stackrel{~}{e}_X`$
$`+`$ $`h.c.`$
$``$ $`\overline{u}_iP_L\stackrel{~}{\chi }_A^+\stackrel{~}{d}_X\left[g\left\{{\displaystyle \frac{m_{u_i}}{\sqrt{2}M_\mathrm{W}\mathrm{sin}\beta }}(O_R)_{A2}(U_{\stackrel{~}{d}}^{})_{Xi}\right\}\right]`$
$`+`$ $`\overline{u}_iP_R\stackrel{~}{\chi }_A^+\stackrel{~}{d}_X\left[g\left\{(O_L)_{A1}(U_{\stackrel{~}{d}}^{})_{Xi}+{\displaystyle \frac{m_{d_k}}{\sqrt{2}M_\mathrm{W}\mathrm{cos}\beta }}(V_{\mathrm{CKM}})_{kj}(V_{\mathrm{CKM}}^{})_{ki}(O_L)_{A2}(U_{\stackrel{~}{d}}^{})_{X,j+3}\right\}\right]`$
$`+`$ $`\overline{d}_iP_L\stackrel{~}{\chi }_A^{}\stackrel{~}{u}_X\left[g\left\{{\displaystyle \frac{m_{d_i}}{\sqrt{2}M_\mathrm{W}\mathrm{cos}\beta }}(O_L^{})_{A2}(U_{\stackrel{~}{u}}^{})_{Xj}\right\}(V_{\mathrm{CKM}})_{ij}\right]`$
$`+`$ $`\overline{d}_iP_R\stackrel{~}{\chi }_A^{}\stackrel{~}{u}_X\left[g\left\{(O_R^{})_{A1}(U_{\stackrel{~}{u}}^{})_{Xj}{\displaystyle \frac{m_{u_j}}{\sqrt{2}M_\mathrm{W}\mathrm{sin}\beta }}(O_R^{})_{A2}(U_{\stackrel{~}{u}}^{})_{X,j+3}\right\}(V_{\mathrm{CKM}}^{})_{ij}\right]`$
$`+`$ $`\overline{\nu }_iP_R\stackrel{~}{\chi }_A^+\stackrel{~}{e}_X\left[g\left\{(O_L)_{A1}(U_{\stackrel{~}{e}}^{})_{Xj}+{\displaystyle \frac{m_{e_i}}{\sqrt{2}M_\mathrm{W}\mathrm{cos}\beta }}(O_L)_{A2}(U_{\stackrel{~}{e}}^{})_{X,j+3}\right\}(U_{\mathrm{MNS}}^{})_{ij}\right]`$
$`+`$ $`\overline{e}_iP_L\stackrel{~}{\chi }_A^{}\stackrel{~}{\nu }_X\left[g\left\{{\displaystyle \frac{m_{e_i}}{\sqrt{2}M_\mathrm{W}\mathrm{cos}\beta }}(O_L^{})_{A2}(U_{\stackrel{~}{\nu }}^{})_{Xi}\right\}\right]`$
$`+`$ $`\overline{e}_iP_R\stackrel{~}{\chi }_A^{}\stackrel{~}{\nu }_X\left[g\left\{(O_R^{})_{A1}(U_{\stackrel{~}{\nu }}^{})_{Xi}\right\}\right]`$
$`+`$ $`\overline{u}_iP_L\stackrel{~}{\chi }_A^0\stackrel{~}{u}_X\left[{\displaystyle \frac{g}{\sqrt{2}}}\left\{{\displaystyle \frac{m_{u_i}}{M_\mathrm{W}\mathrm{sin}\beta }}(O_N^{})_{A4}(U_{\stackrel{~}{u}}^{})_{Xi}{\displaystyle \frac{4}{3}}\mathrm{tan}\theta _W(O_N^{})_{A1}(U_{\stackrel{~}{u}}^{})_{X,i+3}\right\}\right]`$
$`+`$ $`\overline{u}_iP_R\stackrel{~}{\chi }_A^0\stackrel{~}{u}_X\left[{\displaystyle \frac{g}{\sqrt{2}}}\left\{{\displaystyle \frac{m_{u_i}}{M_\mathrm{W}\mathrm{sin}\beta }}(O_N)_{A4}(U_{\stackrel{~}{u}}^{})_{X,i+3}+\left[(O_N)_{A2}+{\displaystyle \frac{1}{3}}\mathrm{tan}\theta _W(O_N)_{A1}\right](U_{\stackrel{~}{u}}^{})_{Xi}\right\}\right]`$
$`+`$ $`\overline{d}_iP_L\stackrel{~}{\chi }_A^0\stackrel{~}{d}_X\left[{\displaystyle \frac{g}{\sqrt{2}}}\left\{{\displaystyle \frac{m_{d_i}}{M_\mathrm{W}\mathrm{cos}\beta }}(O_N^{})_{A3}(U_{\stackrel{~}{d}}^{})_{Xj}+{\displaystyle \frac{2}{3}}\mathrm{tan}\theta _W(O_N^{})_{A1}(U_{\stackrel{~}{d}}^{})_{X,j+3}\right\}(V_{\mathrm{CKM}})_{ij}\right]`$
$`+`$ $`\overline{d}_iP_R\stackrel{~}{\chi }_A^0\stackrel{~}{d}_X[{\displaystyle \frac{g}{\sqrt{2}}}\{{\displaystyle \frac{m_{d_i}}{M_\mathrm{W}\mathrm{cos}\beta }}(O_N)_{A3}(U_{\stackrel{~}{d}}^{})_{X,j+3}`$
$`+[(O_N)_{A2}+{\displaystyle \frac{1}{3}}\mathrm{tan}\theta _W(O_N)_{A1}](U_{\stackrel{~}{d}}^{})_{Xj}\}(V_{\mathrm{CKM}})_{ij}]`$
$`+`$ $`\overline{\nu }_iP_R\stackrel{~}{\chi }_A^0\stackrel{~}{\nu }_X\left[{\displaystyle \frac{g}{\sqrt{2}}}\left[(O_N)_{A2}\mathrm{tan}\theta _W(O_N)_{A1}\right](U_{\stackrel{~}{\nu }}^{})_{X,j}(U_{\mathrm{MNS}})_{ij}\right]`$
$`+`$ $`\overline{e}_iP_L\stackrel{~}{\chi }_A^0\stackrel{~}{e}_X\left[{\displaystyle \frac{g}{\sqrt{2}}}\left\{{\displaystyle \frac{m_{e_i}}{M_\mathrm{W}\mathrm{cos}\beta }}(O_N^{})_{A3}(U_{\stackrel{~}{e}}^{})_{Xi}+2\mathrm{tan}\theta _W(O_N^{})_{A1}(U_{\stackrel{~}{e}}^{})_{X,i+3}\right\}\right]`$
$`+`$ $`\overline{e}_iP_R\stackrel{~}{\chi }_A^0\stackrel{~}{e}_X\left[{\displaystyle \frac{g}{\sqrt{2}}}\left\{{\displaystyle \frac{m_{e_i}}{M_\mathrm{W}\mathrm{cos}\beta }}(O_N)_{A3}(U_{\stackrel{~}{e}}^{})_{X,i+3}+\left[(O_N)_{A2}\mathrm{tan}\theta _W(O_N)_{A1}\right](U_{\stackrel{~}{e}}^{})_{Xi}\right\}\right]`$
$`+`$ $`h.c.`$
## Appendix C Trilinear interactions of fermions or bosons with $`Z^0`$-boson or photon
The interactions of $`Z^0`$-boson or photon with any fermion or any boson follow from $`SU(2)_L\times U(1)_Y`$ gauge symmetry.
$`_f`$ $`=`$ $`g{\displaystyle \underset{f}{}}\overline{f}\left[\begin{array}{c}\text{/}\\ Z\hfill \end{array}\left(I_3^f{\displaystyle \frac{1}{c_W}}Q^f{\displaystyle \frac{s_W^2}{c_W}}\right)+\begin{array}{c}\text{/}\\ A\hfill \end{array}\left(s_WQ^f\right)\right]f,`$ (165)
$`_b`$ $`=`$ $`g{\displaystyle \underset{b}{}}\left(b^{}i\underset{\mu }{\overset{}{}}b\right)\left[Z^\mu \left(I_3^f{\displaystyle \frac{1}{c_W}}Q^f{\displaystyle \frac{s_W^2}{c_W}}\right)+A^\mu \left(s_WQ^f\right)\right].`$ (166)
The interaction Lagrangians are flavour-diagonal in the weak basis. In the mass basis the interactions with photon remain diagonal, because all mixed states $`A`$ must have the same charge, $`Q_A`$. On the other side the interactions with $`Z`$-boson, which depend on charge ($`Q_A`$) and third component of the weak isospin ($`I_{3A}`$) are not in general flavour-diagonal in the mass basis, because the mixed states may have different $`I_{3A}`$ values.
For LFV processes the interaction of photon and $`Z`$-boson with charginos, neutralinos and sfermion fields is needed. The Lagrangians for the corresponding weak-basis fields is easily written knowing the charges and $`I_3`$-s of these fields. After transformation from the weak-basis to the mass-basis the following interaction Lagrangians are obtained
$`_\chi ^{}`$ $`=`$ $`g\overline{\stackrel{~}{\chi }_A^{}}\begin{array}{c}\text{/}\\ Z\hfill \end{array}\left[(E_{AB}^{L(\chi ^{})}P_L+E_{AB}^{R(\chi ^{})}P_R)+{\displaystyle \frac{s_W^2}{c_W}}\delta _{BA}\right]\stackrel{~}{\chi }_B^{}+e\overline{\stackrel{~}{\chi }_A^{}}\begin{array}{c}\text{/}\\ A\hfill \end{array}\delta _{AB}\stackrel{~}{\chi }_B^{},`$ (171)
$`_{\chi ^0}`$ $`=`$ $`g\overline{\stackrel{~}{\chi }_A^0}\begin{array}{c}\text{/}\\ Z\hfill \end{array}\left[E_{AB}^{L(\chi ^0)}P_L+E_{AB}^{R(\chi ^0)}P_L\right]\stackrel{~}{\chi }_B^0,`$ (174)
$`_{\stackrel{~}{f}}`$ $`=`$ $`g\stackrel{~}{f}_X^{}i\underset{\mu }{\overset{}{}}\stackrel{~}{f}_YZ^\mu [D_{XY}^{\stackrel{~}{f}}]e\stackrel{~}{f}_X^{}i\underset{\mu }{\overset{}{}}\stackrel{~}{f}_YA^\mu [Q^{\stackrel{~}{f}}],`$ (175)
where
$`E_{AB}^{L(\chi ^{})}`$ $`=`$ $`{\displaystyle \frac{1}{c_W}}\left((O_L)_{A1}(O_L^{})_{B1}+{\displaystyle \frac{1}{2}}(O_L)_{A2}(O_L^{})_{B2}\right),`$ (176)
$`E_{AB}^{R(\chi ^{})}`$ $`=`$ $`{\displaystyle \frac{1}{c_W}}\left((O_R)_{A1}(O_R^{})_{B1}+{\displaystyle \frac{1}{2}}(O_R)_{A2}(O_R^{})_{B2}\right),`$ (177)
$`E_{AB}^{L(\chi ^0)}`$ $`=`$ $`{\displaystyle \frac{1}{c_W}}\left({\displaystyle \frac{1}{2}}(O_N)_{A3}(O_N^{})_{B3}{\displaystyle \frac{1}{2}}(O_N)_{A4}(O_N^{})_{B4}\right),`$ (178)
$`E_{AB}^{R(\chi ^0)}`$ $`=`$ $`E_{AB}^{L(\chi ^0)},`$ (179)
$`D_{XY}^{\stackrel{~}{f_L}}`$ $`=`$ $`I_3^{\stackrel{~}{f_L}}(U_{\stackrel{~}{f}})_{Xi}(U_{\stackrel{~}{f}}^{})_{Yi}Q_3^{\stackrel{~}{f_L}}{\displaystyle \frac{s_W^2}{c_W}}.`$ (180)
Here, in definition of $`D_{XY}^{\stackrel{~}{f}}`$ index $`i`$ is summed over generation indices $`1,2,3`$. $`Q^{\stackrel{~}{f_L}}`$ is a charge of the sfermion $`\stackrel{~}{f_L}`$, and $`I_3^{\stackrel{~}{f_L}}`$ is its third component of isospin in the weak basis. The charges and third components of the isospin of the weak-basis fields for charginos and neutralinos are explicitly written in the definitions of constants $`E_{AB}^{L,R(\chi ^{})}`$ and $`E_{AB}^{L,R(\chi ^0)}`$.
## Appendix D Loop functions
In this Appendix the loop functions appearing in the $`Z`$-boson amplitude and the box amplitudes are listed.
## Z-boson loop functions
The $`Z`$-boson amplitude comprises two-loop functions from the triangle-diagram part of the amplitude, $`F_1(a,b,c)`$ and $`F_2(a,b,c)`$, and two-loop functions from the self-energy part of the amplitude, $`f_1(a,b)`$ and $`f_2(a,b)`$ are
$`F_1(a,b,c)`$ $`=`$ $`{\displaystyle \frac{1}{bc}}\left[{\displaystyle \frac{a\mathrm{ln}ab\mathrm{ln}b}{ab}}{\displaystyle \frac{a\mathrm{ln}ac\mathrm{ln}c}{ac}}\right],`$ (181)
$`F_2(a,b,c)`$ $`=`$ $`{\displaystyle \frac{3}{8}}{\displaystyle \frac{1}{4}}{\displaystyle \frac{1}{bc}}\left[{\displaystyle \frac{a^2\mathrm{ln}ab^2\mathrm{ln}b}{ab}}{\displaystyle \frac{a^2\mathrm{ln}ac^2\mathrm{ln}c}{ac}}\right],`$ (182)
$`f_1(a,b)`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \frac{\mathrm{ln}a}{2}}+{\displaystyle \frac{a^2+b^2+2b^2(\mathrm{ln}a\mathrm{ln}b)}{4(ab)^2}},`$ (183)
$`f_2(a,b)`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \frac{\mathrm{ln}b}{2}}+{\displaystyle \frac{a^2b^2+2a^2(\mathrm{ln}b\mathrm{ln}a)}{4(ab)^2}}.`$ (184)
These loop functions are evaluated by neglecting the momenta of incoming and outgoing particles. The functions $`F_1`$ and $`F_2`$ are symmetric regarding replacement of their arguments ($`a`$, $`b`$, $`c`$). In the limit of two equal argument the functions $`F_1`$ and $`F_2`$ have the following form,
$`F_1(a,b,b)`$ $`=`$ $`{\displaystyle \frac{aba\mathrm{ln}a+a\mathrm{ln}b}{(ab)^2}},`$ (185)
$`F_2(a,b,b)`$ $`=`$ $`{\displaystyle \frac{1}{4}}{\displaystyle \frac{\mathrm{ln}b}{4}}+{\displaystyle \frac{a^2b^2+2a^2(\mathrm{ln}b\mathrm{ln}a)}{8(ab)^2}}.`$ (186)
The arguments of logarithms appearing in the $`F_1`$ and $`F_2`$ can be divided by any constant, what can be used to redefine these functions as functions of two variables, for instance $`b/a`$ and $`c/a`$. The unpleasant $`\mathrm{ln}b`$ term in $`f_2`$ function can be replaced with $`\mathrm{ln}(b/a)`$ because of unitarity cancellations in the sum $`N_{jBY}^{R(e)}N_{iAX}^{R(e)}`$, and therefore $`f_1`$ and $`f_2`$ can be expressed in terms of one variable ($`b/a`$) only.
## Box-loop functions
Box amplitude contains two-loop functions, $`d_0`$ and $`d_2`$.
$`d_0(x,y,z,w)`$ $`=`$ $`{\displaystyle \frac{x\mathrm{ln}x}{(yx)(zx)(wx)}}+{\displaystyle \frac{y\mathrm{ln}y}{(xy)(zy)(wy)}}`$ (187)
$`+`$ $`{\displaystyle \frac{z\mathrm{ln}z}{(xz)(yz)(wz)}}+{\displaystyle \frac{w\mathrm{ln}w}{(xw)(yw)(zw)}},`$
$`d_2(x,y,z,w)`$ $`=`$ $`{\displaystyle \frac{1}{4}}\{{\displaystyle \frac{x^2\mathrm{ln}x}{(yx)(zx)(wx)}}+{\displaystyle \frac{y^2\mathrm{ln}y}{(xy)(zy)(wy)}}`$ (188)
$`+`$ $`{\displaystyle \frac{z^2\mathrm{ln}z}{(xz)(yz)(wz)}}+{\displaystyle \frac{w^2\mathrm{ln}w}{(xw)(yw)(zw)}}\}.`$
As mentioned before, they are also by evaluated neglecting the momenta of incoming and outgoing particles.
## Appendix E Meson states and quark currents
Meson states are assumed to contain valence quarks only. The quark-antiquark ($`q_aq_b^c`$) content of the pseudoscalar-meson states is given in the table below. The quark-antiquark content of the vector-meson states is obtained replacing the fields $`K^+`$, $`K^0`$, $`\pi ^+`$, $`\pi ^0`$, $`\pi ^{}`$, $`\overline{K}^0`$ $`K^{}`$, $`\eta _8`$, $`\eta _1`$, $`\eta `$ and $`\eta ^{}`$ by fields $`K^+`$, $`K^0`$, $`\rho ^+`$, $`\rho ^0`$, $`\rho ^{}`$, $`\overline{K^0}`$ $`\overline{K^{}}`$, $`\varphi _8`$, $`\varphi _1`$, $`\varphi `$ and $`\omega `$, and the angle $`\theta _P`$ by the angle $`\theta _V`$.
From the quark content of the meson fields one can find the meson content of e.g. axial vector (A) and vector (V) quark currents, (factor of proportionality, Lorentz indices and spinor structures are neglected),
$`(\overline{u}u)_A`$ $``$ $`\left({\displaystyle \frac{c_P}{\sqrt{6}}}{\displaystyle \frac{s_P}{\sqrt{3}}}\right)\eta ^{}+\left({\displaystyle \frac{s_P}{\sqrt{6}}}+{\displaystyle \frac{c_P}{\sqrt{3}}}\right)\eta ^{{}_{}{}^{}}+{\displaystyle \frac{1}{\sqrt{2}}}\pi ^0k_{\overline{u}u}^\eta \eta ^{}+k_{\overline{u}u}^\eta ^{}\eta ^{}+k_{\overline{u}u}^{\pi ^0}\pi ^0`$
$`(\overline{d}d)_A`$ $``$ $`\left({\displaystyle \frac{c_P}{\sqrt{6}}}{\displaystyle \frac{s_P}{\sqrt{3}}}\right)\eta ^{}+\left({\displaystyle \frac{s_P}{\sqrt{6}}}+{\displaystyle \frac{c_P}{\sqrt{3}}}\right)\eta ^{}{\displaystyle \frac{1}{\sqrt{2}}}\rho ^0k_{\overline{d}d}^\eta \eta ^{}+k_{\overline{d}d}^\eta ^{}\eta ^{}+k_{\overline{d}d}^{\pi ^0}\pi ^0`$
$`(\overline{s}s)_A`$ $``$ $`\left({\displaystyle \frac{2c_P}{\sqrt{6}}}{\displaystyle \frac{s_P}{\sqrt{3}}}\right)\eta ^{}+\left({\displaystyle \frac{2s_P}{\sqrt{6}}}+{\displaystyle \frac{c_P}{\sqrt{3}}}\right)\eta ^{}k_{\overline{s}s}^\eta \eta ^{}+k_{\overline{s}s}^\eta ^{}\eta ^{}`$ (189)
$`(\overline{u}u)_V`$ $``$ $`\left({\displaystyle \frac{c_V}{\sqrt{6}}}{\displaystyle \frac{s_V}{\sqrt{3}}}\right)\varphi ^{}+\left({\displaystyle \frac{s_V}{\sqrt{6}}}+{\displaystyle \frac{c_V}{\sqrt{3}}}\right)\omega ^{}+{\displaystyle \frac{1}{\sqrt{2}}}\rho ^0k_{\overline{u}u}^\varphi \varphi ^{}+k_{\overline{u}u}^\omega \omega ^{}+k_{\overline{u}u}^{\rho ^0}\rho ^0`$
$`(\overline{d}d)_V`$ $``$ $`\left({\displaystyle \frac{c_V}{\sqrt{6}}}{\displaystyle \frac{s_V}{\sqrt{3}}}\right)\varphi ^{}+\left({\displaystyle \frac{s_V}{\sqrt{6}}}+{\displaystyle \frac{c_V}{\sqrt{3}}}\right)\omega ^{}{\displaystyle \frac{1}{\sqrt{2}}}\rho ^0k_{\overline{d}d}^\varphi \varphi ^{}+k_{\overline{d}d}^\omega \omega ^{}+k_{\overline{d}d}^{\rho ^0}\rho ^0`$
$`(\overline{s}s)_V`$ $``$ $`\left({\displaystyle \frac{2c_V}{\sqrt{6}}}{\displaystyle \frac{s_V}{\sqrt{3}}}\right)\varphi ^{}+\left({\displaystyle \frac{2s_V}{\sqrt{6}}}+{\displaystyle \frac{c_V}{\sqrt{3}}}\right)\omega ^{}k_{\overline{s}s}^\varphi \varphi ^{}+k_{\overline{s}s}^\omega \omega ^{}`$ (190)
Here $`s_V=\mathrm{sin}\theta _V`$, and $`c_V=\mathrm{cos}\theta _V`$, and the numerical factors are normalizatons of mesons expressed in terms quark fields. The combinations of constants $`k_{\overline{q}_aq_b}^{P,V}`$ are contained in expressions for all quark currents (from the scalar-quark current to the tensor-quark current) and we introduce them to abbreviate the expressions for box form factors.
The photon-penguin and the Z-boson–penguin amplitudes contain combinations of constants $`k_{\overline{q}_aq_b}^{P,V}`$ given in Table 3. For instance, $`k_\gamma ^{\rho ^0}=Q_uk_{\overline{u}u}^{\rho _0}+Q_dk_{\overline{d}d}^{\rho _0}+Q_sk_{\overline{s}s}^{\rho _0}`$.
| $`k`$ | $`\rho ^0`$ | $`\varphi `$ | $`\omega `$ |
| --- | --- | --- | --- |
| $`k_\gamma ^V`$ | $`\frac{1}{\sqrt{2}}`$ | $`\frac{1}{\sqrt{6}}c_V`$ | $`\frac{1}{\sqrt{6}}s_V`$ |
| $`k_Z^V`$ | $`\frac{1}{\sqrt{2}}c_{2W}`$ | $`\frac{c_V}{\sqrt{6}}c_{2W}+\frac{s_V}{2\sqrt{3}}`$ | $`\frac{s_V}{\sqrt{6}}c_{2W}\frac{c_V}{2\sqrt{3}}`$ |
| $`k`$ | $`\pi ^0`$ | $`\eta `$ | $`\eta ^{}`$ |
| $`k_Z^P`$ | $`\frac{1}{\sqrt{2}}`$ | $`\frac{c_P}{\sqrt{6}}+\frac{s_P}{2\sqrt{3}}`$ | $`\frac{s_P}{\sqrt{6}}\frac{c_P}{2\sqrt{3}}`$ |
Table 3. Combinations of $`k_{\overline{q}_aq_b}^{P,V}`$ constants appearing in photon-penguin and the Z-boson–penguin $`\mathrm{}\mathrm{}^{}P(V)`$ amplitudes.
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# 1 Introduction
## 1 Introduction
It is well known since the ancient times that two concepts of space as absolute or relative are possible; for example Democritos believed in the existence of both atoms and empty space, while Aristotle with his opinion that space only epitomizes the place of the objects comes close to a relational concept of space. Consequently, although Mach was not the first who formulated mechanics in terms of purely relational quantities his profound critique of foundations of Newtonian mechanics played a key role in Einsteins development of general relativity. It has been difficult to identify spacetime geometries uniquely from statement of Machs principle since it is difficult to translate it into some precise mathematical language. Moreover, each statement of Machs principle should be accompanied by a declaration of the theoretical framework in which it is applied.
Machs principle has also guided us to developments of multidimensional gravitation theories. The idea of the ordinary spacetime is viewed as a hypersurface embedded in a higher dimensional space have been considered in several different contexts, such as strings , D-branes , Randall-Sundrum models , and non-compactified versions of Kaluza-Klein theories . A Campbell-Magaards theorem , and its extension , has acquired fundamental relevance for granting the mathematical consistency of multidimensional embedding theories and also has been applied to investigate how lower-dimensional theories could be related to (3+1)-dimensional Einstein gravity.
The dimensional reduction in order to explain the property of nature is now widely used in theoretical physics , , . From the pioneering paper by Jordan physicists have explored the odd, however apparently fruitful thought, that Kaluza-Klein theories with special concern for the idea that the new 5-dimensional metric component, a spacetime scalar, might play the role of a varying gravitational constant. However, Jordan and his colleagues took the next step of separating the scalar field from the original 5-dimensional metric context unified gravitational-electromagnetic framework.
On the other hand, many physicists have expressed reservations about the 4-dimensional formalism. Though Lorentz recognized Einstein’s hypothesis of geometrization of gravity as one of the basic principles, his enthusiasm for it gradually waned soon afterwards. In 1910 Lorentz said ”Die Vorstellung (die auch Redner nur ungen aufgeben würde), dass Raum und Zeit etwas vollig Verschiedenes seien und dass es eine wahre Zeit gebe (die Gleichzeitigkeit würde denn unabhängig vom Orte bestehen).
An investigation into how lower and higher-dimensional theories of gravity are related to 3-dimensional Universe is therefore well motivated from a physical point of view. A potential bridge between gravitational theories of different dimensionality may be found at least two distinct ways by employing the embedding or foliation approaches . In the embedding approach, the geometry is determined in terms of quantities which are defined exclusively in lower-dimensional. In the foliation approach, the lower-dimensional space-time geometry is determined by inducing the metric on the leaves, so that in this case the metric tensor depends on the extra coordinate. According to this world-view, our perception of normal 3-dimensional space arises because we leave on a domain wall (time-brane) in a bulk space. This kind of infinitely thin branes models can, however, be only treated as approximation. One can describe a thick brane in framework of 4-dimensional general relativity as domain wall separating two different states past and future and to an observer in ordinary space feels the thickness of brane as present.
This paper details the new approach to the problem of separate the temporal scalar field from the 4-dimensional metric context of general relativity which where outlined in previous letter . Section 2 develops the formalism related to the correspondence between the ordinary 4-dimensional gravity with 3-dimensional theories with temporal scalar field. In Section 3 the exact solutions for the static sphere has been obtained. Finally, Section 4 briefly discussed some aspects of these theories. We choose units such that G<sub>N</sub> = c = 1, and let Latin indices run 1-4 and Greek indices run 1-3.
## 2 Dimensional reduction approach
Embedding theorems of differential geometry provide a natural framework for relating higher and lower- dimensional theories of gravity. Procedure introduced by Campbell- Magaards theorems to embed a Riemannian manifold into Einstein space has been applied to investigate how low- dimensional theories could be related to (3+1) dimensional Einstein gravity. Such an embedding has a namber of applications. For example, employing classical dimensional reduction techniques , a 3-dimensional theory containing temporal scalar field , can be obtained. Such a possibility naturally appears within the frame of multidimensional models type of Kaluza-Klein theory.
The claim that, in the spirit of Newtonian theory, the space and time are treated in completely different ways can be obtained by separate the temporal scalar field from the original 4-dimensional metric context of general relativity.
A second tenet of this theory is that all classical macroscopic physical quantities, such as matter density and pressure, could be given a geometrical interpretation. In this way, it is proposed that the classical energy-momentum tensor, which enters the right-hand side of the 3-dimensional Einstein equations, can be generated by pure geometrical means. In other words, it is claimed that geometrical curvature induces matter in three dimensions, and to an observer in the ordinary space the extra dimensions (time) appear as the matter source for gravity. However, induced stress-energy tensor does not in general uniquely determine the matter content and the interpretation can lead to quite different kinematic quantities.
One can use the local an isotropic embedding of the 3-dinensional, Riemannian manifold with line element
$$ds^2=g_{\alpha \beta }(x^\mu )dx^\alpha dx^\beta ,$$
(1)
in (3+1)-dimensional manifold defined by the metric
$$ds^2=h_{\alpha \beta }dx^\alpha dx^\beta +\varphi ^2dt^2,$$
(2)
where h<sub>αβ</sub> = h<sub>αβ</sub>(x<sup>μ</sup>,t) and $`\varphi `$=$`\varphi `$(x<sup>μ</sup>,t) are functions of the (3+1) variables {x<sup>μ</sup>,t} . The claim that any energy-momentum can be generated by an embedding mechanism may be translated in geometrical language as saying that Riemannian 3-dimensional manifold is embeddable into a 4-dimensional Ricci manifold. In the coordinates (2), the field equations take the form
$$\begin{array}{c}\hfill {}_{}{}^{(4)}R_{\alpha \beta }^{}={}_{}{}^{(3)}R_{\alpha \beta }^{}\frac{\varphi _{;\alpha ;\beta }}{\varphi }\\ \hfill \frac{1}{2\varphi ^2}(\frac{\stackrel{}{\varphi }}{\varphi }\underset{\alpha \beta }{\overset{}{h}}\underset{\alpha \beta }{\overset{}{h}}\frac{1}{2}h^{\gamma \delta }\underset{\gamma \delta }{\overset{}{h}}\underset{\alpha \beta }{\overset{}{h}}+h^{\gamma \delta }\underset{\alpha \gamma }{\overset{}{h}}\underset{\beta \delta }{\overset{}{h}}),\end{array}$$
(3)
$${}_{}{}^{(4)}R_{\alpha y}^{}=\frac{\varphi }{2}_\beta (\frac{1}{\varphi }h^{\beta \gamma }\underset{\gamma \alpha }{\overset{}{h}}\frac{1}{\varphi }\delta _\beta ^\alpha h^{\gamma \delta }\underset{\gamma \delta }{\overset{}{h}}),$$
(4)
$${}_{}{}^{(4)}R_{yy}^{}=\varphi \varphi _{;\alpha ;\alpha }\frac{1}{2}h^{\gamma \delta }\underset{\gamma \delta }{\overset{}{h}}\frac{1}{2}\stackrel{\gamma \delta }{\stackrel{}{h}}\underset{\gamma \delta }{\overset{}{h}}+\frac{1}{2}h^{\gamma \delta }\underset{\gamma \delta }{\overset{}{h}}\frac{\stackrel{}{\varphi }}{\varphi }\frac{1}{4}h^{\gamma \beta }h^{\delta \alpha }\underset{\delta \beta }{\overset{}{h}}\underset{\gamma \alpha }{\overset{}{h}}.$$
(5)
where overdots denote normal time derivatives. As in classical general relativity the metric is determined by Einstein’s equations.
$${}_{}{}^{(3)}R_{\alpha \beta }^{}\frac{1}{2}g_{\alpha \beta }{}_{}{}^{(3)}R=\kappa {}_{}{}^{(3)}T_{\alpha \beta }^{},$$
(6)
where $`\kappa `$ is the Einstein constant and $`{}_{}{}^{(3)}T_{\alpha \beta }^{}`$ energy-momentum tensor, which is a function of all matter fields and metric. These give the components of the induced energy-momentum tensor since Einsteins equations (6) hold.
$$\begin{array}{c}\hfill {}_{}{}^{(3)}T_{\alpha \beta }^{}=\frac{\varphi _{;\alpha ;\beta }}{\varphi }\frac{1}{2\varphi ^2}[\frac{\stackrel{}{\varphi }}{\varphi }\underset{\alpha \beta }{\overset{}{h}}\underset{\alpha \beta }{\overset{}{h}}+h^{\lambda \mu }\underset{\alpha \lambda }{\overset{}{h}}\underset{\beta \mu }{\overset{}{h}}\frac{1}{2}h^{\mu \nu }\underset{\mu \nu }{\overset{}{h}}\underset{\alpha \beta }{\overset{}{h}}]+\\ \hfill +\frac{h_{\alpha \beta }}{8\varphi ^2}[\stackrel{\mu \nu }{\stackrel{}{h}}\underset{\mu \nu }{\overset{}{h}}+\left(h^{\mu \nu }\underset{\mu \nu }{\overset{}{h}}\right)^2].\end{array}$$
(7)
It can be shown that the field equations of this theory can be obtained from the reduced action that are formally similar to the Jordan, Brans-Dicke action with $`\omega `$ = 0 . The Jordan, Brans-Dicke theory of gravity , represents a natural extension of general relativity, where a nonminimally-coupled scalar field parametrizes the space-time dependence of Newtons constant. It is expected on the general ground that dimensional reduction of multi-dimensional theory yields a (generalized) Jordan, Brans-Dicke theory. It is well known that the Jordan, Brans-Dicke theory accommodates both Machs principle and Diracs large number hypothesis , . Note that this conclusion depends crucially on the particular formulation of Machs principle used . Newton’s gravitational constant G<sub>N</sub> replaced by dynamical scalar field acts as the source of the (local) gravitational coupling with G<sub>N</sub> $``$ $`\varphi ^1`$ and consequently is determined by the total matter in the universe through auxiliary scalar field equations.
By the way the presented 3-dimensional effective theory is derived by the integration of the action with respect to the extra dimension as it customary in Kaluza-Klein compactification. In contrast with standard Jordan, Brans-Dicke theory this temporal scalar field don’t play the role of a varying gravitational ”constant”. Take it into account one can show that such 3-dimensional theory introduce a temporal scalar field, which will (locally and approximately) play the role of Newtonian gravitational potential. The next step is to suppose that it is possible to find a special coordinate system, such that we can make the following split of 4 -dimensional metric
$$g_{ab}=\left(\begin{array}{cc}g_{\alpha \beta }+\varphi ^2A_\alpha A_\beta & \varphi ^2A_\alpha \\ \varphi ^2A_\beta & \varphi ^2\end{array}\right)$$
(8)
where $`\varphi `$ may be regarded as a temporal scalar field, and A<sub>α</sub> is 3-dimensional vector, whose role is to be determined later.
## 3 Three dimensional gravity theories with temporal scalar field.
Alternative ways to construct the temporal scalar field theory is to consider the more general 3-dimensional theory governed by the action where gravity is coupled to temporary scalar field. From the formal viewpoint, scalar-tensor models are purely metric theories including a nonminimal coupling between a temporal scalar field and the curvature scalar R of the metric g. The most general action can be submitted by allowing $`\omega `$($`\varphi `$) to depend on the scalar field $`\varphi `$ and introducing a cosmological function $`\lambda `$($`\varphi `$). The action for these theories in the Jordan frame is:
$`S={\displaystyle \left(f\left(\varphi \right){}_{}{}^{(3)}R+\frac{\omega \left(\varphi \right)}{\varphi }g^{\mu \nu }\varphi _{,\mu }\varphi _{,\nu }2\varphi \lambda \left(\varphi \right)+L_{matt}\right)\sqrt{\left|g\right|}d^3x𝑑t},`$ (9)
The dynamics of the scalar field $`\varphi `$ depends on the functions $`f\left(\varphi \right),\omega `$($`\varphi `$) and $`\lambda `$($`\varphi `$). The matter part of the action $`L_{matt}`$ depends on the material fields and the metric but does not involve the scalar field $`\varphi .`$ This ensures a satisfaction of the ”weak equivalence principle”, that is, the statement that the paths of test particles in gravitational field are independent of their masses. We may demand that, in a spacetime consistent with Mach’s principle, the metric tensor should be completely determined by the mass distribution in the universe. Different choices of function $`\omega `$($`\varphi `$) and $`\lambda `$($`\varphi `$) give different temporal scalar tensor theories. A simple example elaborates on further possibilities. One can neglect the cosmological function assuming $`\lambda `$($`\omega `$) = 0 and use $`f\left(\varphi \right)=\varphi ,`$ $`\omega `$($`\varphi `$) = const. Then the variation of (9) with respect to g<sup>μν</sup> and $`\varphi `$ gives, respectively, the field equations
$$\begin{array}{c}\hfill {}_{}{}^{(3)}R_{\mu \nu }^{}\frac{1}{2}g_{\mu \nu }{}_{}{}^{(3)}R=\frac{8\pi }{\varphi }{}_{}{}^{\left(3\right)}T_{\mu \nu }^{}\text{ }\\ \hfill \frac{\omega }{\varphi ^2}\left(\varphi _{;\mu }\varphi _{;\nu }\frac{1}{2}g_{\mu \nu }\varphi _{;\rho }\varphi ^{;\rho }\right)\frac{1}{\varphi }\left(\varphi _{;\mu ;\nu }g_{\mu \nu }\varphi _{;\rho }^{;\rho }\right),\end{array}$$
(10)
$`\varphi _{;\mu }^{;\mu }={\displaystyle \frac{8\pi }{3+2\omega }}T,`$ (11)
where $`{}_{}{}^{(3)}T_{\mu \nu }^{}`$ is the energy-momentum tensor of matter fields.
The system of equations (10), (11) have the vacuum solution g<sup>μν</sup> = $`\eta ^{\mu \nu }`$ and $`\varphi `$ = const, where $`\eta ^{\mu \nu }`$ is the metric tensor of a flat space. Obviously, this solution identical to Minkowski’s space of general relativity and this solution fixes the arbitrariness in the coordinate system. Notice that this explanation viable in a more general theories when $`f\left(\varphi \right)\varphi `$ and $`\omega `$($`\varphi `$) $``$ const. It is necessary to remark that, according to Mach, the inertial forces observed locally in accelerated laboratory may be interpreted as gravitational effects having their origin in distant matter accelerated relative to laboratory. Consequently, there must be energy density of matter everywhere in machian universe. In completely empty space there can not be no gravitational field at all; in this case it would not be possible neither distribution of light, nor existence of scales and hours . In this case we must choose the solution g<sup>μν</sup> = $`0`$ and $`\varphi `$ = 0. On the other hand if we choose the case g<sup>μν</sup> = $`\eta ^{\mu \nu }`$ and $`\varphi `$ = 0 then we came to ”Newtonian” world.
The most appropriate coordinates in which to study the static spherically symmetric field are isotropic. Hence, the metric will be assumed to have the form
$`ds^2=e^\beta \left(dr^2+r^2\left(d\theta ^2+\mathrm{sin}^2\theta d\phi ^2\right)\right),`$ (12)
where, $`\beta `$ and $`\varphi `$ will be assumed to be function of r only. Expressing the line element in isotropic form equations (10), (11) gives:
$`\begin{array}{c}\frac{\beta ^{}\varphi ^{}}{\varphi }\frac{2\beta ^{}}{r}2\omega \left(\frac{\varphi ^{}}{\varphi }\right)^22\beta ^{\prime \prime }\frac{2\varphi ^{\prime \prime }}{\varphi }=0\\ \\ \frac{3\beta ^{}}{r}+\frac{\beta ^{}^2}{2}+\beta ^{\prime \prime }+\left(\frac{2}{r}+\beta ^{}\right)\frac{\varphi ^{}}{\varphi }=0\\ \\ \frac{4}{r}+\beta ^{}+\frac{2\varphi ^{\prime \prime }}{\varphi ^{}}=0\end{array}`$ (18)
The solution of equations (18) for the static sphericaly symmetric configuration can be easily obtained. This gives
$`\begin{array}{c}\beta =\mathrm{ln}\left(\frac{r^4\left(\frac{r+\mu }{r\mu }\right)^{\sqrt{\frac{8}{2+\omega }}}}{\left(2+\omega \right)\left(r^2\mu ^2\right)^2}\right)+c_1\\ \\ \varphi =c_2\left(\frac{r+\mu }{r\mu }\right)^{\sqrt{\frac{2}{2+\omega }}}\end{array},`$ (22)
where $`\mu ,`$c<sub>1</sub> and c<sub>2</sub> arbitrary constant.
A point to be noted is that under choices the case $`\omega `$=0 the solution (22) becomes identical to Schwarzschild solution of general relativity. It is well known that the exterior Schwarzschild solution of general relativity does not incorporate Mach’s principle. Moreover, the anti-Machian character of the general relativity field equations also shows up in solutions which allow for a curved spacetime even in the absence of any matter . The point is that the field equations admit curved vacuum solutions and also solutions that are asymptotically flat. However, a single mass point have inertia but in a consistent relativity theory there cannot be inertia relative to ”space” but in a Mach’s philosophy the inertia of a body gets determined by the presence of all other bodies in the universe.
However, that despite the mathematical equivalence of the Schwarzschild and (22) with $`\omega `$=0 solutions, the fact that they come from different physical theories makes them conceptually distinct. It is surprising that in 3- dimensional temporal scalar tensor theories there is not single matter particle solution. Thus in this sense the fourth dimension (time) generates effective matter sources and in turn act to curve 3-dimensional space. In other words, existence of even a sole particle in the space fills the entire universe with matter. Consequently, according to the Machian hypothesis there are the matter ”everywhere” in this universe.
## 4 Conclusion
There has been a lot of activity on describes our world as a 4-dimensional surface (brane) world embedded in higher dimensional space-time (bulk) , , , , . This concept leads to a great variety of specific models both in cosmological context and in the description of local self-gravitating objects. Modifying the usual picture, where the our Universe represent as 4-dimensional differentiable manifolds, these our developments are based on the idea that matter fields could be confined to 3-dimensional world, while time could live in a higher dimensional space. So, like many authors, we try to describe N-dimensional gravitational models with extra fields were obtained from some multidimensional model by dimensional reduction based on decomposition of the manifold . Most of theories built on this idea represent the 4-dimensional space-time and its associated extra space by differentiable manifolds hence such a gravity theory is equivalent to the Jordan, Brans-Dicke theory with vanishing parameter $`\omega `$ and a potential term , . On the other hand after the conceptual foundations of non-Euclidian geometry had been laid by Lobachevsky, Bolyai, Gauss and the examinations on mechanics in non-Euclidean spaces had been clarified (at the end of XIX centuries , ), it was natural to consider how Newton’s law should fit into the new framework. In this scheme, in the spirit of Newtonian theory, space and time are treated in completely different ways.
We have presented here a new point of view to the General relativity theory. In present article the possibility of description of 3-dimensional gravity theories in terms of dimensional reduction is analyzed. By the way the idea that the matter in three dimensions can be explained from a 4-dimensional Riemannian manifold is a consequence of this point of view. This theory represents some modification of general relativity can be interpreted as a dynamical theory of evolution of 3- dimensional Riemannian geometry. A related question in whether plausible brane physics can generically be recovered and whether there are testable correction to Newtonian gravity. It would be interesting to consider these question further.
If the brane world is real, one may find some evidences of higher-dimensions. The present can be viewed as confined to a time-brane, a 3-dimensional hyper-surface embedded in space-time with time extra dimensions. It is necessary to remark that, the man feels the present (the thickness of brane) as 0.5 seconds. The field equations of this theory formally have a structure similar to the Jordan, Brans-Dicke theory in three dimension with the free parameter $`\omega `$ = 0. It is expected on the general ground that dimensional reduction of a multi-dimensional theory in the 3-dimensional scheme yields a temporal scalar field theory. Consequently, these ideas may be continued by suggested considering the more general case leading to the new theory with an additional temporal scalar field. These models based on temporal scalars have the theoretically appealing property that the same scalar field that participates in the gravity sector is enhanced by potential. Introducing the temporal scalar field a different actions (9) with different choices of functions $`\lambda `$($`\omega `$), $`f\left(\varphi \right)`$and $`\omega `$($`\varphi `$) has been proposed.
Finally an important question with regards to three dimensional gravity theories with temporal scalar field theories is how to choose reasonable forms of the field equations and the energy-momentum tensor. An examination based upon case of static spherically symmetric configuration with restriction to the functions $`\lambda `$($`\omega `$) = 0, $`f\left(\varphi \right)=\varphi `$and $`\omega `$($`\varphi `$) = const., show that there is solution of fields equations (18) identical to the Schwarzschild solution in general relativity. However, the geometrical curvature induces matter in three dimensions, and the extra dimensions (time) appear as the additional effective matter source for gravity.
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# The Complex Wind Torus and Jets of PSR B1706-44
## 1 Introduction
PSR B1706$``$44, discovered by Johnston et al. (1992) is among the most interesting pulsars for study at high energies. It is one of a handful of pulsars detected by EGRET in GeV $`\gamma `$-rays. It is quite similar to the Vela pulsar with a characteristic age $`\tau _c=P/(2\dot{P})=1.7\times 10^4`$ yr and a spindown luminosity of $`\dot{E}4\times 10^{36}`$ erg/s, but is $`10\times `$ more distant at $`d=3d_3`$ kpc. Early Chandra HRC/ACIS data provided a first detection of X-ray pulsations and showed a compact $`10^{\prime \prime }`$ surrounding pulsar wind nebula (PWN) (Gotthelf, Halpern, & Dodson, 2002; Dodson & Golap, 2002). More recent XMM-Newton spectroscopy (McGowan et al., 2004) has provided improved measurements of the X-ray spectrum and pulsations. Early claims that the PWN is detected in TeV $`\gamma `$-rays (Kifune et al., 1995; Chadwick et al., 1997) have not been supported by recent HESS observations (Aharonian et al., 2005).
PSR B1706$``$44 is superposed on a radio-bright spur of the supernova remnant G343.1$``$2.3, which has a similar, albeit unreliable, $`\mathrm{\Sigma }D`$ distance of $`3`$ kpc (McAdam, Osborne & Parkinson, 1993). The pulsar DM gives a distance of $`2.3\pm 0.3`$ kpc in the Cordes & Lazio (2002) model. Dodson & Golap (2002) have argued for an association. In particular, they found a faint southern extension of the SNR, which would place the pulsar within the full SNR boundary. They also noted an approximately N-S elongation of the X-ray PWN, pointing roughly back to the SNR center and argued that this would represent a trailed nebula. The required velocity for travel from the approximate geometric center of the SNR, about $`12^{}`$ away, was $`1000d_3/\tau _4`$ km/s where $`\tau _4`$ is the pulsar age in units of $`10^4`$ yr. There are, however, some challenges to this SNR association. Koribalski et al. (1995) in an HI absorption study of the pulsar found velocity components setting lower and upper bounds for the distance of $`d_{min}=2.4\pm 0.6`$ kpc and $`d_{max}=3.2\pm 0.4`$ kpc. However a prominent HI emission feature seen in the bright limb of G343.1$``$2.3 at $`32`$km/s is not seen by Koribalski et al. in the absorption spectrum of the pulsar, suggesting that it lies in front of the SNR. Also scintillation studies (Nicastro, Johnston & Koribalski, 1996; Johnston, Nicastro & Koribalski, 1998) suggest a low transverse velocity for the pulsar, $`v89`$ km/s. This estimate has been supported by more recent scintillation measurements (Johnston, priv. comm). Thus the distances of the pulsar and SNR are still fairly uncertain. We adopt here a generic distance of 3 kpc in the discussion that follows, but carry through the scaling to show the distance dependence.
Ng & Romani (2004) re-examined the AO1 50 ks HRC and $`15`$ ks ACIS-S exposures, and found that the compact PWN could be well modeled as an equatorial torus + polar jets, rather similar to the structures seen around the young Crab and Vela pulsars. They found that the brightest arc of PWN emission lay behind the pulsar’s inferred motion from the SNR center, while a polar jet extended well in front of the pulsar position. This is difficult to reconcile with a bow shock/trail interpretation. On the other hand, the PWN symmetry axis did indeed point back to the SNR center, suggesting that a more careful evaluation of the connection was in order. This is particularly interesting, since comparison of the PWN symmetry (= pulsar spin) and proper motion axes can constrain the origin of pulsar birth kicks (Spruit & Phinney, 1998; Lai, Chernoff & Cordes, 2001; Romani, 2004).
We have obtained a deeper 100 ks ACIS-I exposure of PSR B1706$``$44 and surroundings. The X-ray exposure coverage is compared to the overall geometry of G343.1$``$2.3 in Figure 1. Together with new ATCA radio continuum imaging we are able to study the rich structure in this PWN and further constrain its connections with the SNR.
## 2 Observations and Data Analysis
PSR B1706$``$44 was observed with the Chandra ACIS-I array (4 ACIS-I chips along with the S3 and S4 chips) on February 1-2, 2004 with standard imaging (3.2s TE) exposures. The CCD array was operated in ‘very faint’ (VF) mode, allowing improved rejection of particle backgrounds. The total live-time was 98.8 ks and no episodes of strong background flaring were observed. Hence all data are included in our analysis. The pulsar was positioned near the standard aimpoint of the I3 chip and all observing conditions were normal. We have also compared our new exposure with the archival (February 3, 2001) 14.3 ks ACIS-S3 exposure (obsID 0757). As usual the backside illuminated S3 chip suffered more from particle background and after cutting out periods of background flares, 11 ks of clean exposure remained. All analysis was performed using CIAO 3.2 and CALDB 3.0.0, including automatic correction for the ACIS QE degradation. These data were nearly free of pile-up; the maximum pixel counts at the pulsar position indicate only $`2.5`$% pile-up while the best fit model for the point source has an expected pile up fraction of $`3.5`$%. For sources with low pile-up we can maximize the spatial resolution of the ACIS image by removing the standard pixel randomization and applying an algorithm correcting the position of split pixel events (Mori et al., 2001). This decreases the on-axis PSF width in our data set by $`10`$%. These data are compared with radio observations of the PWN.
### 2.1 Radio Imaging and Astrometry
Data for the radio maps shown here were collected at the Australia Telescope Compact Array (ATCA) in Narrabri (latitude -30.3) (Frater, Brooks & Whiteoak, 1992). For the 1.4 GHz map in Figures 1 and 2a, the data acquisition and analysis are described in Dodson & Golap (2002). For the image contours in Figure 3a, the data first presented in Dodson & Golap (2002) were re-imaged including the 6km baselines and uniform weighting to highlight the high resolution features. The restoring beam size is $`9\stackrel{}{\mathrm{.}}0\times 7\stackrel{}{\mathrm{.}}8`$. To show the nebular structure, an 11mJy point source PSF has been subtracted at the position of the pulsar. Two maxima appear flanking the pulsar position. These are unlikely to be artifacts due to pulsar variability, as diffractive scintillation for this pulsar is particularly weak (Johnston, Nicastro & Koribalski, 1998). Since the data were collected in five sessions, spread over more than a year, it is in principal possible for slow refractive scintillation to change the pulsar flux between epochs and distort its PSF. However, each epoch used $``$12h of integration, so any residual epoch PSF should be close to circularly symmetric, in contrast to the structure near the pulsar which is clearly bipolar. Further observations, with pulsar binning, have been requested to confirm this result.
For the 4.8 GHz map in Figure 2b, observations were made at 4.8- and 8.6 GHz with the array in the standard configurations 0.75A, 1.5A and 6A on 06 Jan, 16 Feb and 11 Apr 2002. The maximum and minimum baselines for the 4.8-GHz data were 1 and 100 k$`\lambda `$ (angular resolutions of $`3.4^{^{}}`$ to $`2\stackrel{}{\mathrm{.}}1`$) for a total of 26 hours observation. In all cases we observed the two frequencies with bandwidths of 128 MHz. We used the ATNF correlator mode that divides each integration’s data into separate phase bins spanning the pulsar period. This firstly allowed the strongly pulsed point source flux to be excluded from the image and secondly allowed us to self-calibrate using the relatively strong point source flux from the pulsar. After data editing and calibrating we inverted the image with a uv-taper of $`20^{^{\prime \prime }}`$ and deconvolved it with the full polarization maximum entropy task PMOSMEM in MIRIAD.
The most important test of the SNR association would, of course, be a direct astrometric proper motion. With a 1.4 GHz flux of $``$11 mJy, PSR B1706$``$44 is relatively bright. As such it is suitable for phase referenced VLBI astrometry, if an in-beam reference could be found. Unfortunately searches for phase references adequate for Australian Long Baseline Array (LBA) and US VLBA experiments have not detected comparison sources with compact fluxes greater than $`1`$mJy. Attempts were made at external phase reference VLBA astrometry. However at 1.4 GHz, the nearest known reference source (2.5 away) was scatter broadened to $`50`$ mas. With the strong ionosphere at such low elevation, the next nearest known source (10 away) is too distant for effective calibration. Since the pulsar spectrum is steep, an attempt at VLBA astrometry at 5 GHz was also unsuccessful; at this low elevation the system temperature was $`45\times `$ nominal and only six VLBA antennae could be used, reducing the sensitivity to $`15\%`$ of nominal. So unfortunately we have only tied-array astrometry at present. Even if the pulsar does travel from the geometric center of G343.1$``$2.3, the expected proper motion is only $`40`$ mas/yr; the existing time base of VLA/ATCA imaging does not yet allow a serious constraint on this motion. We must conclude that a direct proper motion measurement awaits substantially increased (SKA or EVLA) capabilities and a long-duration, large base-line experiment.
### 2.2 X-ray Spatial Analysis
To show the diffuse emission surrounding PSR B1706$``$44 we plot (Fig. 2a) a 1-7keV image with point sources removed (except the pulsar). These data are exposure corrected to minimize the chip gaps and heavily smoothed on a $`20^{\prime \prime }`$ scale. The diffuse emission is an edge-brightened, radius $`110^{\prime \prime }`$ cavity surrounding the pulsar with a faint extension to the west. Contours of the 1.38 GHz radio map show good correlation with the radio emission in the bar crossing G343.1$``$2.3 (Fig. 1). We will refer to this structure as the ‘nebula’.
Moving in toward smaller scale, in 2b we show a 1-7keV image, smoothed with a $`1\stackrel{}{\mathrm{.}}5`$ Gaussian. Point sources have not been removed. This shows that the cross structure fit by Ng & Romani (2004) extends across $`1^{}`$. Narrow X-ray jets, which we refer to here as the ‘outer jet’ (extending south) and ‘outer counter jet’ (extending north) start $`10^{\prime \prime }`$ from the pulsar and continue to $`30^{\prime \prime }`$. Bracketing these is faint diffuse X-ray emission which we will call the ‘equatorial PWN’. For comparison we draw contours of a 4.8 GHz ATCA image with $`21^{\prime \prime }\times 18^{\prime \prime }`$ restoring beam. These observations have the pulsar ‘gated out’ and show that the radio PWN has a hollow center bracketing the equatorial PWN. Diffuse radio peaks are, in fact, seen just east and west of this X-ray structure.
Finally, we show in Figure 3 a lightly smoothed image of the central region of the PWN, stretched to bring out the faint outer jets. The contours are drawn from a 1.38 GHz ATCA image, where the 6-km baselines have been weighted to produce a $`9\stackrel{}{\mathrm{.}}0\times 7\stackrel{}{\mathrm{.}}8`$ restoring beam. A point source PSF has been subtracted at the pulsar position. Two local radio maxima with peak fluxes $`2`$mJy and $`2.5`$mJy bracket the ‘torus’ structure. The radio then shows a sub-luminous zone surrounding the ‘equatorial PWN’; beyond $`30^{\prime \prime }`$ the radio brightens again, as in figure 2b. No emission appears along the ‘outer jets’. Indeed there appear to be evacuated channels in the radio emission, but improved S/N and resolution are needed to probe this sub-mJy structure. The frames on the right show the innermost region of the PWN with the best-fit torus + inner jet model (§2.3).
The overall geometry of the PWN is strongly reminiscent of that surrounding the Vela pulsar. In particular Pavlov et al. (2003) have described a series of ACIS images of the Vela nebula which show a torus-like structure, an inner jet and counter jet and a faint narrow outer jet system. This imaging sequence showed that the Vela outer jet, which is patchy and strongly bent, varies dramatically on timescales of days to weeks. Apparent motion of blobs within the jets suggests mildly relativistic bulk velocities and strong instabilities. For PSR B1706$``$44 our single sensitive image does not let us comment on variability. However we argue that the relatively straight and narrow jets, $`3\times `$ longer than those of Vela, and symmetric PWN structure are a consequence of a static uniform external environment and a low pulsar velocity. At 1.4-8.5 GHz Dodson et al. (2003) have found that the Vela PWN has two bright patches bracketing the X-ray torus and jets, in a structure quite similar to that in Figure 2b. Polarization imaging of the Vela radio structure suggests that these two patches represent the limbs of a toroidal B field structure. This implies that the rotation axis controls the PWN symmetry to large radii.
### 2.3 Nebula Structure Fits
Following Ng & Romani (2004) we have fitted our new ACIS image to a point source PSF, Doppler boosted equatorial torus, polar jets and uniform background. The fitting minimizes residuals using a Poisson-based likelihood function. Monte Carlo simulations of Poisson realizations of the best-fit model are in turn re-fitted to generate statistical errors and their co-variance matrices. Table 1 contains the best-fit values. The torus radius and axis inclination and position angles are $`r`$, $`\zeta `$ and $`\mathrm{\Psi }`$, respectively. See Ng & Romani (2004) for the definition of the other parameters and the details of the fitting technique. In Table 1 the inner jet/counter-jet are constrained to lie along the torus axis in the fits.
In addition to the statistical errors, there are certainly systematic errors, in particular induced by unmodeled PWN components. For example, it is clear that there are counts in excess of the torus+jet model in a cap surrounding the inner counter-jet. Interestingly similar structure is seen in the Crab PWN. We have made an attempt to constrain the systematic biases by modifying the fitting model. For example, allowing the (inner) jet and counter-jet to have a free position and amplitude shifts the best-fit position angle to $`\mathrm{\Psi }=165\pm 0.5^{}`$ and the inclination to $`\zeta =56.7\pm 1.0^{}`$. We therefore infer systematic errors about $`3\times `$ larger than our rather small statistical errors.
We have also measured the outer jet/counter-jet system. Minimizing the residual to a 1-D line passing through the pulsar, the two jets together lie at $`\mathrm{\Psi }_{\mathrm{outer}}=169.4\pm 0.15^{}`$. If the jets are fitted separately, we obtain $`\mathrm{\Psi }_{\mathrm{outer}}=168.4\pm 0.2^{}`$ and $`170.9\pm 0.2^{}`$ for the outer jet and counter-jet, respectively. Thus the two jets are mis-aligned at the $`8\sigma `$ level. A fit to the count distribution about the best fit axis shows that the narrow outer counter-jet has a Gaussian FWHM across the jet of $`2\stackrel{}{\mathrm{.}}3\pm 0\stackrel{}{\mathrm{.}}2`$. The outer jet appears broader at the base with an initial width of $`4\stackrel{}{\mathrm{.}}9\pm 0\stackrel{}{\mathrm{.}}5`$, continuing at FWHM=$`2\stackrel{}{\mathrm{.}}7\pm 0\stackrel{}{\mathrm{.}}3`$ for its outer half. These estimates have been corrected for the telescope PSF, which is quite uniform this close to the aimpoint.
It is important to note that at the observation roll angle, the read-out direction lies at $`\mathrm{\Psi }=168.7^{}`$! Due consideration, however, shows that the jet structure cannot be produced by the read-out trail. First the jets cover only $`1^{}`$; the read-out excess should cover the full I3 chip. Second the pulsar provides only $`2900`$ $`17`$ keV counts. The read-out trail (out-of-time) image of this source should contribute only 36 counts over the full $`8.3^{}`$ strip across I3 and $`2.5`$ counts in the ‘jet regions’; the outer jet and counter-jet have $`92`$ and $`93`$ $`17`$ keV counts, respectively. Finally the outer jets are much harder than the soft X-ray emission from the pulsar; indeed with a mean detected photon energy $`2`$ keV these are the hardest extended features in the image.
The over-all system, showing an asymmetric torus, broad inner jets and narrow outer jets is, of course, very similar to the Vela PWN as studied with Chandra by Pavlov et al. (2003). We will discuss the comparison with the Vela system in §3, highlighting the differences. We interpret these as suggesting that the PSR B1706$``$44 PWN has developed from a low velocity pulsar.
### 2.4 Spectral Analysis
For the best possible constraints on the source spectrum, we have reprocessed both the $`11`$ ks cleaned ACIS-S data set and our new $`100`$ ks ACIS-I data set with the new time-dependent gain adjustment and CTI correction available in CIAO 3.2. The updated RMFs should in particular improve the low energy calibration, important for obtaining the best estimates of $`N_H`$. As noted above, in these data sets the pile-up was negligible at $`3\%`$. To model the aperture corrections, 10 PSFs with monochromatic energies from 0.5 to 9.5keV were simulated using the Chandra Ray-Tracer program, ChaRT. The enclosed energy fraction as a function of radius was fitted to a linear function of energy and this was used to correct the ARFs used in the spectral fit. In extracting the pulsar spectrum, an aperture of radius 1<sup>′′</sup> was used to minimize nebular contamination. Results from the combined fits of the ACIS-I and ACIS-S pulsar data sets are listed in Table 1; the spectral fits are substantially better for composite models with both thermal and power-law components. Spectral parameter errors are projected multi-dimensional $`1\sigma `$ values. We quote both absorbed and unabsorbed fluxes. As is often the case with low statistics X-ray spectra, projected (multi-dimensional) errors on the fluxes are very large due to spectral parameter uncertainties. Thus, we follow other authors in quoting flux errors as $`1\sigma `$ single parameter values.
To get the best constraints on the point source spectrum, Table 2 gives fits with $`N_H`$ held to the value from the power law fits to the extended emission. We have also compared our results with the XMM-Newton fitting of McGowan et al. (2004), by fitting counts in the $`20^{\prime \prime }`$ aperture used in that observation. Our parameters and fluxes for the thermal component are generally in very good agreement. However, since this aperture contains much of the torus and central PWN, XMM-Newton substantially overestimates the non-thermal flux for the point source. Their fit power law flux corresponds to $`7.5\times 10^{13}\mathrm{erg}/\mathrm{cm}^2/\mathrm{s}`$ (0.5-8keV, unabsorbed). In the $`20^{\prime \prime }`$ aperture we find $`9.6\times 10^{13}\mathrm{erg}/\mathrm{cm}^2/\mathrm{s}`$ (0.5-8keV); while the small Chandra point source aperture gives $`2.3\times 10^{13}\mathrm{erg}/\mathrm{cm}^2/\mathrm{s}`$ (0.5-8keV) for the power law component. We find a similar $`3\times `$ excess in the power-law \+ atmosphere flux for the fit to the large XMM-Newton aperture. Conversion of the power law flux observed in our small point source aperture to the XMM-Newton band shows that the expected PN+MOS (0.2-10 keV) count rate is 18% of the total (power law + thermal) counts in the $`20^{\prime \prime }`$ aperture. In soft (0.2-1.35 keV) and hard (1.35-10 keV) bands the predicted fraction of the counts from the power law are 12% and 21%, respectively. However the light curves of McGowan et al. (2004) show that the pulse fractions are 21% (soft), 12% (hard) and 11% (total). Since the small aperture power law produces only 12% of the soft counts but 21% are pulsed, there must a thermal pulse component. Conversely, since the power law produces 21% of the hard band flux, but this only has a pulse fraction of 12%, some of the power law counts must be unpulsed. Extrapolation of the PWN count excess above the point source PSF in the sub-luminous zone at 2$`3^{\prime \prime }`$ produces $`1`$% of the point source aperture counts. Thus the larger scale torus emission does not contribute significantly to the point source power law and cannot account for its unpulsed component. This suggests that part of the magnetospheric emission is nearly isotropic or that there is a very compact ($`1^{\prime \prime }`$) PWN component at the pulsar position.
For the thermal component, the fit flux gives an emitting area (effective radius) as a function of distance. Our fit to a pure blackbody gives $`R_{eff}=2.8d_3`$ km. Thus for reasonable distances, this flux represents hot $`T2\times 10^6`$ K emission from a small fraction of the stellar surface ($`4.5`$% for an $`R^{\mathrm{}}=13.1`$ km star). The light element neutron star atmosphere models, such as the pure H $`10^{12}`$ G model grid used here (Zavlin et al., 1996), have large Wien excesses. When fit they give lower $`T_{eff}`$. Also, the black body departures allow one, in principle, to fit both the surface redshift and radius. In practice, these are typically highly degenerate in CCD-quality data. We assume here a generic surface radius of $`R_s=10`$ km, corresponding to $`R_{\mathrm{}}=R_s(12GM/R_sc^2)^{1/2}=13.1`$ km. With $`N_H`$ free (giving $`5.9\pm 0.9\times 10^{21}\mathrm{cm}^2`$) our thermal flux normalization gives a radiating radius of $`R_{\mathrm{}}=27.4d_3`$km, which is difficult to reconcile with expected neutron star radii for any $`d>1.8`$ kpc. However, when $`N_H`$ is fixed at the nebular value of $`5\times 10^{21}\mathrm{cm}^2`$, we get an effective radius of $`R_{\mathrm{}}=16.1d_3`$km, which is tolerable even at our nominal 3 kpc distance.
Analysis of the low signal-to-noise, extended flux depends critically on the background subtraction. Given the limited statistics, only simple absorbed power-law fits were attempted for all non-thermal sources. The results are listed in Table 3. For consistency, all fits are to the 0.5-8 keV range and we quote both absorbed and unabsorbed fluxes.
Note that there is significant softening of the extended emission as one progresses to larger scales and that the jet components appear to be the hardest of all. This is certainly consistent with the idea that the central pulsar supplies fresh energetic electrons and that synchrotron burn-off increasingly softens the spectrum as older populations are viewed in the outer PWN. Again this trend is common in the well-measured young PWNe. Allowing the photon index to vary for the different nebula components, the best fit to a global absorption value for the extended emission gives us our fiducial $`N_H=5\times 10^{21}\mathrm{cm}^2`$. This is consistent with free-fit values for the point source, but given complexities of the composite thermal+power law model, we consider the nebular fit value more robust. Note that with DM=75.7 $`\mathrm{cm}^3\mathrm{pc}`$, the $`H/n_e21`$ for this sight-line is large, but not unprecedented for low $`|b|`$ pulsars. This is also consistent with the $`H_I`$ absorption measurements and a fiducial SNR distance $`3`$ kpc, given the appreciable uncertainties.
## 3 Interpretation and Conclusions
A number of authors have discussed the evolution of a PWN within an expanding supernova remnant. For example, van der Swaluw (2001) and Chevalier (2005) describe the early evolution when the supernova ejecta are in free expansion. Later, after the remnant interior is heated by the passage of the reverse shock, the PWN evolves within the Sedov phase supernova remnant whose radius is $`R_{\mathrm{SNR}}=1.17(E_0/\rho )^{1/5}t^{2/5}`$ for an explosion energy $`E_0`$ in a $`\gamma =5/3`$ medium of density $`\rho `$. PSR B1706$``$44 has a characteristic age $`10^4\tau _4`$ yr with $`\tau _41.7`$, so G343.1$``$2.3 should be safely in the Sedov phase with an expected angular size
$$\theta _{\mathrm{SNR}}16^{}(E_{51}/n_0)^{1/5}t_4^{2/5}/d_3$$
$`(1)`$
for a supernova releasing energy $`E_0=10^{51}E_{51}\mathrm{erg}`$ in an external medium density $`n_0\mathrm{cm}^3`$, at a true age $`10^4t_4`$y at a distance $`3d_3`$ kpc. The observed size then implies $`E_{51}11n_0t_4^2d_3^5`$, requiring a fairly energetic explosion for $`d>2`$ kpc. During the Sedov phase the interior pressure is
$$P_{SNR}10^9E_{51}^{2/5}n_0^{3/5}t_4^{6/5}\mathrm{g}/\mathrm{cm}/\mathrm{s}^2$$
$`(2)`$
and is relatively constant away from the SNR limb.
The pulsar blows a wind bubble within this SNR interior, whose radius is $`R_{\mathrm{PWN}}(E_{}/E_0)^{1/3}R_{\mathrm{SNR}}`$ for a PWN bubble energy $`E_{}=f\dot{E}\tau _c`$ (van der Swaluw & Wu, 2001). Although the accuracy of this dependence of PWN radius on pulsar injection energy has been questioned (Blondin, Chevalier & Frierson, 2001), we adopt it for the following estimates. With the observed ratio of radii, $`R_{\mathrm{PWN}}/R_{\mathrm{SNR}}=1.8^{}/25^{}`$, we obtain $`E_{}=3.7\times 10^4E_0`$, i.e. this PWN has quite low internal energy. This is also reflected in the low radio and X-ray fluxes. Together we use these estimates, the observed size of the SNR, equation (1) and the measured $`\dot{E}_{36}=3.4`$ and $`\tau _c=1.75\times 10^4`$y to write $`f(R_{\mathrm{PWN}}/R_{\mathrm{SNR}})^3E_0/(\dot{E}\tau _c)=2.1n_0t_4^2d_3^5`$. Now, if the PWN is adiabatic and we assume spindown with constant $`B`$ and braking index $`n=3`$ from an initial period $`P_0`$, we find that the total energy in the plerion is $`[(P/P_0)^21]\dot{E}\tau _c`$. Then, setting $`f=(P/P_0)^21`$ and eliminating the true age $`t`$ using $`t=\tau _c[1(P_0/P)^2]`$ for magnetic dipole spindown, we obtain a constraint on the initial spin period
$$[1(P_0/P)^2]^3/(P_0/P)^2=0.68n_0d_3^5$$
$`(3)`$
which has a solution of $`P_0=0.61P=62`$ ms for d=3 kpc, and $`P_0=0.79P=80`$ ms for d=2 kpc. The corresponding true ages are 0.61$`\tau _c`$ ($`1.1\times 10^4`$y) and 0.38$`\tau _c`$ ($`0.67\times 10^4`$y), respectively. These numerical values are for $`n_0=1`$ and the density dependence from Equation (3) is quite weak. van der Swaluw & Wu (2001) present a similar sum for $`P_0`$, assuming a known $`E_0`$; the above formulation emphasizes the sensitivity to the poorly known $`d`$. Note that with the large implied initial period, the integrated PWN energy is quite comparable to the present spin energy, with $`f1.7`$ at d=3 kpc and $`f0.62`$ at d=2 kpc. So the spindown luminosity is roughly constant in the adiabatic phase and the PWN growth is closer to $`t^{11/15}`$ than to the $`t^{3/10}`$ law appropriate for impulsive energy injection (van der Swaluw, 2001).
Inside this wind bubble, the Sedov interior pressure confines the PWN, giving rise to a termination shock at
$$\theta _{\mathrm{WS}}(\dot{E}/4\pi cP_{\mathrm{SNR}})^{1/2}/d.$$
$`(4)`$
which results in $`\theta _{\mathrm{WS}}1\stackrel{}{\mathrm{.}}2\dot{E}_{36}^{1/2}E_{51}^{1/5}n_0^{3/10}t_4^{3/5}d_3^1`$. If we apply the SNR estimate for $`E_0`$ above, this becomes $`\theta _{\mathrm{WS}}0\stackrel{}{\mathrm{.}}72\dot{E}_{36}^{1/2}n_0^{1/2}t_4d_3^2`$. Then, using $`\dot{E}_{36}4`$ and applying the age estimate following Equation (3) we get $`\theta _{\mathrm{WS}}1\stackrel{}{\mathrm{.}}5n_0^{1/2}`$ (d=3 kpc) or $`\theta _{\mathrm{WS}}2\stackrel{}{\mathrm{.}}1n_0^{1/2}`$ (d=2 kpc). These estimates are reasonably consistent with the observed $`3^{\prime \prime }`$ torus radius, especially since an equatorially concentrated flow should have a stand-off distance 1.5-2$`\times `$ this spherical scale. The polar jets can have an initial shock at somewhat larger angle, with the resulting pitch angle scattering illuminating the jets somewhat further from the pulsar.
Of course, this bubble is offset from the center of G343.1$``$2.3 at $`R=0.5R_{SNR}`$ (figure 1). This is inside the $`0.68R_{SNR}`$ where van der Swaluw, Downes & Keegan (2004) note that the increasing density causes the pulsar to be supersonic, so a bow shock should not have yet formed. These authors however compute numerical models of a fast moving pulsar in a SNR interior. As the pulsar moves, the PWN should become highly asymmetric with a ‘relic PWN’ at the SNR center and the pulsar placed near the leading edge of the PWN; see van der Swaluw, Downes & Keegan (2004) figures 7 and 8. We see no PWN structure near the geometric center of G343.1$``$2.3 and, if the ‘bubble nebula’ is identified with the shocked pulsar wind, the pulsar is certainly not offset from its center along the proper motion axis (away from the SNR center). So these models are an inadequate description of G343.1$``$2.3. From Figure 2, the pulsar is well centered in the bubble nebula, with any offset from its center along the axis to the SNR substantially less than $`30\mathrm{\Delta }_{30}`$ arcsec. Thus
$$v<40\mathrm{\Delta }_{30}d_3/t_4\mathrm{km}/\mathrm{s},$$
$`(5)`$
and the pulsar cannot have moved far from the explosion center. This is, of course, consistent with the scintillation results.
We can reconcile the symmetric PWN with the offset SNR shell if we assume that the pulsar progenitor exploded toward the edge of a quasi-spherical cavity. One scenario (also posited by Gvaramadze 2002, Bock & Gvaramadze 2002) that can associate the low velocity pulsar with G343.1$``$2.3 is to assume that the progenitor star had a stellar wind of mass loss rate $`\dot{M}_810^8M_{}`$/yr and wind speed $`10^8v_{\mathrm{w8}}`$ cm/s over $`tt_710^7`$y, typical of the $`10M_{}`$ stars that dominate the pulsar progenitors (Maeder, 1981). This evacuates a stellar wind bubble of size
$$\theta _{SW}=46^{}\left(\dot{M}_8v_{\mathrm{w8}}^2/n_0\right)^{1/5}t_7^{3/5}/d_3.$$
$`(6)`$
During the main sequence lifetime, the star moving at $`10v_6`$ km/s travels $`2^{}v_6t_7/d_3`$ and so it can easily traverse its wind bubble. Thus, one can imagine an off-center supernova in a nearly symmetric stellar wind bubble of radius $`25^{}`$: the supernova blast wave expands to fill the bubble, passing to the Sedov phase near its present radius. The supernova produces a neutron star with little or no kick, placing the pulsar near its present position. This has the added advantage of accommodating the rather large SNR size with a more modest energy of a few$`\times 10^{51}`$ erg. The PWN energy and size estimates above would then be somewhat amended; this would require a careful numerical simulation.
For the reasons detailed in the introduction, it is not yet clear that PSR B1706$``$44 and G343.1$``$2.3 are associated. So for completeness we can consider the case when the shocked pulsar wind blows an adiabatic bubble in a static, low $`P_{\mathrm{ext}}`$ external medium (Castor, McCray & Weaver, 1975). If we assume that the pulsar was born (sans SNR) or entered a confining region of the ISM $`10^4`$y ago and that since then it has been spinning down at the present energy loss rate, we find that it will blow a bubble of angular size
$$\theta _{\mathrm{BN}}0.76(\dot{E}t^3/\rho )^{1/5}/d120^{\prime \prime }(\dot{E}_{36}/n_0)^{1/5}\tau _4^{3/5}/d_3.$$
$`(7)`$
These estimates change somewhat for a pulsar born at $`P_0P`$; since we are not making the association with the SNR G343.1$``$2.3, we can make no estimate of the initial spin period. As first noted by Dodson & Golap (2002), the $`4^{}`$ wide radio spur across the face of G343.1$``$2.3 has the approximate scale of such a ‘bubble nebula’. If the PWN stays unmixed (relativistic) then the interior of the bubble will have a pressure $`P_{\mathrm{BN}}\dot{E}t/(4\pi R^3)1.6\times 10^{10}\left(n_0^3\dot{E}_{36}^2/\tau _4^4\right)^{1/5}\mathrm{g}/\mathrm{cm}/\mathrm{s}^2`$. In turn, the torus termination shock in this medium is at
$$\theta _{\mathrm{WS}}2\stackrel{}{\mathrm{.}}9(\dot{E}_{36}/n_0)^{3/10}\tau _4^{2/5}/d_3.$$
$`(8)`$
If (e.g. through Rayleigh-Taylor instabilities) the pulsar wind is well mixed with the swept up gas, the adiabatic thermal pressure would be $`2\times `$ larger. Interestingly, the angular scales for this scenario are also reasonably compatible with the observed torus and bubble nebula size. Of course this scenario leaves open the question of the pulsar origin. Again the pulsar would need to have a quite low velocity to produced the observed symmetry.
Turning to the spectral results, we note that Possenti et al. (2002) fit a correlation between spindown energy and the PSR+PWN luminosity: $`L(210\mathrm{keV})=1.8\times 10^{38}\dot{E}_{40}^{1.34}\mathrm{erg}/\mathrm{s}`$. For the PSR B1706$``$44 parameters this predicts a flux $`f(210\mathrm{keV})=4.7\times 10^{12}d_3^2\mathrm{erg}/\mathrm{cm}^2/\mathrm{s}`$. The observed 2-10 keV flux is in fact $`1.7\times 10^{12}\mathrm{erg}/\mathrm{cm}^2/\mathrm{s}`$, even including the outer ‘bubble nebula’; without this component it is half as large. These correlations are not very accurate, but this does imply that the PSR B1706$``$44 PWN is substantially under-luminous for any distance less than 3 kpc. Gotthelf (2003) has derived correlations between the pulsar spindown power and the pulsar/PWN spectral indices. His relation predicts $`\mathrm{\Gamma }_{\mathrm{PSR}}=0.63\pm 0.17`$ (substantially smaller than our power law index $`1.6\pm 0.2`$) and $`\mathrm{\Gamma }_{\mathrm{PWN}}=1.3\pm 0.3`$ (not inconsistent with the values measured for the torus and equatorial PWN).
As described in §2.4, the spectrum softens appreciably from the central torus to the outer bubble nebula. This suggests increased aging of the synchrotron population. Figures 2 and 3 show that the bulk of the radio emission lies in the ‘bubble nebula’ region. So we can take the radio flux and spectral index from Giacani et al. (2001) and compare with our nebula X-ray flux (Fig. 4). Comparing the radio spectral index $`\alpha _R=0.3`$ with the best fit X-ray index $`\alpha _X=0.77`$, shows a break quite close to the $`\mathrm{\Delta }\alpha =0.5`$ expected from synchrotron cooling. The extrapolated intersection of these power laws gives a break frequency of Log\[$`\nu _B(\mathrm{Hz})]=12.2^{+0.9}_{1.1}`$. For the fiducial pulsar age of $`1.7\times 10^4`$y, this corresponds to a nebula field of $`1.4_{0.6}^{+2.1}\times 10^4`$ G. Note that the magnetic pressure from this (photon flux-weighted) average field is $`8\times 10^{10}\mathrm{g}/\mathrm{cm}/\mathrm{s}^2`$, somewhat larger than the nebula pressure estimated from its radius. This may indicate field compression in the nebula limb. In general, if the mean nebula field is $`10^4B_4`$ G for a nebula of angular radius $`100^{\prime \prime }\theta _{100}`$, the total nebula field energy is $`E_B1.5\times 10^{47}B_4^2(\theta _{100}d_3)^3`$ erg. This is comfortably less than the present spin energy $`E_{\mathrm{PSR}}2\times 10^{48}`$ erg, so the nebula can be easily powered even if the pulsar was born close to its present spin period. We find that this cooling break field is substantially larger than the equipartition field of $`1015\mu `$G inferred for the radio and X-ray emitting populations (also the minimum equipartition nebula energy $`9\times 10^{45}`$ erg is substantially smaller). The cooling break field can also be compared to that expected from simple radial evolution of the pulsar surface field: if this field $`B_{}=3\times 10^{12}`$ G falls off as $`r^3`$ to the light cylinder, then as $`1/r`$ to the wind shock where it is compressed we get $`B_{\mathrm{WS}}3B_{}r_{}^3/(r_{\mathrm{LC}}^2r_{\mathrm{WS}})1`$mG. If it continues to fall off as $`1/r`$ beyond this we get a field at the limb of the bubble nebula of $`30\mu `$G. So the best we can do is to infer a mean nebular field $`1030\times `$ the equipartition value, with some generation of new field beyond the torus wind shock. The energetic requirements for this field, required to match the $`\nu _B`$ cooling break, are comfortably less than the energy available from PSR B1706$``$44. These field estimates are consistent with the non-detection of TeV ICS flux from this source (Aharonian et al., 2005).
The narrow outer jets also have a power-law spectrum and are almost certainly synchrotron-emitting. For a reasonable $`0.1B_4`$mG field, the observed X-rays of $`E_\gamma =1.5B_4\mathrm{\Gamma }_{7.5}^2`$ keV require substantial $`e^\pm `$ energies, with $`\mathrm{\Gamma }_e=3\times 10^7\mathrm{\Gamma }_{7.5}`$ near the radiation-reaction limited primary Lorentz factor inferred for many polar cap (Muslimov & Harding, 2003) and outer magnetosphere (Romani, 1996) pulsar models. Since the jet is narrow, confinement of these pairs imposes a (not very restrictive) lower bound on the jet field $`B_4>0.075E_5^{1/3}/(d_3\theta _w)`$ where the maximum observed jet photon has $`E5E_5`$ keV and the observed outer jet half-width is $`\theta _w`$arcsec. A more restrictive upper limit on the mean jet field comes from the observation that the jets do not soften noticeably before their end $`30\theta _{30}`$arcsec from the pulsar. If we assume a jet bulk speed $`\beta `$c, then arguing that the flow time is shorter than the synchrotron cooling time gives us the limit $`B_4<8.5\left(\frac{\beta }{d_3\theta _{30}}\right)^{2/3}E_5^{1/3}.`$
When the observed jet spectrum has an energy index of $`\alpha =\mathrm{\Gamma }10.3`$, we infer a power-law spectrum of $`e^\pm `$ in the jet $`N(\mathrm{\Gamma }_e)d\mathrm{\Gamma }_e=K\mathrm{\Gamma }_e^pd\mathrm{\Gamma }_e`$, with $`p=2\alpha +11.6`$. We can then make an estimate of the minimum jet luminosity, i.e. at ‘equipartition’ when $`B^2=6\pi m_ec^2\mathrm{\Gamma }_eN(\mathrm{\Gamma }_e)𝑑\mathrm{\Gamma }_e`$. Given the observed combined outer jet luminosity (0.5-8 keV) $`L=4\pi d^2f_{\mathrm{oj}}2.7\times 10^{31}\mathrm{erg}/\mathrm{s}`$ and emitting volume $`V2\times \theta _L\times \pi \theta _wd^31.1\times 10^{52}(\theta _L/20)\theta _w^2d_3^3\mathrm{cm}^3`$ with the angles in arcsec, we can estimate the equipartition field for an isotropic plasma as
$$B_{eq}=\left[\frac{18\pi }{\sigma _T}\left(\frac{2\pi m_ec}{E_{\mathrm{max}}}\right)^{1/2}\frac{2(1\alpha )}{12\alpha }\frac{1(E_{\mathrm{min}}/E_{\mathrm{max}})^{(12\alpha )/2}}{1(E_{\mathrm{min}}/E_{\mathrm{max}})^{(1\alpha )}}\frac{L}{V}\right]^{2/7}$$
$`(9)`$
where the observed photon spectrum runs from $`E_{\mathrm{min}}`$ to $`E_{\mathrm{max}}`$. For the observed flux this gives
$$B_{\mathrm{eq}}0.25\times 10^4q(\alpha )[\theta _L\theta _w^2d_3/20]^{2/7}\mathrm{G},$$
$`(10)`$
where $`q(\alpha =0.3)=1`$ is a weak function of $`\alpha `$. The corresponding minimum energy flux for the outer jet is
$$L_{\mathrm{oj}}=8\times 10^{33}\beta d_3^{12/7}\theta _w^{10/7}(\theta _L/20)^{2/7}\mathrm{erg}/\mathrm{s},$$
$`(11)`$
where the jet bulk velocity is $`\beta `$c. This is $`10^3\dot{E}`$ per jet and will, of course, be larger if the jet flow includes ions. Interestingly, if the pulsar couples roughly isotropically to the PWN, then the corresponding fraction of the outflow should subtend a half angle of $`5^{}`$. This is somewhat smaller than the angle subtended by the inner jets, but $`3\times `$ larger than the $`1^{\prime \prime }`$ width of the ends of the jet – there is substantial collimation of the jet energy flux.
We have argued that a low PSR velocity can explain the symmetry of the PWN. The central location of the pulsar and spherical post-shock flow may also allow the equatorial toroidal structure and polar jets to propagate undisturbed to large radii. We do, however, measure a small misalignment of the outer jets, corresponding to a deflection of $`\theta _{\mathrm{de}}=1.3\pm 0.15^{}`$ for each. If we imagine a pressure acting along the jet’s $`30^{\prime \prime }`$ length, then the required perturbation is $`\delta PL_{\mathrm{oj}}\mathrm{tan}\theta _{\mathrm{de}}/(\beta cA_{\mathrm{oj}})5\times 10^{14}L_{34}/(\beta \theta _{30}\theta _wd_3^2)\mathrm{g}/\mathrm{cm}/\mathrm{s}^2`$, where $`A_{\mathrm{oj}}`$ is the jet’s cross sectional area. This is only $`10^3`$ of the total pressure in the nebula. It could be due to ram pressure if the shocked nebular medium flows to the west at $`v1.7n_{\mathrm{neb}}^{1/2}`$ km/s.
Our X-ray measurements have established the PWN symmetry axis, presumably reflecting the pulsar spin axis, to very high precision. Unfortunately our original goal of relating this to the proper motion axis remains unfulfilled. It is true that the torus symmetry axis points roughly toward the center of G343.1$``$2.3, confirming the estimates from earlier Chandra data. However, the PWN symmetry about the pulsar and the low scintillation velocity suggest a very low transverse speed $`40`$km/s, which would preclude a birth site as distant as the SNR center. This low speed makes a direct proper motion challenging, but allows latitudinal asymmetries in the PWN flow to propagate undisturbed to fairly large radius, where they can be imaged with Chandra. Thus, study of this PWN offers some good opportunities to probe outflow dynamics and jet collimation. Study of this, and similar, PWNe may prove useful electrodynamic analogs of the $`10^6\times `$ more powerful AGN jets. A viable scenario for maintaining the G343.1$``$2.3/PSR B1706$``$44 association posits a supernova event near the present pulsar site, with the remnant inflating a pre-existing off-center cavity. However, the residual (small) proper motion could then have any direction. In fact, the faint extension of the PWN (bubble nebula) to the west and the increased radio surface brightness to the east might suggest a rather slow pulsar motion at PA $`80^{}`$. This would be nearly orthogonal to the torus symmetry axis. With the large $`P_0`$ estimated here, this could be construed as suggesting poor rotational averaging of a birth kick (Ng & Romani, 2004). So the PWN/SNR geometry offers both aligned and orthogonal axes. Only a sensitive astrometric campaign can detect or limit the pulsar motion and resolve this ambiguity.
This work was supported in part by NASA grants SAO G04-5060X and NAG5-13344. We thank the referee for several careful readings, which resulted in substantial changes. We also thank R.N. Manchester who provided a current pulsar ephemeris for the VLBI experiment.
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# Third order spectral branch points in Krein space related setups: limit-from{𝑃𝑇}-symmetric matrix toy model, MHD limit-from{𝛼²}-dynamo, and extended Squire equation 11footnote 1Presented at the 3rd International Workshop ”Pseudo-Hermitian Hamiltonians in Quantum Physics”, Istanbul, Turkey, June 20-22, 2005.
## 1 Krein space related physical setups and spectral phase transitions
Some basic spectral properties of the Hamiltonians of $`PT`$symmetric Quantum Mechanics (PTSQM) can be easily explained from the fact that these Hamiltonians are self-adjoint operators in Krein spaces — Hilbert spaces with an indefinite metric structure. In contrast to the purely real spectra of self-adjoint operators in ”usual” Hilbert spaces (with positive definite metric structure), the spectrum of self-adjoint operators in Krein spaces splits into real sectors and sectors with pair-wise complex conjugate eigenvalues. In physical terms these two types of sectors are equivalent to phases of exact $`PT`$symmetry and spontaneously broken $`PT`$symmetry . The reality of the spectrum of a PTSQM Hamiltonian, e.g., with complex potential $`ix^3`$, means that this spectrum is located solely in a real sector and that the operator is quasi-Hermitian<sup>2</sup><sup>2</sup>2Because quasi-Hermitian operators form only a restricted subclass of self-adjoint operators in Krein spaces (pseudo-Hermitian operators in the sense of Ref. ), the question for the reality of the spectrum seems up to now only partially solved. Apart from the requirement for existing $`PT`$symmetry of the differential expression of the operator, a subtle interplay between the operator domain and some, in general not yet sufficiently clearly identified, additional structural aspects of its differential expression seems to be responsible for the quasi-Hermiticity of the operator (exact $`PT`$symmetry of the corresponding PTSQM setup). in the sense of Ref. . In general, the spectrum of a self-adjoint operator in a Krein space is spreading over both sectors.
Considering the spectrum as a Riemann surface depending on model parameters (not necessarily moduli) from a space $`M`$, the boundaries between real and complex spectral sectors can be described as algebraic variety ((multi-component) hypersurface) $`\mathrm{{\rm Y}}`$ in $`M`$, $`\mathrm{{\rm Y}}M`$, where the Riemann surface has exceptional points of branching type . The most generic varieties of this type are those of codimension one, $`\text{codim}(\mathrm{{\rm Y}})=\text{dim}(M)\text{dim}(\mathrm{{\rm Y}})=1`$, what can be easily read off, e.g., from a pseudo-Hermitian $`2\times 2`$matrix model
$$H=\left(\begin{array}{cc}x+y& w\\ w^{}& xy\end{array}\right),H\psi =\lambda \psi $$
(1)
where $`\mathrm{{\rm Y}}`$ is simply a double cone $`y^2=|w|^2M(y,\mathrm{}w,\mathrm{}w)`$. Such codimension-one varieties correspond to a pairwise real-to-complex transition of two eigenvalues what in Riemann surface terms is equivalent to a square root branching of two of its sheets (two spectral branches). At the same time the two (geometric) eigenvectors coalesce into a single (geometric) eigenvector, and an additional associated vector (algebraic eigenvector) appears. Instead of two eigenvalues of geometric and algebraic multiplicity one (diagonal block), a single eigenvalue of geometric multiplicity one, $`(m(\lambda )=1)`$, and algebraic multiplicity two $`(n(\lambda )=2)`$, ($`2\times 2`$Jordan block) forms . When the parameters on this codimension-one variety are further tuned to $`y=\mathrm{}w=\mathrm{}w=0`$ one arrives at a diabolic point where the single eigenvalue gets geometric and algebraic multiplicity two ($`m(\lambda )=n(\lambda )=2`$, diagonal $`2\times 2`$ block). In a (more) general setting such configurations correspond to codimension-three varieties .
Apart from these generic real-to-complex transitions on codimension-one varieties, there may occur higher order intersections of more than two Riemann sheets simultaneously — on varieties $`\mathrm{{\rm Y}}`$ of higher codimension, $`\text{codim}(\mathrm{{\rm Y}})2`$, and with larger Jordan blocks in the spectral decomposition .
Below we present explicit examples for third order intersections in three Krein space related physical models. First, we sketch a few aspects of such intersections semi-analytically (algebraically) for a maximally simplified $`PT`$symmetric $`4\times 4`$matrix toy model. Afterwards, we show numerically that such intersections also occur in the operator spectra of the spherically symmetric $`\alpha ^2`$dynamo of magnetohydrodynamics (MHD) and of the recently analyzed $`PT`$symmetric interpolation model of Ref. (which can be understood as a $`PT`$symmetrically extended, rescaled and Wick-rotated version of the Squire equation of hydrodynamics).
## 2 Spectral triple points in a pseudo-Hermitian $`4\times 4`$matrix toy model
Similar to the two different multiplicity contents (Jordan structures) of two-fold degenerate eigenvalues (see, e.g., also ), one has to distinguish the following three types of Jordan structures/geometric multiplicities $`m(\lambda )`$ corresponding to a spectral triple point:
* type I: $`m(\lambda )=1`$, $`n(\lambda )=3`$; $`3\times 3`$Jordan block; e.g., coalescence of two square-root branch points connected by a purely real spectral segment (see Eq. (12) and Figs. 1 \- 3 below),
* type II: $`m(\lambda )=2`$, $`n(\lambda )=3`$; one $`2\times 2`$Jordan block $`+`$ one simple eigenvalue; e.g., a square-root branch point accidentally coincides with a simple eigenvalue,
* type III: $`m(\lambda )=3`$, $`n(\lambda )=3`$; diagonal matrix $`\lambda I_3`$; e.g., three accidentally coinciding eigenvalues (generalized diabolic point).
Subsequently, we concentrate on the physically most interesting case of type I triple points which correspond to pair-wise coalescing square-root branch points.
As maximally simplified toy model we choose a $`PT`$symmetric (pseudo-Hermitian) $`4\times 4`$matrix setup in a diagonal representation of the parity (involution) operator $`P`$
$$H=PH^{}P=\left(\begin{array}{cc}H_{++}& H_+\\ H_+& H_{}\end{array}\right),P=\left(\begin{array}{cc}I_2& 0\\ 0& I_2\end{array}\right),$$
(2)
where the $`2\times 2`$blocks satisfy the conditions $`H_{\pm \pm }=H_{\pm \pm }^{},H_+=H_+^{}`$. In order to keep the demonstration as simple as possible, it suffices to consider a model based on two real (truncated) elementary pseudo-Hermitian $`2\times 2`$matrices
$$H_{1,2}=\left(\begin{array}{cc}x_{1,2}+y_{1,2}& w_{1,2}\\ w_{1,2}& x_{1,2}y_{1,2}\end{array}\right),x_{1,2},y_{1,2},w_{1,2}$$
(3)
as subsystems, which we embed $`PT`$symmetrically into the also highly reduced<sup>3</sup><sup>3</sup>3Nine of the 16 effective free real parameters of the pseudo-Hermitian $`4\times 4`$matrix $`H`$ are set to zero. and real pseudo-Hermitian $`4\times 4`$matrix $`H`$ with block structure (2)
$$H_1,H_2H=\left(\begin{array}{cccc}x_1+y_1& 0& w_1& z\\ 0& x_2+y_2& 0& w_2\\ w_1& 0& x_1y_1& 0\\ z& w_2& 0& x_2y_2\end{array}\right).$$
(4)
The interaction between the subsystems $`H_1`$ and $`H_2`$ is controlled by the coupling parameter $`z`$. For vanishing coupling, $`z=0`$, the characteristic equation $`\mathrm{\Delta }(\lambda )=det(H\lambda I_4)=0`$ factors as $`\mathrm{\Delta }(\lambda )=det(H_1\lambda I_2)det(H_2\lambda I_2)=0`$, and the four eigenvalues of the matrix $`H`$ are defined by the eigenvalues of $`H_1`$ and $`H_2`$
$$\lambda _{1,\pm }=x_1\pm \sqrt{y_1^2w_1^2},\lambda _{2,\pm }=x_2\pm \sqrt{y_2^2w_2^2}.$$
(5)
The corresponding two square root branch points $`\lambda _1=x_1,\lambda _2=x_2`$ with $`\mathrm{\Delta }(\lambda )=0,_\lambda \mathrm{\Delta }(\lambda )=0`$ are located on the two (reduced) double cones<sup>4</sup><sup>4</sup>4The two diabolic points (with eigenvalues of geometric and algebraic multiplicity two) of the subsystems $`H_1`$, $`H_2`$ are located at $`y_1=w_1=0`$ and $`y_2=w_2=0`$, respectively. (crossed lines) defined by
$$y_1^2=w_1^2,y_2^2=w_2^2.$$
(6)
For non-vanishing interaction $`z0`$, the eigenvalues of $`H`$ are given as roots of the quartic equation
$$\mathrm{\Delta }(\lambda )=\lambda ^4+a_3\lambda ^3+a_2\lambda ^2+a_1\lambda +a_0=0$$
(7)
with coefficients<sup>5</sup><sup>5</sup>5The reason for restricting to the maximally simplified (truncated) $`4\times 4`$matrix setup (4) was in keeping these coefficients $`a_k,k=0,\mathrm{},3`$ in a form sufficiently simple for quick inspection.
$`a_3`$ $`=`$ $`2(x_1+x_2)`$
$`a_2`$ $`=`$ $`(y_1^2w_1^2)(y_2^2w_2^2)+(x_1+x_2)^2+2x_1x_2+z^2`$
$`a_1`$ $`=`$ $`2\left[x_1(y_2^2w_2^2x_2^2)+x_2(y_1^2w_1^2x_1^2)\right]z^2(x_1y_1+x_2+y_2)`$
$`a_0`$ $`=`$ $`(y_1^2w_1^2x_1^2)(y_2^2w_2^2x_2^2)+z^2(x_1y_1)(x_2+y_2).`$ (8)
A smooth branch point coalescence can be easily arranged by $`PT`$symmetrically blowing up a triple root
$$\lambda _1=\lambda _2=\lambda _3=:\beta _{(c)}\lambda _{4(c)}$$
(9)
(”inflection point” configuration $`\mathrm{\Delta }(\lambda )=_\lambda \mathrm{\Delta }(\lambda )=_\lambda ^2\mathrm{\Delta }(\lambda )=0`$) of the quartic equation (7).
For our illustration purpose, it suffices to explicitly derive a suitable parametrization for an arbitrarily chosen triple root<sup>6</sup><sup>6</sup>6A detailed discriminant-based case analysis and complete classification of possible intersection scenarios for the four roots $`\lambda _k,k=1,\mathrm{},4`$ will be presented elsewhere. satisfying the additional branch point conditions (6) for the subsystems. Using these conditions together with (2), (9) in Newton’s identities (Vieta’s formulas)
$`\lambda _1+\lambda _2+\lambda _3+\lambda _4`$ $`=`$ $`a_3`$
$`\lambda _1\lambda _2+\lambda _1\lambda _3+\lambda _1\lambda _4+\lambda _2\lambda _3+\lambda _2\lambda _4+\lambda _3\lambda _4`$ $`=`$ $`a_2`$
$`\lambda _1\lambda _2\lambda _3+\lambda _2\lambda _3\lambda _4+\lambda _1\lambda _2\lambda _4+\lambda _1\lambda _3\lambda _4`$ $`=`$ $`a_1`$
$`\lambda _1\lambda _2\lambda _3\lambda _4`$ $`=`$ $`a_0`$ (10)
and setting for convenience $`\beta _{(c)}=z_{(c)}=1,x_{1(c)}=0`$ we easily obtain the remaining root $`\lambda _{4(c)}`$ and the three parameters $`x_{2(c)},y_{1,2(c)}`$ as
$`\lambda _{4(c)ϵ}`$ $`=`$ $`3+2^{3/2}ϵ>0,ϵ:=\pm 1,`$
$`x_{2(c)ϵ}`$ $`=`$ $`{\displaystyle \frac{1}{2}}(\lambda _{4(c)ϵ}+3)=3+2^{1/2}ϵ,`$
$`y_{1(c)ϵ,\delta }`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left[(3\lambda _{4(c)ϵ}+1)+\delta \sqrt{9\lambda _{4(c)ϵ}^2+2\lambda _{4(c)ϵ}+1}\right],\delta :=\pm 1,`$
$`y_{2(c)ϵ,\delta }`$ $`=`$ $`\lambda _{4(c)ϵ}1+{\displaystyle \frac{\delta }{2}}\sqrt{9\lambda _{4(c)ϵ}^2+2\lambda _{4(c)ϵ}+1}.`$ (11)
A computer algebraic test shows that the Jordan normal form of $`H`$ for this triple root configuration reads
$$H_{(c)}=SH_{(c)J}S^1,H_{(c)J}=\left(\begin{array}{cccc}1& 1& 0& 0\\ 0& 1& 1& 0\\ 0& 0& 1& 0\\ 0& 0& 0& \lambda _{4(c)ϵ}\end{array}\right),$$
(12)
so that the eigenvalue $`\beta _{(c)}=1`$ is indeed a type I triple point with geometric multiplicity one and algebraic multiplicity three. For the $`2D`$illustration of the coalescing branch points in Fig. 1 we have chosen a one-parameter matrix parametrization of the type $`H(t)=H_{(c)}+h(t),h(t0)0`$ providing the blowing-up of the triple root for $`t0`$ with
$`x_1=x_{1(c)},`$ $`x_2=x_{2(c)}+t,`$ (13)
$`y_1=y_{1(c)++}+2t^2`$ $`y_2=y_{2(c)++}t^3,`$
$`w_1=y_{1(c)++}t`$ $`w_2=y_{2(c)++}+3t^2.`$
A characteristic feature of the present coalescing branch point setup consists in the merging of two originally separated complex spectral sectors (of broken $`PT`$symmetry) into a larger single sector and the simultaneous disappearance of the real segment between the two square-root branch points. This is in obvious contrast to coalescing branch point setups in simpler pseudo-Hermitian $`2\times 2`$matrix models, where the branch points are connected by two real or complex branches which disappear upon branch point coalescence.
## 3 Third order spectral transitions of the $`\alpha ^2`$dynamo and the extended Squire setup
Based on the information of the previous section, it is easy to identify configurations with coalescing branch points in the spectra of the MHD $`\alpha ^2`$dynamo operator and of the operator of the extended Squire setup.
In case of the operator of the spherically symmetric $`\alpha ^2`$dynamo
$$\widehat{H}_l[\alpha ]=\left(\begin{array}{cc}Q[1]& \alpha \\ Q[\alpha ]& Q[1]\end{array}\right),Q[\alpha ]:=(_r+1/r)\alpha (r)(_r+1/r)+\alpha (r)\frac{l(l+1)}{r^2}$$
(14)
with $`\alpha `$profile<sup>7</sup><sup>7</sup>7The rather special numerical coefficients in the $`\alpha `$profile are adopted from the recently studied field-reversal scenario for $`\alpha ^2`$dynamos .
$`\alpha (r)`$ $`=`$ $`C[(21.465+2.467\zeta )+(426.412+167.928\zeta )r^2`$ (15)
$`(806.729+436.289\zeta )r^3+(392.276+272.991\zeta )r^4]`$
and physical (realistic) boundary conditions a triple point transition in the decay-mode sector $`(\mathrm{}\lambda <0)`$ is depicted in Fig. 2. Up to now it is still an open question whether such a transition in the sector of growing modes $`(\mathrm{}\lambda >0)`$ will have any physical significance, e.g., in dynamo experiments .
A different situation occurs in the case of the extended Squire setup of Ref.
$$\left[_y^2+gy^2(iy)^\nu \right]\psi (y)=E\psi (y),\psi (y=\pm b)=0$$
(16)
which corresponds to the Bender-Boettcher problem over a finite interval $`[b,b]`$. There a coalescence of two square-root branch points occurs close to the real-to-complex transition point in the associated Herbst-box model (Herbst model over a finite interval) located at the low-energy end of a web-like branch structure. An example for the corresponding triple point transition (which supplements the qualitative and graphical analyzes of Ref. ) is depicted in Fig. 3.
This work has been supported by DFG grant GE 682/12-2.
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# Corrections to Universal Fluctuations in Correlated Systems: the 2D XY-model
## Appendix: Normal ordered procedure
We introduce the normalized Gausssian measure with covariance $`TG`$
$$d\mu _{TG}(\varphi )=\frac{1}{Z}_0\mathrm{exp}\left\{\frac{\beta }{2}|_\mu \varphi _x|^2\right\}D\varphi $$
(43)
It has the characteristic function
$$𝑑\mu _{TG}(\varphi )e^{ij,\varphi }=e^{\frac{T}{2}j,Gj}.$$
(44)
If $`A(\varphi )`$ is a function of $`\varphi `$ which is integrable with respect to the Gaussian measure, it can be written in normal ordered form GlimmJaffe . The normal ordering operation is typically indicated by ::, but it depends on the covariance of the Gaussian measure.
$`A(\varphi )`$ $`=`$ $`:B(\varphi ):,`$ (45)
$`B(\varphi )`$ $`=`$ $`{\displaystyle 𝑑\mu _{TG}(\xi )A(\varphi +\xi )}.`$ (46)
In particular
$$e^{i\alpha \varphi _x}=e^{\frac{1}{2}\alpha ^2TG(0)}:e^{i\alpha \varphi _x}:$$
(47)
It follows that
$`\mathrm{cos}\varphi _x=e^{\frac{1}{2}TG(0)}:\mathrm{cos}\varphi _x:`$ (48)
$`{\displaystyle 𝑑\mu _{TG}(\varphi )}:\mathrm{cos}\varphi _x:=1.`$ (49)
Under general conditions, $`B`$ is a holomorphic function of $`\varphi `$ which can be expanded into a power series. In particular
$$:\mathrm{cos}\varphi _x:=1\frac{1}{2}:\varphi _x^2:+\frac{1}{4!}:\varphi _x^4:+\mathrm{}$$
(50)
It is well known that the normal products $`:\varphi _x^p:`$ are always eigenmodes of the linearized renormalization group at the Gaussian fix point determined by $`TG`$. In two dimensions, with fix point Hamiltonian $`\beta _{SW}=\frac{1}{2}\beta |_\mu \varphi |^2`$ they all have dimension $`0`$.
Consider now the partition function
$$\stackrel{~}{Z}(q)=e^{iqM}Z(q).$$
(51)
It follows from eq.(2) that the moments
$$M^p=\left(i\frac{}{q}\right)^p\stackrel{~}{Z}(q)|_{q=0}.$$
(52)
We restrict our attention to the spin wave approximation, replacing $``$ by $`_{SW}`$. Inserting eq.(48) into the definition (3) of $`Z(q)`$, we obtain
$$\stackrel{~}{Z}(q)=d\mu _{TG}(\varphi )\mathrm{exp}\{\frac{iq}{N}e^{\frac{T}{2}G(0)}\underset{x}{}:\mathrm{cos}\varphi _x:\}$$
(53)
From this and eq.(49) we recover the known result for the magnetization in spin wave approximation,
$$M=e^{\frac{T}{2}G(0)}$$
(54)
Therefore
$$\stackrel{~}{Z}(q)=\widehat{Z}(Mq)$$
(55)
with a new partition function
$$\widehat{Z}(q)=d\mu _{TG}(\varphi )\mathrm{exp}\{\frac{iq}{N}\underset{x}{}:\mathrm{cos}\varphi _x:\},$$
(56)
and
$$M^p/M^p=\left(i\frac{}{q}\right)^p\widehat{Z}(q)|_{q=0}$$
(57)
The partition function $`\widehat{Z}(q)`$ is exactly like $`\stackrel{~}{Z}(q)`$ except for the normal ordering of $`\mathrm{cos}\varphi _x`$. In a diagrammatic expansion, normal ordering eliminates tadpoles. The diagrams which contribute to $`M^p/M^p`$ are therefore exactly the diagrams without tadpoles which contribute to $`M^p`$.
It is more convenient to consider truncated expectation values
$$M^p^c/M^p=\left(i\frac{}{q}\right)^p\mathrm{ln}\widehat{Z}(q)|_{q=0}.$$
(58)
They have expansions in connected diagrams. The n-th order in the loop expansion for $`\mathrm{ln}\widehat{Z}`$ sums connected diagrams with $`n`$ loops. In 1-loop approximation only the term $`1\frac{1}{2}:\varphi _x^2:`$ in the expansion (50) of $`:\mathrm{cos}\varphi _x:`$ contributes, but the $`\frac{1}{4!}:\varphi _x^4:`$-term begins to contribute in 2-loop order, and so on.
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# Transfer Matrices for the Partition Function of the Potts Model on Toroidal Lattice Strips
## I Introduction
The $`q`$-state Potts model has served as a valuable model for the study of phase transitions and critical phenomena potts ; wurev . On a lattice, or, more generally, on a (connected) graph $`G`$, at temperature $`T`$, this model is defined by the partition function
$$Z(G,q,v)=\underset{\{\sigma _n\}}{}e^\beta $$
(1)
with the (zero-field) Hamiltonian
$$=J\underset{ij}{}\delta _{\sigma _i\sigma _j}$$
(2)
where $`\sigma _i=1,\mathrm{},q`$ are the spin variables on each vertex (site) $`iG`$; $`\beta =(k_BT)^1`$; and $`ij`$ denotes pairs of adjacent vertices. The graph $`G=G(V,E)`$ is defined by its vertex set $`V`$ and its edge set $`E`$; we denote the number of vertices of $`G`$ as $`n=n(G)=|V|`$ and the number of edges of $`G`$ as $`e(G)=|E|`$. We use the notation
$$K=\beta J,v=e^K1$$
(3)
so that the physical ranges are $`v0`$ for the Potts ferromagnet, and $`1v0`$ for Potts antiferromagnet, corresponding to $`0T\mathrm{}`$. One defines the (reduced) free energy per site $`f=\beta F`$, where $`F`$ is the actual free energy, via $`f(\{G\},q,v)=lim_n\mathrm{}\mathrm{ln}[Z(G,q,v)^{1/n}]`$, where we use the symbol $`\{G\}`$ to denote the formal limit $`lim_n\mathrm{}G`$ for a given family of graphs. In the present context, this $`n\mathrm{}`$ limit corresponds to the limit of infinite length for a strip graph of fixed width and some prescribed boundary conditions.
In this paper we shall present transfer matrices for the $`q`$-state Potts model partition functions $`Z(G,q,v)`$, for arbitrary $`q`$ and temperature variable $`v`$, on toroidal and Klein bottle strip graphs $`G`$ of the square, triangular, and honeycomb lattices of width $`L_y`$ vertices and of arbitrarily great length $`L_x`$ vertices. We label the lattice type as $`\mathrm{\Lambda }`$ and abbreviate the three respective types as $`sq`$, $`tri`$, and $`hc`$. Each strip involves a longitudinal repetition of $`m`$ copies of a particular subgraph. For the square-lattice strips, this is a column of squares. It is convenient to represent the strip of the triangular lattice as obtained from the corresponding strip of the square lattice with additional diagonal edges connecting, say, the upper-left to lower-right vertices in each square. In both these cases, the length is $`L_x=m`$ vertices. We represent the strip of the honeycomb lattice in the form of bricks oriented horizontally. In this case, since there are two vertices in 1-1 correspondence with each horizontal side of a brick, $`L_x=2m`$ vertices. Summarizing for all of three lattices, the relation between the number of vertices and the number of repeated copies is
$$L_x=\{\begin{array}{cc}m\hfill & \text{if }\mathrm{\Lambda }=sq\text{ or }tri\hfill \\ 2m\hfill & \text{if }\mathrm{\Lambda }=hc\hfill \end{array}.$$
(4)
For the toroidal case the partition function has the general form a ; bcc
$$Z(\mathrm{\Lambda },L_y\times L_x,tor.,q,v)=\underset{d=0}{\overset{L_y}{}}\underset{j}{}b_j^{(d)}(\lambda _{Z,\mathrm{\Lambda },L_y,d,j})^m$$
(5)
where $`\lambda _{Z,\mathrm{\Lambda },L_y,d,j}`$, are eigenvalues of the transfer matrices in the degree-$`d`$ subspace $`T_{Z,\mathrm{\Lambda },L_y,d}`$ and they are independent of the length of the strip length $`L_x`$. The size of the transfer matrices $`T_{Z,\mathrm{\Lambda },L_y,d}`$ is equal to $`d!n_{Z,tor}(L_y,d)`$ that will be discussed in the following section. Here because the dimensions of $`T_{Z,\mathrm{\Lambda },L_y,d}`$ are the same for all three lattices $`\mathrm{\Lambda }=sq,tri,hc`$, we omit $`\mathrm{\Lambda }`$ in the notation. In contrast to cyclic and Möbius strips, there can be more than one coefficient, denoted as $`b_j^{(d)}`$, for each degree $`d`$. The $`b_j^{(d)}`$’s are polynomials of degree $`d`$ in $`q`$ and play a role analogous to multiplicities of eigenvalues $`\lambda _{Z,\mathrm{\Lambda },L_y,d,j}`$, although this identification is formal, since $`b_j^{(d)}`$ may be zero or negative for the small physical values of $`q`$. We exhibit our method for calculating $`T_{Z,\mathrm{\Lambda },L_y,d}`$ for arbitrary $`L_y`$, and we shall construct an analogous formula for Klein bottle strips. Explicit results for arbitrary $`L_y`$ are given for (i) $`T_{Z,\mathrm{\Lambda },L_y,d}`$ with $`d=L_y`$, (ii) the determinant $`det(T_{Z,\mathrm{\Lambda },L_y,d})`$, and (iii) the trace $`Tr(T_{Z,\mathrm{\Lambda },L_y})`$. The results in (i) and (iii) apply for $`\mathrm{\Lambda }=sq,tri,hc`$, while the results in (ii) apply for $`\mathrm{\Lambda }=sq,hc`$. Corresponding results are given for the equivalent Tutte polynomials for these lattice strips and illustrative examples are included. We have calculated the transfer matrices up to widths $`L_y=4`$ for the square and honeycomb lattices and up to $`L_y=3`$ for the triangular lattice. Since the dimensions of these matrices increase rapidly with strip width (e.g., $`dim(T_{Z,\mathrm{\Lambda },L_y,d})=14,35,56,48`$ for $`L_y=4`$ and $`0d3`$), it is not feasible to present many of the explicit results here; instead, we concentrate on general methods and results that hold for arbitrary $`L_y`$.
Various special cases of the Potts model partition function $`Z(G,q,v)`$ are of interest. For example, if one considers the case of antiferromagnetic spin-spin coupling, $`J<0`$ and takes the temperature to zero, so that $`K=\mathrm{}`$ and $`v=1`$, then
$$Z(G,q,1)=P(G,q)$$
(6)
where $`P(G,q)`$ is the chromatic polynomial (in $`q`$) expressing the number of ways of coloring the vertices of the graph $`G`$ with $`q`$ colors such that no two adjacent vertices have the same color bbook ; rtrev . The minimum number of colors necessary for such a coloring of $`G`$ is called the chromatic number, $`\chi (G)`$. The $`q`$-state Potts antiferromagnet (AF) potts ; wurev exhibits nonzero ground state entropy, $`S_0>0`$ (without frustration) for sufficiently large $`q`$ on a given lattice $`\mathrm{\Lambda }`$ or, more generally, on a graph $`G=(V,E)`$. This is equivalent to a ground state degeneracy per site $`W>1`$, since $`S_0=k_B\mathrm{ln}W`$. Thus
$$W(\{G\},q)=\underset{n\mathrm{}}{lim}P(G,q)^{1/n}.$$
(7)
For a strip graph $`G_s`$ of a lattice $`\mathrm{\Lambda }`$ with given boundary conditions, following our earlier notation a we denote the sum of coefficients (generalized multiplicities) $`c_{Z,G,j}`$ of the $`\lambda _{Z,G,j}`$ as
$$C_{Z,G}=\underset{j=1}{\overset{N_{Z,G,\lambda }}{}}c_{Z,G,j}$$
(8)
while for the chromatic polynomial as
$$C_{P,G}=\underset{j=1}{\overset{N_{P,G,\lambda }}{}}c_{P,G,j}.$$
(9)
These sums are independent of the length $`m`$ of the strip. General results for these sums will be given below for the strips of interest.
We recall some previous related work. The partition function $`Z(G,q,v)`$ for the Potts model was calculated for arbitrary $`q`$ and $`v`$ on strips of the lattice $`\mathrm{\Lambda }`$ with toroidal or Klein bottle longitudinal boundary conditions was calculated for (i) $`\mathrm{\Lambda }=sq`$, $`L_y=2,3`$ in s3a , and (ii) for $`L_y=3`$ on the square-lattice with next-nearest-neighbor spin-spin couplings in ka3 . Matrix methods for calculating chromatic polynomials were developed and used in bds ; bm ; baxter86 ; baxter87 and more recently in matmeth ; matmeth2 ; matmeth3 . Ref. sqtran ; cyltran ; tritran developed transfer matrix methods for both $`Z(G,q,v)`$ and the special case $`v=1`$ of chromatic polynomials on strips of the square and triangular lattices with free longitudinal boundary conditions and used them to calculate the latter polynomials for a large variety of widths. These have been termed transfer matrices in the Fortuin-Kasteleyn representation (see eq. (10) below). These methods were applied to calculate the full Potts model partition function for strips of the square and triangular lattices with free boundary conditions and a number of widths in Refs. ts ; tt , and then were generalized for strips of the square, triangular and honeycomb lattices with cyclic and Möbious strips zt ; pt . Here we extend these transfer matrix methods for calculating Potts model partition functions on toroidal and Klein bottle lattice strips and present general results for these lattice strips of arbitrary width. Clearly, new exact results on the Potts model are of value in their own right.
Let $`G^{}=(V,E^{})`$ be a spanning subgraph of $`G`$, i.e. a subgraph having the same vertex set $`V`$ and an edge set $`E^{}E`$. $`Z(G,q,v)`$ can be written as the sum kf ; fk
$$Z(G,q,v)=\underset{G^{}G}{}q^{k(G^{})}v^{|E^{}|}$$
(10)
where $`k(G^{})`$ denotes the number of connected components of $`G^{}`$. Since we only consider connected graphs $`G`$, we have $`k(G)=1`$. The formula (10) enables one to generalize $`q`$ from $`_+`$ to $`_+`$ for physical ferromagnetic $`v`$. More generally, eq. (10) allows one to generalize both $`q`$ and $`v`$ to complex values, as is necessary when studying zeros of the partition function in the complex $`q`$ and $`v`$ planes.
The Potts model partition function $`Z(G,q,v)`$ is equivalent to an object of considerable current interest in mathematical graph theory, the Tutte polynomial, $`T(G,x,y)`$, given by tutte1 ; tutte2 ; tutte3
$$T(G,x,y)=\underset{G^{}G}{}(x1)^{k(G^{})k(G)}(y1)^{c(G^{})}$$
(11)
where $`c(G^{})=|E^{}|+k(G^{})|V|`$ is the number of independent circuits in $`G^{}`$. Now let
$$x=1+\frac{q}{v},y=v+1$$
(12)
so that
$$q=(x1)(y1).$$
(13)
Then the equivalence between the Potts model partition function and the Tutte polynomial for a graph $`G`$ is
$$Z(G,q,v)=(x1)^{k(G)}(y1)^{n(G)}T(G,x,y).$$
(14)
Given this equivalence, we can express results either in Potts or Tutte form. We will use both, since each has its own particular advantages. The Potts model form, involving the variables $`q`$ and $`v`$ is convenient for physical applications, since $`q`$ specifies the number of states and determines the universality class of the transition, and $`v`$ is the temperature variable. The Tutte form has the advantage that many expressions are simpler when written in terms of the Tutte variables $`x`$ and $`y`$.
From (14), one can write the Tutte polynomial as
$$T(\mathrm{\Lambda },L_y\times L_x,tor.,x,y)=\frac{1}{x1}\underset{d=0}{\overset{L_y}{}}\underset{j}{}b_j^{(d)}(\lambda _{T,\mathrm{\Lambda },L_y,d,j})^m$$
(15)
where $`m`$ is given in terms of $`L_x`$ by eq. (4) and it is convenient to factor out a factor of $`1/(x1)`$. (This factor is always cancelled, since the Tutte polynomial is a polynomial in $`x`$ as well as $`y`$.)
From eq. (14) it follows that a ; hca
$$\lambda _{Z,\mathrm{\Lambda },L_y,d,j}=v^{pL_y}\lambda _{T,\mathrm{\Lambda },L_y,d,j}$$
(16)
and
$$T_{Z,\mathrm{\Lambda },L_y,d}=v^{pL_y}T_{T,\mathrm{\Lambda },L_y,d}$$
(17)
where
$$p=\{\begin{array}{cc}1\hfill & \text{if }\mathrm{\Lambda }=sq\text{ or }tri\hfill \\ 2\hfill & \text{if }\mathrm{\Lambda }=hc\hfill \end{array}.$$
(18)
Note that the factor of $`(x1)`$ in eq. (14) cancels the factor $`1/(x1)`$ in eq. (15).
## II Transfer Matrix Method
The chromatic polynomials for strips with periodic boundary condition in the longitudinal direction were discussed in Refs. matmeth ; matmeth2 ; matmeth3 in terms of a compatibility matrix, and the bases of the matrix are given by $`[X|\xi ]`$, where $`X`$ is a subset of $`L_y`$ vertices and $`\xi `$ is an injection from $`X`$ to $`\{1,2,\mathrm{},q\}`$. We now use the same bases for the transfer matrix of the full Potts model partition function. The degree $`d=0`$ subspace of the transfer matrix (equivalently called the “level 0” subspace in matmeth3 ) corresponds to the empty set $`X=\mathrm{}`$, and the bases of the transfer matrix are all of the possible non-crossing partitions of $`L_y`$ vertices. (For the zero-temperature antiferromagnetic Potts model, a non-nearest-neighbor requirement is imposed; this will be discussed further in the next section.) The dimension of this matrix is $`n_{Z,tor}(L_y,0)=C_{L_y}`$, the Catalan number sqtran , as in eq. (97) below. The eigenvalues of the transfer matrix for a cylindrical strip are a subset of the eigenvalues of the transfer matrix $`T_{Z,\mathrm{\Lambda },L_y,d=0}`$ in this degree $`d=0`$ subspace for the corresponding toroidal strip. The degree $`d=1`$ subspace of the transfer matrix is given by all of the possible non-crossing partitions with a color assignment to one vertex (with possible connections with other vertices); i.e., $`X`$ contains one vertex, and the multiplicity is $`q1b^{(1)}`$. This follows because there are $`q`$ possible ways of making this color assignment, but one of these has to be subtracted, since the effect of all the possible color assignments is equivalent to the choice of no specific color assignment, which has been taken into account in the level 0 subspace. In this derivation and subsequent ones we assume that $`q`$ is a sufficiently large integer to begin with, so that the multiplicities are positive-definite; we then analytically continue them downward to apply in the region of small $`q`$ where the coefficients can be zero or negative. For the next subspace we consider all of the non-crossing partitions with two-color assignments to two separated vertices (with possible connections with other vertices). Now the multiplicity can be understood by the sieve formula of matmeth ; matmeth2 ; matmeth3 . Since the two assigned colors should be different, there are $`q(q1)`$ ways of making these assignments. This includes the $`q`$ possible color assignments for each of the two vertices that have been considered in level 1 and hence these must be subtracted. In doing this, the no-color assignment was subtracted twice, and one of these has to be added back. Therefore, the multiplicity is
$$q(q1)2q+1=q^23q+1b^{(2)}.$$
(19)
Recall that a slice of cyclic strips is a tree graph, but a slice of toroidal strips is a circuit graph. For toroidal strips, these coefficients $`b^{(d)}`$ with $`0d2`$ are the same as those for cyclic strips $`c^{(d)}`$ in Ref. cf (see also saleur ) because when the number of vertices is 1 or 2 a tree graph is equivalent to a circuit graph modulo multiple edges. For toroidal strips, the multiplicity for three-color assignment is
$$q(q1)(q2)3q(q1)+3q1=q^36q^2+8q1b^{(3)}.$$
(20)
By similar reasoning, the multiplicity for four-color assignment is
$`q(q1)(q^23q+3)4q(q1)^2+4q(q1)+2q^24q+1`$ (21)
$`=`$ $`q^48q^3+20q^215q+1b^{(4)}.`$ (23)
Further, we find
$`b^{(5)}`$ $`=`$ $`q^510q^4+35q^350q^2+24q1`$ (24)
$`b^{(6)}`$ $`=`$ $`q^612q^5+54q^4112q^3+105q^235q+1`$ (26)
$`b^{(7)}`$ $`=`$ $`q^714q^6+77q^5210q^4+294q^3196q^2+48q1`$ (28)
$`b^{(8)}`$ $`=`$ $`(q^23q+1)(q^613q^5+64q^4147q^3+155q^260q+1)`$ (30)
$`b^{(9)}`$ $`=`$ $`q^918q^8+135q^7546q^6+1287q^51782q^4+1386q^3540q^2+80q1`$ (32)
$`b^{(10)}`$ $`=`$ $`q^{10}20q^9+170q^8800q^7+2275q^64004q^5+4290q^42640q^3+825q^299q+1`$ (34)
and so forth for higher values of $`d`$. We show the calculation for these multiplicities pictorially for $`2d4`$ in Fig. 1. Let us define
$$q=2+t+t^1=2+2\mathrm{cos}\theta =4\mathrm{cos}^2(\theta /2).$$
(37)
Then, in terms of these variables, we find that $`b^{(d)}`$ can be written as
$$b^{(0)}=1$$
(38)
$$b^{(1)}=t+1+t^1=1+2\mathrm{cos}\theta $$
(39)
and, for $`d2`$,
$$b^{(d)}=(1)^d(t+1+t^1)+t^d+t^d=(1)^d(1+2\mathrm{cos}\theta )+2\mathrm{cos}(d\theta )$$
(40)
Therefore, we have the recursion relation
$$b^{(d+2)}+(1)^{d+1}b^{(1)}=(t+\frac{1}{t})(b^{(d+1)}+(1)^db^{(1)})(b^{(d)}+(1)^{d1}b^{(1)})\text{for}d2,$$
(41)
that is
$$b^{(d+2)}=(q2)b^{(d+1)}b^{(d)}+q(1)^db^{(1)}\text{for}d2.$$
(42)
Now the coefficient $`c^{(d)}`$ in cf is given by
$$c^{(d)}=U_{2d}(\frac{\sqrt{q}}{2})=\underset{j=0}{\overset{d}{}}(1)^j\left(\genfrac{}{}{0pt}{}{2dj}{j}\right)q^{dj}$$
(43)
where $`U_n(x)`$ is the Chebyshev polynomial of the second kind. Equivalently, in terms of the angle $`\theta `$ in eq. (37) saleur ; cf ,
$$c^{(d)}=\frac{\mathrm{sin}\left((2d+1)\frac{\theta }{2}\right)}{\mathrm{sin}\left(\frac{\theta }{2}\right)}$$
(44)
In general, we find that the relation between $`b^{(d)}`$ and $`c^{(d)}`$ is as follows,
$`c^{(d)}`$ $`=`$ $`b^{(d)}\text{for}0d2`$ (45)
$`c^{(d)}`$ $`=`$ $`{\displaystyle \underset{d^{}=1}{\overset{d}{}}}b^{(d^{})}\text{for odd}d`$ (47)
$`c^{(d)}`$ $`=`$ $`{\displaystyle \underset{d^{}=2}{\overset{d}{}}}b^{(d^{})}\text{for even}d,`$ (49)
and
$`b^{(d)}`$ $`=`$ $`c^{(d)}c^{(d1)}+(1)^dc^{(1)}`$ (50)
$`=`$ $`U_{2d}({\displaystyle \frac{\sqrt{q}}{2}})U_{2d2}({\displaystyle \frac{\sqrt{q}}{2}})+(1)^d(q1)`$ (52)
$`=`$ $`q^d+{\displaystyle \underset{j=1}{\overset{d}{}}}(1)^j{\displaystyle \frac{2d}{2dj}}\left({\displaystyle \genfrac{}{}{0pt}{}{2dj}{j}}\right)q^{dj}+(1)^d(q1)\text{for}d2`$ (54)
At the value $`q=0`$, $`b^{(d)}`$ has the property that
$$b^{(d)}(q=0)=(1)^d.$$
(55)
We also find the following properties of the $`b^{(d)}`$ polynomials. For $`q=1`$ and $`d2`$, $`b^{(d)}`$ is equal to 2 if $`d=0`$ mod 3 and $`1`$ if $`d=1`$ or $`d=2`$ mod 3. For $`q=2`$ and $`d2`$, $`b^{(d)}`$ is equal to 3 if $`d=0`$ mod 4 and $`1`$ if $`d=1,2,3`$ mod 4. One can derive similar relations for other values of $`q`$. For $`q=3`$, $`b^{(d)}`$ takes on the pattern of values $`4,1,1,4,1,1`$ for $`d=3,4,5,6,7,8`$, and then this pattern repeats for higher integral values of $`d`$. For $`q=4`$ and $`d2`$, $`b^{(d)}=5`$ if $`d`$ is even and $`b^{(d)}=1`$ if $`d`$ is odd. For larger values of $`q`$, $`b^{(d)}`$ does not cycle through a fixed pattern of values as $`d`$ increases, but instead increases. The fixed patterns of values of $`b^{(d)}`$ for integer $`q`$ from 0 to 4 are related to special sum rules for Potts model partition functions on lattice strips with toroidal boundary conditions, similar to those that we exhibited in Refs. ka3 and s3a .
The partitions for toroidal strips are more involved than those for cyclic strips because of the rotational symmetry of the slice of toroidal strips. We list graphically all the possible partitions for strips of widths $`L_y=2`$ and $`L_y=3`$ in Figs. 2 and 3, respectively, where white circles are the original $`L_y`$ vertices and each black circle corresponds to a specific color assignment. In the following discussion, we will simply use the names white and black circles with these meanings understood implicitly. The connections between the black circles and white circles should obey the non-crossing restriction. The apparent crossings in some partitions in Figs. 2 and 3 do not violate this rule when white circles are rotated appropriately, owing to the rotational symmetry for toroidal strips. For example, the crossing in the last partition of $`d=1`$ in Fig. 3 can be removed by moving the bottom white circle to the top so that it becomes equivalent to the fourth partition of $`d=1`$, or moving the top white circle to the bottom so that it becomes equivalent to the eighth partition of $`d=1`$ in the figure. Another example is the crossing in the eleventh partition of $`d=2`$ in Fig. 3 that can be removed by the rotational symmetry of the white circles. In addition, it is necessary to include all the arrangements of black circles when the transfer matrix is constructed. For example, the two partitions of $`d=2`$ in Fig. 2 correspond to exchanging the black circles, i.e. the assignments of the colors. We arrange the partitions of $`d=2`$ in Fig. 3 into pairs such that every two partitions are equivalent when the two black circles are exchanged. We show the six partitions of $`d=3`$ in Fig. 3 which correspond to all the possible arrangements of black circles.
We denote the partitions $`𝒫_{L_y,d}`$ for $`2L_y4`$ as follows:
$$𝒫_{2,0}=\{I;12\},𝒫_{2,1}=\{\overline{2};\overline{1};\overline{12}\},𝒫_{2,2}=\{\overline{1},\widehat{2};\widehat{1},\overline{2}\}$$
(56)
$`𝒫_{3,0}`$ $`=`$ $`\{I;12;13;23;123\},𝒫_{3,1}=\{\overline{3};\overline{2};\overline{1};12,\overline{3};\overline{12};\overline{13};\overline{23};23,\overline{1};\overline{123};13,\overline{2}\},`$ (57)
$`𝒫_{3,2}`$ $`=`$ $`\{\overline{2},\widehat{3};\widehat{2},\overline{3};\overline{1},\widehat{3};\widehat{1},\overline{3};\overline{1},\widehat{2};\widehat{1},\overline{2};\overline{12},\widehat{3};\widehat{12},\overline{3};\overline{1},\widehat{23};\widehat{1},\overline{23};\overline{13},\widehat{2};\widehat{13},\overline{2}\},`$ (59)
$`𝒫_{3,3}`$ $`=`$ $`\{\overline{1},\widehat{2},\stackrel{~}{3};\widehat{1},\overline{2},\stackrel{~}{3};\overline{1},\stackrel{~}{2},\widehat{3};\widehat{1},\stackrel{~}{2},\overline{3};\stackrel{~}{1},\overline{2},\widehat{3};\stackrel{~}{1},\widehat{2},\overline{3}\}`$ (61)
$`𝒫_{4,0}`$ $`=`$ $`\{I;12;13;14;23;24;34;12,34;14,23;123;124;134;234;1234\},`$ (62)
$`𝒫_{4,1}`$ $`=`$ $`\{\overline{4};\overline{3};\overline{2};\overline{1};12,\overline{4};12,\overline{3};\overline{12};13,\overline{4};\overline{13};\overline{14};23,\overline{4};\overline{23};23,\overline{1};\overline{24};24,\overline{1};\overline{34};34,\overline{2};34,\overline{1};12,\overline{34};`$ (68)
$`34,\overline{12};23,\overline{14};123,\overline{4};\overline{123};\overline{124};\overline{134};\overline{234};234,\overline{1};\overline{1234};13,\overline{2};14,\overline{3};14,\overline{2};24,\overline{3};14,\overline{23};`$
$`124,\overline{3};134,\overline{2}\},`$
$`𝒫_{4,2}`$ $`=`$ $`\{\overline{3},\widehat{4};\widehat{3},\overline{4};\overline{2},\widehat{4};\widehat{2},\overline{4};\overline{1},\widehat{4};\widehat{1},\overline{4};\overline{2},\widehat{3};\widehat{2},\overline{3};\overline{1},\widehat{3};\widehat{1},\overline{3};\overline{1},\widehat{2};\widehat{1},\overline{2};12,\overline{3},\widehat{4};12,\widehat{3},\overline{4};\overline{12},\widehat{4};\widehat{12},\overline{4};`$ (78)
$`\overline{12},\widehat{3};\widehat{12},\overline{3};\overline{13},\widehat{4};\widehat{13},\overline{4};\overline{13},\widehat{2};\widehat{13},\overline{2};\overline{14},\widehat{3};\widehat{14},\overline{3};\overline{14},\widehat{2};\widehat{14},\overline{2};14,\overline{2},\widehat{3};14,\widehat{2},\overline{3};\overline{23},\widehat{4};\widehat{23},\overline{4};`$
$`23,\overline{1},\widehat{4};23,\widehat{1},\overline{4};\overline{1},\widehat{23};\widehat{1},\overline{23};\overline{1},\widehat{24};\widehat{1},\overline{24};\overline{24},\widehat{3};\widehat{24},\overline{3};\overline{2},\widehat{34};\widehat{2},\overline{34};\overline{1},\widehat{34};\widehat{1},\overline{34};34,\overline{1},\widehat{2};`$
$`34,\widehat{1},\overline{2};\overline{12},\widehat{34};\widehat{12},\overline{34};\overline{14},\widehat{23};\widehat{14},\overline{23};\overline{123},\widehat{4};\widehat{123},\overline{4};\overline{124},\widehat{3};\widehat{124},\overline{3};\overline{134},\widehat{2};\widehat{134},\overline{2};\overline{1},\widehat{234};`$
$`\widehat{1},\overline{234}\},`$
$`𝒫_{4,3}`$ $`=`$ $`\{\overline{2},\widehat{3},\stackrel{~}{4};\widehat{2},\overline{3},\stackrel{~}{4};\overline{2},\stackrel{~}{3},\widehat{4};\widehat{2},\stackrel{~}{3},\overline{4};\stackrel{~}{2},\overline{3},\widehat{4};\stackrel{~}{2},\widehat{3},\overline{4};\overline{1},\widehat{3},\stackrel{~}{4};\widehat{1},\overline{3},\stackrel{~}{4};\overline{1},\stackrel{~}{3},\widehat{4};\widehat{1},\stackrel{~}{3},\overline{4};\stackrel{~}{1},\overline{3},\widehat{4};\stackrel{~}{1},\widehat{3},\overline{4};`$ (88)
$`\overline{1},\widehat{2},\stackrel{~}{4};\widehat{1},\overline{2},\stackrel{~}{4};\overline{1},\stackrel{~}{2},\widehat{4};\widehat{1},\stackrel{~}{2},\overline{4};\stackrel{~}{1},\overline{2},\widehat{4};\stackrel{~}{1},\widehat{2},\overline{4};\overline{1},\widehat{2},\stackrel{~}{3};\widehat{1},\overline{2},\stackrel{~}{3};\overline{1},\stackrel{~}{2},\widehat{3};\widehat{1},\stackrel{~}{2},\overline{3};\stackrel{~}{1},\overline{2},\widehat{3};\stackrel{~}{1},\widehat{2},\overline{3};`$
$`\overline{12},\widehat{3},\stackrel{~}{4};\widehat{12},\overline{3},\stackrel{~}{4};\overline{12},\stackrel{~}{3},\widehat{4};\widehat{12},\stackrel{~}{3},\overline{4};\stackrel{~}{12},\overline{3},\widehat{4};\stackrel{~}{12},\widehat{3},\overline{4};\overline{14},\widehat{2},\stackrel{~}{3};\widehat{14},\overline{2},\stackrel{~}{3};\overline{14},\stackrel{~}{2},\widehat{3};\widehat{14},\stackrel{~}{2},\overline{3};`$
$`\stackrel{~}{14},\overline{2},\widehat{3};\stackrel{~}{14},\widehat{2},\overline{3};\overline{1},\widehat{23},\stackrel{~}{4};\widehat{1},\overline{23},\stackrel{~}{4};\overline{1},\stackrel{~}{23},\widehat{4};\widehat{1},\stackrel{~}{23},\overline{4};\stackrel{~}{1},\overline{23},\widehat{4};\stackrel{~}{1},\widehat{23},\overline{4};\overline{1},\widehat{2},\stackrel{~}{34};\widehat{1},\overline{2},\stackrel{~}{34};`$
$`\overline{1},\stackrel{~}{2},\widehat{34};\widehat{1},\stackrel{~}{2},\overline{34};\stackrel{~}{1},\overline{2},\widehat{34};\stackrel{~}{1},\widehat{2},\overline{34}\},`$
$`𝒫_{4,4}`$ $`=`$ $`\{\overline{1},\widehat{2},\stackrel{ˇ}{3},\stackrel{~}{4};\widehat{1},\overline{2},\stackrel{ˇ}{3},\stackrel{~}{4};\overline{1},\stackrel{ˇ}{2},\widehat{3},\stackrel{~}{4};\widehat{1},\stackrel{ˇ}{2},\overline{3},\stackrel{~}{4};\stackrel{ˇ}{1},\overline{2},\widehat{3},\stackrel{~}{4};\stackrel{ˇ}{1},\widehat{2},\overline{3},\stackrel{~}{4};\overline{1},\widehat{2},\stackrel{~}{3},\stackrel{ˇ}{4};\widehat{1},\overline{2},\stackrel{~}{3},\stackrel{ˇ}{4};\overline{1},\stackrel{ˇ}{2},\stackrel{~}{3},\widehat{4};`$ (94)
$`\widehat{1},\stackrel{ˇ}{2},\stackrel{~}{3},\overline{4};\stackrel{ˇ}{1},\overline{2},\stackrel{~}{3},\widehat{4};\stackrel{ˇ}{1},\widehat{2},\stackrel{~}{3},\overline{4};\overline{1},\stackrel{~}{2},\widehat{3},\stackrel{ˇ}{4};\widehat{1},\stackrel{~}{2},\overline{3},\stackrel{ˇ}{4};\overline{1},\stackrel{~}{2},\stackrel{ˇ}{3},\widehat{4};\widehat{1},\stackrel{~}{2},\stackrel{ˇ}{3},\overline{4};\stackrel{ˇ}{1},\stackrel{~}{2},\overline{3},\widehat{4};\stackrel{ˇ}{1},\stackrel{~}{2},\widehat{3},\overline{4};`$
$`\stackrel{~}{1},\overline{2},\widehat{3},\stackrel{ˇ}{4};\stackrel{~}{1},\widehat{2},\overline{3},\stackrel{ˇ}{4};\stackrel{~}{1},\overline{2},\stackrel{ˇ}{3},\widehat{4};\stackrel{~}{1},\widehat{2},\stackrel{ˇ}{3},\overline{4};\stackrel{~}{1},\stackrel{ˇ}{2},\overline{3},\widehat{4};\stackrel{~}{1},\stackrel{ˇ}{2},\widehat{3},\overline{4}\}`$
where partitions are separated by a colon, and overline, hat, check and tilde in each partition correspond to different color assignments. It follows that the size of the transfer matrix $`T_{Z,\mathrm{\Lambda },L_y,d}`$ always has the factor $`d!`$. Let us denote the reduced size of the transfer matrix as $`n_{Z,tor}(L_y,d)`$, i.e., we neglect the permutation of black circles. As the derivation of $`b^{(d)}`$ does not take into account permutations, i.e. a set of $`d!`$ eigenvalues should have their coefficients summed to be equal to $`b^{(d)}`$, we find that the sum of all coefficients for a toroidal strip graph, which is equal to the dimension of the total transfer matrix, i.e., the number of ways to color $`L_y`$ vertices, is
$$C_{Z,\mathrm{\Lambda },L_y}=dim(T_{Z,\mathrm{\Lambda },L_y})=\underset{d=0}{\overset{L_y}{}}n_{Z,tor}(L_y,d)b^{(d)}=q^{L_y}\mathrm{for}\mathrm{\Lambda }=sq,tri,hc.$$
(95)
By substituting the expression for $`b^{(d)}`$ in eq. (95), we determine the $`n_{Z,tor}(L_y,d)`$, as follows:
$$n_{Z,tor}(L_y,d)=0\text{for}d>L_y$$
(96)
$$n_{Z,tor}(L_y,0)=C_{L_y}$$
(97)
$$n_{Z,tor}(L_y,1)=\left(\genfrac{}{}{0pt}{}{2L_y1}{L_y1}\right)$$
(98)
$$n_{Z,tor}(L_y,d)=\left(\genfrac{}{}{0pt}{}{2L_y}{L_yd}\right)\text{for}2dL_y$$
(99)
where $`C_n`$ is the Catalan number which occurs in combinatorics and is defined by
$$C_n=\frac{1}{(n+1)}\left(\genfrac{}{}{0pt}{}{2n}{n}\right).$$
(100)
The first few Catalan numbers are $`C_1=1`$, $`C_2=2`$, $`C_3=5`$, $`C_4=14`$, and $`C_5=42`$. Special cases of $`n_{Z,tor}(L_y,d)`$ that are of interest here include
$$n_{Z,tor}(L_y,L_y)=1$$
(101)
$$n_{Z,tor}(L_y,L_y1)=2L_y\text{for}L_y3$$
(102)
In Table 1 we list the first few numbers $`n_{Z,tor}(L_y,d)`$. We notice the following recursion relations:
$`n_{Z,tor}(L_y,1)`$ $`=`$ $`n_{Z,tor}(L_y,0)n_{Z,tor}(L_y1,0)+2n_{Z,tor}(L_y1,1)+n_{Z,tor}(L_y1,2)`$ (105)
$`\text{for}L_y2`$
$`n_{Z,tor}(L_y,2)`$ $`=`$ $`n_{Z,tor}(L_y1,0)+2n_{Z,tor}(L_y1,1)+2n_{Z,tor}(L_y1,2)+n_{Z,tor}(L_y1,3)`$ (109)
$`\text{for}L_y2`$
$`n_{Z,tor}(L_y,d)`$ $`=`$ $`n_{Z,tor}(L_y1,d1)+2n_{Z,tor}(L_y1,d)+n_{Z,tor}(L_y1,d+1)`$ (113)
$`\text{for}L_y3,3dL_y.`$
The last line here is the same as the recursion relation for $`n_{Z,cyc}(L_y,d)`$ when $`1dL_y`$ for cyclic strips (which was denoted as $`n_Z(L_y,d)`$ in cf ). By substituting $`q=0`$ into eq. (95) and using eq. (55), we have the relation
$$\underset{d=0}{\overset{L_y}{}}n_{Z,tor}(L_y,d)(1)^d=0\mathrm{for}\mathrm{\Lambda }=sq,tri,hc.$$
(114)
Another way to realize the number $`n_{Z,tor}(L_y,d)`$ is as follows. For $`d=0`$, i.e. no color assignment, the partitions are identical to those for cyclic strips, and the number of these partitions is $`n_{Z,tor}(L_y,0)=n_{Z,cyc}(L_y,0)=C_{L_y}`$ given by eq. (97). Recall that the number of partitions of size $`n`$ with $`k`$ components is given by the Narayana number $`N(n,k)`$, which is sequence A001263 in sl (see also Ref. Flajolet99 and the references therein), given by
$$N(n,k)=\frac{1}{k}\left(\genfrac{}{}{0pt}{}{n1}{k1}\right)\left(\genfrac{}{}{0pt}{}{n}{k1}\right)=\frac{1}{n}\left(\genfrac{}{}{0pt}{}{n}{k}\right)\left(\genfrac{}{}{0pt}{}{n}{k1}\right).$$
(115)
The first few rows of the triangle of Narayana numbers, or Catalan triangle, are shown in Fig. 4. It is clear that the sum of terms in each row in the triangle is equal to the corresponding Catalan number:
$$\underset{k=1}{\overset{L_y}{}}N(L_y,k)=\frac{1}{L_y}\underset{k=0}{\overset{L_y1}{}}\left(\genfrac{}{}{0pt}{}{L_y}{k+1}\right)\left(\genfrac{}{}{0pt}{}{L_y}{k}\right)=\frac{1}{L_y}\left(\genfrac{}{}{0pt}{}{2L_y}{L_y1}\right)=C_{L_y}=n_{Z,tor}(L_y,0).$$
(116)
For $`d=1`$, we need to assign a color for each component of the level 0 partitions. Therefore, we have
$`n_{Z,tor}(L_y,1)`$ $`=`$ $`{\displaystyle \underset{k=1}{\overset{L_y}{}}}kN(L_y,k)={\displaystyle \underset{k=1}{\overset{L_y}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{L_y1}{k1}}\right)\left({\displaystyle \genfrac{}{}{0pt}{}{L_y}{k1}}\right)=\left({\displaystyle \genfrac{}{}{0pt}{}{2L_y1}{L_y1}}\right)`$ (117)
$`=`$ $`{\displaystyle \frac{1}{2}}\left({\displaystyle \genfrac{}{}{0pt}{}{2L_y}{L_y}}\right)={\displaystyle \frac{L_y+1}{2}}C_{L_y}={\displaystyle \frac{L_y+1}{2}}n_{Z,tor}(L_y,0).`$ (119)
Let us compare $`n_{Z,tor}(L_y,d)`$ with the corresponding number $`n_{Z,cyc}(L_y,d)`$ given in Table 3 of cf . In general, $`n_{Z,tor}(L_y,d)`$ is larger than $`n_{Z,cyc}(L_y,d)`$, except for $`d=0`$ where they are the same. This is a consequence of the fact that certain forbidden partitions for cyclic strips for $`d1`$ are allowed for toroidal strips with rotational translation. We list the first few $`n_{Z,tor}(L_y,d)n_{Z,cyc}(L_y,d)`$ in Table 2. Comparing Table 2 with Table 1, we infer the relation
$$n_{Z,tor}(L_y,d)n_{Z,cyc}(L_y,d)=n_{Z,tor}(L_y,d+1)\text{for}d2.$$
(120)
This can be understood as follows. The extra partitions for toroidal strips, as contrast with cyclic strips, are those with removable crossings. That is, for any connection between two non-adjacent white circles, all the color assignments should be inside or outside the connection while these two white circles can be either color-assigned or not. Therefore, there are three possible correspondences from the partitions for $`n_{Z,tor}(L_y,d+1)`$ to the partitions for $`n_{Z,tor}(L_y,d)n_{Z,cyc}(L_y,d)`$ with $`d2`$: (i) if the top color-assigned vertex and the bottom color-assigned vertex are not within a connection, then connect these two vertices (so that only one color is assigned to them); (ii) if all color-assigned vertices are within a connection and the connection is also assigned a color, then remove the assignment for the connection; (iii) if all color-assigned vertices are within a connection and the connection is not assigned a color, then connect the top and bottom color-assigned vertices (within that connection). These account for all of the extra partitions that toroidal strips have, relative to cyclic strips. As an example, consider the case of width $`L_y=4`$ and degree $`d=2`$ in eq. (94). For this case there are eight partitions that are allowed for toroidal strips but not for cyclic strips. These are $`\overline{13},\overline{2}`$; $`\overline{14},\overline{3}`$; $`\overline{14},\overline{2}`$; $`14,\overline{2},\overline{3}`$; $`\overline{24},\overline{3}`$; $`\overline{14},\overline{23}`$; $`\overline{124},\overline{3}`$; $`\overline{134},\overline{2}`$. Here we only use overline to indicate the color assignment, since permutations of black circles are not considered for $`n_{Z,tor}(L_y,d)`$. These partitions are in one-to-one correspondence with the partitions for $`L_y=4`$ and $`d=3`$ (without color permutation), namely $`\overline{1},\overline{2},\overline{3}`$; $`\overline{1},\overline{3},\overline{4}`$; $`\overline{1},\overline{2},\overline{4}`$; $`\overline{14},\overline{2},\overline{3}`$; $`\overline{2},\overline{3},\overline{4}`$; $`\overline{23},\overline{1},\overline{4}`$; $`\overline{12},\overline{3},\overline{4}`$; $`\overline{34},\overline{1},\overline{2}`$. Eq. (120) is equivalent to the relation
$$n_{Z,tor}(L_y,d)=\underset{d^{}=d}{\overset{L_y}{}}n_{Z,cyc}(L_y,d^{})\text{for}d2.$$
(121)
The construction of the transfer matrix for each level (= degree) $`d`$ can be carried out by methods similar to those for cyclic strips zt . Notice that for the honeycomb lattice, the number of vertices in the transverse direction, $`L_y`$, should be an even number, and the smallest value without degeneracy is $`L_y=4`$. Using the bases described above (e.g. eqs. (56)-(94) for $`2L_y4`$), we define $`J_{L_y,d,i,i+1}`$ as the join operator between vertices $`i`$ and $`i+1`$ and $`D_{L_y,d,i}`$ as the detach operator on vertex $`i`$ for each subspace $`d`$. As compared with the situation for cyclic strips, the corresponding toroidal strips with the same $`L_y`$ and $`d`$ have one more join operator, namely $`J_{L_y,d,L_y,1}J_{L_y,d,L_y,L_y+1}`$. The transfer matrix $`T_{Z,\mathrm{\Lambda },L_y,d}`$ for each $`d`$ is the product of the transverse and longitudinal parts, $`H_{Z,\mathrm{\Lambda },L_y,d}`$ and $`V_{Z,\mathrm{\Lambda },L_y,d}`$, which can be expressed as
$`H_{Z,sq,L_y,d}`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{L_y}{}}}(I+vJ_{L_y,d,i,i+1})`$ (122)
$`H_{Z,tri,L_y,d}`$ $`=`$ $`J_{L_y+1,d,L_y+1,1}{\displaystyle \underset{i=1}{\overset{L_y}{}}}(I+vJ_{L_y+1,d,i,i+1})`$ (124)
$`H_{Z,hc,L_y,d,1}`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{L_y/2}{}}}(I+vJ_{L_y,d,2i1,2i}),H_{Z,hc,L_y,d,2}={\displaystyle \underset{i=1}{\overset{L_y/2}{}}}(I+vJ_{L_y,d,2i,2i+1})`$ (126)
$`V_{Z,sq,L_y,d}`$ $`=`$ $`V_{Z,hc,L_y,d}={\displaystyle \underset{i=1}{\overset{L_y}{}}}(vI+D_{L_y,d,i})`$ (128)
$`V_{Z,tri,L_y,d}`$ $`=`$ $`D_{L_y+1,d,1}{\displaystyle \underset{i=1}{\overset{L_y}{}}}[(I+vJ_{L_y+1,d,i,i+1})(vI+D_{L_y+1,d,i+1})],`$ (130)
where $`[\nu ]`$ denotes the integral part of $`\nu `$, and $`I`$ is the identity matrix. Notice that for triangular strips with periodic transverse boundary conditions, one has to work with width-$`(L_y+1)`$ matrices and then identify the two end vertices, as shown in sqtran . We have
$`T_{Z,sq,L_y,d}`$ $`=`$ $`V_{Z,sq,L_y,d}H_{Z,sq,L_y,d},T_{Z,tri,L_y,d}=V_{Z,tri,L_y,d}H_{Z,tri,L_y,d}`$ (131)
$`T_{Z,hc,L_y,d}`$ $`=`$ $`(V_{Z,hc,L_y,d}H_{Z,hc,L_y,d,2})(V_{Z,hc,L_y,d}H_{Z,hc,L_y,d,1})T_{Z,hc,L_y,d,2}T_{Z,hc,L_y,d,1}.`$ (133)
As discussed above, the definition of $`b^{(d)}`$ does not include the possible permutations of black circles. In principle, there can be $`d!`$ eigenvalues with their coefficients $`b_j^{(d)}`$ added up to $`b^{(d)}`$. In our explicit calculation for toroidal strips, $`b^{(d)}`$ decomposes into two $`b_j^{(d)}`$ (with possible multiplication by an even integer) for the square and honeycomb lattices when $`1<d<L_y`$ up to $`L_y=4`$ and for the triangular lattice when $`1<dL_y`$ up to $`L_y=3`$. As discussed above, we know that there is only one coefficient for both $`d=0`$ and $`d=1`$, i.e., $`b^{(0)}=1`$, $`b^{(1)}=q1`$. For $`d=2,3`$, the coefficients are
$`b_1^{(2)}`$ $`=`$ $`{\displaystyle \frac{1}{2}}q(q3)`$ (134)
$`b_2^{(2)}`$ $`=`$ $`{\displaystyle \frac{1}{2}}(q1)(q2)`$ (136)
$`b_1^{(3)}`$ $`=`$ $`{\displaystyle \frac{1}{3}}(q1)(q^25q+3)`$ (138)
$`=`$ $`{\displaystyle \frac{1}{6}}q(q1)(q5)+{\displaystyle \frac{1}{6}}(q1)(q2)(q3)`$ (140)
$`b_2^{(3)}`$ $`=`$ $`{\displaystyle \frac{1}{6}}q(q2)(q4)`$ (142)
with possible multiplication. In contrast to cyclic and Möbius strips, certain eigenvalues for toroidal strips may appear in more than one level, so that their coefficients are the summation of the corresponding $`b_j^{(d)}`$’s.
For the Klein bottle strip of the square and triangular lattices or the honeycomb lattice with $`L_y`$ even, the sum of coefficients is the same as the sum for the corresponding strip with Möbius longitudinal boundary conditions,
$$C_{Z,L_y,Kb}\underset{j=1}{\overset{N_{Z,L_y,Kb,\lambda }}{}}c_{Z,L_y,Kb,j}=\{\begin{array}{cc}q^{L_y/2}\hfill & \text{for even }L_y\hfill \\ q^{(L_y+1)/2}\hfill & \text{for odd }L_y\hfill \end{array}.$$
(143)
Consider the coefficients for Klein bottle strips. For the square and honeycomb lattices, when the longitudinal boundary condition is changed from toroidal to Klein bottle, we observe the following changes of coefficients:
$`b^{(0)}`$ $``$ $`\pm b^{(0)}`$ (144)
$`b^{(1)}`$ $``$ $`\pm b^{(1)}`$ (146)
$`b_1^{(2)}`$ $``$ $`\pm b_1^{(2)}`$ (148)
$`b_2^{(2)}`$ $``$ $`\pm b_2^{(2)}`$ (150)
$`b_1^{(3)}`$ $``$ $`\pm b^{(1)}.`$ (152)
We find that certain coefficients become zero when the boundary condition is changed from toroidal to Klein bottle, and therefore the number of eigenvalues for Klein bottle strips always appears to be less than the number for the corresponding toroidal strips. This has been observed before in Ref. tk ; tor4 . The partition functions for Klein bottle strips have the same expressions as eqs. (5) and (15) with $`b_j^{(d)}`$ replaced by appropriate coefficients, as in eq. (152). Illustrative calculations will be given below.
## III Zero-temperature Potts antiferromagnet
For the chromatic polynomials, it is necessary that adjacent vertices are not assigned the same color, i.e., in the partition diagrams, adjacent vertices cannot be connected by color-connection lines. For the square and triangular lattices, a transverse slice is a circuit graph; while for the honeycomb lattice, a transverse slice is a set of $`L_y/2`$ two-vertex tree graphs.
### III.1 Square and triangular lattices
For the square and triangular lattices, the second partition of the degree $`d=0`$ subspace and the third partition of the degree $`d=1`$ subspace for width $`L_y=2`$ in Fig. 2 are not allowed. For $`L_y=3`$ in Fig. 3, we keep the first partition in the $`d=0`$ subspace, the first three partitions of the $`d=1`$ subspace, and the first six partitions of the $`d=2`$ subspace. Let us denote the reduced size of the transfer matrix for the square and triangular lattices as $`n_{P,tor}(L_y,d)`$, excluding permutations of black circles. The sum of all coefficients is the dimension of the total transfer matrix (independent of $`L_y`$), which is the chromatic polynomial of the circuit graph given by Theorem 10 of cf , so that for $`L_y2`$,
$$C_{P,\mathrm{\Lambda },L_y}=\underset{d=0}{\overset{L_y}{}}n_{P,tor}(L_y,d)b^{(d)}=(q1)^{L_y}+(q1)(1)^{L_y}\text{for}\mathrm{\Lambda }=sq,tri.$$
(153)
By substituting the expression for $`b^{(d)}`$ in eq. (95), we determine the $`n_{P,tor}(L_y,d)`$. We list the first few numbers $`n_{P,tor}(L_y,d)`$ in Table 3. The strip with $`L_y=1`$ is obtained by identifying the free boundaries of the cyclic strip with $`L_y=2`$, namely, each vertex is attached by a loop such that its chromatic polynomial is identically zero. The column $`n_{P,tor}(L_y,0)`$ in Table 3 is the number of non-crossing non-nearest-neighbor partitions of $`L_y`$ vertices on a circle, and is given by the Riordan number $`R_{L_y}`$ rpstanley ; bernhart99 . If we compare $`n_{P,tor}(L_y,0)`$ with the corresponding number $`n_{P,cyc}(L_y,0)`$ (which was denoted as $`n_P(L_y,d)`$) for cyclic strips, given in Table 1 of cf , we find the relation
$$n_{P,cyc}(L_y,0)n_{P,tor}(L_y,0)=n_{P,tor}(L_y1,0)\text{for}L_y2.$$
(154)
This can be understood as follows. For toroidal strips we have additional bonds connecting the top and bottom vertices of the corresponding cyclic strips, so these vertices cannot have the same color. That is, the $`d=0`$ partitions for cyclic strips with the top and bottom vertices connected (having the same color) is not allowed for the corresponding toroidal strips. The number of these partitions is given by $`n_{P,tor}(L_y1,0)`$. This is clear, since we know $`n_{P,cyc}(L_y,0)=M_{L_y1}`$ cf , where $`M_n`$ is the Motzkin number, and the relation $`M_n=R_n+R_{n+1}`$. We notice the following recursion relations:
$`n_{P,tor}(L_y,1)`$ $`=`$ $`n_{P,tor}(L_y,0)+n_{P,tor}(L_y1,1)+n_{P,tor}(L_y1,2)+(1)^{L_y}\text{for}L_y2`$ (155)
$`n_{P,tor}(L_y,2)`$ $`=`$ $`n_{P,tor}(L_y1,0)+2n_{P,tor}(L_y1,1)+n_{P,tor}(L_y1,2)`$ (159)
$`+n_{P,tor}(L_y1,3)+(1)^{L_y}\text{for}L_y2`$
$`n_{P,tor}(L_y,d)`$ $`=`$ $`n_{P,tor}(L_y1,d1)+n_{P,tor}(L_y1,d)+n_{P,tor}(L_y1,d+1)`$ (163)
$`\text{for}L_y3,3dL_y.`$
The last line here is the same as the recursion relation for $`n_{P,cyc}(L_y,d)`$ when $`1dL_y`$ for cyclic strips cf . By substituting $`q=0`$ into eq. (153) and using eq. (55), we have
$$\underset{d=0}{\overset{L_y}{}}n_{P,tor}(L_y,d)(1)^d=0\mathrm{for}\mathrm{\Lambda }=sq,tri.$$
(164)
Let us denote the number of diagonal dissections of a convex $`n`$-gon into $`k`$ regions as $`D(n,k)`$, which is sequence A033282 in sl . (See also Ref. Flajolet99 .)
$$D(n,k)=\frac{1}{k}\left(\genfrac{}{}{0pt}{}{n3}{k1}\right)\left(\genfrac{}{}{0pt}{}{n+k2}{k1}\right)\text{for}n3,1kn2.$$
(165)
The first few numbers are $`D(3,1)=1`$, $`D(4,1)=1`$, $`D(4,2)=2`$, $`D(5,1)=1`$, $`D(5,2)=5`$, $`D(5,3)=5`$. We also know that the Riordan number for $`n>1`$ is given by the number of dissections of a convex polygon by nonintersecting diagonals with $`n+1`$ edges. To relate $`D(n,k)`$ with the Riordan number, define $`D^{}(n,k)`$ as
$$D^{}(n,k)=D(n+2k,k)=\frac{1}{k}\left(\genfrac{}{}{0pt}{}{n1k}{k1}\right)\left(\genfrac{}{}{0pt}{}{n}{k1}\right)\text{for}n2.$$
(166)
It is clear that the entries of $`D^{}(n,k)`$ are zero for $`[n/2]+1kn`$. Let us also define $`D^{}(1,1)0=n_{P,tor}(1,0)`$, and list the first few rows of the triangle of $`D^{}(n,k)`$ for $`n1,1kn`$ in Fig. 5. The summation of each row in the triangle gives the Riordan number:
$$\underset{k=1}{\overset{[L_y/2]}{}}D^{}(L_y,k)=R_{L_y}=n_{P,tor}(L_y,0).$$
(167)
For $`d=1`$, we need to assign a color for each component of level 0 partitions, so that
$$n_{P,tor}(L_y,1)=\underset{k=1}{\overset{[L_y/2]}{}}(L_y+1k)D^{}(L_y,k)=\underset{k=1}{\overset{[L_y/2]}{}}\left(\genfrac{}{}{0pt}{}{L_y}{k}\right)\left(\genfrac{}{}{0pt}{}{L_y1k}{k1}\right).$$
(168)
For entries $`n_{P,tor}(L_y,d)`$ in Table 3 with $`d>1`$, we find the following relation with known number sequences. Motivated by eq. (154), let us define $`n_{P,tor}^{}(L_y,d)`$ such that $`n_{P,tor}^{}(1,d)=n_{P,tor}(1,d)`$ and
$$n_{P,tor}^{}(L_y,d)=n_{P,tor}(L_y,d)+n_{P,tor}(L_y1,d)\text{for}L_y2.$$
(169)
We list the first few numbers $`n_{P,tor}^{}(L_y,d)`$ in Table 4. Comparing these with the triangular array given by sequence A025564 in sl (which result from pairwise sums of entries in the trinomial array (sequence A027907)), we find that $`n_{P,tor}^{}(L_y,d)`$ for $`d2`$ are columns in the array. The first few rows of this triangular array are listed in Fig. 6. Let us denote a particular entry as $`T(n,k)`$; then $`T(n,k)=T(n1,k2)+T(n1,k1)+T(n1,k)`$ starting with rows $`[1]`$,$`[1,2,1]`$,$`[1,3,4,3,1]`$. We find the relation
$$n_{P,tor}^{}(L_y,d)=T(L_y+1,L_y+1d)=T(L_y+1,L_y+1+d)\text{for}d2.$$
(170)
We list the first few $`n_{P,tor}^{}(L_y,d)n_{P,cyc}(L_y,d)`$ for $`L_y2`$ in Table 5. Because of eq. (154), we have $`n_{P,tor}^{}(L_y,0)n_{P,cyc}(L_y,0)=0`$. For $`d=1`$, we find $`n_{P,tor}^{}(L_y,1)n_{P,cyc}(L_y,1)`$ is given by sequence A005775 in sl , the elements of which are the numbers of compact rooted directed animals of size $`L_y3`$ having 3 source points. Combining the results in Table 5 and Table 4, we infer the relation
$$n_{P,tor}^{}(L_y,d)n_{P,cyc}(L_y,d)=n_{P,tor}^{}(L_y,d+1)\text{for}d2,$$
(171)
or equivalently,
$$n_{P,tor}^{}(L_y,d)=\underset{d^{}=d}{\overset{L_y}{}}n_{P,cyc}(L_y,d^{})\text{for}d2.$$
(172)
We recall that the sum of the coefficients in (9) for a Klein bottle strip of the square or triangular lattice of width $`L_y`$ is (independent of $`L_x`$), given as Theorem 11 in cf , is
$$C_{P,L_y,Kb}=\underset{j=1}{\overset{N_{P,L_y,Kb,\lambda }}{}}c_{P,L_y,Kb,j}=0$$
(173)
### III.2 Honeycomb lattice
For a strip of the honeycomb lattice with toroidal boundary conditions, the width $`L_y`$ must be even. In each transverse slice of such a strip, the $`L_y`$ vertices are connected in a pairwise manner, just as for the strips with cyclic boundary conditions. Let us denote the reduced size of the transfer matrix as $`n_{P,tor,hc}(L_y,d)`$, without considering permutations of black circles in the partition diagram. The sum of all coefficients is the dimension of the total transfer matrix (independent of $`L_x`$), which is the same as eq. (6.19) of hca for cyclic strips, namely
$$C_{P,hc,L_y}=\underset{d=0}{\overset{L_y}{}}n_{P,tor,hc}(L_y,d)b^{(d)}=\left(q(q1)\right)^{L_y/2}.$$
(174)
By substituting the expression for $`b^{(d)}`$ in eq. (95), we determine the $`n_{P,tor,hc}(L_y,d)`$. We list the first few numbers $`n_{P,tor,hc}(L_y,d)`$ in Table 6. We infer the following recursion relations for $`L_y4`$:
$`n_{P,tor,hc}(L_y,0)`$ $`=`$ $`n_{P,tor,hc}(L_y2,0)+4n_{P,tor,hc}(L_y2,1)n_{P,tor,hc}(L_y2,2)`$ (177)
$`n_{P,tor,hc}(L_y2,3)`$
$`n_{P,tor,hc}(L_y,1)`$ $`=`$ $`3n_{P,tor,hc}(L_y2,0)+10n_{P,tor,hc}(L_y2,1)+n_{P,tor,hc}(L_y2,2)`$ (179)
$`n_{P,tor,hc}(L_y,2)`$ $`=`$ $`3n_{P,tor,hc}(L_y2,0)+8n_{P,tor,hc}(L_y2,1)+4n_{P,tor,hc}(L_y2,2)`$ (183)
$`+3n_{P,tor,hc}(L_y2,3)+n_{P,tor,hc}(L_y2,4)`$
$`n_{P,tor,hc}(L_y,3)`$ $`=`$ $`n_{P,tor,hc}(L_y2,0)+2n_{P,tor,hc}(L_y2,1)+3n_{P,tor,hc}(L_y2,2)`$ (187)
$`+4n_{P,tor,hc}(L_y2,3)+3n_{P,tor,hc}(L_y2,4)+n_{P,tor,hc}(L_y2,5)`$
$`n_{P,tor,hc}(L_y,d)`$ $`=`$ $`n_{P,tor,hc}(L_y2,d2)+3n_{P,tor,hc}(L_y2,d1)+4n_{P,tor,hc}(L_y2,d)`$ (191)
$`+3n_{P,tor,hc}(L_y2,d+1)+n_{P,tor,hc}(L_y2,d+2)\text{for}4dL_y.`$
By substituting $`q=0`$ into eq. (174) and using eq. (55), we have
$$\underset{d=0}{\overset{L_y}{}}n_{P,tor,hc}(L_y,d)(1)^d=0.$$
(192)
It is clear that the number of degree $`d=0`$ partitions $`n_{P,tor,hc}(L_y,0)`$ is the same as the corresponding number $`n_{P,cyc,hc}(L_y,0)`$ for the corresponding cyclic strips (which was denoted as $`n_P(hc,L_y,d)`$ in hca ). Combining these results with those in Table 2 of hca , we observe that
$$n_{P,tor,hc}(L_y,1)=\frac{N_{P,cyc,hc,L_y,\lambda }}{2}$$
(193)
where $`N_{P,cyc,hc,L_y,\lambda }`$ is the total number of $`\lambda `$ for cyclic strips. We list the first few $`n_{P,tor,hc}(L_y,d)n_{P,cyc,hc}(L_y,d)`$ in Table 7. Compare Table 7 with Table 6, it is easy to see the relation
$$n_{P,tor,hc}(L_y,d)n_{P,cyc,hc}(L_y,d)=n_{P,tor,hc}(L_y,d+1)\text{for}d2,$$
(194)
or equivalently,
$$n_{P,tor,hc}(L_y,d)=\underset{d^{}=d}{\overset{L_y}{}}n_{P,cyc,hc}(L_y,d^{})\text{for}d2.$$
(195)
We also have
$$n_{P,tor,hc}(L_y,1)n_{P,cyc,hc}(L_y,1)=\underset{d^{}=3}{\overset{L_y}{}}(1)^{d^{}+1}n_{P,tor,hc}(L_y,d^{})=\underset{\mathrm{odd}d^{}=3}{\overset{L_y1}{}}n_{P,cyc,hc}(L_y,d^{}).$$
(196)
## IV Properties of Transfer Matrices at Special Values of Parameters
In this section we derive some properties of the transfer matrices $`T_{Z,\mathrm{\Lambda },L_y,d}`$ at special values of $`q`$ and $`v`$, and, correspondingly, properties of $`T_{T,\mathrm{\Lambda },L_y,d}`$ at special values of $`x`$ and $`y`$.
### IV.1 $`v=0`$
From (1) or (10) it follows that for any graph $`G`$, the Potts model partition function $`Z(G,q,v)`$ satisfies
$$Z(G,q,0)=q^{n(G)}.$$
(197)
Since this holds for arbitrary values of $`q`$, in the context of the lattice strips considered here, it implies that
$$(T_{Z,\mathrm{\Lambda },L_y,d})_{v=0}=0\mathrm{for}1dL_y,$$
(198)
i.e. these are zero matrices. Secondly, restricting to toroidal strips for simplicity, and using the basic results $`n=L_yL_x=L_ym`$ for $`\mathrm{\Lambda }=sq,tri`$ and $`n=2L_ym`$ for $`\mathrm{\Lambda }=hc`$, eq. (197) implies that
$$Tr[(T_{Z,\mathrm{\Lambda },L_y})^m]_{v=0}=\{\begin{array}{cc}q^{L_ym}\hfill & \text{ for }\mathrm{\Lambda }=sq,tri\hfill \\ q^{2L_ym}\hfill & \text{ for }\mathrm{\Lambda }=hc\hfill \end{array}.$$
(199)
With our explicit calculations, we have
$$(T_{Z,\mathrm{\Lambda },L_y,0})_{jk}=0\mathrm{for}v=0,\mathrm{and}j2,$$
(200)
i.e., all rows of these matrices except the first vanish, and
$$(T_{Z,\mathrm{\Lambda },L_y,0})_{11}=q^{pL_y}\mathrm{for}v=0,$$
(201)
where $`p`$ was given in eq. (18). For this $`v=0`$ case, since all of the rows except the first are zero, the elements $`(T_{Z,\mathrm{\Lambda },L_y,0})_{1k}`$ for $`k2`$ do not enter into $`Tr[(T_{Z,\mathrm{\Lambda },L_y,0})^m]`$, which just reduces to the $`m`$’th power of the (1,1) element:
$$Tr[(T_{Z,\mathrm{\Lambda },L_y,0})^m]=[(T_{Z,\mathrm{\Lambda },L_y,0})_{11}]^m.$$
(202)
Corresponding to this, all of the eigenvalues $`\lambda _{Z,\mathrm{\Lambda },L_y,d,j}`$ vanish except for one, which is equal to $`(T_{Z,\mathrm{\Lambda },L_y,0})_{11}`$ in eq. (201). As will be seen, this is reflected in the property that $`det(T_{Z,\mathrm{\Lambda },L_y,d})`$ has a nonzero power of $`v`$ as a factor for all of the lattice-$`\mathrm{\Lambda }`$ strips considered here. Note that the condition $`v=0`$ is equivalent to the Tutte variable condition $`y=1`$.
### IV.2 $`q=0`$
Another fundamental relation that follows, e.g., by setting $`q=0`$ in eq. (10), is
$$Z(G,0,v)=0.$$
(203)
Now the coefficients $`b^{(d)}`$ evaluated at $`q=0`$ satisfy eq.(55). Although $`b^{(d)}`$ decomposes into two $`b_j^{(d)}`$ for $`d>1`$, one of them is equal to zero and the other is $`(1)^d`$ at $`q=0`$ for the cases we have, namely $`b_1^{(2)}=0`$, $`b_2^{(2)}=1`$, $`b_1^{(3)}=1`$, $`b_2^{(3)}=0`$. Hence, in terms of transfer matrices, we derive the sum rule
$$\underset{0\mathrm{even}dL_y}{}\underset{j}{}(\lambda _{Z,\mathrm{\Lambda },L_y,d,j}(0,v))^m\underset{1\mathrm{odd}dL_y}{}\underset{j}{}(\lambda _{Z,\mathrm{\Lambda },L_y,d,j}(0,v))^m=0.$$
(204)
This is similar to a sum rule that we obtained in Ref. s5 . Since it applies for arbitrary $`m`$, it implies that there must be a pairwise cancellation between various eigenvalues in different degree-$`d`$ subspaces, which, in turn, implies that at $`q=0`$ there are equalities between these eigenvalues. For example, for the $`L_y=2`$ strip of the square lattice and $`q=0`$, two of the eigenvalues of $`T_{Z,sq,2,1}`$ become equal to the eigenvalues of $`T_{Z,sq,2,0}`$, while the third eigenvalue of $`T_{Z,sq,2,1}`$ becomes equal to $`\lambda _{Z,sq,2,2}=v^2`$. Note that setting $`q=0`$ does not, in general, lead to any additional vanishing eigenvalues for the $`T_{Z,\mathrm{\Lambda },L_y,d}`$ of the $`\mathrm{\Lambda }=sq,hc`$ lattices, and hence our formulas below for $`T_{Z,\mathrm{\Lambda },L_y,d}`$ do not contain overall factors of $`q`$. However, for the triangular lattice, certain eigenvalues do vanish at $`q=0`$.
In terms of the variables $`x`$ and $`y`$ in the Tutte polynomial, the value $`q=0`$ is equivalent to the value $`x=1`$ (unless $`v=0`$). In contrast to the vanishing of $`Z(G,q,v)`$ at $`q=0`$, the Tutte polynomial $`T(G,1,y)`$ is nonzero for general $`y`$. This different behavior can be traced to the feature that in eq. (14), $`Z(G,q,v)`$ is proportional to $`T(G,x,y)`$ multiplied by the factor $`(x1)^{k(G)}=(x1)`$, so at $`x=1`$, $`Z(G,0,v)=0`$ even if $`T(G,1,y)0`$.
### IV.3 $`v=1`$, i.e., $`y=0`$
As discussed above, the special value $`v=1`$, i.e., $`y=0`$, corresponds to the zero-temperature Potts antiferromagnet, and in this case the Potts model partition function reduces to the chromatic polynomial, as indicated in eq. (6). In the previous two sections we have determined how the dimensions $`n_{Z,tor}(L_y,d)`$ for the square, triangular, and honeycomb lattice strips reduce to the dimensions $`n_{P,tor}(L_y,d)`$. In all cases, in each degree-$`d`$ subspace for $`0dL_y1`$ for $`\mathrm{\Lambda }=sq,tri,hc`$, some eigenvalues vanish, so that $`n_{P,tor,\mathrm{\Lambda }}(L_y,d)<n_{Z,tor,\mathrm{\Lambda }}(L_y,d)`$ and hence $`det(T_{Z,\mathrm{\Lambda },L_y,d})=0`$ at $`v=1`$. This property is reflected in the powers of $`(v+1)`$ and $`y`$ that appear, respectively, in our formulas below for $`det(T_{Z,\mathrm{\Lambda },L_y,d})`$ and $`det(T_{T,\mathrm{\Lambda },L_y,d})`$.
### IV.4 $`v=q`$, i.e., $`x=0`$
For the graph $`G=G(V,E)`$, setting $`x=0`$ and $`y=1q`$ in the Tutte polynomial $`T(G,x,y)`$ yields the flow polynomial $`F(G,q)`$, which counts the number of nowhere-0 $`q`$-flows (without sinks or sources) that there are on $`G`$ boll :
$$F(G,q)=(1)^{e(G)n(G)+1}T(G,0,1q).$$
(205)
Therefore, the flow polynomial for toroidal lattice strips has the form
$$F(\mathrm{\Lambda },L_y\times L_x,tor.,q)=\underset{d=0}{\overset{L_y}{}}\underset{j}{}b_j^{(d)}(\lambda _{F,\mathrm{\Lambda },L_y,d,j})^m.$$
(206)
We found that for all $`L_y1`$ and $`0dL_y1`$ for $`\mathrm{\Lambda }=sq,hc`$, when one sets $`x=0`$ in the Tutte polynomial, or equivalently, $`q=v`$ in the Potts model partition function, some of the eigenvalues in each degree-$`d`$ subspace vanish, and hence $`det(T_{Z,\mathrm{\Lambda },L_y,d})`$ and $`det(T_{T,\mathrm{\Lambda },L_y,d})`$ vanish. This is reflected in the powers of $`(q+v)`$ and $`x`$ that appear, respectively, in our formulas below for $`det(T_{Z,\mathrm{\Lambda },L_y,d})`$ and $`det(T_{T,\mathrm{\Lambda },L_y,d})`$. The expressions for the determinants for strips of the triangular lattice are more complicated because of the fact alluded to above that one has to work with width-$`(L_y+1)`$ matrices and then identify the two end vertices.
## V General Results for Toroidal Strips
In this section we present general results that we have obtained for transfer matrices of toroidal strips and their properties. These are valid for arbitrarily large strip widths $`L_y2`$ (as well as arbitrarily great lengths).
### V.1 Determinants
For the square lattice, we find
$$det(T_{T,sq,L_y,d})=(xy)^{d!L_yn_{Z,tor}(L_y1,d)}$$
(207)
where $`n_{Z,tor}(L_y,d)`$ was given by eq. (97)-(99). This result applies for all $`d`$, i.e., $`0dL_y`$ since $`n_{Z,tor}(L_y1,d)=0`$ for $`d>L_y1`$ (c.f. eq. (96)). This is equivalent to the somewhat more complicated expression for $`det(T_{Z,sq,L_y,d})`$:
$$det(T_{Z,sq,L_y,d})=v^{d!L_yn_{Z,tor}(L_y,d)}\left[\left(1+\frac{q}{v}\right)(1+v)\right]^{d!L_yn_{Z,tor}(L_y1,d)}.$$
(208)
These determinant formulas can be explained as follows. By the arrangement of the partitions, as shown in Figs. 2 and 3, the matrix $`J_{L_y,d,i,i+1}`$ has the lower triangular form and the matrix $`D_{L_y,d,i}`$ has the upper triangular form. From the definition of the transfer matrix in eq. (133), the determinant of $`T_{Z,sq,L_y,d}`$ is the product of the diagonal elements of $`H_{Z,sq,L_y,d}`$ and $`V_{Z,sq,L_y,d}`$. The diagonal elements of $`H_{Z,sq,L_y,d}`$ have the form $`(1+v)^r`$, where $`r`$ is the number of edges in the corresponding partition, and the diagonal elements of $`V_{Z,sq,L_y,d}`$ have the form $`v^{L_y}(1+q/v)^s`$, where $`s`$ is the number of vertices which do not connect to any other vertex in the corresponding partition. Therefore, the power of $`(1+v)`$ in eq. (208) is the sum of the number of nearest-neighbor edges of all the $`(L_y,d)`$-partitions. Let us compare the $`(L_y,d)`$-partitions and $`(L_y1,d)`$ partitions. Since each edge of a $`(L_y,d)`$-partition corresponds to adding a new vertex in a $`(L_y1,d)`$-partition (at $`L_y`$ possible places) and connecting it with the vertex above it, the power of $`(1+v)`$ is $`d!L_yn_{Z,tor}(L_y1,d)`$, where the factor $`d!`$ comes from the permutation of black circles. It is clear that the power of $`v`$ is $`d!L_yn_{Z,tor}(L_y,d)`$. The power of $`(1+q/v)`$ in eq. (208) is the sum of the number of unconnected vertices of all the $`(L_y,d)`$-partitions. Now consider the $`(L_y,d)`$-partitions and $`(L_y1,d)`$ partitions again. Since each unconnected vertex of a $`(L_y,d)`$-partition corresponds to adding a unconnected vertex to a $`(L_y1,d)`$-partition in $`L_y`$ possible ways, the power of $`(1+q/v)`$ is again $`d!L_yn_{Z,tor}(L_y1,d)`$.
Next, taking into account that the generalized multiplicity is $`b^{(d)}`$, we have, for the total determinant
$`det(T_{Z,sq,L_y})`$ $`=`$ $`{\displaystyle \underset{d=0}{\overset{L_y}{}}}(y1)^{L_yn_{Z,tor}(L_y,d)b^{(d)}}(xy)^{L_yn_{Z,tor}(L_y1,d)b^{(d)}}`$ (209)
$`=`$ $`(y1)^{L_y_{d=0}^{L_y}n_{Z,tor}(L_y,d)b^{(d)}}(xy)^{L_y_{d=0}^{L_y}n_{Z,tor}(L_y1,d)b^{(d)}}.`$ (211)
Using eq. (95) together with eq. (96) so that $`_{d=0}^{L_y}n_{Z,tor}(L_y1,d)b^{(d)}`$=$`_{d=0}^{L_y1}n_{Z,tor}(L_y1,d)b^{(d)}`$, we have, finally,
$$det(T_{Z,sq,L_y})=(y1)^{L_yq^{L_y}}(xy)^{L_yq^{L_y1}}.$$
(212)
This agrees with the conjecture given as eq. (3.70) of our earlier Ref. s3a for the determinant of the transfer matrix of the toroidal strip of the square lattice with arbitrary width $`L_y`$.
For the honeycomb lattice, we find
$$det(T_{T,hc,L_y,d})=(x^2y)^{d!L_yn_{Z,tor}(L_y1,d)}.$$
(213)
Equivalently,
$$det(T_{Z,hc,L_y,d})=(v^2)^{d!L_yn_{Z,tor}(L_y,d)}\left[\left(1+\frac{q}{v}\right)^2(1+v)\right]^{d!L_yn_{Z,tor}(L_y1,d)}.$$
(214)
This can be understood as follows: by an argument similar to that given before, the power of $`(1+v)`$ is the same as for the square lattice case. Comparing $`T_{Z,sq,L_y,d}`$ and $`T_{Z,hc,L_y,d}`$ in eq. (133), one sees that $`V_{Z,hc,L_y,d}=V_{Z,sq,L_y,d}`$ has been multiplied twice for the honeycomb lattice, so that the powers of $`v`$ and $`(1+q/v)`$ become twice of the corresponding powers for the square lattice.
Taking into account that the generalized multiplicity is $`b^{(d)}`$, the total determinant for the $`hc`$ lattice is given by
$$det(T_{T,hc,L_y})=(x^2y)^{L_yq^{L_y1}}.$$
(215)
Equivalently,
$$det(T_{Z,hc,L_y})=(v^2)^{L_yq^{L_y}}\left[\left(1+\frac{q}{v}\right)^2(1+v)\right]^{L_yq^{L_y1}}.$$
(216)
It is clear that $`det(T_{T,hc,L_y,d})`$ is related to $`det(T_{T,sq,L_y,d})`$ by the replacement $`xx^2`$ (holding $`y`$ fixed). This, together with the fact that $`n_{Z,tor}(L_y,d)`$ is the same for these lattices means that the total determinants $`det(T_{T,hc,L_y})`$ is related to $`det(T_{T,sq,L_y})`$ by the same respective replacements. Correspondingly, $`det(T_{Z,hc,L_y,d})`$ is related to $`det(T_{Z,sq,L_y,d})`$ by the replacements of the respective factors $`v`$ by $`v^2`$ (cf. eq. (18)) and $`(1+q/v)`$ by $`(1+q/v)^2`$.
### V.2 Traces
The trace of the total transfer matrix is the $`m=1`$ case in eq. (5) or (15). For strips of the square lattice, this corresponds to a $`L_y`$-vertex circuit with a loop attached to each vertex, as illustrated in Fig. 7. We have $`Tr(T_{T,sq,L_y})`$ given by the corresponding Tutte polynomial,
$$Tr(T_{T,sq,L_y})=y^{L_y}\left(y1+\frac{x^{L_y}1}{x1}\right).$$
(217)
Equivalently, we find
$$Tr(T_{Z,sq,L_y})=(1+v)^{L_y}[(v+q)^{L_y}+(q1)v^{L_y}].$$
(218)
In principle, we can also consider the $`m=1`$ case for the Klein bottle strips, but the result is not as simple as that listed here.
For the triangular lattice, the total trace corresponds to a $`L_y`$-vertex circuit with each edge doubled and with a loop attached to each vertex as illustrated in Fig. 8. Therefore, the trace is given by the Tutte polynomial of this graph
$$Tr(T_{T,tri,L_y})=y^{L_y}\left((y1)(y+1)^{L_y}+\underset{j=0}{\overset{L_y1}{}}(y+1)^j(x+y)^{L_y1j}\right).$$
(219)
Equivalently,
$$Tr(T_{Z,tri,L_y})=q(v+1)^{L_y}\left((v^2+2v)^{L_y}+\underset{j=0}{\overset{L_y1}{}}(v^2+2v)^j(v^2+2v+q)^{L_y1j}\right).$$
(220)
For the honeycomb lattice, the total trace corresponds to a $`2L_y`$-vertex circuit with every other edge doubled as shown in Fig. 9. Hence, the trace is given by the Tutte polynomial of this graph
$$Tr(T_{T,hc,L_y})=(x+y)^{L_y}\left(\frac{x^{L_y}1}{x1}\right)+(y1)(y+1)^{L_y}+\underset{j=0}{\overset{L_y1}{}}(y+1)^j(x+y)^{L_y1j}.$$
(221)
Equivalently,
$`Tr(T_{Z,hc,L_y})`$ $`=`$ $`(v^2+2v+q)^{L_y}\left((v+q)^{L_y}v^{L_y}\right)`$ (224)
$`+qv^{L_y}\left((v^2+2v)^{L_y}+{\displaystyle \underset{j=0}{\overset{L_y1}{}}}(v^2+2v)^j(v^2+2v+q)^{L_y1j}\right).`$
### V.3 Eigenvalues for $`d=L_y`$ and $`d=L_y1`$ for $`\mathrm{\Lambda }=sq,hc`$
It was shown earlier tk ; tor4 that the $`\lambda `$’s for a strip with Klein bottle boundary conditions are a subset of the $`\lambda `$’s for the same strip with torus boundary conditions. From eq. (99) one knows that there is only one $`\lambda `$, denoted as $`\lambda _{Z,\mathrm{\Lambda },L_y,L_y}`$, for degree $`d=L_y`$ for strips of the square and honeycomb lattices because all the permutation of black circles are equivalent. That is, for this value of $`d`$, the transfer matrix reduces to $`1\times 1`$, i.e. is a scalar. We found that for a toroidal or Klein bottle strip of the square, or honeycomb lattice with width $`L_y`$,
$$\lambda _{T,\mathrm{\Lambda },L_y,L_y}=1.$$
(225)
Equivalently, in terms of Potts model variables,
$$\lambda _{Z,sq,L_y,L_y}=v^{L_y}$$
(226)
$$\lambda _{Z,hc,L_y,L_y}=v^{2L_y}.$$
(227)
For the width $`L_y`$ triangular strips, because one has to start by constructing the width-$`(L_y+1)`$ transfer matrices, there can be more than one $`\lambda `$ even for degree $`d=L_y`$.
Concerning the $`\lambda `$’s for degree $`d=L_y1`$, we find that one of these is the same, independent of $`L_y`$, for the toroidal square strips we have calculated, namely
$$\lambda _{T,sq,L_y,L_y1,1}=x.$$
(228)
This also appears for degree $`d=L_y1`$ in the case of the cyclic square-lattice strips, and we conjecture that all the toroidal square-lattice strips have this eigenvalue. For the strips of the square lattice with $`L_y=3,4`$, there is another common $`\lambda `$ for degree $`d=L_y1`$, namely
$$\lambda _{T,sq,L_y,L_y1,2}=y.$$
(229)
For the toroidal honeycomb strip with $`L_y=4`$, one of the eigenvalues for degree $`d=L_y1`$, say that for $`j=1`$, is given by
$$\lambda _{T,hc,L_y,d=L_y1,j=1}=x^2$$
(230)
as for the corresponding cyclic strip. We conjecture that all of the toroidal strips of the honeycomb lattice have this eigenvalue. In terms of Potts model quantities these results are
$`\lambda _{Z,sq,L_y,L_y1,1}`$ $`=`$ $`v^{L_y1}(v+q)`$ (231)
$`\lambda _{Z,sq,L_y,L_y1,2}`$ $`=`$ $`v^{L_y}(v+1)\text{except for}L_y=2`$ (233)
$`\lambda _{Z,hc,L_y,L_y1,1}`$ $`=`$ $`v^{2(L_y1)}(v+q)^2.`$ (235)
The expressions for the other eigenvalues are, in general, more complicated.
## VI Some Illustrative Calculations
### VI.1 Square-Lattice Strip, $`L_y=2`$
The toroidal strip of the square lattice with width $`L_y=2`$ is equivalent to the ladder graph with all the transverse edges doubled. The Potts model partition function $`Z(sq,L_y\times m,q,v)`$ and Tutte polynomial $`T(sq,L_y\times m,x,y)`$ were calculated for the toroidal and Klein bottle strips of the square lattice with width $`L_y=2`$ in Ref. s3a . We express the results here in terms of transfer matrices $`T_{Z,sq,2,d}`$. For $`d=0`$, we have
$$T_{Z,sq,2,0}=\left(\begin{array}{cc}q^2+4qv+qv^2+5v^2+2v^3& (q+2v)(1+v)^2\\ v^3(2+v)& v^2(1+v)^2\end{array}\right).$$
(236)
The eigenvalues of $`T_{Z,sq,2,0}`$ are the same as $`\lambda _{Z,s2t,(5,6)}`$ given as eqs. (3.16) and (3.17) of s3a . For $`d=1`$, we have
$$T_{Z,sq,2,1}=v\left(\begin{array}{ccc}q+3v+v^2& v(2+v)& (1+v)^2\\ v(2+v)& q+3v+v^2& (1+v)^2\\ v^2(2+v)& v^2(2+v)& v(1+v)^2\end{array}\right).$$
(237)
The eigenvalues of $`T_{Z,sq,2,1}`$ are the same as $`\lambda _{Z,s2t,(2,3,4)}`$ given as eqs. (3.13) to (3.15) of s3a . The matrix $`T_{T,sq,2,2}=v^2`$ has been given above.
### VI.2 Square-Lattice Strip, $`L_y=3`$
For the $`L_y=3`$ toroidal and Klein bottle strips of the square lattice, the Potts model partition function $`Z(sq,L_y\times m,q,v)`$ and Tutte polynomial $`T(sq,L_y\times m,x,y)`$ were calculated in Ref. s3a . We express the results here in terms of transfer matrices $`T_{Z,sq,3,d}`$. For $`d=0`$, we have
$$T_{Z,sq,3,0}=\left(\begin{array}{ccccc}s_1& v_1s_2& v_1s_2& v_1s_2& v_1^3s_5\\ v^3s_3& v^2v_1s_4& v^3v_1v_2& v^3v_1v_2& v^2v_1^3\\ v^3s_3& v^3v_1v_2& v^2v_1s_4& v^3v_1v_2& v^2v_1^3\\ v^3s_3& v^3v_1v_2& v^3v_1v_2& v^2v_1s_4& v^2v_1^3\\ v^5v_3& v^4v_1v_2& v^4v_1v_2& v^4v_1v_2& v^3v_1^3\end{array}\right)$$
(238)
where we use the notations $`v_1=1+v`$, $`v_2=2+v`$, $`v_3=3+v`$ and
$`s_1`$ $`=`$ $`q^3+6q^2v+15qv^2+qv^3+16v^3+3v^4`$ (239)
$`s_2`$ $`=`$ $`q^2+5qv+qv^2+8v^2+3v^3`$ (240)
$`s_3`$ $`=`$ $`q+4v+v^2`$ (241)
$`s_4`$ $`=`$ $`q+3v+v^2`$ (242)
$`s_5`$ $`=`$ $`q+3v.`$ (243)
The eigenvalues of $`T_{Z,sq,3,0}`$ consist of a linear term $`v^2(v+q)(v+1)`$ and roots of a cubic equation. The linear term, denoted as $`\lambda _{Z,s3t,10}`$ in s3a , has multiplicity two, so the corresponding coefficient is $`2b^{(0)}=2`$. The cubic equation is the same as eqs. (3.48) to (3.51) of s3a . For $`d=1`$, we have
$$T_{Z,sq,3,1}=v\left(\begin{array}{cccccccccc}s_6& vs_3& vs_3& v_1s_3& vv_1v_2& v_1s_4& v_1s_4& vv_1v_2& v_1^3& vv_1v_2\\ vs_3& s_6& vs_3& vv_1v_2& v_1s_4& vv_1v_2& v_1s_4& vv_1v_2& v_1^3& v_1s_3\\ vs_3& vs_3& s_6& vv_1v_2& v_1s_4& v_1s_4& vv_1v_2& v_1s_3& v_1^3& vv_1v_2\\ v^3& 0& 0& v^2v_1& 0& 0& 0& 0& 0& 0\\ v^3v_3& v^2s_3& v^2s_3& v^2v_1v_2& vv_1s_4& v^2v_1v_2& v^2v_1v_2& v^2v_1v_2& vv_1^3& v^2v_1v_2\\ v^2s_3& v^3v_3& v^2s_3& v^2v_1v_2& v^2v_1v_2& vv_1s_4& v^2v_1v_2& v^2v_1v_2& vv_1^3& v^2v_1v_2\\ v^2s_3& v^2s_3& v^3v_3& v^2v_1v_2& v^2v_1v_2& v^2v_1v_2& vv_1s_4& v^2v_1v_2& vv_1^3& v^2v_1v_2\\ 0& 0& v^3& 0& 0& 0& 0& v^2v_1& 0& 0\\ v^4v_3& v^4v_3& v^4v_3& v^3v_1v_2& v^3v_1v_2& v^3v_1v_2& v^3v_1v_2& v^3v_1v_2& v^2v_1^3& v^3v_1v_2\\ 0& v^3& 0& 0& 0& 0& 0& 0& 0& v^2v_1\end{array}\right)$$
(244)
where
$$s_6=q^2+5qv+8v^2+v^3.$$
(245)
The eigenvalues of $`T_{Z,sq,3,1}`$ consist of the roots of a cubic equation, which enter with multiplicity two, and the roots of a quartic equation. The cubic equation is given by eqs. (2.23) to (2.26) in s3a , and since it has multiplicity two, its roots have coefficients $`2b^{(1)}=2(q1)`$. The quartic equation is the same as that given in eqs. (3.43) to (3.47) of s3a . For $`d=2`$, we have
$$T_{Z,sq,3,2}=v^2\left(\begin{array}{cccccccccccc}s_5& 0& v& 0& 0& v& v_1& 0& 0& 0& v_1& 0\\ 0& s_5& 0& v& v& 0& 0& v_1& 0& 0& 0& v_1\\ v& 0& s_5& 0& v& 0& v_1& 0& v_1& 0& 0& 0\\ 0& v& 0& s_5& 0& v& 0& v_1& 0& v_1& 0& 0\\ 0& v& v& 0& s_5& 0& 0& 0& v_1& 0& 0& v_1\\ v& 0& 0& v& 0& s_5& 0& 0& 0& v_1& v_1& 0\\ v^2& 0& v^2& 0& 0& 0& vv_1& 0& 0& 0& 0& 0\\ 0& v^2& 0& v^2& 0& 0& 0& vv_1& 0& 0& 0& 0\\ 0& 0& v^2& 0& v^2& 0& 0& 0& vv_1& 0& 0& 0\\ 0& 0& 0& v^2& 0& v^2& 0& 0& 0& vv_1& 0& 0\\ v^2& 0& 0& 0& 0& v^2& 0& 0& 0& 0& vv_1& 0\\ 0& v^2& 0& 0& v^2& 0& 0& 0& 0& 0& 0& vv_1\end{array}\right).$$
(246)
The eigenvalues of $`T_{Z,sq,3,2}`$ consist of two linear terms and roots of three quadratic equations; two of the quadratic equations occur with multiplicity two. Both of the linear terms $`v^2(v+q)`$ and $`v^3(v+1)`$ have the coefficient $`b_2^{(2)}=(q1)(q2)/2`$. One of the quadratic equations is given by eq. (3.39) of s3a with coefficient $`q(q3)/2`$. The other two quadratic equations have multiplicity two, as noted; one of them is given by eq. (2.20) of s3a and its roots have coefficient $`2b_2^{(2)}=(q1)(q2)`$. The other is given by eq. (2.21) of s3a and its roots have coefficient $`2b_1^{(2)}=q(q3)`$. The matrix $`T_{T,sq,3,3}=v^3`$ has been given above and has coefficient $`b^{(3)}`$.
### VI.3 Square-Lattice Strip, $`L_y=4`$
For the $`L_y=4`$ toroidal and Klein bottle strips of the square lattice, the sizes of matrices are $`14\times 14`$, $`35\times 35`$, $`56\times 56`$, and $`48\times 48`$ for levels $`0d3`$, as indicated in Table 1. The eigenvalues of $`T_{Z,sq,4,0}`$ are roots of three quadratic equations and a sixth degree equation. One of the quadratic equations has multiplicity two. The eigenvalues of $`T_{Z,sq,4,1}`$ consist of a linear term $`v^3(1+v)(q+v)`$, roots of a cubic equation, a sixth-degree equation, an eighth-degree equation and a ninth-degree equation. The eighth-degree equation has multiplicity two. The eigenvalues of $`T_{Z,sq,4,2}`$ consist of a linear term, roots of two cubic equations, three quartic equations, two sixth-degree equations and an eighth-degree equation. The linear term is again $`v^3(1+v)(q+v)`$ which has multiplicity four with coefficient $`2(q^23q+1)`$. The eighth-degree equation and one of the sixth-degree equations have multiplicity two. The eigenvalues of $`T_{Z,sq,4,3}`$ consist of two linear terms $`v^4(1+v)`$ and $`v^3(v+q)`$, and roots of four quadratic equations and a quartic equation. The linear terms and one of the quadratic equations have multiplicity two, and the other equations have multiplicity four. The matrix $`T_{T,sq,4,4}=v^4`$ has been given above with coefficient $`b^{(4)}`$. Some of these equations are too lengthy to list here; they are available from the authors. At the special value $`v=1`$, the Potts model partition function reduces to the chromatic polynomial. The non-zero eigenvalues are the same as $`\lambda _{st4,j}`$ for $`1j33`$ from eqs. (2.3) to (2.21) in tor4 .
### VI.4 Triangular-Lattice Strip, $`L_y=2`$
The toroidal strip of the triangular lattice with $`L_y=2`$ is equivalent to the ladder graph with next-nearest-neighbor coupling and all transverse edges doubled. The corresponding Klein bottle strip is the same as the toroidal strip. For $`d=0`$, we have
$$T_{Z,tri,2,0}=\left(\begin{array}{cc}q^2+6qv+qv^2+12v^2+8v^3+2v^4& (q+4v+2v^2)(1+v)^2\\ v^2(2q+12v+13v^2+6v^3+v^4)& v^2(2+v)^2(1+v)^2\end{array}\right).$$
(247)
The eigenvalues of $`T_{Z,tri,2,0}`$, denoted as $`\lambda _{tt2,(1,2)}`$, are roots of the following equation:
$`\xi ^2`$ $``$ $`(v^6+q^2+6v^5+20v^3+15v^4+16v^2+qv^2+6vq)\xi `$ (249)
$`+v^2(1+v)^2(2v^4+4v^3+4v^3q+11qv^2+q^2v^2+4vq+4q^2v+2q^2)=0.`$
For $`d=1`$, we have
$$T_{Z,tri,2,1}=v\left(\begin{array}{ccc}q+6v+4v^2+v^3& q+6v+4v^2+v^3& (2+v)(1+v)^2\\ q+6v+4v^2+v^3& q+6v+4v^2+v^3& (2+v)(1+v)^2\\ v(q+12v+13v^2+6v^3+v^4)& v(q+12v+13v^2+6v^3+v^4)& v(2+v)^2(1+v)^2\end{array}\right).$$
(250)
The eigenvalues of $`T_{Z,tri,2,1}`$ consist of a zero, $`\lambda _{tt2,0}=0`$, and roots of the following quadratic equation:
$$\xi ^2v(v^5+6v^4+15v^3+20v^2+16v+2q)\xi +2v^3(2+v)(1+v)^2(v^2+vq+q)=0$$
(251)
with roots $`\lambda _{tt2,(3,4)}`$. Their coefficients are equal to $`b^{(1)}=q1`$. For $`d=2`$, we have
$$T_{Z,tri,2,2}=v^2\left(\begin{array}{cc}1& 1\\ 1& 1\end{array}\right).$$
(252)
$`T_{Z,tri,2,2}`$ has two eigenvalues. One of them is $`\lambda _{tt2,5}=2v^2`$ with coefficient $`b_1^{(2)}=q(q3)/2`$; while the other eigenvalue is zero with coefficient $`b_2^{(2)}=(q1)(q2)/2`$. Therefore, the Potts model partition function for the triangular lattice strip with $`L_y=2`$ and toroidal boundary condition is given by
$$Z(tri,2\times L_x,tor.,q,v)=b^{(0)}\underset{j=1}{\overset{2}{}}(\lambda _{tt2,j})^{L_x}+b^{(1)}\underset{j=3}{\overset{4}{}}(\lambda _{tt2,j})^{L_x}+b_1^{(2)}(\lambda _{tt2,5})^{L_x}.$$
(253)
### VI.5 Triangular-Lattice Strip, $`L_y=3`$
We illustrate our results for the $`L_y=3`$ toroidal strip of the triangular lattice. For $`d=0`$, we have
$$T_{Z,tri,3,0}=\left(\begin{array}{ccccc}t_1& v_1t_2& v_1t_2& v_1t_2& v_1^3t_5\\ v^2t_3& v^2v_1v_2t_4& v^2v_1v_2t_4& v^2v_1t_6& v^2v_1^3v_2^2\\ v^2t_3& v^2v_1t_6& v^2v_1v_2t_4& v^2v_1v_2t_4& v^2v_1^3v_2^2\\ v^2t_3& v^2v_1v_2t_4& v^2v_1t_6& v^2v_1v_2t_4& v^2v_1^3v_2^2\\ v^4v_3t_7& v^3v_1v_2t_8& v^3v_1v_2t_8& v^3v_1v_2t_8& v^3v_1^3v_2^3\end{array}\right)$$
(254)
where we define
$`t_1`$ $`=`$ $`q^3+9q^2v+33qv^2+4qv^3+50v^3+21v^4+3v^5`$ (255)
$`t_2`$ $`=`$ $`q^2+8qv+2qv^2+20v^2+14v^3+3v^4`$ (256)
$`t_3`$ $`=`$ $`q^2+10qv+2qv^2+30v^2+22v^3+7v^4+v^5`$ (257)
$`t_4`$ $`=`$ $`q+6v+4v^2+v^3`$ (258)
$`t_5`$ $`=`$ $`q+6v+3v^2`$ (259)
$`t_6`$ $`=`$ $`2q+14v+14v^2+6v^3+v^4`$ (260)
$`t_7`$ $`=`$ $`3q+18v+15v^2+6v^3+v^4`$ (261)
$`t_8`$ $`=`$ $`q+12v+13v^2+6v^3+v^4.`$ (262)
The characteristic equation of $`T_{Z,tri,3,0}`$ yields the following cubic and quadratic equations:
$$\xi ^3+f_{31}\xi ^2+f_{32}\xi +f_{33}=0$$
(263)
where
$`f_{31}`$ $`=`$ $`(q^3+v^9+39v^2q+9v^8+36v^7+84v^6+129v^5+137v^4+96v^3+12v^3q+2v^4q`$ (265)
$`+9q^2v)`$
$`f_{32}`$ $`=`$ $`v^2(1+v)(46q^3v+76q^3v^2+246q^2v^5+56v^3q^3+1190v^6q+72q^2v^6+q^3v^6+24v^9q`$ (270)
$`+2v^{10}q+3q^4+35v^{10}+3v^{11}+9v^7q^2+534v^7q+147v^8q+483v^3q^2+213v^2q^2`$
$`+8q^3v^5+28q^3v^4+1272v^4q+2q^4v+414v^3q+794v^5+192v^4+1618v^5q+1091v^7`$
$`+1243v^6+188v^9+580v^8+455v^4q^2)`$
$`f_{33}`$ $`=`$ $`v^5(1+v)^4(22q^3v+72q^3v^2+105q^2v^5+105v^3q^3+72v^6q+24q^2v^6+6q^4+19v^7q`$ (274)
$`+90v^3q^2+30v^2q^2+2v^4q^4+12q^4v^3+24q^4v^2+12q^3v^5+62q^3v^4+60v^4q+19q^4v`$
$`+24v^3q+24v^5+8v^4+90v^5q+22v^7+30v^6+6v^8+148v^4q^2)`$
with roots $`\lambda _{tt3,j}`$ for $`1j3`$ and
$$\xi ^2v^3(q2)(1+v)\xi +v^6(q2)^2(1+v)^2=0$$
(275)
with roots $`\lambda _{tt3,j}`$ for $`j=4,5`$. For $`d=1`$, we have
$`T_{Z,tri,3,1}=v\times `$ (288)
$`\left(\begin{array}{cccccccccc}t_9& t_9& vt_{10}& v_1t_{11}& v_1t_4& v_1t_4& v_1v_2s_3& vv_1v_2^2& v_1^3v_2& v_1t_{11}\\ vt_{10}& t_9& t_9& vv_1v_2^2& v_1v_2s_3& v_1t_4& v_1t_4& v_1t_{11}& v_1^3v_2& v_1t_{11}\\ t_9& vt_{10}& t_9& v_1t_{11}& v_1t_4& v_1v_2s_3& v_1t_4& v_1t_{11}& v_1^3v_2& vv_1v_2^2\\ v^2t_{12}& v^2t_{12}& 0& v^2v_1v_2& 0& 0& v^2v_1v_2& 0& 0& v^2v_1v_2\\ v^2t_{13}& v^2t_{13}& vt_{14}& v^2v_1v_2t_{15}& vv_1v_2t_4& vv_1v_2t_4& vv_1t_8& vv_1t_{16}& vv_1^3v_2^2& v^2v_1v_2t_{15}\\ vt_{14}& v^2t_{13}& v^2t_{13}& vv_1t_{16}& vv_1t_8& vv_1v_2t_4& vv_1v_2t_4& v^2v_1v_2t_{15}& vv_1^3v_2^2& v^2v_1v_2t_{15}\\ v^2t_{13}& vt_{14}& v^2t_{13}& v^2v_1v_2t_{15}& vv_1v_2t_4& vv_1t_8& vv_1v_2t_4& v^2v_1v_2t_{15}& vv_1^3v_2^2& vv_1t_{16}\\ v^2t_{12}& 0& v^2t_{12}& v^2v_1v_2& 0& v^2v_1v_2& 0& v^2v_1v_2& 0& 0\\ v^3v_3t_{17}& v^3v_3t_{17}& v^3v_3t_{17}& v^3v_1v_2v_3t_{18}& v^2v_1v_2t_8& v^2v_1v_2t_8& v^2v_1v_2t_8& v^3v_1v_2v_3t_{18}& v^2v_1^3v_2^3& v^3v_1v_2v_3t_{18}\\ 0& v^2t_{12}& v^2t_{12}& 0& v^2v_1v_2& 0& 0& v^2v_1v_2& 0& v^2v_1v_2\end{array}\right)`$
where
$`t_9`$ $`=`$ $`q^2+8qv+qv^2+20v^2+8v^3+v^4`$ (291)
$`t_{10}`$ $`=`$ $`2q+10v+5v^2+v^3`$ (292)
$`t_{11}`$ $`=`$ $`q+8v+5v^2+v^3`$ (293)
$`t_{12}`$ $`=`$ $`q+6v+2v^2`$ (294)
$`t_{13}`$ $`=`$ $`4q+qv+24v+20v^2+7v^3+v^4`$ (295)
$`t_{14}`$ $`=`$ $`q^2+9qv+2qv^2+30v^2+22v^3+7v^4+v^5`$ (296)
$`t_{15}`$ $`=`$ $`5+4v+v^2`$ (297)
$`t_{16}`$ $`=`$ $`q+14v+14v^2+6v^3+v^4`$ (298)
$`t_{17}`$ $`=`$ $`2q+18v+15v^2+6v^3+v^4`$ (299)
$`t_{18}`$ $`=`$ $`4+3v+v^2.`$ (300)
The characteristic equation of $`T_{Z,tri,3,1}`$ yields the following quartic equation and sixth-degree equation:
$$\xi ^4+f_{41}\xi ^3+f_{42}\xi ^2+f_{43}\xi +f_{44}=0$$
(301)
where
$`f_{41}`$ $`=`$ $`v(84v^5+36v^6+9v^7+2v^3q+9v^2q+2q^2+v^8+140v^3+130v^4+98v^2+23vq)`$ (302)
$`f_{42}`$ $`=`$ $`v^3(1+v)(4v^{10}+46v^9+2v^9q+242v^8+22v^8q+118v^7q+736v^7+383v^6q+2q^2v^6`$ (306)
$`+1388v^6+797v^5q+16q^2v^5+1638v^5+1071v^4q+1164v^4+56v^4q^2+920v^3q+384v^3`$
$`+109v^3q^2+396v^2q+141v^2q^2+94q^2v+3q^3v+6q^3)`$
$`f_{43}`$ $`=`$ $`v^6(v+2)(1+v)^2(33q^3v+48q^3v^2+111q^2v^5+42v^3q^3+834v^6q+19q^2v^6+4v^9q`$ (311)
$`+3v^{10}+262v^7q+46v^8q+348v^3q^2+277v^2q^2+3q^3v^5+18q^3v^4+100q^2v+360v^2q`$
$`+1836v^4q+1248v^3q+1277v^5+812v^4+12q^3+1585v^5q+196v^3+602v^7+1127v^6`$
$`+35v^9+194v^8+268v^4q^2)`$
$`f_{44}`$ $`=`$ $`v^9(v+2)^2(1+v)^5(6v^6+18v^5q+10v^5+34v^4q+18v^4q^2+10v^4+43v^3q^2+4v^3`$ (314)
$`+6v^3q^3+25v^3q+10v^2q+18q^3v^2+34v^2q^2+18q^3v+10q^2v+6q^3)`$
with roots $`\lambda _{tt3,j}`$ for $`6j9`$ and
$$\xi ^6+f_{61}\xi ^5+f_{62}\xi ^4+f_{63}\xi ^3+f_{64}\xi ^2+f_{65}\xi +f_{66}=0$$
(315)
where
$`f_{61}`$ $`=`$ $`v(7v^3+2v^4+12v^2+q^2+v^3q+7vq+3v^2q)`$ (316)
$`f_{62}`$ $`=`$ $`v^2(4v^8+168v^4+208v^5+117v^6+32v^7+3v^7q+q^4+q^2v^6+14q^3v+2v^3q^3+6q^3v^2`$ (319)
$`+73v^2q^2+21v^6q+26v^4q^2+168v^4q+172v^3q+55v^3q^2+6q^2v^5+84v^5q)`$
$`f_{63}`$ $`=`$ $`v^6(1+v)(2v^7q8v^7+11v^6q+q^2v^660v^6+24v^5q+7q^2v^5200v^510v^4q`$ (323)
$`336v^4+27v^4q^2+2v^3q^3+47v^3q^2116v^3q272v^3200v^2q+7q^3v^2+32v^2q^2`$
$`+13q^3v48q^2v+q^44q^3)`$
$`f_{64}`$ $`=`$ $`v^8(1+v)^2(22q^3v+9q^3v^2+38q^2v^5+17v^3q^3+v^8q^2212v^6q+28q^2v^6+2q^4`$ (327)
$`+8v^7q^264v^7q8v^8q14v^3q^2+104v^2q^2+q^4v^2+2q^3v^5+9q^3v^44v^4q+q^4v`$
$`+344v^3q+872v^5+672v^4288v^5q+136v^7+492v^6+16v^8+3v^4q^2)`$
$`f_{65}`$ $`=`$ $`v^{12}(q4)(v+2)(1+v)^3(v^3+4v^2+6v+q)(q+2v)^2`$ (329)
$`f_{66}`$ $`=`$ $`v^{14}(v+2)^2(1+v)^4(q+2v)^4`$ (331)
with roots $`\lambda _{tt3,j}`$ for $`10j15`$. Their coefficients are equal to $`b^{(1)}=q1`$. For $`d=2`$, we have
$$T_{Z,tri,3,2}=v^2\left(\begin{array}{cccccccccccc}t_{19}& v& t_{12}& 0& t_{19}& v& v_1v_2& 0& v_1& v_1& v_1v_2& 0\\ v& t_{19}& 0& t_{12}& v& t_{19}& 0& v_1v_2& v_1& v_1& 0& v_1v_2\\ v& t_{19}& t_{19}& v& t_{12}& 0& v_1& v_1& 0& v_1v_2& v_1v_2& 0\\ t_{19}& v& v& t_{19}& 0& t_{12}& v_1& v_1& v_1v_2& 0& 0& v_1v_2\\ 0& t_{12}& v& t_{19}& t_{19}& v& 0& v_1v_2& 0& v_1v_2& v_1& v_1\\ t_{12}& 0& t_{19}& v& v& t_{19}& v_1v_2& 0& v_1v_2& 0& v_1& v_1\\ v^2v_3& v^2v_3& vt_{12}& 0& vt_{12}& 0& vv_1v_2& 0& 0& vv_1v_2& vv_1v_2& 0\\ v^2v_3& v^2v_3& 0& vt_{12}& 0& vt_{12}& 0& vv_1v_2& vv_1v_2& 0& 0& vv_1v_2\\ vt_{12}& 0& vt_{12}& 0& v^2v_3& v^2v_3& vv_1v_2& 0& vv_1v_2& 0& vv_1v_2& 0\\ 0& vt_{12}& 0& vt_{12}& v^2v_3& v^2v_3& 0& vv_1v_2& 0& vv_1v_2& 0& vv_1v_2\\ 0& vt_{12}& v^2v_3& v^2v_3& vt_{12}& 0& 0& vv_1v_2& 0& vv_1v_2& vv_1v_2& 0\\ vt_{12}& 0& v^2v_3& v^2v_3& 0& vt_{12}& vv_1v_2& 0& vv_1v_2& 0& 0& vv_1v_2\end{array}\right)$$
(332)
where
$$t_{19}=q+5v+v^2.$$
(333)
The characteristic equation $`T_{Z,tri,3,2}`$ yields two linear equations, a quadratic equation and two quartic equations. The eigenvalues which are the solutions of the two linear equations, $`\lambda _{tt3,16}=v^3(v+1)(v+2)`$ and $`\lambda _{tt3,17}=v^2(2v+q)`$, have the coefficient $`b_2^{(2)}=(q1)(q2)/2`$. One of the quartic equations is given by
$`\xi ^4v^2(9v^2+2v^3+14v+2q)\xi ^3`$ (334)
$`+v^4(36v^5+135v^4+246v^3+192v^2+4v^6+7v^3q+33v^2q+54vq+4q^2)\xi ^2`$ (335)
$`v^7(v+2)(1+v)(9v^2+2v^3+14v+2q)(q+2v)\xi +v^{10}(v+2)^2(1+v)^2(q+2v)^2=0`$ (336)
with roots $`\lambda _{tt3,j}`$ for $`18j21`$ that also have $`b_2^{(2)}`$ as their coefficient. The roots of the following two equations have coefficient $`b_1^{(2)}=q(q3)/2`$:
$$\xi ^2v^2(3q+24v+13v^2+3v^3)\xi +v^5(1+v)(6v^2+5vq+6q)=0$$
(338)
with roots $`\lambda _{tt3,(22,23)}`$ and
$$\xi ^4+v^4\xi ^3+v^6(q+vq+v^2)\xi ^2v^{10}q(1+v)\xi +v^{12}q^2(1+v)^2=0$$
(339)
with roots $`\lambda _{tt3,j}`$ for $`24j27`$. For $`d=2`$, we have
$$T_{Z,tri,3,3}=v^3\left(\begin{array}{cccccc}1& 0& 0& 1& 0& 0\\ 0& 1& 1& 0& 0& 0\\ 0& 0& 1& 0& 0& 1\\ 0& 0& 0& 1& 1& 0\\ 1& 0& 0& 0& 1& 0\\ 0& 1& 0& 0& 0& 1\end{array}\right).$$
(340)
All the eigenvalues of $`T_{Z,tri,3,3}`$ have multiplicity two. These eigenvalues are $`\lambda _{tt3,28}=2v^3`$ with coefficient $`b_1^{(3)}=(q1)(q^25q+3)/3`$ and the roots of $`\xi ^2v^3\xi +v^6=0`$, denoted as $`\lambda _{tt3,(29,30)}`$, with coefficient $`2b_2^{(3)}=q(q2)(q4)/3`$. Therefore, the Potts model partition function for the triangular lattice strip with $`L_y=3`$ and toroidal boundary condition is given by
$`Z(tri,3\times L_x,tor.,q,v)`$ $`=`$ $`b^{(0)}{\displaystyle \underset{j=1}{\overset{5}{}}}(\lambda _{tt3,j})^{L_x}+b^{(1)}{\displaystyle \underset{j=6}{\overset{15}{}}}(\lambda _{tt3,j})^{L_x}+b_2^{(2)}{\displaystyle \underset{j=16}{\overset{21}{}}}(\lambda _{tt3,j})^{L_x}`$ (343)
$`+b_1^{(2)}{\displaystyle \underset{j=22}{\overset{27}{}}}(\lambda _{tt3,j})^{L_x}+b_1^{(3)}(\lambda _{tt3,28})^{L_x}+2b_2^{(3)}{\displaystyle \underset{j=29}{\overset{30}{}}}(\lambda _{tt3,j})^{L_x}.`$
For the corresponding Klein bottle strip, the coefficients for $`\lambda _{tt3,j}`$ with $`10j15`$, $`18j21`$, $`24j27`$ and $`j=4,5,29,30`$ become zero. The coefficient for $`\lambda _{tt3,17}`$ changes sign and the coefficient for $`\lambda _{tt3,28}`$ becomes $`b^{(1)}=1q`$. Consequently, the Potts model partition function for the triangular lattice strip with $`L_y=3`$ and Klein bottle boundary condition is given by
$`Z(tri,3\times L_x,Kb.,q,v)`$ $`=`$ $`b^{(0)}{\displaystyle \underset{j=1}{\overset{3}{}}}(\lambda _{tt3,j})^{L_x}+b^{(1)}\left({\displaystyle \underset{j=6}{\overset{9}{}}}(\lambda _{tt3,j})^{L_x}(\lambda _{tt3,28})^{L_x}\right)`$ (346)
$`+b_2^{(2)}\left((\lambda _{tt3,16})^{L_x}(\lambda _{tt3,17})^{L_x}\right)+b_1^{(2)}{\displaystyle \underset{j=22}{\overset{23}{}}}(\lambda _{tt3,j})^{L_x}.`$
At the special value $`v=1`$, the Potts model partition function reduces to the chromatic polynomial. The non-zero eigenvalues are the same as $`\lambda _{tt,j}`$ for $`1j11`$ from eqs. (4.2) to (4.9) in t .
### VI.6 Honeycomb-Lattice Strip, $`L_y=4`$
For the $`L_y=4`$ toroidal and Klein bottle strips of the honeycomb lattice, the linear sizes of the matrices are 14, 35, 56, 48 for levels $`0d3`$, as indicated in Table 1. The eigenvalues of $`T_{Z,hc,4,0}`$ are roots of two quadratic equations and an eighth-degree equation. One of the quadratic equations has multiplicity two. The eigenvalues of $`T_{Z,hc,4,1}`$ are roots of a seventh-degree equation, an eighth-degree equation and a twelfth-degree equation. The eighth-degree equation has multiplicity two. The eigenvalues of $`T_{Z,hc,4,2}`$ are roots of a quartic equation, a fifth-degree equation, a sixth-degree equation, two eighth-degree equations and a eleventh-degree equation. The sixth-degree equation and one of the eighth-degree equations have multiplicity two. The eigenvalues of $`T_{Z,hc,4,3}`$ consist of a linear term $`v^6(v+q)^2`$, roots of a quadratic equation, a cubic equation and two quartic equations. The linear terms and the cubic equation have multiplicity two, and the other equations have multiplicity four. The matrix $`T_{T,hc,4,4}=v^8`$ has been given above with coefficient $`b^{(4)}`$. Some of these equations are too lengthy to list here. They are available from the authors.
At the special value $`v=1`$, the Potts model partition function reduces to the chromatic polynomial. The linear sizes of the matrices reduce to 6, 18, 34, 36 for levels $`0d3`$, as indicated in Table 6. The eigenvalues of $`T_{P,hc,4,0}`$ consist of two linear terms and roots of a quartic equation. They are given by
$$\lambda _{ht4,0}=0$$
(347)
$$\lambda _{ht4,1}=(q1)^2(q3)^2$$
(348)
$`\xi ^4(q^812q^7+66q^6220q^5+496q^4796q^3+922q^2734q+314)\xi ^3`$ (349)
$`+(q^{12}20q^{11}+188q^{10}1092q^9+4344q^812428q^7+26192q^641022q^5+47597q^4`$ (350)
$`40318q^3+24247q^29782q+2169)\xi ^2`$ (351)
$`(4q^{12}80q^{11}+736q^{10}4124q^9+15705q^842922q^7+86535q^6129964q^5`$ (352)
$`+144575q^4116374q^3+64525q^222280q+3680)\xi +4(q1)^4(q2)^4=0`$ (353)
with roots $`\lambda _{ht4,j}`$ for $`2j5`$. $`\lambda _{ht4,0}=0`$ appears as one of the eigenvalues of $`T_{P,hc,4,1}`$ with multiplicity two. The remaining eigenvalues of $`T_{P,hc,4,1}`$ are roots of a cubic equation, a quartic equation and a sixth-degree equation. The cubic equation has multiplicity two. These equations are given by
$`\xi ^3(q^610q^5+44q^4110q^3+167q^2150q+68)\xi ^2`$ (354)
$`+(q^814q^7+90q^6342q^5+832q^41322q^3+1345q^2812q+230)\xi `$ (355)
$`(q1)^2(2q^622q^5+102q^4256q^3+373q^2308q+115)=0`$ (356)
with roots $`\lambda _{ht4,j}`$ for $`j=6,7,8`$.
$`\xi ^4(q^68q^5+29q^464q^3+95q^286q+37)\xi ^3`$ (357)
$`+(q1)^2(q^814q^7+86q^6302q^5+672q^41000q^3+1030q^2724q+280)\xi ^2`$ (358)
$`(q1)^4(q^816q^7+112q^6444q^5+1084q^41664q^3+1581q^2886q+253)\xi `$ (359)
$`+(q3)^2(q1)^6=0`$ (360)
with roots $`\lambda _{ht4,j}`$ for $`9j12`$.
$`\xi ^6(q^612q^5+65q^4204q^3+407q^2510q+321)\xi ^5+(q^{10}20q^9+187q^8`$ (361)
$`1064q^7+4046q^610670q^5+19694q^425280q^3+22278q^213120q+4368)\xi ^4`$ (362)
$`(q^{12}24q^{11}+274q^{10}1948q^9+9541q^833732q^7+87865q^6169306q^5+239311q^4`$ (363)
$`243300q^3+172063q^279202q+19177)\xi ^3+(4q^{12}92q^{11}+992q^{10}6588q^9`$ (364)
$`+29885q^897284q^7+232764q^6412788q^5+540290q^4512128q^3+336720q^2`$ (365)
$`139512q+27993)\xi ^24(q1)^2(q^{10}20q^9+182q^8984q^7+3483q^68402q^5`$ (366)
$`+13956q^415756q^3+11602q^25078q+1032)\xi +16(q2)^2(q1)^6=0`$ (367)
with roots $`\lambda _{ht4,j}`$ for $`13j18`$. $`\lambda _{ht4,0}=0`$ is again one of the eigenvalues of $`T_{P,hc,4,2}`$. The remaining eigenvalues of $`T_{P,hc,4,2}`$ are roots of a quadratic equation, two cubic equations, two fifth-degree equations and a seventh-degree equation. They are given by
$$\lambda _{ht4,(19,20)}=\frac{1}{2}\left[q^24q+7\pm (q3)(q^22q+5)^{1/2}\right]$$
(368)
$`\xi ^5(2q^414q^3+41q^260q+41)\xi ^4+(q^814q^7+91q^6358q^5+932q^41636q^3`$ (369)
$`+1875q^21270q+396)\xi ^3(q^{10}18q^9+148q^8726q^7+2339q^65148q^5+7824q^4`$ (370)
$`8136q^3+5615q^22372q+477)\xi ^2+(q1)^2(q^814q^7+85q^6292q^5+626q^4`$ (371)
$`874q^3+800q^2448q+125)\xi 4(q1)^4=0`$ (372)
with roots $`\lambda _{ht4,j}`$ for $`21j25`$.
$`\xi ^5(2q^415q^3+51q^290q+73)\xi ^4+(q^815q^7+104q^6429q^5+1152q^42064q^3`$ (373)
$`+2415q^21694q+570)\xi ^3(q^{10}19q^9+167q^8881q^7+3070q^67372q^5+12396q^4`$ (374)
$`14520q^3+11459q^25550q+1265)\xi ^2+(q1)^2(q^815q^7+104q^6428q^5+1144q^4`$ (375)
$`2039q^3+2392q^21706q+575)\xi 2q(q1)^6=0`$ (376)
with roots $`\lambda _{ht4,j}`$ for $`26j30`$.
$`\xi ^3(q^47q^3+21q^232q+24)\xi ^2+(q^611q^5+53q^4139q^3+206q^2164q+58)\xi `$ (377)
$`(q1)^2(q^47q^3+18q^224q+15)=0`$ (378)
with roots $`\lambda _{ht4,j}`$ for $`j=31,32,33`$.
$`\xi ^3(q^48q^3+27q^240q+24)\xi ^2+(q1)^2(q^410q^3+44q^290q+72)\xi `$ (379)
$`(q3)^2(q1)^4=0`$ (380)
with roots $`\lambda _{ht4,j}`$ for $`j=34,35,36`$.
$`\xi ^7(3q^422q^3+75q^2140q+126)\xi ^6+(3q^844q^7+304q^61282q^5+3601q^4`$ (381)
$`6890q^3+8830q^27038q+2781)\xi ^5(q^{12}22q^{11}+230q^{10}1504q^9+6839q^8`$ (382)
$`22756q^7+56785q^6107160q^5+152314q^4160280q^3+120302q^259250q+15045)\xi ^4`$ (383)
$`+(q^{14}26q^{13}+318q^{12}2414q^{11}+12685q^{10}48806q^9+141948q^8317702q^7`$ (384)
$`+551665q^6743122q^5+769969q^4601782q^3+341290q^2128886q+25246)\xi ^3`$ (385)
$`(q^{14}24q^{13}+280q^{12}2088q^{11}+11028q^{10}43286q^9+129310q^8297078q^7`$ (386)
$`+525432q^6710236q^5+721717q^4535598q^3+276095q^289634q+14145)\xi ^2`$ (387)
$`+4(q1)^2(q^96q^814q^7+249q^61080q^5+2522q^43564q^3+3069q^21499q`$ (388)
$`+326)\xi 4(q2)^2(q1)^6=0`$ (389)
with roots $`\lambda _{ht4,j}`$ for $`37j43`$. The second fifth- degree equation and the first cubic equation given above have multiplicity two. $`\lambda _{ht4,0}=0`$ is again one of the eigenvalues of $`T_{P,hc,4,3}`$ with multiplicity two. The remaining eigenvalues of $`T_{P,hc,4,3}`$ are two linear terms, roots of a quadratic equation and two cubic equations. They are given by
$$\lambda _{ht4,44}=q^24q+5$$
(390)
$$\lambda _{ht4,45}=(q1)^2$$
(391)
$$\lambda _{ht4,(46,47)}=\frac{1}{2}\left[q^26q+13\pm (q3)(q^26q+17)^{1/2}\right]$$
(392)
$$\xi ^3(2q^28q+11)\xi ^2+(q^48q^3+26q^238q+23)\xi (q1)^2=0$$
(393)
with roots $`\lambda _{ht4,j}`$ for $`j=48,49,50`$.
$$\xi ^3(2q^28q+13)\xi ^2+(q^48q^3+26q^234q+19)\xi 3(q1)^2=0$$
(394)
with roots $`\lambda _{ht4,j}`$ for $`j=51,52,53`$. $`\lambda _{ht4,j}`$ for $`j=45,46,47`$ have multiplicity two, and the other eigenvalues of $`T_{P,hc,4,3}`$ have multiplicity four. Finally, we have $`\lambda _{ht4,54}=1`$ as the eigenvalue of $`T_{P,hc,4,4}`$. The corresponding coefficients are
$`b_{ht4,j}`$ $`=`$ $`1\text{for}1j5`$ (395)
$`b_{ht4,j}`$ $`=`$ $`2(q1)\text{for}6j8`$ (397)
$`b_{ht4,j}`$ $`=`$ $`q1\text{for}9j18`$ (399)
$`b_{ht4,j}`$ $`=`$ $`{\displaystyle \frac{1}{2}}(q1)(q2)\text{for}19j25`$ (401)
$`b_{ht4,j}`$ $`=`$ $`(q1)(q2)\text{for}26j30`$ (403)
$`b_{ht4,j}`$ $`=`$ $`q(q3)\text{for}31j33`$ (405)
$`b_{ht4,j}`$ $`=`$ $`{\displaystyle \frac{1}{2}}q(q3)\text{for}34j43`$ (407)
$`b_{ht4,44}`$ $`=`$ $`{\displaystyle \frac{2}{3}}(q1)(q^25q+3)`$ (409)
$`b_{ht4,j}`$ $`=`$ $`{\displaystyle \frac{1}{3}}(q1)(q^25q+3)\text{for}45j47`$ (411)
$`b_{ht4,j}`$ $`=`$ $`{\displaystyle \frac{2}{3}}q(q2)(q4)\text{for}48j53`$ (413)
$`b_{ht4,54}`$ $`=`$ $`q^48q^3+20q^215q+1.`$ (415)
Although $`\lambda _{ht4,0}=0`$ does not contribute the the chromatic polynomial, we can still calculate its coefficient for the toroidal strip, and we get $`(q1)(2q^27q+12)/6`$. The sum of all the coefficients is equal to $`q^2(q1)^2`$, as dictated by the $`L_y=4`$ special case of the general result (174). Hence, the chromatic polynomial for the honeycomb lattice strip with $`L_y=4`$ and toroidal boundary condition is given by
$`P(hc,4\times L_x,tor.,q)`$ $`=`$ $`b^{(0)}{\displaystyle \underset{j=1}{\overset{5}{}}}(\lambda _{ht4,j})^m+b^{(1)}\left(2{\displaystyle \underset{j=6}{\overset{8}{}}}(\lambda _{ht4,j})^m+{\displaystyle \underset{j=9}{\overset{18}{}}}(\lambda _{ht4,j})^m\right)`$ (424)
$`+b_2^{(2)}\left({\displaystyle \underset{j=19}{\overset{25}{}}}(\lambda _{ht4,j})^m+2{\displaystyle \underset{j=26}{\overset{30}{}}}(\lambda _{ht4,j})^m\right)`$
$`+b_1^{(2)}\left(2{\displaystyle \underset{j=31}{\overset{33}{}}}(\lambda _{ht4,j})^m+{\displaystyle \underset{j=34}{\overset{43}{}}}(\lambda _{ht4,j})^m\right)`$
$`+b_1^{(3)}\left(2(\lambda _{ht4,44})^m+{\displaystyle \underset{j=45}{\overset{47}{}}}(\lambda _{ht4,j})^m\right)`$
$`+4b_2^{(3)}{\displaystyle \underset{j=48}{\overset{53}{}}}(\lambda _{ht4,j})^m+b^{(4)}(\lambda _{ht4,54})^m`$
where $`m=L_x/2`$ as given in eq. (4). For the corresponding Klein bottle strip, $`\lambda _{ht4,j}`$ for $`6j8`$, $`26j33`$, $`48j53`$ and $`j=44`$ do not contribute. The coefficient for $`\lambda _{ht4,j}`$ with $`9j12`$, $`21j25`$, $`34j36`$ and $`j=1`$ changes sign. The coefficient for $`\lambda _{ht4,45}`$ becomes $`b^{(1)}=q1`$, the coefficients for $`\lambda _{ht4,(46,47)}`$ become $`b^{(1)}=1q`$ and the coefficients for $`\lambda _{ht4,54}`$ become one. Therefore, the chromatic polynomial for the honeycomb lattice strip with $`L_y=4`$ and Klein bottle boundary condition is given by
$`P(hc,4\times L_x,Kb.,q)`$ $`=`$ $`b^{(0)}\left((\lambda _{ht4,1})^m+{\displaystyle \underset{j=2}{\overset{5}{}}}(\lambda _{ht4,j})^m+(\lambda _{ht4,54})^m\right)`$ (431)
$`+b^{(1)}\left({\displaystyle \underset{j=9}{\overset{12}{}}}(\lambda _{ht4,j})^m+{\displaystyle \underset{j=13}{\overset{18}{}}}(\lambda _{ht4,j})^m+(\lambda _{ht4,45})^m{\displaystyle \underset{j=46}{\overset{47}{}}}(\lambda _{ht4,j})^m\right)`$
$`+b_2^{(2)}\left({\displaystyle \underset{j=19}{\overset{20}{}}}(\lambda _{ht4,j})^m{\displaystyle \underset{j=21}{\overset{25}{}}}(\lambda _{ht4,j})^m\right)`$
$`+b_1^{(2)}\left({\displaystyle \underset{j=34}{\overset{36}{}}}(\lambda _{ht4,j})^m+{\displaystyle \underset{j=37}{\overset{43}{}}}(\lambda _{ht4,j})^m\right).`$
The singular locus $``$ for the $`L_x\mathrm{}`$ limit of the strip of the honeycomb lattice with $`L_y=4`$ and toroidal boundary conditions is shown in Fig. 10. For comparison, chromatic zeros are calculated and shown for length $`L_x=13`$ (i.e., $`n=104`$ vertices). The locus $``$ crosses the real axis at the points $`q=0`$, $`q=2`$, and at the maximal point $`q=q_c`$, where
$$q_c2.480749\mathrm{for}\{G\}=(hc,4\times \mathrm{},PBC_y,PBC_x).$$
(432)
We note that this is rather close to the value for the infinite 2D honeycomb lattice, namely $`q_c(hc)=(3+\sqrt{5})/22.61803`$. As is evident from Fig. 10, the locus $``$ separates the $`q`$ plane into different regions including the following: (i) region $`R_1`$, containing the semi-infinite intervals $`q>q_c`$ and $`q<0`$ on the real axis and extending outward to infinite $`|q|`$; (ii) $`R_2`$ containing the interval $`2<q<q_c`$; (iii) $`R_3`$ containing the real interval $`0<q<2`$; and (iv) the complex-conjugate pair $`R_4,R_4^{}`$ centered approximately at $`q=2.55\pm 0.6i`$. The (nonzero) densities of chromatic zeros have the smallest values on the curve separating regions $`R_1`$ and $`R_3`$ in the vicinity of the point $`q=0`$ and on the curve separating regions $`R_2`$ and $`R_3`$ in the vicinity of the point $`q=2`$.
It is instructive to compare this locus $``$ in Fig. 10 with the corresponding loci that we obtained previously for strips of the honeycomb lattice with periodic longitudinal boundary conditions and free or periodic transverse boundary conditions. We recall that for a given infinite-length lattice strip with specified transverse boundary conditions, the locus $``$ is the same for periodic longitudinal and twisted periodic longitudinal boundary conditions (i.e., cyclic versus Möbius if the transverse B.C. are free, and toroidal versus Klein bottle if the transverse B.C. are periodic). For the cyclic strip of the honeycomb lattice, $``$ was given as (i) Fig. 1 of Ref. pg for width $`L_y=2`$, and (ii) Figs. 7, 8, and 9 of pt for $`L_y=3,4,5`$ (see also the zeros in Fig. 17 of Ref. hca for $`L_y=3`$). We recall that $`L_y`$ must be even for strips of the honeycomb lattice with toroidal boundary conditions. Comparing our results for $``$ for the width $`L_y=4`$ toroidal strip of the honeycomb lattice in Fig. 10 with our earlier results for the $`L_y=4`$ cyclic strip of this lattice in Fig. 8 of pt , we see that the locus $``$ is significantly simpler for toroidal, versus cyclic, boundary conditions, with fewer regions and, in particular, no evidence of the tiny sliver phases that are present in the latter case. Furthermore, the value of $`q_c2.481`$ that we obtain (eq. (432)) for the width $`L_y=4`$ infinite-length strip of the honeycomb lattice with toroidal boundary conditions is closer to the value of $`q_c(hc)2.618`$ for the infinite 2D than the value $`q_c2.155`$ for the corresponding cyclic $`L_y=4`$ strip. This is in agreement with one’s general expectations, since toroidal B.C. minimize finite-size effects to a greater extent than cyclic B.C. and hence expedite the approach to the $`L_y\mathrm{}`$ limit.
In region $`R_1`$, the dominant eigenvalue is the root of eq. (353) with largest magnitude. Denoting it as $`\lambda _{ht4,2}`$, we have
$$W=(\lambda _{ht4,2})^{1/8},qR_1.$$
(433)
This is the same as $`W`$ for the corresponding $`L_x\mathrm{}`$ limit of the strip of the honeycomb lattice with the same width $`L_y=4`$ and cylindrical $`(PBC_y,FBC_x)`$ boundary conditions, calculated in hca . This equality of the $`W`$ functions for the $`L_x\mathrm{}`$ limit of two strips of a given lattice with the same transverse boundary conditions and different longitudinal boundary conditions in the more restrictive region $`R_1`$ defined by the two boundary conditions is a general result wcy ; bcc . In hca we have listed values of $`W`$ for a range of values of $`q`$ for the $`L_x\mathrm{}`$ limit of various strips of the honeycomb lattice, including $`(hc,4\times \mathrm{},PBC_y,FBC_x)`$. Since $`W`$ is independent of $`BC_x`$ for $`q`$ in the more restrictive region $`R_1`$ defined by $`FBC_x`$ and $`(T)PBC_x`$ (which is the $`R_1`$ defined by $`PBC_x`$ here), it follows, in particular, that
$$W(4\times \mathrm{},PBC_y,(T)PBC_x,q)=W(4\times \mathrm{},PBC_y,FBC_x,q)\mathrm{for}q2.48\mathrm{}$$
(434)
In region $`R_2`$, the largest root of the seventh degree equation (389) is dominant; we label this as $`\lambda _{ht4,37}`$, so that
$$|W|=|\lambda _{ht4,37}|^{1/8},qR_2$$
(435)
(in regions other than $`R_1`$, only $`|W|`$ can be determined unambiguously w ). Thus, $`q_c`$ is the relevant solution of the equation of degeneracy in magnitude $`|\lambda _{ht4,2}|=|\lambda _{ht4,37}|`$. In region $`R_3`$,
$$|W|=|\lambda _{ht4,13}|^{1/8},qR_3$$
(436)
where $`\lambda _{ht4,13}`$ is the largest root of the sixth degree equation (367). In regions $`R_4,R_4^{}`$,
$$|W|=|\lambda _{ht4,9}|^{1/8},qR_4,R_4^{}$$
(437)
where $`\lambda _{ht4,9}`$ is the largest root of the quartic equation (360). We calculate the following specific values: $`|W(q=1)|=1.67290..`$, $`|W(q=2)|=1.35817..`$, and $`|W(q=q_c)|=1.26949..`$ These values may be compared with those listed in Ref. hca . As is evident from Fig. 10, the curve $``$ has support for $`Re(q)<0`$. For strips with finite lengths $`m3`$, certain chromatic zeros have support for $`Re(q)<0`$.
## VII Summary
In summary, in this paper we have presented a method for calculating transfer matrices for the $`q`$-state Potts model partition functions $`Z(G,q,v)`$, for arbitrary $`q`$ and temperature variable $`v`$, on strip graphs $`G`$ of the square, triangular, and honeycomb lattices with width $`L_y`$ and arbitrarily great length $`L_x`$, having toroidal and Klein bottle boundary conditions. Using this method, we have derived a number of general properties of these transfer matrices. In particular, we have found some very simple formulas for the determinant $`det(T_{Z,\mathrm{\Lambda },L_y,d})`$, and trace $`Tr(T_{Z,\mathrm{\Lambda },L_y})`$. A number of explicit exact calculations were given.
Acknowledgments: The research of R.S. was partially supported by the NSF grant PHY-00-98527. The research of S.C.C. was partially supported by the NSC grant NSC-93-2119-M-006-009.
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# Exciton Dissociation Dynamics in Model Donor-Acceptor Polymer Heterojunctions: I. Energetics and Spectra
## I Introduction
Organic semiconducting polymers are currently of broad interest as potential low-cost materials for photovoltaic and light-emitting display applications. While research into the fabrication, chemistry, and fundamental physics of these materials continues to advance, devices based upon organic semiconductors have begun to appear in the market place. Composite materials fabricated by mixing semiconducting polymers with different electron and hole accepting properties have been used to make highly efficient solar cells and photovoltaic cells.Tang (1986); M. Granström et al. (1998); Bach et al. (1998) The advantage gained by mixing materials is that when a conjugated polymer material absorbs light, the primary species produced is an exciton, i.e. an electron/hole pair bound by Coulombic attraction, rather than free charge carriers. As a result, a single layer conventional Schottky-barrier cell of a molecular semiconductor sandwiched between two metal electrodes has poor power conversion and charge collection efficiencies. In composite materials, the driving force for charge separation is due to the mismatch between the HOMO and LUMO levels of the blended materials. This creates a band off-set at the interface between the two materials. A type I heterojunction is when the HOMO and LUMO levels of one of the materials both lie within the HOMO-LUMO gap of the other. A type II heterojunction occurs when both the HOMO and LUMO of one material are simultaneously shifted up or down relative to the HOMO and LUMO levels of the other material. By and large, organic semiconductor heterojunctions fall into the type II category.Shockley (1950) The chemical structure of some of the more important materials for device applications are shown in Fig. 1.
Photocurrent production within a type II material hinges upon its ability to dissociate the bound electron/hole pairs produced by photoexcitation. To a first approximation, the dissociation threshold is determined by the ratio of the exciton binding energy, $`\epsilon _B`$, to the band off-set at the heterojunction, $`\mathrm{\Delta }E`$. If $`\mathrm{\Delta }E>\epsilon _B`$, the exciton will be energetically unstable and can undergo fission to form charge-separated states and eventually free charge carriers. On the other hand, when $`\mathrm{\Delta }E<\epsilon _B`$, the exciton remains the favorable species. For typical organic semiconductors, $`\epsilon _B0.50.`$6 eV. In heterojunctions composed of poly(benzimidazobenzophenanthroline ladder) (BBL) and poly(1,4-phenylenevinylene) (PPV), $`\mathrm{\Delta }E`$ is considerably greater than $`\epsilon _B`$Alam and Jenekhe (2004). As a result, recent reports of devices fabricated from nanoscale layers of BBL and PPV or its derivative poly(2-methoxy-5(2”-ethyl-hexyloxy)-1,4-phenylenevinylene) (MEH-PPV) indicate very efficient photoinduced charge transfer and large photovoltaic power conversion efficiencies under white light illumination.Yu et al. (1995); Halls et al. (1995, 1999)
Surprisingly, type II materials can also be used to produce highly efficient light-emitting diodes (LEDs). In demixed blends of poly(9,9-dioctylfluorene-co-benzothiadiazole) (F8BT) and poly(9,9-dioctylfluorene-co-N-(4-butylphenyl) diphenylamine) (TFB), the exciton is the stable species and fabricated devices show excellent LED performanceMorteani et al. (2003). While chemically similar to TFB, blends of F8BT with poly(9,9-dioctylfluorene-co-bis-N,N-(4 -butylphenyl)-bis-N,N-phenyl-1,4-phenylenediamine) (PFB) exhibit very poor LED performance. Changing the diphenylamine to a phenylenediamine within the co-polymer shifts the position of the HOMO energy level by +0.79 eV relative to the HOMO level of TFB. (Table 1) and destabilizes the exciton.
In this paper we investigate the role of interchain coupling and lattice reorganization in determining the essential states model polymer semiconductor systems consisting of parallel chains of PPV:BBL, TFB:F8BT, and PFB:F8BT. In the case of the PPV:BBL heterojunction, the exciton binding energy is smaller than the band off-set leading to electron transter from the photoexcited PPV to the BBL chain. In the case of TFB:F8BT and PFB:F8BT junctions, the band off-sets are close to the exciton binding energy and we see considerable mixing between intramolecular excitonic states with intermolecular charge-transfer or exciplex states. Moreover, in TFB:F8BT, avoided crossing between Born Oppenhemier potential energy surfaces may open channels for exciton regeneration that are thermally accessible from the exciplex states as recently reported by Morteani et al.Morteani et al. (2004). In a subsequent paper we will report upon our simulations of the charge-injection and photoexcitation dynamics within these systems. Bittner et al. (2005)
## II Theoretical Approach
Our basic description is derived starting from a model for the on-chain electronic excitations of a single conjugated polymer chain.Karabunarliev and Bittner (2003a); Karabunarliev and Bittner (2004); Karabunarliev and Bittner (2003b) This model accounts for the coupling of excitations within the $`\pi `$-orbitals of a conjugated polymer to the lattice phonons using localized valence and conduction band Wannier functions ($`|\overline{h}`$ and $`|p`$) to describe the $`\pi `$ orbitals and two optical phonon branches to describe the bond stretches and torsions of the the polymer skeleton.
$`H`$ $`=`$ $`{\displaystyle \underset{\mathrm{𝐦𝐧}}{}}(F_{\mathrm{𝐦𝐧}}^{}+V_{\mathrm{𝐦𝐧}})A_𝐦^{}A_𝐧`$ (1)
$`+`$ $`{\displaystyle \underset{\mathrm{𝐧𝐦}i\mu }{}}\left({\displaystyle \frac{F_{\mathrm{𝐧𝐦}}^{}}{q_{i\mu }}}\right)A_𝐧^{}A_𝐦q_{i\mu }`$
$`+`$ $`{\displaystyle \underset{i\mu }{}}\omega _\mu ^2(q_{i\mu }^2+\lambda _\mu q_{i\mu }q_{i+1,\mu })+p_{i\mu }^2`$
where $`A_𝐧^{}`$ and $`A_𝐧`$ are Fermion operators that act upon the ground electronic state $`|0`$ to create and destroy electron/hole configurations $`|n=|\overline{h}p`$ with positive hole in the valence band Wannier function localized at $`h`$ and an electron in the conduction band Wannier function $`p`$. Finally, $`q_{i\mu }`$ and $`p_{i\mu }`$ correspond to lattice distortions and momentum components in the $`i`$-th site and $`\mu `$-th optical phonon branch. Except as noted below, our model and approach is identical to what we have described in our earlier publications. Karabunarliev and Bittner (2003a); Karabunarliev and Bittner (2004); Karabunarliev and Bittner (2003c)
Since we are dealing with multiple polymer chains, we consider both interchain and intrachain single particle contributions. For the intrachain terms, we use the hopping terms and site energies derived for isolated polymer chains of a given species, $`t_{i,||}`$, where our notation denotes the parallel hopping term for the $`i`$th chain ($`i=1,2`$). For PPV and similar conjugated polymer species, these are approximately $`0.5`$eV for both valence and conduction $`\pi `$ bands.
We include two intramolecular optical phonon branches which correspond roughly to the high-frequency C=C bond stretching modes within a given repeat unit and a second low-frequency mode, which in the case of PPV are taken to represent the phenylene torsional modes. The electron-phonon couplings are assumed to be transferable between the various chemical species. Since the modes are assumed to be intramolecular, we do not include interchain couplings in the phonon Hessian matrix.
Upon transforming $`H`$ into the diabatic representation by diagonalizing the electronic terms at $`q_{i\mu }=0`$, we obtain a series of vertical excited states $`|a_{}`$ with energies, $`\epsilon _a^{}`$ and normal modes, $`Q_\xi `$ with frequencies, $`\omega _\xi `$. (We will assume that the sum over $`\xi `$ spans all phonon branches).
$`H`$ $`=`$ $`{\displaystyle \underset{a}{}}\epsilon _a^{}|a_{}a_{}|+{\displaystyle \underset{ab\xi }{}}g_{ab\xi }^{}q_\xi (|a_{}b_{}|+|b_{}a_{}|)`$ (2)
$`+`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{\xi }{}}(\omega _\xi ^2Q_\xi ^2+P_\xi ^2).`$
The adiabatic or relaxed states can be determined then by iteratively minimizing $`\epsilon _a(Q_\xi )=a|H|a`$ according to the self-consistent equations
$`{\displaystyle \frac{d\epsilon _a(Q_\xi )}{dQ_\xi }}=g_{aa\xi }+\omega _\xi ^2Q_\xi =0.`$ (3)
Thus, each diabatic potential surface for the nuclear lattice motion is given by
$`\epsilon _a(Q_\xi )=\epsilon _a+{\displaystyle \frac{1}{2}}{\displaystyle \underset{\xi }{}}\omega _\xi ^2(Q_\xi Q_\xi ^{(a)})^2.`$ (4)
While our model accounts for the distortions in the lattice due to electron/phonon coupling, we do not account for any adiabatic change in the phonon force constants within the excited states.
While following along a relaxation pathway, crossing between diabatic states can occur. In order to insure maximum correlation between the relaxed state, $`\psi _b(Q_\xi )`$ and its parent vertical state $`\psi _a(0)`$, we compute the inner product, $`S_{ab}=\psi _a(0)|\psi _b(Q_\xi )`$ and select the state with the greatest overlap with the parent diabatic state. This allows us to follow a given state and correlate the adiabatic states with the parent vertical states. Figures 9 and 11 are examples of this. The dashed lines connect the parent vertical state with the final relaxed adiabatic state. We will comment upon these figures later in this paper.
Since we are dealing with inter-chain couplings we make the following set of assumptions. First, the single-particle coupling between chains is expected to be small compared to the intramolecular coupling. For this, we assume that the perpendicular hopping integral $`t_{}=0.01eV`$. This is consistent with LDF calculations performed by Vogl and Campbell and with the $`t_{}0.15f_1`$ estimate used in an earlier study of interchain excitons by Yu et al.Vogl and Campbell (1990); Yu et al. (1996) Furthermore, we assume that the $`J(r)`$, $`K(r)`$, and $`D(r)`$ two-particle interactions depend only upon the linear distance between two sites, as in the intrachain case. Since these are expected to be weak given that the interchain separation, $`d`$, is taken to be some what greater than the inter-monomer separation. The third approximation that we make takes into account the difference in size of the repeat units of the polymer chains. As indicated in Fig. 1, the monomeric unit in BBL is roughly two times the size as a PPV monomeric unit. Thus, for the PPV:BBL junction we adopt the modified ladder scheme show in Fig. 2.
Finally, the most important assumption that we make is that the site energies for the electrons and holes for the various chemical species can be determined by comparing the relative HOMO and LUMO energies to PPV. These are listed in Table 1. For example, the HOMO energy for PPV (as determined by its ionization potential) is -5.1eV. For BBL, this energy is -5.9 eV. Thus, we assume that the $`\overline{f}_o`$ for a hole on a BBL chain is 0.8 eV lower than $`\overline{f}_o`$ for PPV at -3.55eV. Likewise for the conduction band. The LUMO energy of PPV is -2.7eV and that of BBL is -4.0eV. Thus, we shift the band center of the BBL chain 1.3eV lower than then PPV conduction band center to 1.45eV. For the F8BT, TFB, and PFB chains, we adopt a similar scheme as discussed below. The site energies and transfer integrals used throughout are indicated in Table 1. We believe our model produces a reasonable estimate of the band off-sets in the PN-junctions formed at the interface between these semiconducting polymers.
## III Excitonic vs. Charge Separated States in Donor/Acceptor Polymer Junctions
### III.1 Electronic States of Isolated Chains
Before moving on to consider the electronic states of the heterojunction systems, we make a brief examination of the CI states of the isolated chains. The relevant data for the vertical and adiabatic excitonic and charge-separated states are given in Table 2. Since we have assumed a symmetric band structure, $`f_1=\overline{f}_1`$, the electron/hole wavefunctions will have even or odd symmetry under spatial inversion. Hence the total hamiltonian can be block diagonalized into even and odd parity blocks. The even parity states we term XT states since these are optically coupled to the ground state and are dominated by geminate electron/hole configurations. On the other hand, the odd-parity states are purely charge-separated states (CT) which possess no geminate electron/hole configurations, are optically coupled weakly at best to the ground state. We define the exciton binding energy to be the energy difference between the lowest energy XT and CT states. In each of the model systems considered, we obtain an exciton binding energy between 0.5 and 0.65V.
A uniform site model for F8BT, TFB, and PFB, may be a gross simplification of the physical systems. For example, recent semi-empirical CI calculations by Jespersen, et alJespersen et al. (2004) indicate that the lowest energy singlet excited state of F8BT consists of alternating positive and negative regions corresponding to the electron localized on the benzothiadiazole units and the hole localized on the fluorene units. These are consistent with a previous study by Cornil et al.Cornil et al. (2003) which places the LUMO on the benzothiadiazole units. Cornil et al.Cornil et al. (2003) also report the HOMO and LUMO levels of the isolated fluorene and benzothiadiazole Ref. Cornil et al. (2003) indicating a $`\mathrm{\Delta }=1.56`$ eV off-set between the fluorene and benzothiadiazole LUMO levels and a 0.66 eV off-set between the fluorene and benzothiadiazole HOMO energy levels. Similarly, PM3 level calculations at the optimized geometry indicate a LUMO offset of 1.48eV and a HOMO offset of 0.8eV. The HOMO and LUMO orbitals for FBT (where we replaced the octyl side chains in F8BT with methyl groups) are shown in Fig. 3. This clearly indicates the localization of the HOMO and LUMO wave functions on the co-polymer units.
We can include this alternation into our model by modulating the site energies of the F8BT chainKarabunarliev and Bittner (2004). Thus, in F8BT we include a 0.5 eV modulation of both the valence and conduction band site energies relative to the band center. Table 1. Hence, the fluorene site energies are at 2.42eV and -3.04eV for the conduction and valence band while the benzothiadiazole site energies are 1.42eV and -4.04eV. This results in a shift in the excitation energy to 0.28 eV relative to the unmodulated polymer and a 0.09 eV increase in the exciton binding energy. Furthermore, the absorption spectrum consists of two distinct peaks at 2.14eV (563nm) and 4.4eV (281nm) which are more or less on par with the 2.77eV (448nm) $`S_oS_1`$ and the 4.16eV (298 nm) $`S_oS_9`$ transitions computed by Jespersen et al. Jespersen et al. (2004) and observed at 2.66eV and 3.63eV by Stevens, et al. Stevens et al. (2001).
The charge-separated character of the lowest singlet excited state of the F8BT-mod chain is readily apparent in Fig. 5 A and B. This is the state which give rise to the absorption maxima at 2.14 eV and is the only state present in the emission spectra. The odd-numbered sites are fluorene-sites and the even-numbered sites are benzothiadiazole sites. The alternating maxima in Fig. 5 A and B show that the hole is largely localized on the odd-sites and the electron is localized on the even-sites. The relaxed state is a self-trapped state with the electron largely localized on site 6 and the hole distributed on either side on sites 5 and 7. States C and D are the vertical and relaxed odd-parity state corresponding to a dissociated charge-transfer state. Finally, E and F in this figure correspond to the state responsible for the absorption peak at 4.4eV. In this state, the electron and hole are more or less uniformly delocalized over the entire polymer segment and carries an oscillator strength comparable to the 2.14eV state, again consistent with the semiempirical data. Jespersen et al. (2004)
For PFB and TFB monomers, semiempirical (PM3) calculations at the ground state equilibrium configuration give HOMO energies of -8.00748eV and -7.95994 eV and LUMO energies of -0.01862 eV and 0.05201 eV respectively for isolated phenylenediamine and diphenylamine units. In comparison, the semiempirical HOMO and LUMO levels for the fluorene segment are -8.83304eV and -0.32915eV, respectively. In the combined co-polymers, the LUMO is more or less localized on the fluorene as in the F8BT case as shown in Fig. 6. To account for this modulation within our model, we include an 0.2eV modulation about the band centers listed in Table 1 in the conduction and valence site energies for the PFB and TFB chains. Since the net modulation is on the order of the exciton binding energy, the the electron and hole will be localized on alternate repeat units in the lowest excited state as indicated in Fig. 7-A for the vertical exciton and Fig. 7-B for the adiabatic or self-trapped exciton. Even though the electron and hole are on different lattice sites, they remain a bound pair. Fig. 7-C and -D show the dissociated electron/hole pair as evidenced by the nodal line along the $`e=h`$ diagonal.
Table 2 summarizes the relevant states for the isolated chains discussed in this paper. It is interesting to note that for the unmodulated chains (PPV and BBL) the adiabatic exciton binding energy is 0.1eV higher than in the modulated co-polymer chains (F8BT, TFB, and PFB). Having established models for the electronic states of the isolated chains, we turn our attention towards modeling the electronic states of the polymer heterojunctions.
### III.2 Electronic states of heterojunctions
We consider three model heterojunction systems: PPV:BBL, TFB:F8BT, and PFB:F8BT. As discussed in the introduction, devices fabricated using these polymers are either layered nanostructures or phase-separated polymer blends. The crucial energetic consideration is the exciton destabalization threshold. In an ideal non-interaction electron model, a type II heterojunction will destabilize an exciton present at the interface if the exciton binding energy is less than the band-off set. In organic semiconductors, the Coulombic interaction between the electron and the hole is quite significant with exciton binding energies in excess of 0.5 eV in most materials. When this is on the order of the band off-set, excitons can be stable at the interface giving LED behavior. By selecting materials with different HOMO and LUMO energies, one can effectively tune LED or photovoltaic performance of a heterojunction material.
#### III.2.1 PPV:BBL
The first of these semiconductor heterojunction materials we consider here, PPV:BBL, has been used in the fabrication of solar-cells and photovoltaic devices and a recent review by Alam and Jenehke provide a succinct overview of recent progress towards the fabrication of efficient solar-cell devices using layers of organic polymer donor/acceptor heterojunction materials.Alam and Jenekhe (2004); and John A. Osaheni (1994); M. Granström et al. (1998); Bach et al. (1998); Kamohara et al. (2005); Russell et al. (2002); Jenekhe and Yi (2000); Chen and Jenekhe (1997); Lin et al. (1996); Weng et al. (1995); Xue et al. (2004); Osaheni et al. (1994); Gao and Dane (2004); Gregg et al. (2002); Saito and Sugi (1992); Yu et al. (1995); Jenekhe and Yi (2000); Manoj et al. (2003); Narayan et al. (1995); Nakamura et al. (2004); Schilinsky et al. (2004) By alternating nanoscale bilayers of high electron affinity (EA) acceptor polymers such as BBL with an EA = 4.0eV with a donor such as PPV or its soluble derivative MEH-PPV exhibit rapid and efficient photoinduced charge transfer and large photovoltaic power conversions under white light illumination at 80-100 m/cm<sup>2</sup> (AM1.5).Alam and Jenekhe (2004) Since the band off-set at the heterojunction of the materials is larger than the exciton binding energy in PPV:BBL, photoexcitation of either BBL or PPV will result in the production of a charge-separated interchain species.
The relevant states for the PPV:BBL heterojunction are shown in Fig. 8. The top two figures are electron/hole densities for inter-chain charge-separated state while the bottom two are the vertical and adiabatic densities for the exciton. Fig. 9 indicates the correlation between the vertical and relaxed energy levels. The vertical exciton (Fig. 8-C) corresponds to the state with the largest transition dipole to the $`S_o`$ ground state.
First, we note that the lowest energy excited state in this system is the charge-transfer state in which the hole resides on the PPV chain (sites 1-10 in Fig. 8) and the electron resides on the BBL chain (sites 11-15). There is some mixing between the interchain charge-transfer state and geminate electron/hole configurations on either chain as evidenced by the smallest contour rings in Fig. 8 A and B. Clearly, $`>`$99.9% of the population is in interchain charge-transfer configurations. The vertical excitonic state, Fig. 8 C, is largely localized on the PPV chain with some mixing between the chains. We find that the extent of the interchain mixing is determined largely by the interchain hopping integral, which we set at $`t_{}=0.15t_{||}=0.0804eV`$. In the self-trapped exciton (Fig. 8 D), the interchain mixing is weaker.
The energy level diagram shown in Fig.9 shows the correlation between the vertical excited states (at the ground state equilibrium geometry ($`Q_\mu =0`$) with the adiabatically relaxed excited states, each of which has a unique excited state equilibrium geometry. This correlation diagram indicates a number of intersystem crossings can occur as the system relaxes from the vertical to adiabatic geometries. In particular, notice that the excitonic state (C) intersects another state as it relaxes to D. This state is nearly degenerate with the adiabatic exciton (D) and is predominantly an interchain charge-transfer state. It is interesting to speculate the role that such intersections and close degeneracies will play in the photophysics of this system.
#### III.2.2 PFB:F8BT vs. TFB:F8BT
As discussed earlier, TFB:F8BT and PFB:F8BT sit on either side of the exciton destabalization threshold. In TFB:F8BT, the band off-set is less than the exciton binding energy and these materials exhibit excellent LED performance. On the other hand, devices fabricated from PFB:F8BT where the exciton binding energy is less than the off-set, are very poor LEDs but hold considerable promise for photovoltic devices. In both of these systems, the lowest energy state is assumed to be an interchain exciplex as evidenced by a red-shifted emission about 50-80ns after the initial photoexcitation. Morteani et al. (2003) In the case of TFB:F8BT, the shift is reported to be 140$`\pm `$20 meV and in PFB:F8BT the shift is 360$`\pm `$30 meV relative to the exciton emission, which originates from the F8BT phase. Bearing this in mind, we systematically varied the separation distance between the cofacial chains from $`r=2a5a`$ (where $`a`$ = unit lattice constant) and set $`t_{}=0.01`$ in order to tune the Coulomb and exchange coupling between the chains and calibrate our parameterization. Fig. 10 shows the variation of the lowest few diabatic states with $`r`$ for both heterojunctions. The dashed line in each corresponds to the exciton energy for an isolated F8BT chain. For large interchain separations, the exciton remains localized on the F8BT chain in both cases. As the chains come into contact, dipole-dipole and direct Coulomb couplings become significant and we begin to see the effect of exciton destabilization. For TFB:F8BT, we select and interchain separation of $`r=2.8a`$ giving a 104meV splitting between the vertical exciton and the vertical exciplex and 87.4 meV for the adiabatic states. For PFB:F8BT, we chose $`r=3a`$ giving a vertical exciton-exciplex gap of 310 meV and an adiabatic gap of 233 meV. In both TFB:F8BT and PFB:F8BT, the separation produce interchain exciplex states as the lowest excitations. with energies reasonably close to the experimental shifts.
Fig. 11 compares the vertical and adiabatic energy levels in the TFB:F8BT and PFB:F8BT chains and Figs. 12 and 13 show the vertical and relaxed exciton and charge-separated states for the two systems. Here, sites 1-10 correspond to the TFB or PFB chains and 11-20 correspond to the F8BT chain. The energy levels labeled in Fig. 11 correspond to the states plotted in Figs. 12 and 13. We shall refer to states A and B as the vertical exciplex and vertical exciton and to states C and D as the adiabatic exciplex and adiabatic exciton respectively. Roughly, speaking a pure exciplex state will have the charges completely separated between the chains and will contain no geminate electron/hole configurations. Likewise, strictly speaking, a pure excitonic state will be localized to a single chain and have only geminate electron/hole configurations. Since site energies for the the F8BT chain are modulated to reflect to internal charge-separation in the F8BT co-polymer as discussed above, we take our “exciton” to be the lowest energy state that is localized predominantly along the diagonal in the F8BT “quadrant”.
In the TFB:F8BT junction, the lowest excited state is the exciplex for both the vertical and adiabatic lattice configurations with the hole on the TFB and the electron on the F8BT. (Fig.12A,C) In the vertical case, there appears to be very little coupling between intrachain and interchain configurations. However, in the adiabatic cases there is considerable mixing between intra- and inter-chain configurations. First, this gives the adiabatic exciplex an increased transition dipole moment to the ground state. Secondly, the fact that the adiabatic exciton and exciplex states are only 87 meV apart means that at 300K, about 4% of the total excited state population will be in the adiabatic exciton.
For the PFB:F8BT heterojunction, the band off-set is greater than the exciton binding energy and sits squarely on the other side of the stabilization threshold. Here the lowest energy excited state (Figs. 13A and B) is the interchain charge-separated state with the electron residing on the F8BT (sites 11-20 in the density plots in Fig 13) and the hole on the PFB (sites 1-10). The lowest energy exciton is almost identical to the exciton in the TFB:F8BT case. Remarkably, the relaxed exciton (Fig. 13D) shows slightly more interchain charge-transfer character than the vertical exciton (Fig. 13C). While the system readily absorbs at 2.3eV creating a localized exciton on the F8BT, luminescence is entirely quenched since all population within the excited states is readily transfer to the lower-lying interchain charge-separated states with vanishing transition moments to the ground state.
The mixing between the exciplex and exciton in TFB:F8BT can be seen in the predicted emission spectra for the system. Fig. 14 shows the absorption and emission spectra for the TFB:F8BT and PFB:F8BT computed by taking the transition dipole moment between the vertical and relaxed excited states to the ground state. While the fluorescent emission from TFB:F8BT is weaker than the corresponding absorption, it is not entirely quenched as in the PFB:F8BT case.
## IV Discussion
In this paper, we developed a model for donor-acceptor semiconducting polymer heterojunctions. For the PPV:BBL heterojunction, the band off-sets at the polymer interface is sufficient to dissociate an electron/hole pair created by photoexcitation. Because of this natural internal driving force to create charge-separated states, PPV:BBL blends are of considerable interest to materials chemists and device engineers developing efficient photovoltaic cells for solar energy conversion Alam and Jenekhe (2004). TFB:F8BT and PFB:F8BT are very similar heterojunctions. Within our model, the single parametric difference between TFB and PFB is the location of the band-center for the valence band. In TFB:F8BT, the off-set between the valence bands is (-2.98eV+ 3.54eV) 0.56eV which is right at the exciton binding energy. Hence, we see considerable mixing between the intrachain exciton localized on the F8BT and a charge-separated state. In the PFB-F8BT case, the valence band off-set is 0.79 eV, clearly greater than the exciton binding energy. Hence, the lowest energy excited states are predominantly interchain charge-separated states.
In a forthcoming paper, Bittner et al. (2005), we continue along the lines set forth in this paper by computing the state-to-state relaxation dynamics following both photoexcitation and electron/hole injection. We will compare results obtained within a vertical approximation (whereby transitions occur via phonon creation/annihilation, but the lattice is frozen in the ground-state equilibrium geometry) to those obtained within an adiabatic Marcus-Jortner-Hush approximation where the transitions occur between the equilibrium geometry of each excited state. We will also discuss the role that the avoided crossings and conical intersections of the Born-Oppenheimer potential energy surfaces of the excited states play in the electron/hole capture process, exciplex formation, and in the thermal regeneration of intramolecular excitons in TFB:F8BT.
###### Acknowledgements.
This work was supported in part by the National Science Foundation and the Robert A. Welch Foundation.
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# OVERVIEW OF LHCb
## 1 Introduction
The LHCb experiment $`^{\mathrm{?},\mathrm{?}}`$ has been designed to study CP violation and rare phenomena in $`B`$ meson decays with very high precision. The physics goal being to provide an understanding of quark flavor physics, and possibly to reveal physics beyond the Standard Model. The experiment will be based at CERN and will study the decay of $`b`$ quarks from $`b\overline{b}`$ pairs produced in proton-proton collision. It will operate at an average luminosity of $`2\times 10^{32}`$cm<sup>-2</sup>s<sup>-1</sup>, much lower than the maximum design luminosity of the LHC, in order to have a number of interactions per crossing dominated by single interaction. This facilitates the triggering and reconstruction by assuring low channel occupancy. In addition, the radiation damages are more manageable. To achieve its physics goals, the experiment will have to perform, a high track reconstruction efficiency, a good $`\pi K`$ separation for momenta from few to $``$100 GeV/c, a very good propertime resolution ($``$40 fs), and, a high trigger efficiency (both for final states including leptons and for those with hadrons only). These needs have led to the detector design that is represented on Fig. 1 and that is described in the next section.
## 2 LHCb detector description
LHCb detector is a single-arm spectrometer with a forward angular coverage from 10 mrad to 300 (250) mrad in the bending (non-bending) plane. The choice of the geometry has been motivated by the fact that at high energies both the $`b`$ and $`\overline{b}`$ hadrons are produced at small angles with respect to the beam pipe $`^\mathrm{?}`$. This is illustrated by Fig. 2 obtained from PYTHIA event generator $`^\mathrm{?}`$.
Figure 1 shows a side view of the detector that consists of the beam pipe, the vertex detector (VELO), the dipole magnet, the tracking system (TT, T1-T3), two Ring Imaging Cherenkov detectors (RICH1 and RICH2), the calorimeter system (SPD/PS, ECAL and HCAL) and the muon system (M1-M5). The VELO has to provide precise measurements of track coordinates close to the interaction region in order to get good proper time resolution. It is made of circular silicon stations placed along and perpendicularly to the beam axis. With the excellent momentum resolution achieved by the tracking system, the proper time of reconstructed $`B`$ mesons can be measured with a resolution of $``$40 fs $`^\mathrm{?}`$. The tracking system is composed of the dipole magnet which is a warm magnet and has a field integral of 4 Tm, a Trigger Tracker (TT) located in front of the magnet entrance and three tracking station (T1-T3) placed behind the magnet. TT stations are made of silicon sensors and play two roles. Firstly, it will be used in the Level-1 trigger to assign transverse momentum information to large impact parameter tracks. Secondly, it will be used in the offline analysis to reconstruct tracks, in particular the decay products of long-lived neutral particle that decay outside the VELO. Each T station consists of an Inner Tracker (IT) close to the beam pipe and an Outer Tracker (OT) surrounding the IT. The IT is made of silicon strip detectors and the OT of straw tubes. Charged tracks are reconstructed with a high efficiency of $``$95% with a low rate of wrongly reconstructed tracks, which does not introduce significant additionnal combinatorial background in the reconstructed $`B`$ meson signals. RICH system is composed of two elements located in front of TT and behind T stations. To cover the momentum range from few to $``$100 GeV/c, three different radiators have been chosen : silica aerogel and two fluorocarbon gases, C<sub>4</sub>F<sub>10</sub> and CF<sub>4</sub>. In this momentum range, the average efficiency for kaon identification is 88% for an average pion misidentification rate of 3%. The main purpose of the calorimeter system is to identify electrons and hadrons and to provide measurements of their energy and position. These measurements are used by the trigger system and for the off-line analysis. The structure consists of four elements, a scintillator pad detector (SPD), a preshower (PS), an electromagnetic calorimeter (ECAL) and an hadronic calorimeter (HCAL). These elements employ similar technologies, i.e. scintillators coupled to wavelength-shifting fibers read out by fast photodectors. The electron identification efficiency for tracks in ECAL acceptance is $``$94% with a pion misidentification rate of $``$0.7%. The muon detector is used to identify muons for the trigger and off-line analysis. It consists of five stations, M1 in front of the calorimeter system and M2-M5 behind the calorimeter, interleaved with iron shielding plates. Each station is made of four layers of Multi Wire Proportional Chambers (MWPC) except for M1 which is made of two. For tracks in muon detector acceptance, muon identification efficiency is $``$93% with a pion misidentification rate of $``$1%.
## 3 LHCb trigger
The trigger $`^\mathrm{?}`$ is one of the biggest challenge for the LHCb experiment, the $`b\overline{b}`$ pair creation cross section being less than 1% of the total cross section. It is designed to distinguish events containing $`B`$ mesons from minimum-bias events through the presence of particles with a large transverse momentum ($`p_T`$) and the existence of secondary vertices. The trigger is divided in three levels : Level-0 (L0) implemented in custom electronics, Level-1 (L1) and High Level Trigger (HLT) both executed in farm processors. The L0 trigger will reduce the 40 MHz LHC beam crossing rate to 1 MHz. The events are triggered by requiring at least one lepton or hadron with a $`p_T`$ exceeding 1 to 3 GeV/c. Events can also be rejected based on global event variables such as track multiplicities and number of interactions. The L1 trigger selects events with an output rate of 40 kHz. The L1 algorithm reconstructs tracks in the VELO and matches these tracks to L0 muons or calorimeter clusters to identify them and measure their momenta. The fringe field of the magnet between the VELO and TT is used to determine the momenta of particles with a resolution of 20-40% and events are selected based on tracks with a large $`p_T`$ and significant impact parameter to the primary vertex. Finally the HLT will select the stored events with a frequency of 200 Hz. The HLT algorithm has access to all data and starts by reconstructing the VELO tracks and the primary vertex. A fast pattern recognition program links the VELO tracks to the tracking stations T1-T3. The final selection of interesting events is a combination of the confirmation of the L1 decision with better resolution, and selection cuts dedicated to specific final states. In addition to the events stored by HLT, 1.8 KHz will be stored to get systematics from the data (e.g. a trigger on high di-muon mass will be used to calibrate tracking, inclusive $`b`$ $`\mu `$ events will be selected to calibrate trigger, and $`D^{}`$ events to calibrate the particle identification).
## 4 LHCb expected physics performances
“Toy Monte Carlo” programs have been used to estimate the LHCb sentivities to some $`CP`$ observables. These programs include signal resolution, efficiency, purity, etc. taken from studies with fully-simulated events. Nevertheless, assumption need to be made concerning the properties of background events due to the lack of statistics. In the real analyses, these properties and the systematics effects will be extracted from the data. Here only statistical errors are given.
The LHCb physics program includes many topics. For some of them, expected sensitivities corresponding to one year of data taking (integrated luminosity of 2 fb<sup>-1</sup>) and a $`b\overline{b}`$ pair production cross-section assumed to be 0.5 mb will be given hereafter. $`B_d^0`$ mixing phase will be measured with $`B^0J/\psi K_S^0`$ decay with a sensitivity $``$0.02. $`B_s^0`$ mixing phase $`\varphi _s`$ and decay-width $`\mathrm{\Delta }\mathrm{\Gamma }_s/\mathrm{\Gamma }_s`$ will be assessed with $`B_s^0J/\psi \varphi `$ decay. This channel presents several challenges. Indeed, an angular analysis is needed as $`J/\psi `$ and $`\varphi `$ are vector mesons. In addition, the oscillation frequency $`\mathrm{\Delta }m_s`$ is expected to be large, requiring excellent proper-time resolution. The expected sensitivities to $`\varphi _s`$ and $`\mathrm{\Delta }\mathrm{\Gamma }_s/\mathrm{\Gamma }_s`$ are respectively 0.064 and 0.018. $`\mathrm{\Delta }m_s`$ will be measured with $`B_s^0D_s^{}\pi ^+`$ decay that is flavor-specific and give access to $`\mathrm{\Delta }m_s`$ through flavour asymmetry analyses (see Fig. 3). The sensitivity on $`\mathrm{\Delta }m_s`$ is 0.009 (0.016) for $`\mathrm{\Delta }m_s`$ = 15 ps<sup>-1</sup> (30 ps<sup>-1</sup>). Moreover, oscillations can be observed (5$`\sigma `$) for $`\mathrm{\Delta }m_s`$ values up to 68 ps<sup>-1</sup>.
The measurement of the angle $`\gamma `$ of the unitarity triangle will be done in three different ways. Time-dependant decay asymmetries in $`B_s^0D_s^{}K^\pm `$ decays combined with $`\varphi _s`$ measurement from $`B_s^0J/\psi \varphi `$ decay, will give $`\sigma (\gamma )=14^{}15^{}`$ without theoretical uncertainty. Time-dependant $`CP`$ asymmetries in $`B_d^0\pi ^+\pi ^{}`$ and $`B_s^0K^+K^{}`$ decays, in combination with measurement with $`B^0J/\psi K_S^0`$ and $`B_s^0J/\psi \varphi `$ respectively, will give $`\sigma (\gamma )=4^{}6^{}`$ assuming U-spin symmetry $`^\mathrm{?}`$. Time-integrated rates of $`B^0D^0K^{}`$, $`B^0\overline{D}{}_{}{}^{0}K_{}^{0}`$, and $`B^0D_{CP}^0K^0`$ decays $`^\mathrm{?}`$, will give $`\sigma (\gamma )=7^{}8^{}`$. In the presence of new physics, these three different approaches could allow to identify it. In addition, very rare decays like $`B_s^0\mu ^+\mu ^{}`$ decay, $`b`$ $``$ $`s`$ penguin processes as $`B^0\varphi K_S^0`$, $`B_d^0\mu ^+\mu ^{}K^0`$,$`B_d^0K^0\gamma `$, $`B_s^0\varphi \gamma `$ and $`B_s^0\varphi \varphi `$ decays, $`b`$ $``$ $`d`$ penguin processes with $`B^0\rho \pi `$ decay, $`B_c`$ mesons and $`b`$-baryons will be studied in great detail by LHCb.
## References
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# An Efficient Numerical Scheme for Simulating Particle Acceleration in Evolving Cosmic-Ray Modified Shocks
## 1 Introduction
Collisionless shocks are widely thought to be effective accelerators of energetic, nonthermal particles (hereafter Cosmic-Rays or CRs). Those particles play central roles in many astrophysical problems. The physical basis of the responsible diffusive shock acceleration (DSA) process is now well established through in-situ measurements of heliospheric shocks and through analytic and numerical calculations . While test particle DSA model treatments are relatively well developed; e.g., , it has long been recognized that DSA is an integral part of collisionless shock physics and that there are substantial and highly nonlinear backreactions from the CRs to the bulk flows and to the MHD wave turbulence mediating the CR diffusive transport (see, for example, and references therein). Most critically, the CRs can capture a large fraction of the kinetic energy dissipated across such transitions. As they diffuse upstream the CRs form a pressure gradient that decelerates and compresses the entering flow inside a broad shock precursor. That, in turn, can lead to greatly altered full shock jump conditions, especially if the most energetic CRs, which can have very large scattering lengths, escape the system and carry significant energy with them. Also in response to the momentum dependent scattering lengths and flow speed variations through the shock precursor the CR momentum distribution will take on different forms than in a simple discontinuity. Effective analytic (e.g., ) and numerical (e.g., ) methods have been developed that allow one to compute steady-state modified shock properties given an assumed diffusion behavior. On the other hand, as the CR particle population evolves in time during the formation of such a shock the shock dynamics and the CR-scattering wave turbulence evolve as well. For dynamically evolving phenomena, such as supernova remnants, the time scale for shock modification can be comparable to the dynamical time scales of the problem.
The above factors make it essential to be able to include both nonlinear and time dependent effects in studies of DSA. Generally, numerical simulations are called for. Full plasma simulations offer the most complete time dependent treatments of the associated shock microphysics , but are far too expensive to follow the shock evolution over the time, length and energy scales needed to model astrophysical CR acceleration. The most powerful alternative approach utilizes continuum methods, with a kinetic equation for each CR component combined with suitably modified compressible fluid dynamical equations for the bulk plasma (see §2 below). By extending that equation set to include relevant wave action equations for the wave turbulence that mediates CR transport, a self-consistent, closed system of equations is possible (e.g., )). Continuum DSA simulations of the kind just described are still quite challenging and expensive even with only one spatial dimension. The numerical difficulty derives especially from the very large range of CR momenta that must be followed, which usually extends to hundreds of GeV/c or beyond on the upper end and down to values close to those of the bulk thermal population, with nonrelativistic momenta. The latter are needed in part to account for “injection” of CRs due to incomplete thermalization that is characteristic of collisionless shocks.
One computational constraint comes from the fact that CR resonant scattering lengths from MHD turbulence, $`\lambda (p)`$, are generally expected to be increasing functions of particle rigidity, $`pc/Ze`$. The characteristic length coupling the CRs of a given momentum, $`p`$, to the bulk flow and defining the width of the modified shock precursor is the so-called diffusion length, $`x_d(p)=\frac{1}{3}\lambda \upsilon /u_s`$, where $`\upsilon `$ is the CR particle speed, and $`u_s`$ is the bulk flow speed into the shock. One must spatially resolve the modified shock transition for the entire range of $`x_d(p)`$ in order to capture the physics of the shock formation and the spatial diffusion of the CRs, in particular. The relevant $`x_d(p)`$ typically spans several orders of magnitude, beginning close to the dissipation length of the thermal plasma, which defines the thickness of the classical, “viscous” gas shock, also called the “subshock” in modified structure. This resolution requirement generally leads to very fine spatial grids in comparison to the “outer scale” of the problem, which must exceed the largest $`x_d(p)`$.
Two approaches have been applied successfully so far to manage this constraint in DSA simulations. Berezhko and collaborators developed a method that normalizes the spatial variable by $`x_d(p)`$ at each momentum value of interest during solution of the CR kinetic equation. This approach establishes an spatial grid that varies economically in tune with $`x_d(p)`$. Derived CR distribution properties at different momenta can be combined to estimate feed-back on the bulk flow at appropriate scales. The method was designed for use with CR diffusion properties known a priori. It is not readily applied to CR diffusion behaviors incorporating arbitrary, nonlinear feedback between the waves and the CRs. As an alternative that can accommodate those latter diffusion properties, Kang et al. have implemented diffusive CR transport into a multi-level adaptive mesh refinement (AMR) environment. The benefit of AMR in this context comes from the feature that the highest resolutions are only necessary very close to the subshock, which can still be treated as a discontinuity satisfying standard Rankine-Hugoniot relations. By efficient use of spatial gridding both of these computational strategies can greatly reduce the cost of time dependent DSA simulations.
On the other hand, the above methods do not directly address the principal computational cost in such simulations, so they remain much more costly compared to purely hydrodynamic or MHD simulations. This is because the dependence of $`f`$ on CR momentum, $`p`$, adds a physical dimension to the problem. In practice, the spatial evolution of the kinetic equation for each CR constituent must be updated over the entire spatial grid at multiple momentum values; say, $`N_p`$. The value of $`N_p`$ is usually large, since the spanned range of CR momentum is typically several orders of magnitude. Physically, CRs propagate in momentum space during DSA in response to adiabatic compression in the bulk flow, sometimes by momentum diffusion (see, for example, equation 1 below), or because of various irreversible energy loss mechanisms, such as Coulomb or radiative losses. The associated evolution rates for $`f`$ depend on the process, but generally depend on $`f/\mathrm{ln}p`$. The conventional approach to evolving $`f`$ approximates $`f/\mathrm{ln}p`$ through low order finite differences in $`\mathrm{\Delta }\mathrm{ln}p`$ (e.g., ). Experience has shown that converged solutions of $`f(p)`$ using such methods require $`\mathrm{\Delta }\mathrm{ln}p<0.1`$ . In that case, for example, a mere five decades of momentum coverage requires more than 100 grid points in $`\mathrm{ln}p`$. Since spatial update of the kinetic equation at each momentum grid point requires computational effort comparable to that for any of the accompanying hydrodynamical equations (e.g., the mass continuity equation), CR transport then dominates the computational effort by a very large factor, commonly exceeding an order of magnitude.
An attractive alternative approach to evolving the kinetic equation replaces $`f`$ by its integral moments over a discrete set of finite momentum volumes, in which case $`f/\mathrm{ln}p`$ is replaced by $`f`$ evaluated at the boundaries of those volumes. The method we outline here follows that strategy. Because $`f(\mathrm{ln}p)`$ is relatively smooth, simple subvolume models can effectively be applied over moderately large momentum volumes. We have found this method to give accurate solutions to the evolution of $`f`$ with an order of magnitude fewer momentum bins than needed in our previous finite difference calculations. The computational effort to evolve the CR population is thereby reduced to a level comparable to that for the hydrodynamics. In recognition of its distinctive features we refer to the method as “Coarse-Grained Momentum finite Volume” or “CGMV”. It extends related ideas introduced in , and for test particle CR transport. Those previous presentations, while satisfactorily following CR transport in many large-scale, smooth flows, did not include spatial or momentum diffusion, so could not explicitly follow evolution of $`f`$ during DSA. Instead, analytic, test-particle solutions for $`f(p)`$ were applied at shock jumps. Here we extend the CGMV method so that it can be applied to the treatment of fully nonlinear CR modified shocks.
We outline the basic CGMV method and its implementation in Eulerian hydrodynamics codes in §2. Several tests are discussed in §3, and our conclusions are presented in §4.
## 2 The Method
### 2.1 Basic Equation
The standard diffusion-convection form of the kinetic equation describing the evolution of the isotropic CR distribution function, $`f(x,p,t)`$, can be written in one spatial dimension as (e.g., ).
$$\frac{f}{t}+u\frac{f}{x}=\frac{1}{3}\left(\frac{u}{x}\right)\frac{f}{y}+\frac{}{x}\left(\kappa \frac{f}{x}\right)+\frac{1}{p^3}\frac{}{y}\left(pD\frac{f}{y}\right)+S,$$
(1)
where $`u`$ is the bulk flow speed, $`y=\mathrm{ln}p`$, $`\kappa `$ is the spatial diffusion coefficient, $`D`$ is the momentum diffusion coefficient, and $`S`$ is a representative source term. We henceforth express particle momentum in units of $`mc`$, where $`m`$ is the particle mass. The first RHS term in equation (1) represents “momentum advection” in response to adiabatic compression or expansion. For simplicity of presentation equation 1 neglects for now propagation of the scattering turbulence with respect to the bulk plasma, which can be a significant influence when the sonic and Alfvénic Mach numbers of the flow are comparable. Although it is numerically straightforward to include this effect, the details are somewhat complex, so we defer that to a follow-up work focussed on CR transport in MHD shocks.
Full solution of the problem at hand requires simultaneous evolution of the hydrodynamical flow, as well as the diffusion coefficients, $`\kappa `$ and $`D`$. Again postponing full MHD, the added equations to be solved are the standard gasdynamic equations with CR pressure included. Expressed in conservative, Eulerian formulation for one dimensional plane-parallel geometry, they are
$$\frac{\rho }{t}+\frac{(u\rho )}{x}=0,$$
(2)
$$\frac{(\rho u)}{t}+\frac{(\rho u^2+P_g+P_c)}{x}=0,$$
(3)
$$\frac{(\rho e_g)}{t}+\frac{(\rho e_gu+P_gu+P_cu)}{x}=u\frac{P_c}{x}L(x,t),$$
(4)
where $`P_\mathrm{g}`$ and $`P_\mathrm{c}`$ are the isotropic gas and the CR pressure, respectively, $`e_\mathrm{g}=P_\mathrm{g}/\rho (\gamma _\mathrm{g}1)+u^2/2`$ is the total energy density of the gas per unit mass and the rest of the variables have their usual meanings. The injection energy loss term, $`L(x,t)`$, accounts for the energy of the suprathermal particles transferred at low energy to the CRs. As usual, CR inertia is neglected in such computations, since the mass fraction of the CRs is generally tiny. We note for completeness that $`P_c`$ can be computed from $`f`$ using the expression
$$P_c=\frac{4\pi }{3}mc^2_{p_{min}}^{p_{max}}p^4f\frac{dp}{\sqrt{1+p^2}}.$$
(5)
In the simulations described below we set the particle mass, $`m=1`$, for convenience .
### 2.2 The CGVM Method
As mentioned in §1, the momentum advection and diffusion terms in equation 1 typically require $`\mathrm{\Delta }y<0.1`$ when using low order finite difference methods in the momentum coordinate . The resulting large number of grid points in $`y`$ makes finding the solution of equation 1 the dominate effort in simulations of DSA. On the other hand, previous studies of DSA as well as direct observations of CRs in different environments have shown that $`f(p)`$ is commonly well described by the form $`fp^{q(p)}`$, where $`q(p)=\frac{\mathrm{ln}f}{y}`$, is a slowly varying function of $`y`$. Thus, we may expect a piecewise powerlaw form to provide an efficient and accurate, two-parameter subgrid model for $`f(y)`$. Two moments of $`f(p)`$ are sufficient to recover the subgrid model parameters. We find it convenient to use
$$n_i=_{p_i}^{p_{i+1}}p^2f(p)𝑑p=_{y_i}^{y_{i+1}}\mathrm{exp}(3y)f(y)𝑑y,$$
(6)
and
$$g_i=_{p_i}^{p_{i+1}}p^3f(p)𝑑p=_{y_i}^{y_{i+1}}\mathrm{exp}(4y)f(y)𝑑y.$$
(7)
The first of these moments, $`n_i`$, is proportional to the spatial number density of CRs in the momentum bin $`\mathrm{\Delta }y_i=[y_i,y_{i+1}]`$, while for relativistic CRs, $`g_i`$ is proportional to the energy density or pressure contribution of CRs in the bin. Then, for example,
$$n_i=\frac{f_ip_i^3}{q_i3}\left[1d_i^{3q_i}\right],$$
(8)
where $`f_i=f(p_i)=(p_{i+1}/p_i)^{q_i}f_{i+1}`$, and $`d_i=p_{i+1}/p_i`$ with obvious extension to $`g_i`$. Either of these moments, plus their ratio, $`g_i/p_in_i`$, can be used in straightforward fashion (e.g., iteration) to find both $`f_i`$ and the intrabin index, $`q_i`$.
To evolve $`n_i`$ and $`g_i`$ we need the associated moments of equation 1 over the finite momentum volume bounded by $`\mathrm{\Delta }y_i`$. The result for $`n_i`$ is
$$\frac{n_i}{t}+u\frac{n_i}{x}=F_{n_i}F_{n_{i+1}}n_i\frac{u}{x}+\frac{}{x}\left(K_{n_i}\frac{n_i}{x}\right)+S_{n_i},$$
(9)
where $`F_{n_i}=\{\dot{p}_i+q(p_i)D(p_i)/p_i)\}p^2_if(p_i)`$ is a flux in momentum space, with $`\dot{p}=p\frac{1}{3}\frac{u}{x}`$, and where $`K_{ni}`$ and $`S_{ni}`$ are averaged over the momentum interval, according to
$$K_{n_i}=\frac{_{p_i}^{p_{i+1}}p^2\kappa \frac{f}{x}𝑑p}{_{p_i}^{p_{i+1}}p^2\frac{f}{x}𝑑p}\frac{_{p_i}^{p_{i+1}}\kappa fp^2𝑑p}{n_i},$$
(10)
and
$$S_{n_i}=_{p_i}^{p_{i+1}}p^2S(p)𝑑p.$$
(11)
Extension of the momentum flux term $`F_{n_i}`$ to include other processes such as radiative or Coulomb losses is obvious (e.g., ). In practice these fluxes should be upwinded according to the signs of $`\dot{p}`$ and $`q`$. Evaluation of $`F_{n_i}`$ and $`q(p_i)`$ at the boundaries of the included momentum range requires application of suitable boundary conditions, of course. We usually have set $`n_{N_p+1}=0`$. In most cases we pick a sufficiently large maximum momentum, $`p_{N_p+1}`$, that this condition is important only late in the simulation, if at all. Appropriate conditions at the lowest momenta can be more involved, depending on how one intends to connect the CR particle distribution to the thermal particle distribution, as in the injection models discussed below.
The $`g_i`$-associated moment of equation 1 is
$$\frac{g_i}{t}+u\frac{g_i}{x}=F_{g_i}F_{g_{i+1}}+g_i\left(4\frac{\dot{p}_i}{p_i}+q_i\frac{D}{p^2}_i\right)+\frac{}{x}\left(K_{g_i}\frac{g_i}{x}\right)+S_{g_i},$$
(12)
where $`F_{g_i}=p_iF_{n_i}`$,
$$K_{g_i}=\frac{_{p_i}^{p_{i+1}}\kappa \frac{f}{x}p^3𝑑p}{_{p_i}^{p_{i+1}}\frac{f}{x}p^3𝑑p}\frac{_{p_i}^{p_{i+1}}\kappa fp^3𝑑p}{g_i},$$
(13)
$`g_iD/p^2_i=_{p_i}^{p_{i+1}}pDf𝑑p`$, and $`S_{g_i}`$ is given by an analogous expression to equation 11.
We note that momentum binning in the CGMV scheme is quite flexible, so that it can be easily adapted to either uniform or nonuniform momentum bin sizes, or to a momentum range that evolves during the simulation. We have successfully implemented both nonuniform and evolving momentum bin structures, although, for brevity we do not illustrate them here.
To update the distribution function we simultaneously integrate equations (9) and (12) over the timestep $`\mathrm{\Delta }t^k`$. Our implementation applies these methods in Eulerian hydrodynamical codes, so we split the update into two parts. Spatial advection is carried out by a second order, explicit van Leer scheme. The remaining terms are evolved using a second-order, semi-implicit Crank Nicholson scheme. For example, given values $`n_i^k`$ at timestep $`t^k`$, and assuming for illustration a uniform spatial grid, the implicit part of the solution is given by the tridiagonal system
$$A_i^+n_{i+1}^{k+1}+A_i^0n_i^{k+1}+A_i^{}n_{i1}^{k+1}=C_i^0,$$
(14)
where
$`A_i^+=\delta K_{i+1/2},`$
$`A_i^0=1+\delta (K_{i+1/2}+K_{i1/2}),`$
$`A_i^{}=\delta K_{i1/2},`$
$`C_i^0=n_i^k\left[1\delta (K_{i+1/2}+K_{i1/2})\right]+n_{i+1}^k\delta K_{i+1/2}+n_{i1}^k\delta K_{i1/2}`$
$`+\mathrm{\Delta }t^k\left[\overline{F}_{n_i}^k\overline{F}_{n_{i1}}^k+3n_i^k({\displaystyle \frac{\dot{p}_i}{p_i}})+S_{n_i}\right],`$
with $`\delta =(1/2)\mathrm{\Delta }t^k/(\mathrm{\Delta }x)^2`$, and $`K_{i+1/2}=(1/2)(K_{i+1}+K_i)`$.
The coefficients in equation 14 are obtained with the aid of the solutions to equations 2 \- 4, which are updated prior to solution of equations 9 and 12. We note again that similar methods can be applied to follow the evolution of the wave turbulence that resonantly scatters CRs and that defines the spatial and momentum diffusion coefficients. In that case one begins with the wave action equation for the appropriate waves rather than the particle kinetic equation (e.g., ).
The solution of equation (14) for either $`n_i`$ or $`g_i`$ is quite analogous to our previous FD methods. Thus, since the CGMV method evolves two quantities rather than one, the relative effort required for a given $`N_p`$ is roughly twice in the CGMV scheme that required in the FD scheme. Our tests confirm this expected scaling. On the other hand, since $`N_p`$ can be dramatically reduced in a CGMV simulation the method can still be more efficient by a large factor.
### 2.3 A Flux Fraction Injection model
The most common source term represented by $`S`$ in equation (1) is injection at the shock of low energy CRs from the thermal plasma. There is presently no generally accepted theory for that process. However, we have implemented two commonly used models successfully into the CGMV scheme. For completeness we outline those here.
The simplest and one of the most frequently applied injection models assumes that a small, fixed fraction of the thermal particle flux through the gas subshock, $`ϵ_{inj}`$, is injected at a momentum $`p_{inj}=\alpha c_{s_2}`$, where $`\alpha `$ is a constant greater than unity and $`c_{s_2}`$ is the plasma sound speed immediately downstream of the subshock (e.g., ). This gives $`S(p)=ϵ_{inj}(\rho _1/\mu m_p)u_sw(xx_s)\delta (pp_{inj})`$, where $`\rho _1`$ is the plasma mass density just upstream of the subshock, $`\mu m_p`$ is the plasma mean particle mass, $`u_s`$ is the subshock speed with respect to the plasma immediately upstream and $`w`$ is a normalized weight function that allows the injection to be distributed across the numerical shock structure. In this case $`S_{n_i}=\frac{1}{4\pi }ϵ_{inj}(\rho _1/\mu m_p)u_sw(xx_s)`$, while $`S_{g_i}=p_{inj}S_{n_i}`$ in the momentum bin with $`p_i<p_{inj}<p_{i+1}`$. Both $`S_{n_i}`$ and $`S_{g_i}`$ are zero, otherwise. The energy extracted from the thermal plasma is simply $`L=\frac{1}{2}ϵ_{inj}w(x)\alpha ^2c_{s_2}^2\rho _1u_s`$. For convenience we call this injection model the “flux fraction” or “FF” model.
### 2.4 A Thermal Leakage Injection model
A more sophisticated approach to injection physics includes models of the physical processes moderating particle orbits in the post shock flow region in order to estimate the probability that particles of a given speed will be able to escape back upstream, across the subshock. In such “thermal leakage” (TL) models for CR injection at shocks, most of the downstream thermal protons are locally confined by nonlinear MHD waves and only particles well into the tail of the postshock Maxwellian distribution can leak upstream across the subshock. In particular, “leaking” particles not only must have velocities large enough to swim against the downstream flow in order to return across the shock, they must also avoid being scattered during that passage by the MHD waves that mediate the plasma subshock. To model TL injection we utilize a “transparency function”, $`\tau _{\mathrm{esc}}`$, expressing the probability that supra-thermal particles at a given velocity can leak upstream from behind the subshock (see for details). In particular we set
$$\tau _{esc}(\upsilon ,u_2)=H\left[\stackrel{~}{\upsilon }(1+ϵ_B)\right]\frac{(1\frac{1}{\stackrel{~}{\upsilon }})}{(1\frac{u_2}{\upsilon })}\mathrm{exp}\left\{\left[\stackrel{~}{\upsilon }(1+ϵ_B)\right]^2\right\},$$
(15)
where $`u_2`$ is the postshock flow speed in the subshock frame, $`H`$ is the Heaviside step function, and the particle velocity is normalized to $`\stackrel{~}{\upsilon }=vϵ_B/u_2`$. The parameter, $`ϵ_B=B_0/B_{}`$, measures the ratio of the amplitude of the postshock MHD wave turbulence $`B_{}`$ to the general magnetic field aligned with the shock normal, $`B_0`$. Both hybrid simulations and theory suggest that $`0.25ϵ_B0.35`$ , so that the model is well constrained. With this $`\tau _{esc}`$ the shock is completely “opaque” to particles with momenta less than $`p_1`$, i.e., $`\tau _{esc}=0`$ for $`p<p_1`$, where $`p_1=m_pu_2(1+ϵ_B)/ϵ_B`$. So $`p_1(t)`$ is the lowest momentum of the first momentum bin in the TL model and changes in time with the postshock flow speed. For $`ϵ_B0.3`$, $`p_145m_pu_2`$. The shock becomes virtually transparent to particles with momenta two to three times greater than $`p_1`$. For strong, unmodified shocks $`p_1`$ in the TL model and $`p_{inj}`$ in the FF model are similar when $`ϵ_B0.3`$ and $`\alpha 2`$. Under those circumstances the initial injection rates will be roughly similar, although differences in model physics lead to different behaviors as such shocks become modified (see, e.g., , , ).
The TL model is implemented in the CGMV scheme by the following numerical approach. After solution of equations (9) and (12) the net changes in $`n_i`$ and $`g_i`$ are corrected (reduced) in the upstream region by application of the transparency function as follows:
$$n_i^{k+1}=n_i^k+_{p_i}^{p_{i+1}}\tau _{esc}(p)(\stackrel{~}{f}_i^{k+1}f_i^k)p^2𝑑p$$
(16)
and
$$g_i^{k+1}=g_i^k+_{p_i}^{p_{i+1}}\tau _{esc}(p)(\stackrel{~}{f}_i^{k+1}f_i^k)p^3𝑑p$$
(17)
where $`\stackrel{~}{f}_i^{k+1}`$, found using equation (8), is the CR distribution updated with equations (9) and (12). The energy loss rate of the bulk plasma to injection into the $`i`$-th CR momentum bin can be approximated by
$$L_i(x,t)\frac{4\pi m_pc^2}{3}(\frac{u}{x})_{p_i}^{p_{i+1}}\frac{\tau _{esc}}{p}p^3(\sqrt{p^2+1}1)f_i^n𝑑p$$
(18)
(see ).
With the piece-wise power-law subgrid model ($`f_i(p)=f_i(p/p_i)^{q_i}`$) the integrals in equations (16)-(18) can be written:
$$_{p_i}^{p_{i+1}}\tau _{esc}(p)f_i(p)p^2𝑑p=\frac{n_i(q_i3)p_i^{q_i3}}{(1d_i^{3q_i})}_{p_i}^{p_{i+1}}\tau _{esc}p^{(3q_i)}𝑑y,$$
(19)
$$_{p_i}^{p_{i+1}}\tau _{esc}(p)f_i(p)p^3𝑑p=\frac{g_i(q_i4)p_i^{q_i4}}{(1d_i^{4q_i})}_{p_i}^{p_{i+1}}\tau _{esc}p^{(4q_i)}𝑑y,$$
(20)
and
$$L_i(x,t)\frac{4\pi m_pc^2}{3}(\frac{u}{x})\frac{g_ip_i^{q_i4}(q_i4)}{(1d_i^{4q_i})}_{p_i}^{p_{i+1}}\frac{\tau _{esc}}{p}p^{4q_i}(\sqrt{p^2+1}1)𝑑y$$
(21)
In practice, this leakage step is significant only for the lowest few momentum bins, so that this correction need not be applied to all bins.
## 3 Discussion
In order to test the performance of the new CGMV scheme we have installed it into two distinct one-dimensional Eulerian hydrodynamic (HD) codes that we have previously applied to studies of CR-modified shocks using conventional finite difference (FD) methods to solve the diffusion convection equation. In this section we briefly discuss the results and compare the CGMV and FD behaviors. Both of the host HD codes are constructed from high order, conservative Riemann solver-based schemes designed to capture shocks sharply. First we describe results from the CGMV scheme installed in a second order “Total Variation Diminished” (TVD) HD code based on the finite difference scheme of Harten . This is the HD version of the MHD-CR code used by us in a previous study of CR modified shocks . Gas subshocks in the TVD scheme typically spread over 2-3 numerical zones. An outline of the code mechanics and the FD CR scheme can be found there, in and references cited in those papers. The FD solver employed a Crank-Nicholson routine originally introduced in for evolving $`p^4f(p)`$ that is similar to equation (14). For the TVD tests we applied the FF injection model.
In addition we present CGMV tests carried out with our Cosmic Ray Adaptive refinement SHock (CRASH) code. CRASH is based on the high order Godunov-like shock tracking algorithm of LeVeque . The hydrodynamics routine in that code employs a nonlinear Riemann solver to follow shock discontinuities within the zones of an initially uniform grid. Thus, gas subshocks in CR-modified shocks remain discontinuous throughout a simulation, allowing CR transport to be modeled down almost to the scale of the physical shock thickness. CRASH also employs adaptive mesh refinement (AMR) around shocks in order to reduce the computational effort on the spatial grid. Refinement is centered on the subshock and each level spans 100 zones with a resolution twice as fine as the level above it. The number of refinement levels depends on what is required to capture diffusion of the lowest energy CRs. The standard, previously documented version of CRASH uses the same FD methods as the TVD code to solve the diffusion convection equation for CR transport. It is described in and . For the CRASH tests discussed in this paper we employed the TL injection model.
### 3.1 TVD-CR Tests
We first examine some results obtained using the TVD-CR code with both FD and CGMV schemes used to model the evolution of a strong CR-modified shock. Fig. 1 and Fig. 2 illustrate the evolution of shocks formed by the reflection of a Mach 30 flow (adiabatic index, $`\gamma =5/3`$) off the left grid boundary. The resulting piston-driven shock initially has a Mach number, $`M_s40`$. The density and sound speed of plasma entering from the right boundary were set to $`\rho =1`$ and $`c_s=\frac{1}{30}`$, respectively, so that the inflow speed was unity. The time unit for the calculations is also set by these scalings. In order to relate hydrodynamical variables to CR momenta it is necessary to fix the unit flow speed (the inflow speed in all the simulations discussed in this paper) with respect to the speed of light; i.e., $`c=1/\beta `$. We set $`\beta =10^2`$ in the TVD-CR simulations. Time steps were fixed by a standard Courant condition, $`\mathrm{\Delta }t=0.8\mathrm{\Delta }x/\mathrm{max}(u\pm c_s)`$.
For the simulation illustrated in Fig. 1 evolution of the CR distribution is followed over the momentum range $`[p_{min},p_{max}]=[2\times 10^4,1.6\times 10^3]`$ ($`y_{max}y_{min}16`$). The simulation represented in Fig. 2 included the momentum range $`[p_{min},p_{max}]=[2\times 10^4,2.4\times 10^5]`$ ($`y_{max}y_{min}21`$).
The CR diffusion coefficient, $`\kappa `$ is spatially uniform and set to $`\kappa (p)=0.1p^{0.51}`$. In a quest for a reasonably generalized behavior that required minimal computational effort, this choice was motivated by results from Malkov, who found self-similar analytic steady-state solutions for strong CR-modified shocks that apply to all powerlaw forms of $`\kappa (p)`$, so long as the powerlaw index is steeper that $`0.5`$ . Thus, our $`\kappa `$ choice leads to fairly general shock behaviors in a way that minimizes the width of the precursor. That width, which determines the minimum space that must be simulated ahead of the subshock, is set by $`x_d(p_{cutoff})`$, where $`p_{cutoff}`$ represents the maximum momentum contained. The TVD-CR simulations utilize a uniform, fixed grid, so, for example, the Bohm diffusion form modeled in the CRASH simulations below would lead to excessive costs for the TVD-CR FD test simulations presented in this section. The spatial resolution required for the calculations is set by $`x_d(p_{min})`$, since accurate solutions of the diffusion-convection equation require good structural information in the diffusive shock precursor upstream of the subshock. Previous convergence tests have shown that $`\mathrm{\Delta }x<0.1x_d`$ is desirable (e.g., ). For the problems illustrated in Fig. 1 and Fig. 2 these considerations led us to set $`\mathrm{\Delta }x=3.8\times 10^4`$ for both the FD and CGMV simulations. By varying this resolution, we verified that the shock evolution is reasonably converged with respect to $`\mathrm{\Delta }x`$.
The simulation followed in Fig. 1 assumes a pre-existing CR population, $`f(p)p^{4.5}`$, corresponding to an upstream CR pressure, $`P_c=P_g`$. No fresh injection is included at the shock; i.e., $`ϵ_{inj}=0`$. This test then provides a simple and direct comparison between the CGMV and FD schemes for solving the diffusion-convection equation, since it omits any complications related to the injection model.
This simulation pair evolves the shocks until $`t=30`$, which is sufficient to accelerate CRs to ultrarelativistic momenta. The spatial grid spans the interval , which is sufficient to contain the leading edge of the CR precursor to the end of the simulation.
The CR-modified shock spatial structures and the CR momentum distributions at the subshock are shown in Fig. 1 at times $`t=2,10,20,\mathrm{and}30`$. Before comparing the solutions obtained with the two different methods, it is useful to summarize briefly the physics captured during the shock evolution. All the behaviors described here have been reported previously by multiple authors. The figure shows the well-known property of strongly modified shocks that the CRs extract most of the energy flux into the structure. That leads to a substantial drop in the postshock gas pressure, $`P_g`$, and a large increase in the postshock density, $`\rho `$. Together those indicate a strongly reduced postshock gas temperature. The decreased temperature is evident in the $`p^4f`$ plot at the subshock, which is dominated at low momenta by the Maxwellian distribution of the bulk plasma. As CRs diffuse upstream against the inflowing gas they compress the flow within the precursor, preheating the gas (adiabatically in these simulations). Initially, while the CR pressure is relatively small compared to the incoming momentum flux, the gas subshock remains strong enough to produce a full four times density compression on top of the precompression. However, the subshock weakens once $`P_c\rho u^2`$, reducing the subshock compression in this case to a factor $`2.6`$, corresponding to a subshock Mach number near 2.3. That evolution explains the well-known, transitory “density spike” in the shock structures seen after $`t=2`$. We note that since energy extracted from the flow by CRs becomes increasingly spread upstream and downstream due to CR diffusion, the total compression in such an evolving modified shock would always exceed the factor of seven one would predict for a strong, fully relativistic gas shock. For this simulation no significant CR energy escapes the spatial grid through upstream diffusion. However, at late times ($`t>20`$) the partial pressure due to CRs just below $`p_{max}`$ is sufficient that escape across the upper momentum boundary is significant. This contributes to the slow decrease in $`P_c`$ behind the shock and the increasing total compression through the shock transition that is visible in Fig. 1.
In the early evolutionary stages of this flow, while shock modification is modest, the CR momentum distribution resembles the powerlaw form, $`p^4f(p)=const`$, predicted by test particle theory for a strong shock (e.g., ). With the spatial diffusion coefficient used in this simulation the high momentum cutoff to the distribution increases with time approximately as $`p_{cutoff}t^2`$ (see, e.g., ). As shock modification intensifies, most of the flow compression shifts from the subshock to the precursor. Then DSA of high momentum CRs occurs predominantly within the precursor rather than near the less important subshock. Consequently the CR distribution develops the familiar upwards-concave form resulting from the momentum dependent CR diffusion length. CRs of higher momentum experience a greater velocity jump within the precursor, so gain more energy each time they are reflected within the shock structure. That flattens the distribution, $`f(p)`$ at momenta below $`p_{cutoff}`$. The result is a bump in the distribution of $`p^4f(p)`$. On the other hand, CRs with momenta only a little above the injection range remain trapped close to the subshock. Their distribution closely approaches the steady state, powerlaw, test particle form appropriate to the weakened subshock. That feature extends upwards in momentum as the bump near $`p_{cutoff}`$ moves upwards.
Looking finally to compare the two methods used to evolve the shock evolution displayed in Fig. 1 we see results from the FD scheme with $`\mathrm{\Delta }y=0.11`$ and the CGMV scheme with $`\mathrm{\Delta }y=1.0`$. The agreement is generally very good. All the dynamical quantities, including shock jumps and the CR momentum distributions show excellent agreement. The curves representing the FD and CGMV distributions of $`\rho `$, $`P_g`$ and $`P_c`$ and virtually indistinguishable in the plots. Most notably, all features formed in the FD evolution of the CR momentum distribution are faithfully reproduced by the much coarser CGMV distribution.
Given the excellent comparison in this strongly modified flow it is satisfying to note that the execution time required for the CGMV solution was a little less than 20% of that for the FD solution, demonstrating the significantly higher efficiency of the former method. The speed-up observed in our implementations of the two methods is roughly in accordance with what we would predict for a given reduction factor in the number of momentum values used, since the CGMV method requires one to evolve two distributions $`n_i`$ and $`g_i`$ for each momentum bin. We address convergence with respect to momentum resolution in our discussion of a second shock.
Fig.2 shows a Mach 30 flow similar to that in Fig. 1. In this case FF injection is included with commonly assumed values, $`ϵ_{inj}=10^3`$ and $`\alpha =2.0`$ (see §2.3), the upper momentum bound is increased to $`p_{max}=2.4\times 10^5`$ and the spatial grid extends farther from the piston to $`x_{max}=25`$. A negligible pre-existing CR population is included to avoid numerical issues coming from the fact that our CGMV scheme requires computation of the ratio $`g_i/p_in_i`$ over the entire grid.
The simulations evolved the shock until $`t=70`$, which leads to $`p_{cutoff}10^3`$. The spatial grid, spanning the interval , is sufficient to contain the CR leading edge of the precursor almost to the end of the simulations. However, after $`t50`$ some CR energy escapes through the right boundary, due to diffusion upstream, mimicking the behavior of a “free escape” boundary (FEB) (e.g., and references therein). Just as for the shock simulated in Fig.1, this energy loss amounts to a cooling process, so that the total shock compression increases with time as the simulation ends.
Again the agreement between FD and CGMV simulations is very good. Very early in the evolution of the shock, when the CRs are dominated by freshly injected nonrelativistic particles, the shock evolution is slightly faster in the CGMV scheme. That influence becomes insignificant later on, so that the modified shock structures found by the two schemes are almost identical as is the distribution $`p^4f(p)`$ at the shock. There is a small residual effect that the position of the subshock in the CGMV simulation lags slightly behind that of the FD simulation, and that $`P_c`$ is slightly higher in the CGMV simulation near the piston, where the shock first formed.
Since the efficiency of the CGMV method comes from its ability to cover the momentum range coarsely, it is important to evaluate how broad the momentum bins can be and still faithfully model the evolution of the shock. Fig. 3 illustrates convergence of the CGMV scheme with respect to momentum bin size, $`\mathrm{\Delta }y`$, at $`t=70`$ for the flow modeled in Fig. 2. The upper panel plots the spatial $`P_c`$ distributions, while the lower panel shows the particle momentum distributions at the gas subshock. For reference the corresponding FD solution ($`\mathrm{\Delta }y=0.11`$) is shown by the dotted red curves. Solutions from the CGMV scheme are plotted for $`\mathrm{\Delta }y=1`$, $`\mathrm{\Delta }y=1.40`$ and $`\mathrm{\Delta }y=2.1`$. Even the coarsest of the CGMV solutions is in basic agreement with the other CGMV solutions and with the FD solution. Fine details in the momentum distribution are naturally obscured as the CGMV bin size increases. The largest bins with $`\mathrm{\Delta }y=2.3`$ span a decade in CR momentum ($`p_{i+1}/p_i=10.2`$), but still capture the basic dynamical properties of the CRs correctly.
However, the quality of the CGMV solutions deteriorates for still larger momentum bins in these experiments, once the simple subgrid model for the momentum distribution becomes inadequate. As illustrated in the lower panel of Fig. 3 already momentum bins larger than roughly $`\mathrm{\Delta }y1.5`$ cannot closely follow sharper structures in the momentum distibution that develops at the ends of the CR distribution. That enhances $`P_c`$ upstream of the subshock, where the flow is both cold and strongly compressed as it approaches the subshock. When the errors become excessive, for $`\mathrm{\Delta }y>2.3`$, in this case, $`P_c`$-induced overcompression in this region can cause the Riemann solver in our TVD code to perform poorly or even to fail in high Mach number flows. In general the largest allowed momentum bin size, $`\mathrm{\Delta }_{max}`$ should depend on the strongest curvature of the CR momentum distribution function as well as the degree of shock modification.
### 3.2 CRASH Tests
For a third test example we illustrate in Figs. 4 and 5 simulation results using the CRASH code with the TL injection model applied to a Mach 10 flow reflecting off the left computational boundary. The initial gas shock Mach number is approximately 13. As for the previous test, the upstream gas density and flow speed are set to unity, with upstream sound speed, $`c_s=\frac{1}{10}`$, and $`\beta =5\times 10^3`$. The time unit is defined accordingly. CR momenta are tracked over the range $`[p_1,10^4]`$, where $`p_1(t)`$ is the smallest momentum that can leak upstream (see equation 15).
In this case a Bohm-type diffusion model with $`\kappa _B=p^2/\sqrt{p^2+1}`$, is adopted and the TL injection parameter, $`ϵ_B=0.2`$ is used. The CRASH test was significantly more computationally demanding than the TVD-CR tests. Note first that in the CRASH simulation the value of $`\kappa _B(p=1)=1/\sqrt{2}`$, while $`\kappa (p=1)=0.1`$ in the previous examples shown in Figs. 1-3. Consequently, $`x_d(p=1)`$ is about seven times greater in the current case, and the nominal physical scale of the precursor and its formation timescale are similarly lengthened. In addition, the stronger momentum dependence of Bohm diffusion coefficient means that the precursor width expands more strongly as $`p_{cutoff}`$ increases. The associated time rate of increase in $`p_{cutoff}`$ is, however, slower, so that the shock must evolve longer to reach a given $`p_{cutoff}`$. These factors substantially increase the size of the physical domain needed to reach a given $`p_{cutoff}`$.
Fig. 4 shows the early evolution of this CR-modified shock for $`t20`$ as computed with both the FD and the CGMV methods. The spatial domain for this simulation is . The base spatial grid included $`10^4`$ zones, giving $`\mathrm{\Delta }x_0=2\times 10^3`$. Since it is necessary to resolve structures near the subshock on scales of the diffusion length for freshly injected, suprathermal CRs, the AMR feature of the CRASH code is utilized. The FD simulation is carried out with 7 refined grid levels; four levels of refined grid are applied in the CGMV simulation. 240 momentum points ($`\mathrm{\Delta }y=0.058`$)are used in the FD simulation, while the CGMV simulation includes 20 momentum bins ($`\mathrm{\Delta }y=0.72`$). The time step for each refinement level, $`\mathrm{\Delta }t_{l_g}`$, is determined by a standard Courant condition, that is, $`\mathrm{\Delta }t_{l_g}=0.3\mathrm{\Delta }x_{l_g}/\mathrm{max}(u\pm c_s)`$. Although the Crank-Nicholson scheme is stable with an arbitrary time step, the diffusion convection equation is solved with the time step smaller than $`\mathrm{\Delta }t_{DC,l_g}2\mathrm{\Delta }y(\mathrm{\Delta }x_{l_g}/u_s`$) to maintain good accuracy in the momentum space advection (i.e., $`dy/dt=\frac{1}{3}\frac{u}{x}`$). With $`\mathrm{\Delta }y=0.058`$, the required time step is smaller by a factor of three or so than the hydrodynamic time step in the FD simulation. Consequently, the FD diffusion convection solver is typically subcycled about 3 times with $`\mathrm{\Delta }t_{DC,l_g}`$ for each hydrodynamic time step. Because of the much larger $`\mathrm{\Delta }y`$, subcycling is not necessary in the CGMV simulation. That adds another relative economy to the CGMV calculation.
At the end of this simulation, $`t=20`$, the modified shock structure is approaching a dynamical equilibrium in the sense that the postshock values of $`\rho `$, $`P_g`$ and $`P_c`$ will not change much at later times. Since this shock is weaker than the Mach 40 shocks examined earlier modifications are more moderate. On the other hand, as expected from the stronger momentum dependence of $`x_d(p)`$, the shock precursor broadens much more quickly in the present case. The cutoff in the CR distribution has reached roughly $`p_{cutoff}10`$ by $`t=20`$. Longer term evolution of this shock will be addressed below.
The agreement between the FD and CGMV solutions shown in Fig. 4 is good, although not as close as it was in the examples illustrated in Fig. 1 and Fig. 2. The more apparent distinctions between the two solutions in the present case come from effective differences in the application of the TL injection model with Bohm diffusion in FD and CGMV methods. Recall that the CGMV scheme applies the diffusion coefficients averaged across the momentum bins (see equations 10, 13). The Bohm diffusion model has a very steep momentum dependence for nonrelativistic particles; namely, $`\kappa p^2`$. At low momenta where injection takes place the averaging increases the effective diffusion coefficient, and, thus, the leakage flux of suprathermal particles, leading to higher injection rate compared to the FD scheme for the same TL model parameters. Consequently, the distribution function in the second bin at $`p_21.4\times 10^2`$ is slightly higher in the CGMV scheme, as evident in Figs. 4-5. Note that $`f(p_1)`$ is anchored on the tail of Maxwellian distribution. The CGMV solutions accordingly show slightly more efficient CR acceleration than the FD solutions at early times. In this test $`P_c`$ is about 5 % greater in the CGMV simulation at $`t=20`$. Since the CGMV scheme can be implemented with nonuniform momentum bins, such differences could be reduced by making the momentum bins smaller at low momentum in instances where the details relating to the injection rate were important.
We show in Fig. 5 the evolution of this same shock extended to $`t=10^3`$, as computed with the CGMV method. This simulation is computed on the domain , spanned by a base spatial grid of $`2\times 10^4`$ zones, giving $`\mathrm{\Delta }x_0=4\times 10^2`$. We also included 7 refined grid levels at the subshock, giving $`\mathrm{\Delta }x_7=3.1\times 10^4`$. This grid spacing is insufficient for convergence at the injection momentum, $`p_{inj}10^2`$, so that the very early evolution is somewhat slower than in the simulations shown in Fig. 4. However, once shock modification becomes strong evolution becomes roughly self-similar, as pointed out previously . The time asymptotic states do not depend sensitively on the early injection history. The self-similar behavior results with Bohm diffusion from a match between the upstream and downstream extensions of the CR population. One also sees from the form of the distribution function in Fig. 5 that the postshock gas temperature has stabilized, while the previously-explained concave form to the CR distribution is better developed than it was at earlier times.
This simulation illustrates nicely the relative efficiency of the CGMV scheme. The equivalent FD simulation would be very much more expensive, because this model requires a long execution time and a large spatial domain. With Bohm diffusion $`\kappa p`$ for ultrarelativistic CRs, so that the scale of the precursor, $`x_dp_{cutoff}`$. At the same time the peak in the CR momentum distribution extends relatively slowly, with $`p_{cutoff}t`$. The required spatial grid is, thus, 40 times longer than for the shorter simulation illustrated in Fig. 4. The simulated time interval in the extended simulation was 50 times longer. Together those increase the total computational time by a factor 2000. The FD calculation with $`x=[0,20]`$ to $`t=20`$ took about 2 CPU days on our fastest available processor, so the extended simulation would have been unrealistic using the FD method. The extended CGMV simulation, however, required only about 10 times the effort of the shorter FD simulation, clearly demonstrating the efficiency of the CGMV scheme. This speed-up is a result of combination of several factors: 20 times larger grid spacing, no need for subcycling for the diffusion convection solver, and, of course, a smaller number of momentum bins.
## 4 Conclusions
Detailed time dependent simulations of nonlinear CR shock evolution are very expensive if one allows for inclusion of arbitrary, self-consistent and possibly time dependent spatial diffusion, as well as various other momentum dependent transport processes. The principal computational cost in such calculations is typically the CR transport itself, and, in self-consistent calculations, the analogous transport of the MHD wave turbulence that mediates CR transport. Tracking these behaviors requires adding at least one physical dimension to the simulations compared to the associated hydrodynamical calculations, since the collisionless media involved are sensitive to the phase space configurations of the particles and waves. Particle kinetic equations (commonly the so-called diffusion convection equation) provide a straightforward approach to addressing this problem and can be coupled conveniently with hydrodynamical equations that track mass and bulk momentum and energy effectively.
Momentum derivatives of the CR distribution function in the diffusion convection equation are most frequently handled by finite differences. Although it is simple, that approach requires moderately fine resolution in momentum space. That is a primary reason that such calculations are costly. Here we introduce a new scheme to solve the diffusion convection equation based on finite volumes in momentum space with a momentum bin spacing as much as an order of magnitude larger than that of the usual finite difference scheme. We demonstrate that this Coarse Grained Momentum finite Volume ( CGMV) method can be used successfully to model the evolution of strong, CR-modified shocks at much lower computational cost than the finite difference approach. The computation efficiency is greatly increased, not only because the number of momentum bins is smaller, but also because the required spatial grid spacing is less demanding due to the coarse-grained averaging of the diffusion coefficient used in the CGMV method. In addition, larger momentum bin size can eliminate the need of subcycling of the diffusion convection solver that can be necessary in some instances using finite differences in momentum. Thus, the combination of the CGMV scheme with AMR techniques as developed in our CRASH code, for example, should allow more detailed modeling of the diffusive shock acceleration process with a strongly momentum dependent diffusion model such as Bohm diffusion, or self-consistent treatments of CR diffusion and wave turbulence transport.
## Acknowledgments
TWJ is supported by NSF grant AST03-07600, by NASA grants NAG5-10774, NNG05GF57G and by the University of Minnesota Supercomputing Institute. HK was supported by KOSEF through the Astrophysical Research Center for the Structure and Evolution of Cosmos (ARCSEC).
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# Suppressed fusion cross section for neutron halo nuclei
## Abstract
Fusion reactions of neutron-halo nuclei are investigated theoretically with a three-body model. The time-dependent wave-packet method is used to solve the three-body Schrödinger equation. The halo neutron behaves as a spectator during the Coulomb dissociation process of the projectile. The fusion cross sections of <sup>11</sup>Be - <sup>209</sup>Bi and <sup>6</sup>He - <sup>238</sup>U are calculated and are compared with measurements. Our calculation indicates that the fusion cross section is slightly hindered by the presence of weakly bound neutrons.
NUCLEAR REACTION, <sup>209</sup>Bi(<sup>10</sup>Be,X), (<sup>11</sup>Be,X), <sup>238</sup>U(<sup>4</sup>He,X), (<sup>6</sup>He,X), Calculated fusion cross section, Three-body model, Time-dependent wave-packet method preprint: nucl-th/0506073
Physics at the dripline is one of current subjects in nuclear physics. Adding neutrons to a nucleus as many as possible, we often find that some neutrons are bound very weakly. These neutrons form a spatially extended neutron cloud which is called the neutron halo. The halo nuclei have been attracting many researchers since its discovery in the secondary beam experiment Tanihata85 . The neutron-halo nucleus exhibits many properties different from normal nuclei. For instance, it violates the nuclear saturation of density and binding energy. Experimentally, these nuclei are mostly studied with reactions using the secondary beam.
Since the halo nuclei are weakly bound and easily break up, developments of the reaction theories capable of describing coupling to the continuum channels are required. For reactions in medium and high incident energies, the Glauber and the eikonal theories have been successful Suzuki03 ; Bertsch89 ; Ogawa92 ; Yabana92 . In these theories, the wave function of the halo neutron in the projectile is frozen during the reaction. The success of the eikonal theories rests on the fact that the dynamics can be well separated into fast and slow motions. Namely, the relative motion of projectile and target is much faster than that of halo neutrons. There is, however, no such clear distinction in collisions at low energy. In spite of many experimental and theoretical efforts in the last decade, the reaction mechanism of halo nuclei at low incident energies is still an open question.
The fusion of halo nuclei shows an example of the controversial subjects. In early stages, a simple and intuitive theoretical argument has been proposed Takigawa91 ; Hussein92 : The fusion probability was expected to be enhanced because of the spatially extended density of the halo neutron. A further enhancement was predicted by the coupling to soft modes inherent to the halo structure. Breakup effects were predicted to be small, sustaining a large enhancement of fusion cross section especially at sub-barrier energies Hussein92 ; Takigawa93 . These early studies utilized a specific reaction model. We have developed a time-dependent wave packet method for fully quantum calculations using a three-body model Yabana95 ; Yabana97 ; Yabana03 ; Yabana04 ; Nakatsukasa04 . Our results have disagreed with the fusion enhancement. In fact, the calculation has even suggested a slight suppression of the fusion probability because of the presence of the halo neutron. However, recent coupled-channel calculations using a similar three-body model, again, indicate a substantial enhancement in the fusion probability at sub-barrier energies Hagino00 ; Diaz02 .
Experimental results also show a rather confusing status. In early measurements, an enhancement of the fusion cross section at sub-barrier energies was reported for reactions using a two-neutron-halo nucleus <sup>6</sup>He Kolata98 ; Trotta00 . Recently, however, Ref. Raabe04 have reported that the enhancement previously obtained in Ref. Trotta00 is not due to the fusion but to the transfer-induced fission. In case of a one-neutron-halo nucleus, <sup>11</sup>Be, the experimental results did not present an evident conclusion either Signorini98 . A recent comparison between <sup>10</sup>Be and <sup>11</sup>Be fusion cross sections suggests no enhancement for <sup>11</sup>Be Signorini04 . Thus, we may say, at least, that the present experimental status indicates no evidence for a fusion enhancement of halo nuclei.
In this letter, we report calculations of the fusion cross sections in the three-body model. We previously reported fusion probabilities of halo nuclei mostly for the case of total angular momentum $`J=0`$ which corresponds to the head-on collision Yabana03 ; Yabana04 ; Nakatsukasa04 . We now include a full range of impact parameter, for $`0J30`$, then, for the first time, present results of fusion cross sections calculated with the time-dependent wave-packet method. Calculations are performed for fusion cross sections of <sup>11</sup>Be - <sup>209</sup>Bi and <sup>6</sup>He - <sup>238</sup>U. Results are compared with experimental data Raabe04 ; Signorini04 .
In our three-body model, the projectile is described as a weakly bound two-body system of the core and the halo neutron(s). The projectile ($`P`$), which is composed of the core ($`C`$) plus neutron(s) ($`n`$), and the target ($`T`$) constitute the three-body model. For <sup>11</sup>Be, we assume a <sup>10</sup>Be$`+n`$ structure in which the halo neutron occupies a $`2s`$ orbital in the projectile ground state. For <sup>6</sup>He, a di-neutron model of $`\alpha +2n`$ is assumed. The time-dependent Schrödinger equation for the three-body model is given as
$`i\mathrm{}{\displaystyle \frac{}{t}}\mathrm{\Psi }(𝐑,𝐫,t)=\{{\displaystyle \frac{\mathrm{}^2}{2\mu }}_𝐑^2{\displaystyle \frac{\mathrm{}^2}{2m}}_𝐫^2`$
$`+V_{nC}(r)+V_{CT}(R_{CT})+V_{nT}(r_{nT})\}\mathrm{\Psi }(𝐑,𝐫,t),`$ (1)
where we denote the relative $`n`$-$`C`$ coordinate as $`𝐫`$ and the relative $`P`$-$`T`$ coordinate as $`𝐑`$. The reduced masses of $`n`$-$`C`$ and $`P`$-$`T`$ motions are $`m`$ and $`\mu `$, respectively.
The Hamiltonian in Eq. (1) includes interaction potentials among constituents. The $`n`$-$`C`$ potential, $`V_{nC}(r)`$, is a real potential. This potential should produce a halo structure of the projectile. The core-target potential, $`V_{CT}(R_{CT})`$, consists of the nuclear and Coulomb potentials. The nuclear potential is complex and its short-ranged imaginary part is inside the Coulomb barrier between the core and the target nuclei. The $`n`$-$`T`$ potential, $`V_{nT}(r_{nT})`$, is taken to be real. We treat the constituents as spin-less point particles, and ignore the non-central forces. The nuclear potentials of $`V_{nC}`$, $`V_{CT}`$, and $`V_{nT}`$ are taken to be of a Woods-Saxon shape. For <sup>11</sup>Be - <sup>209</sup>Bi reaction, we take $`V_0=53`$ MeV, $`r_0=1.3`$ fm, $`a=0.7`$ fm for $`V_{nC}`$, $`V_0=65`$ MeV, $`r_0=1.16`$ fm, $`a=0.8`$ fm for the real part of $`V_{CT}`$, $`V_0=60`$ MeV, $`r_0=0.6`$ fm, $`a=0.4`$ fm for the imaginary part of $`V_{CT}`$, and $`V_0=44.019`$ MeV, $`r_0=1.27`$ fm, $`a=0.67`$ fm for $`V_{nT}`$. For <sup>6</sup>He - <sup>238</sup>U reaction, we adopt $`V_0=69.305`$ MeV, $`r_0=0.667`$ fm, $`a=0.65`$ fm ($`V_{nC}`$), $`V_0=95`$ MeV, $`r_0=1.1101`$ fm, $`a=0.5834`$ fm (real $`V_{CT}`$), $`V_0=10`$ MeV, $`r_0=0.6`$ fm, $`a=0.7`$ fm (imaginary $`V_{CT}`$), and $`V_0=93.4`$ MeV, $`r_0=0.986`$ fm, $`a=1.184`$ fm ($`V_{nT}`$). The $`V_{nC}`$ provides the $`2s`$ bound orbital at $`0.6`$ MeV for <sup>11</sup>Be, and $`0.97`$ MeV for <sup>6</sup>He, which are close to the experimental separation energies of one and two neutrons, respectively. The adopted parameter set for $`V_{nC}`$ is close to that used in the studies of the $`\alpha `$+<sup>6</sup>He and $`p`$+<sup>6</sup>He scatterings by Rusek et al. 6heopt . This parameter set gives a large value of $`B(E2;g.s.2^+)`$ for <sup>6</sup>He as pointed out in Ref. 6heopt .
We define the fusion cross section in terms of the loss of flux because of the imaginary potential in $`V_{CT}`$. The use of the imaginary potential is approximately equivalent to the incoming boundary condition inside the barrier. Once the core and the target nuclei are in the range of the imaginary $`V_{CT}`$ inside the Coulomb barrier, the total flux decreases by the absorption, no matter whether the neutron is captured by the target or not. Therefore, this includes not only complete but also a part of the incomplete fusion in which the charged core (<sup>10</sup>Be for <sup>11</sup>Be and <sup>4</sup>He for <sup>6</sup>He) and the target fuse while the neutron escapes.
The wave-packet calculation is performed with the partial wave expansion. Here, we use the partial wave expansion in the body-fixed frame Pack74 ; Imanishi87 . Although this is equivalent to the one in the space-fixed frame, the body-fixed frame is computationally superior for states at finite total angular momentum, $`J0`$. In the body-fixed frame, the channels are specified by the projection of $`J`$ on the body-fixed $`z`$-axis, $`\mathrm{\Omega }`$, and the magnitude of the relative angular momentum conjugate to the angle between $`𝐫_{nC}`$ and $`𝐑_{PT}`$ which we denote as $`l`$. The coupling between different $`\mathrm{\Omega }`$ channels is present only for $`\mathrm{\Delta }\mathrm{\Omega }=\pm 1`$, given by the Coriolis term. This sparse coupling term makes computations much easier in the body-fixed channel. Another advantage in the body-fixed frame is a natural truncation scheme set by the maximum value of $`\mathrm{\Omega }`$.
We adopt the model space of the radial coordinates, $`0<R<50`$ fm and $`0<r<60`$ fm, and discretize them with steps $`(\mathrm{\Delta }R,\mathrm{\Delta }r)=(0.2,1.0)`$ fm $`(\mathrm{\Delta }R,\mathrm{\Delta }r)=(0.3,0.6)`$ fm for <sup>11</sup>Be and and <sup>6</sup>He, respectively. In the case of <sup>11</sup>Be, we also performed a calculation with a finer mesh $`\mathrm{\Delta }r=0.6`$ fm and confirmed that the fusion probabilities coincide with the present results in a good accuracy. The maximum $`n`$-$`C`$ relative angular momentum is taken as $`l_{\mathrm{max}}=70`$ to obtain converged results. If the $`n`$-$`T`$ potential, $`V_{nT}`$ is switched off, one can use a much smaller cutoff, typically $`l_{\mathrm{max}}=4`$. The discrete-variable representation is used for the second differential operator Colbert92 . The time evolution of the wave packet is calculated employing the fourth order Taylor expansion method.
The initial wave packet in the incident channel is a spatially localized Gaussian wave packet multiplied with an incoming wave for the P-T relative motion while the $`n`$-$`C`$ wave function is a bound orbital corresponding to the halo structure. The average P-T distance of the wave packet is set as 22 fm and the average incident energy approximately equal to the barrier top energy. The time evolution is calculated until the wave packet consists only of outgoing waves and the average P-T distance reaches 30 fm.
A single wave packet solution contains information for a certain range of energy. We extract the fusion probability for a fixed incident energy by the energy projection procedure Yabana97 .
Since the ground state of <sup>11</sup>Be is known to be well described by a single neutron halo model, our calculation for <sup>11</sup>Be is expected to be realistic. On the other hand, a di-neutron model for <sup>6</sup>He may be too crude. In order to include two-neutron correlation, we need a three-body $`\alpha +n+n`$ treatment for <sup>6</sup>He. Then, this requires a four-body description for the reaction, which is beyond the present computational ability. We consider the present di-neutron treatment may provide a qualitative description for <sup>6</sup>He.
Now, let us show the calculated results. In Figs. 1 and 2, we show calculated fusion probabilities at the head-on collision ($`J=0`$) as a function of the incident energy. In Fig. 1, the fusion probability of <sup>11</sup>Be - <sup>209</sup>Bi (solid curve) is compared with that of <sup>10</sup>Be - <sup>209</sup>Bi (dotted). The latter is a result of a simple two-body calculation with the $`V_{CT}`$ potential. We also show, by a dashed line, a fusion probability of <sup>11</sup>Be calculated without the neutron-target potential, $`V_{nT}`$. Figure 1 indicates that the fusion probability is slightly suppressed by the presence of the halo neutron. The suppression is observed in both calculations with and without $`V_{nT}`$. This fact suggests that the suppression originates from the core-target Coulomb field which induces the dissociation of the halo neutron. The fusion probability of <sup>6</sup>He - <sup>238</sup>U system also shows a behavior similar to <sup>11</sup>Be (Fig. 2) The suppression of the fusion probability is seen to be even stronger than that of <sup>11</sup>Be. This may be due to a difference of the effective charges of the halo neutron: $`(1/11)(4e)`$ for <sup>11</sup>Be and $`(2/6)(2e)`$ for <sup>6</sup>He. The Coulomb dissociation effect is more significant for <sup>6</sup>He. There is only a minor difference between the solid and dashed curves for both cases, <sup>11</sup>Be and <sup>6</sup>He. We have confirmed that the results do not depend on choice of $`V_{nT}`$. Therefore, we conclude that reaction mechanism associated with $`V_{nT}`$ does not play a major role in the fusion process of halo nuclei. It may be rather surprising that a long-ranged attractive potential of a halo nucleus, which is expected in a folding picture of $`V_{CT}`$ and $`V_{nT}`$, does not contribute much to the fusion probability. This means that the reaction dynamics are very different from those naturally expected by a simple picture.
In order to understand why, we should look in detail at the time evolution of the wave packet. Then, we find the following picture for the dynamics of halo nuclei Yabana03 : When the core nucleus is decelerated by the target Coulomb field, the halo neutrons keep their incident velocity, leaving the core nucleus. This yields the Coulomb breakup of the projectile. After the breakup, the fusion takes place between two charged particles, the core and the target nuclei. The halo neutrons behave like a spectator. Based on this picture, we propose a mechanism that may explain the reason for the slight suppression of the fusion probability. The incident energy of the projectile is initially carried by the core and the halo neutron. After the Coulomb breakup, since the dissociated neutrons, as a spectator, keep their shared energy, the energy of the core nucleus is smaller than the total incident energy of the projectile. This practical decrease in the incident energy makes the fusion probability suppressed for a neutron-halo nucleus in comparison with the nucleus without halo neutrons. In this spectator picture, one may naively expect that effect of the breakup simply leads to an energy shift in curves of the fusion probability. However, our calculation indicates that the fusion probability above the barrier is suppressed, while the probability below the barrier is almost unchanged. Thus the effect is not that simple and depends on the incident energy.
In Refs.Hagino00 ; Diaz02 , a strongly enhanced sub-barrier fusion probability was reported for halo nuclei by solving a similar three-body problem. This apparently disagrees with our results which indicate the suppressed fusion probability. We found that this discrepancy originates from a slow convergence of the calculated result with respect to $`l_{\mathrm{max}}`$ Yabana03 . The treatment of the imaginary potential is also different : we include it between the core and the target, while it is in the coordinate of the center of mass of the whole projectile relative to the target in Ref. Hagino00 ; Diaz02 . We examined both treatments of the imaginary potential and found that the results are identical. In the coupled-discretized continuum-channels (CDCC) approach which is adopted in Refs.Hagino00 ; Diaz02 , a small value of $`l_{\mathrm{max}}`$ is applied for the $`n`$-$`C`$ relative angular momentum, typically $`l_{\mathrm{max}}=4`$. In Fig. 3, we show the fusion probability for different cutoff values of $`l_{\mathrm{max}}`$. If we adopt a small $`l_{\mathrm{max}}`$ value, a strong enhancement of the sub-barrier fusion probability is obtained in our three-body calculations. This is consistent to results of the CDCC calculation Hagino00 ; Diaz02 . However, increasing $`l_{\mathrm{max}}`$, the fusion probability gradually decreases, and eventually becomes smaller than the values without a halo neutron (dashed line). The result of $`l_{\mathrm{max}}=30`$ cannot be distinguished from that of $`l_{\mathrm{max}}=50`$. To get a converged result, therefore, one needs to adopt $`l_{\mathrm{max}}30`$. We confirmed that this slow convergence with respect to $`l_{\mathrm{max}}`$ is due to the neutron-target potential, $`V_{nT}`$. If we ignore $`V_{nT}`$ and includes only the Coulomb breakup induced by $`V_{CT}`$, the convergence is much faster, usually $`l_{\mathrm{max}}=4`$ is enough.
Extending the calculations to finite $`J`$ up to $`J=30`$, we achieve the total fusion cross section. In addition to $`l_{\mathrm{max}}`$, we must also specify a maximum value of $`\mathrm{\Omega }`$ which we denote as $`\mathrm{\Omega }_{\mathrm{max}}`$. We compare fusion probabilities with $`\mathrm{\Omega }_{\mathrm{max}}=0,5,`$ and $`10`$ for a fixed $`J`$, and find that the results with $`\mathrm{\Omega }_{\mathrm{max}}=5`$ and 10 are identical. Furthermore, a small difference between results of $`\mathrm{\Omega }_{\mathrm{max}}=0`$ and $`\mathrm{\Omega }_{\mathrm{max}}=5`$ is visible only for energies well above the barrier, $`E>42`$ MeV. Thus, we here adopt the no-Coriolis approximation ($`\mathrm{\Omega }_{\mathrm{max}}=0`$) to calculate the cross sections.
We show calculated fusion cross sections of <sup>10,11</sup>Be - <sup>209</sup>Bi in Fig. 4. The $`V_{CT}`$ for <sup>10</sup>Be - <sup>209</sup>Bi is constructed so that the fusion cross section of this system is reproduced by a simple potential picture. Adding a halo neutron, the fusion cross section decreases slightly for an entire energy region. This is because the suppression of the fusion probability observed in the $`J=0`$ also appears in finite $`J`$.
The measurement indicates that the fusion cross sections of <sup>11</sup>Be do not differ from those of <sup>10</sup>Be within the experimental error. At this stage, it is difficult to judge whether or not our calculated results of slight suppression in fusion cross section conflict with the measurements. In our calculation, a part of the halo neutron proceeds to the forward direction keeping the incident velocity after fusion. Therefore, the observation of these forward neutrons in the incomplete fusion process may support our picture. In Ref.Signorini04 , the forward neutron emission accompanying the fusion is discussed.
We next consider the fusion cross section of <sup>6</sup>He - <sup>238</sup>U. In Fig. 5, we show the calculated fusion cross sections of <sup>6</sup>He and <sup>4</sup>He on <sup>238</sup>U. The experimental data in Ref.Raabe04 indicate, again, that there is no difference between <sup>6</sup>He and <sup>4</sup>He fusion cross section within a statistical error. Our calculation indicates that the cross section for <sup>6</sup>He is suppressed from the one for <sup>4</sup>He at energies above the barrier. The magnitude of suppression in <sup>6</sup>He is larger than that in <sup>11</sup>Be at high energies. On the other hand, at sub-barrier energies, the calculated fusion cross sections of <sup>4</sup>He and <sup>6</sup>He are almost identical to each other. We again need to wait for measurements with high statistics to make a definite conclusion whether our calculation agrees with the experimental data.
In summary, we calculated the fusion cross section for nuclei with neutron halo structure. The reaction is described as the three-body model, and the three-body Schrödinger equation is solved exactly with the time-dependent wave-packet method. We compare calculated cross sections with measurements for <sup>11,10</sup>Be - <sup>209</sup>Bi and <sup>6,4</sup>He - <sup>238</sup>U collisions. Recent measurements of these systems suggest that there is no evidence for an enhancement in the fusion cross section of halo nuclei. Our calculation indicates that the fusion cross section is even suppressed by adding a halo neutron. Further measurements with increased statistics are desirable to obtain definite conclusion whether the fusion cross section is increased or suppressed by the presence of the halo neutron.
This work has been supported by the Grant-in-Aid for Scientific Research in Japan (Nos. 14540369 and 17540231). A part of the numerical calculations have been performed at RCNP, Osaka University and at SIPC, University of Tsukuba.
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# Correlated particle dynamics in concentrated quasi-two-dimensional suspensions
## 1 Introduction
The Brownian motion of colloid particles and macromolecules is correlated through hydrodynamic interactions, i.e., flows that the motion induces in the host liquid . In an unconfined suspension these correlations decay with inter-particle distance $`r`$ as $`1/r`$. They are positive, i.e., particles drag one another in the same direction. The long range of the interaction leads to strong many-body effects, manifest in an appreciable dependence of transport coefficients on concentration. A concentration effect of particular interest is hydrodynamic screening , which sets in over distances much larger than the inter-particle distance and renormalizes the prefactor of the $`1/r`$ pair interaction.
Colloids may sometimes be spatially confined by rigid boundaries, as in porous media, biological constrictions, or microfluidic devices. Attention has been turned recently toward the dynamics of such confined suspensions due to the emergence of microfluidic applications and the development of new visualization and manipulation techniques (digital video microscopy and optical tweezers ), allowing to study dynamics at the single- and few-particle level. The dynamics of confined suspensions have been studied also by computer simulations and have renewed the interest in related hydrodynamic problems .
We have recently demonstrated the dramatic effect that confinement between two flat surfaces has on the hydrodynamic pair interaction at large inter-particle distances . Such a quasi-two-dimensional (Q2D) suspension is portrayed in figure 1. Despite the confinement, the hydrodynamic interaction is still long-range, decaying as $`1/r^2`$ rather than $`1/r`$. There is “anti-drag” in the transverse direction, i.e., the correlation between two particles moving perpendicular to their connecting line is negative. Arguably the most striking finding, however, is that the concentration of the suspension has no effect on the large-distance correlation, i.e., there is no hydrodynamic screening. These three properties are in stark contrast with their unconfined counterparts mentioned above.
In the current paper we review these results and then extend the analysis to shorter inter-particle distances where the static pair-correlation of the suspension plays a crucial role. Throughout the paper we present the theoretical analysis along with the corresponding experimental results. Section 2 introduces the model system and the corresponding terminology, and section 3 describes the experimental setup and methods. In section 4 we briefly discuss the dynamics of a single particle and in section 5 address the hydrodynamic interaction between two isolated particles. Section 6, which constitutes the main part of the current work, deals with the three-body correction to the pair-interaction at finite concentration. Finally, in section 7 we conclude and discuss the results.
## 2 Model system
The geometry considered in this work is depicted in figure 1(b). Identical, spherical particles of radius $`a`$ are suspended in a liquid of viscosity $`\eta `$ and temperature $`T`$, confined in a slab of width $`w`$ between two planar solid surfaces. The $`\widehat{x}`$ and $`\widehat{y}`$ axes are taken parallel, and the $`\widehat{z}`$ axis perpendicular, to the surfaces, where $`z=0`$ is the mid-plane. The particles behave as hard spheres with no additional equilibrium interaction. For simplicity we consider cases where particle motion is restricted to two dimensions, i.e., to a monolayer lying at the mid-plane. We use the notation $`𝐫(𝝆,z)`$ for three-dimensional position vectors, where $`𝝆(x,y)`$ is a two-dimensional position vector in the monolayer. The area fraction occupied by particles is denoted by $`\varphi `$. The Reynolds number is very low, and the hydrodynamics, therefore, are well described by viscous Stokes flows . The confining boundaries are impermeable and rigid, imposing no-slip boundary conditions on the flow.
We characterize the dynamic pair correlation between two particles by the coupling mobilities $`B_{\mathrm{L},\mathrm{T}}^\mathrm{c}(\rho )`$ as functions of the inter-particle distance $`\rho `$. These are the off-diagonal terms in the mobility tensor of a particle pair, i.e., the proportionality coefficients relating the force acting on one particle with the change in velocity of the other. The two independent coefficients, $`B_\mathrm{L}^\mathrm{c}`$ and $`B_\mathrm{T}^\mathrm{c}`$, correspond, respectively, to the coupling along and transverse to the line connecting the pair. The dynamic correlation between two Brownian particles is similarly characterized by two coupling diffusion coefficients, $`D_{\mathrm{L},\mathrm{T}}^\mathrm{c}(\rho )`$, which, due to the Einstein relation, are simply related to the coupling mobility coefficients via the thermal energy, $`D_{\mathrm{L},\mathrm{T}}^\mathrm{c}(\rho )=k_\mathrm{B}TB_{\mathrm{L},\mathrm{T}}^\mathrm{c}(\rho )`$. \[In four coefficients were considered, $`D_{\mathrm{L},\mathrm{T}}^\pm `$, whose relation with the ones considered here is $`D_{\mathrm{L},\mathrm{T}}^\mathrm{c}=(D_{\mathrm{L},\mathrm{T}}^+D_{\mathrm{L},\mathrm{T}}^{})/2`$.\]
## 3 Experimental setup
The experimental system consists of an aqueous suspension of monodisperse silica spheres (diameter $`2a=1.58\pm 0.04`$ $`\mu `$m, density 2.2 g/cm<sup>3</sup>, Duke Scientific), undergoing Brownian motion while being tightly confined between two parallel glass plates in a sealed thin cell (figure 1). The inter-plate separation is $`w=1.76\pm 0.05`$ $`\mu `$m, i.e., slightly larger than the sphere diameter, $`2a/w0.90`$. Digital video microscopy and subsequent data analysis are used to locate the centres of the spheres in the field of view and then extract time-dependent two-dimensional trajectories. Details of the setup and measurement methods can be found elsewhere . Measurements were made at four values of area fraction, $`\varphi =`$ 0.254, 0.338, 0.547, $`0.619\pm 0.001`$. (Larger area fractions could not be checked because the suspension began to crystallize .) From equilibrium studies of this system we infer that, for the purpose of this study, the particles can be regarded as hard spheres.
The coupling diffusion coefficients as functions of the inter-particle distance are directly measured from the tracked trajectories as
$$D_\mathrm{L}^\mathrm{c}(\rho )=x_1(t)x_2(t)_\rho /(2t),D_\mathrm{T}^\mathrm{c}(\rho )=y_1(t)y_2(t)_\rho /(2t),$$
(1)
where $`x_i(t)`$ and $`y_i(t)`$ are the displacements of particle $`i`$ of the pair during a time interval $`t`$ along and transverse to their connecting line, respectively. The average $`_\rho `$ is taken over all pairs whose mutual distance falls in a narrow range ($`\pm 0.09`$ $`\mu `$m) around $`\rho `$.
## 4 Single particle
Consider a single particle whose centre, lying on the mid-plane between the two confining surfaces, is defined as the origin. A force $`𝐟_1`$ is applied to the particle in the $`i`$ direction parallel to the surfaces ($`i=x,y`$). As a result, the particle moves with velocity
$$u_{1i}=B_\mathrm{s}f_{1i}=B_0[(a/w)\mathrm{\Delta }_\mathrm{s}(a/w)]f_{1i},$$
(2)
where $`B_\mathrm{s}`$ is the self-mobility of the particle in the given geometry and $`B_0=(6\pi \eta a)^1`$ is its self-mobility in an unconfined liquid. Alternatively, we may consider a free Brownian particle. Its mean-square displacement during a time interval $`t`$ will be
$$\rho ^2(t)=4D_0[(a/w)\mathrm{\Delta }_\mathrm{s}(a/w)]t,$$
(3)
where $`D_0=k_\mathrm{B}TB_0`$. The dimensionless factor $`(a/w)\mathrm{\Delta }_\mathrm{s}(a/w)`$, representing the effect of confinement, becomes unity in the limit $`a/w1`$. Approximate expressions for this factor for larger values of the confinement ratio $`a/w`$ were the subject of many previous works (see and references therein), and their validity has been confirmed in recent experiments .
The particle motion makes the surrounding liquid flow. At distances much larger than $`w`$ the flow velocity is given by
$`v_i(𝐫)`$ $`=`$ $`(a/w)B_0\mathrm{\Delta }_{ij}(𝐫)f_{1j}`$
$`\mathrm{\Delta }_{ij}(𝝆,z)`$ $`=`$ $`\lambda _0{\displaystyle \frac{w^2}{\rho ^2}}\left(\delta _{ij}{\displaystyle \frac{2\rho _i\rho _j}{\rho ^2}}\right)H(z/w),`$ (4)
where $`i,j=x,y`$, the vertical component $`v_z`$ is exponentially small in $`\rho /w`$, and $`\lambda _0`$ is a dimensionless prefactor of order 1 dependent only on the confinement ratio $`a/w`$. The flow field (4) can be viewed as produced by an effective two-dimensional mass dipole (source doublet), located at the particle centre and oriented in the $`j`$ direction . The flow has a perpendicular profile $`H(z/w)`$ which vanishes on the two confining surfaces, $`H(1/2)=H(1/2)=0`$, and is normalized to 1 at the mid-plane, $`H(0)=1`$.
As argued in , the far flow has the dipolar shape (4) regardless of particle size. Changing the confinement ratio $`a/w`$ merely modifies the effective mass-dipole strength, i.e., the prefactor $`\lambda _0(a/w)`$. In the limit $`a/w0`$ one finds $`\lambda _0=9/16`$ . As will be shown below, we find for $`a/w0.45`$ $`\lambda _00.49`$. Since the maximum possible confinement ratio is $`a/w=1/2`$, it is inferred that the function $`\lambda _0(a/w)`$ actually changes very little with $`a/w`$. This weak dependence on the confinement ratio can be understood in terms of mutual compensation of two opposing effects: when $`a/w`$ is larger, on one hand, the particle displaces a larger liquid volume as it moves, yet, on the other hand, its self-mobility decreases \[i.e., $`(a/w)\mathrm{\Delta }_\mathrm{s}(a/w)`$ gets smaller\].
## 5 Pair interaction
Let us introduce a second particle whose centre lies at $`𝐫=(𝝆,0)`$. The particle is torque-free. It is also force-free in the $`xy`$ plane. Being confined to the mid-plane, it evidently cannot be assumed force-free in the $`z`$ direction. However, the flow (4) does not exert any force in that direction on a mid-plane-placed particle. Hence, as long as we restrict the discussion to such flows, the particle can be regarded as force-free. It is thus entrained by the flow with velocity
$$u_{2i}=\frac{1}{4\pi a^2}_A𝐫^{\prime \prime }v_i(𝐫^{\prime \prime })=\left(1\frac{4a^2}{3w^2}\right)v_i(𝝆,0),$$
(5)
where the integration is over the particle surface $`A`$. In obtaining the last equality we have assumed a parabolic (Poiseuille) vertical profile, $`H(z/w)=14z^2/w^2`$ . Equation (5) is, in fact, a manifestation of Faxen’s first law as applied to the Q2D geometry. For a force-free particle, Faxen’s law yields $`𝐮=𝐯+(a^2/6)^2𝐯`$. The dipolar flow (4) has $`(_{xx}+_{yy})𝐯=0`$ and $`_{zz}|_{z=0}𝐯=[w^2H^{\prime \prime }(0)/H(0)]𝐯`$, from which equation (5) readily follows.
Substituting equation (4) in equation (5), we obtain a relation between the force $`𝐟_1`$ exerted on the first particle and the resulting velocity change $`𝐮_2`$ of the second particle,
$$u_{2i}=[14a^2/(3w^2)](a/w)B_0\mathrm{\Delta }_{ij}(𝝆,0)f_{1j},$$
(6)
thus identifying the off-diagonal coupling terms of the pair-mobility tensor as $`B_{ij}^\mathrm{c}(𝝆)=[14a^2/(3w^2)](a/w)B_0\mathrm{\Delta }_{ij}(𝝆,0)`$. To get the coupling mobility along the line connecting the pair, we consider the relation between $`f_{1x}`$ and $`u_{2x}`$ for $`𝝆=\rho \widehat{x}`$. Using equation (4) we get
$`B_\mathrm{L}^\mathrm{c}(\rho )`$ $`=`$ $`[14a^2/(3w^2)](a/w)B_0\mathrm{\Delta }_{xx}(\rho \widehat{x},0)=(a/w)B_0\mathrm{\Delta }_\mathrm{L}(\rho )`$
$`\mathrm{\Delta }_\mathrm{L}(\rho )`$ $`=`$ $`\lambda w^2/\rho ^2.`$ (7)
Similarly, for $`𝝆=\rho \widehat{y}`$ the relation between $`f_{1x}`$ and $`u_{2x}`$ yields the coupling mobility transverse to the connecting line,
$`B_\mathrm{T}^\mathrm{c}(\rho )`$ $`=`$ $`[14a^2/(3w^2)](a/w)B_0\mathrm{\Delta }_{xx}(\rho \widehat{y},0)=(a/w)B_0\mathrm{\Delta }_\mathrm{T}(\rho )`$
$`\mathrm{\Delta }_\mathrm{T}(\rho )`$ $`=`$ $`\lambda w^2/\rho ^2.`$ (8)
In equations (7) and (8) we have defined
$$\lambda (a/w)=[14a^2/(3w^2)]\lambda _0(a/w).$$
(9)
This is a refinement of our previous analysis . In a Q2D geometry the so-called stokeslet approximation, equating the coupling mobility with the flow velocity per unit force, strictly holds only when $`a`$ is much smaller than both the inter-particle distance $`\rho `$ and the slab width $`w`$, whereupon $`\lambda \lambda _09/16`$. If $`a`$ is much smaller than $`\rho `$ but comparable to $`w`$ then, no matter how large the inter-particle distance may be, we have $`\lambda <\lambda _0`$. This is due to the second term in Faxen’s law, $`a^2^2𝐯`$. In an unconfined liquid it contributes a negligible correction to the stokeslet limit, $`O(a^2/r^2)`$, whereas in Q2D the newly introduced length $`w`$ leads to a significant correction of $`O(a^2/w^2)`$.
As found in equations (7) and (8), the pair hydrodynamic interaction in Q2D is very different from its unconfined counterpart. The decay with distance is faster, $`1/r^2`$ instead of $`1/r`$, yet the interaction is still long-range . (Its decay, in fact, is slower than near a single surface, where the interaction falls off as $`1/r^3`$ .) The transverse coupling is negative, i.e., particles exert “anti-drag” on one another as they move perpendicular to their connecting line. (In the unconfined case one has $`\mathrm{\Delta }_\mathrm{T}=\mathrm{\Delta }_\mathrm{L}/2`$, both coefficients being positive.) These properties are direct consequences of the far flow field (4). The $`1/\rho ^2`$ decay is that of the flow due to a 2D mass dipole. The negative transverse coupling is a result of circulation flows in the dipolar field .
For a pair of force-free Brownian particles, the hydrodynamic coupling derived above implies that their motions be correlated according to equation (1), with $`D_{\mathrm{L},\mathrm{T}}^\mathrm{c}(\rho )=(a/w)D_0\mathrm{\Delta }_{\mathrm{L},\mathrm{T}}(\rho )`$. Thus, the dimensionless couplings $`\mathrm{\Delta }_{\mathrm{L},\mathrm{T}}(\rho )`$ can be directly measured from the statistics of particle trajectories. These measurements for different area fractions $`\varphi `$ are presented in figure 2. The large-distance behaviour for all $`\varphi `$ values fits well the $`\pm \lambda w^2/\rho ^2`$ dependence of equations (7) and (8) with $`\lambda 0.36`$. The negative sign of $`\mathrm{\Delta }_\mathrm{T}`$ confirms the predicted “anti-drag” between particles located transverse to the direction of motion. Substituting $`\lambda 0.36`$ and $`a/w0.45`$ in equation (9) we find $`\lambda _0(0.45)0.49`$, which is only slightly smaller than the value for a vanishing confinement ratio, $`\lambda _0(0)=9/160.56`$.
## 6 Concentration effect
As the area fraction $`\varphi `$ is increased, the pair hydrodynamic interaction should become affected by the presence of other particles. In unconfined suspensions hydrodynamic screening sets in at distances much larger than the typical inter-particle distance and renormalizes the interaction by a concentration-dependent prefactor . Similarly, one expects the pair interaction in Q2D to have the form of equations (7) and (8) yet with a modified, $`\varphi `$-dependent prefactor. We find, however, that this is not the case, as is clearly demonstrated in figure 2(e).
Consider again particle 1, located at the origin and exerting a force $`𝐟_1`$ parallel to the confining surfaces. This creates a far flow velocity $`𝐯`$ according to equation (4) which, in the absence of any other particle, entrains particle 2 at $`𝐫=(𝝆,0)`$ with velocity $`𝐮_2𝐯(𝐫)`$ given by equation (6). We now introduce particle 3 at $`𝐫^{}=(𝝆^{},0)`$. It obstructs the flow and thus modifies $`𝐮_2`$.
The particle being force-free, the lowest force-distribution moment it can exert is a force dipole,
$$S_{ij}(𝐫^{})=\lambda _s(a/wB_0)^1\frac{1}{4\pi a^2}_A𝑑𝐫^{\prime \prime }(𝐫^{\prime \prime }𝐫^{})_iv_j(𝐫^{\prime \prime }),$$
(10)
where the integration is over the surface of particle 3, and $`\lambda _s`$ is a dimensionless factor of order 1. The force dipole, located at $`𝐫^{}`$, changes the flow velocity at $`𝐫`$ by
$$\delta v_i^s(𝐫,𝐫^{})=S_{jk}(𝐫^{})(a/wB_0)_j\mathrm{\Delta }_{ki}(𝐫𝐫^{}).$$
(11)
This holds regardless of confinement. In an unconfined system the far flows decay as $`1/r`$, leading to $`S(r^{})^2`$ and $`\delta v^s(r^{})^2|𝐫𝐫^{}|^2`$. Then, upon integration over the third-particle position $`𝐫^{}`$, one gets a correction $`1/r`$, which renormalizes the coefficient of the bare $`1/r`$ interaction. In Q2D, however, the situation is different. First, we realize that, due to the vanishing of $`v_z`$ and the symmetry of the mid-plane, only the terms $`i,j=x,y`$ of $`S_{ij}`$ in equation (10) are non-zero. Then, since the far flows decay as $`1/\rho ^2`$, one gets $`\delta v^s(\rho ^{})^3|𝝆𝝆^{}|^3`$ which, upon integration over $`𝝆^{}`$, yields a correction $`1/\rho ^4`$. At large distances this is much smaller than the bare pair interaction ($`1/\rho ^2`$) and will be neglected.
Since the force dipole created by particle 3 has a negligible far-field effect in Q2D, we proceed to the third moment of the force distribution,
$$T_{ijk}(𝐫^{})=\lambda _t(a/wB_0)^1\frac{1}{4\pi a^2}_A𝑑𝐫^{\prime \prime }(𝐫^{\prime \prime }𝐫^{})_i(𝐫^{\prime \prime }𝐫^{})_jv_k(𝐫^{\prime \prime }),$$
(12)
where $`\lambda _t`$ is another dimensionless prefactor depending only on the confinement ratio $`a/w`$. In principle, $`\lambda _t`$ (and $`\lambda _s`$) could be calculated using the short-range hydrodynamics in a Q2D geometry , yet we shall not pursue this technically complicated calculation here. The vanishing of $`v_z`$ and the mid-plane symmetry make all terms $`T_{ijk}`$, for which one or three of the indices are $`z`$, vanish. In addition, the terms with $`(i,j,k)=(x,y)`$ are negligible at large distances compared to those with two $`z`$ indices, since they involve two additional powers of $`1/\rho ^{}`$. We are thus left with $`T_{zzi}`$, $`i=(x,y)`$, yielding through equations (4) and (12)
$$T_{zzi}(𝝆^{})=\frac{1}{3}\lambda _ta^2\left(1\frac{12a^2}{5w^2}\right)\mathrm{\Delta }_{ij}(𝝆^{},0)f_{1j}.$$
(13)
In the integration we have assumed again a parabolic vertical profile, $`H(z/w)=14z^2/w^2`$ . Note that within this assumption the similar leading terms in $`1/\rho ^{}`$ from all higher moments of the force distribution, involving higher $`z`$ derivatives, vanish. The moment $`T(𝝆^{})`$ changes the flow velocity at $`𝝆`$ by
$$\delta v_i^t(𝐫,𝐫^{})=T_{zzj}(𝝆^{})(a/wB_0)_{zz}|_{z=0}\mathrm{\Delta }_{ji}(𝝆𝝆^{},0).$$
(14)
Combining equations (4), (5), (13) and (14), we find the correction to the velocity of particle 2 due to particle 3,
$`\delta u_{2i}`$ $`=`$ $`(a/wB_0)C_1\mathrm{\Delta }_{ij}(𝝆𝝆^{})\mathrm{\Delta }_{jk}(𝝆^{})f_{1k}`$ (15)
$`C_1`$ $`=`$ $`\lambda _t{\displaystyle \frac{8a^2}{3w^2}}\left(1{\displaystyle \frac{4a^2}{3w^2}}\right)\left(1{\displaystyle \frac{12a^2}{5w^2}}\right),`$
where the reference to $`z=0`$ is hereafter omitted for brevity.
Equation (15) yields the correction to $`𝐮_2`$ given a certain position $`𝝆^{}`$ of particle 3. We now wish to average this correction over all possible $`𝝆^{}`$,
$$\delta 𝐮_2=d^2\rho ^{}p(𝝆,𝝆^{})\delta 𝐮_2(𝝆,𝝆^{}),$$
(16)
where $`p(𝝆,𝝆^{})`$ is the probability density of finding particle 3 at $`𝝆^{}`$ given that particle 2 is at $`𝝆`$ and particle 1 is at the origin. Employing the superposition approximation, we take this probability as
$$p(𝝆,𝝆^{})\frac{\varphi }{\pi a^2}g(\rho ^{})g(|𝝆𝝆^{}|)\frac{\varphi }{\pi a^2}[1+h(\rho ^{})+h(|𝝆𝝆^{}|)],$$
(17)
where $`g(\rho )`$ is the pair correlation function of the Q2D suspension (normalized to 1 at $`\rho \mathrm{}`$), and $`h(\rho )=g(\rho )1`$. We have assumed that the monolayer of particles is a disordered 2D liquid, and thus the pair correlation has no angular dependence.
Using equations (15)–(17), and specializing to the relation between $`\delta u_{2x}`$ and $`f_{1x}`$, we obtain the average corrections to the coupling mobilities,
$`\delta B_{\mathrm{L},\mathrm{T}}^\mathrm{c}(\rho )`$ $`=`$ $`(a/w)B_0\delta \mathrm{\Delta }_{\mathrm{L},\mathrm{T}}(\rho )`$
$`\delta \mathrm{\Delta }_{\mathrm{L},\mathrm{T}}(\rho )`$ $`=`$ $`(\pi a^2)^1C_1\varphi {\displaystyle }d^2\rho ^{}[1+2h(\rho ^{})][\mathrm{\Delta }_{xx}(𝝆𝝆^{})\mathrm{\Delta }_{xx}(𝝆^{})+`$ (18)
$`\mathrm{\Delta }_{xy}(𝝆𝝆^{})\mathrm{\Delta }_{xy}(𝝆^{})],`$
where the longitudinal and transverse couplings are obtained, as in section 5, by taking $`𝝆=\rho \widehat{x}`$ and $`𝝆=\rho \widehat{y}`$, respectively. \[Note that equation (18), for $`h(\rho )=0`$, has exactly the same form as the one used in based on a less detailed analysis.\]
A brief examination of equation (18) raises two expectations. First, the integrand scales as $`(\rho ^{})^2|𝝆𝝆^{}|^2`$ which, upon integration over $`𝝆^{}`$, is expected to yield a correction $`1/\rho ^2`$, i.e., a renormalization of the prefactor of the bare coupling in accord with hydrodynamic screening. Second, since the fields $`\mathrm{\Delta }_{ij}`$ decay slowly with distance, one expects features (oscillations) in the static correlation function $`h(\rho )`$ to be smoothed by integration and thus be manifest only weakly in the dynamic coupling. As is shown below, due to the unique shape of the far flow in Q2D, both of these expectations turn out to be false.
The convolution integral of equation (18) is worked out in the Appendix. The result is
$$\delta \mathrm{\Delta }_{\mathrm{L},\mathrm{T}}(\rho )=C\varphi w^2\left[\frac{h(\rho )}{\rho ^2}2_\rho ^{\mathrm{}}𝑑\xi \frac{h(\xi )}{\xi ^3}+O(a^2/\rho ^4)\right],$$
(19)
where $`C=2(w^2/a^2)\lambda _0^2C_1`$, and the results for the longitudinal and transverse couplings are identical. Equation (19) is the main result of the current work. It gives the leading three-body correction to the pair hydrodynamic interaction in Q2D. If static pair correlations are neglected, $`h(\rho )=0`$, this correction vanishes, as was presented in . Moreover, the correction for a nonzero $`h(\rho )`$, although finite, is short-range: since the suspension is disordered, the static pair correlation $`h(\rho )`$ decays exponentially beyond an equilibrium correlation length. Thus, the correction in equation (19) decreases with distance much faster than the $`1/\rho ^2`$ decay of the bare interaction. This remarkable disagreement with the usual notion of hydrodynamic screening (i.e., concentration-dependent change of the prefactor at large distances), known from unconfined systems, is fully confirmed by our experimental measurements as is seen in figure 2(e).
The first term in equation (19) represents the leading correction to equations (7) and (8) as the inter-particle distance $`\rho `$ decreases. It scales as $`h(\rho )/\rho ^2`$ and, therefore, should be in phase with the oscillations of the static pair correlation. The second term is smaller, of order $`ah(\rho )/\rho ^3`$. Higher-order corrections, which have not been treated here, come from several sources. Corrections of $`O(a^2/\rho ^4)`$ arise from the integration in equation (18) (see Appendix), as well as the force-dipole terms ($`\delta v^s`$) discussed above. Terms of $`O(h^2/\rho ^2)`$ come from the approximation employed in equation (17). Finally, as $`\rho `$ becomes comparable to $`w`$, short-distance hydrodynamics set in and the far flow (4), on which our entire analysis has been relying, should be corrected . Thus, the expression formulated in equation (19) is relevant to suspensions of sufficiently high concentration, such that the static pair correlation $`h(\rho )`$ is significant at distances $`\rho >w`$ where the far flow holds.
To compare equation (19) with experiment we need the static pair correlation function for the various values of area fraction. These functions are directly measurable from snapshots of the Q2D suspensions such as the one shown in figure 1(a). (For more details, see .) The results are presented in figure 3.
The measured $`h(\rho )`$ are subsequently used in equation (19) to calculate the correction to the pair interaction. The results are shown in figure 4. With the coefficient $`\lambda `$ of the bare interaction already found, there is only one unknown coefficient, $`C`$, fitted as $`C0.85`$. Working back the definitions of the various coefficients introduced during the calculation, we find the coefficient of the third force moment for our experimental system ($`a/w0.45`$) to be $`\lambda _t1.77`$. Comparing the fits in figure 4 with those in figure 2, we see that the introduction of the concentration-dependent correction leads to a good agreement between the theory and the measured longitudinal interaction not only at asymptotically large distances (as was presented in ) but also at intermediate distances, $`\rho >2w`$. As anticipated, equation (19) is particularly successful for the high-concentration systems \[figure 4(c,d)\], where oscillations related to $`h(\rho )`$ are clearly observed (compare with figure 3). The discrepancies seen in figure 4 between the theory and experiment for $`\rho 2w`$ can be attributed to several factors, as discussed above, which have not been included in the current analysis.
Our calculation has yielded an identical correction to the transverse pair interaction, $`\delta \mathrm{\Delta }_\mathrm{T}=\delta \mathrm{\Delta }_\mathrm{L}`$. This does not agree with the measured transverse interaction, which does not exhibit a wiggly behaviour similar to that of $`\mathrm{\Delta }_\mathrm{L}`$ (see figure 2). Thus, we cannot currently offer an accurate analysis of the concentration effect on the transverse interaction at intermediate and short distances. A possible reason for the difference between the longitudinal and transverse modes may lie in the higher sensitivity of the latter to short-range hydrodynamics. To demonstrate this difference we show in figure 5 the deviations of the longitudinal and transverse bare interactions from their large-distance behaviour in the limit of very small particles ($`a/w1`$), as obtained from an exact solution . The departure of the transverse interaction from its asymptotic behaviour is found to be more pronounced than that of the longitudinal one.
## 7 Discussion
The correlated dynamics of particles in Q2D suspensions are very different from those in unconfined systems. Confinement affects the decay of the dynamic correlation with distance ($`1/r^2`$ instead of $`1/r`$) as well as its sign (the correlation transverse to the line connecting the particles becomes negative). In the current work we have focused on the effect of concentration, i.e., the presence of additional particles, on the pair interaction. Unlike unconfined suspensions, where the large-distance coupling changes with concentration, in Q2D suspensions increasing the concentration has no effect at large inter-particle distances (no hydrodynamic screening). The qualitatively different behaviour of Q2D suspensions has been theoretically and experimentally corroborated. It stems from the dipolar shape of the far flow induced by particle motion. As argued in , the key property of this flow is that it is governed by the displacement of liquid mass rather than the diffusion of liquid momentum.
Concentration-dependent effects on the pair correlation do set in in Q2D suspensions, yet only at intermediate and short distances. The range of these corrections is determined by the larger of two lengths: the range of the static pair correlation $`h(\rho )`$ of the suspension and the confinement width $`w`$. In the former case, which is valid for sufficiently high concentration, the spatial dependence of the dynamic coupling reflects the features of the static pair correlation. In this case we have been able to provide a quantitative account for the concentration effect on the longitudinal interaction at intermediate distances, which is in good agreement with the experimental measurements. In the latter case (range of $`h`$ shorter than $`w`$), the correction depends on short-range hydrodynamics whose analysis lies beyond the scope of the current work.
We have not treated the effect of particle motion perpendicular to the bounding surfaces. Such fluctuations must exist in practice, yet our results suggest that they have a minor effect at large and intermediate distances. This is expected since the flows produced by perpendicular fluctuations decay exponentially with distance and thus have only a short-range effect .
This research was supported by the Israel Science Foundation (77/03), the National Science Foundation (CTS-021774 and CHE-9977841) and the NSF-funded MRSEC at The University of Chicago. H.D. acknowledges additional support from the Israeli Council of Higher Education (Alon Fellowship).
## Appendix
We need to calculate the integral appearing in equation (18):
$$I(𝝆)=d^2\rho ^{}[1+2h(\rho ^{})][\mathrm{\Delta }_{xx}(𝝆^{})\mathrm{\Delta }_{xx}(𝝆𝝆^{})+\mathrm{\Delta }_{xy}(𝝆^{})\mathrm{\Delta }_{xy}(𝝆𝝆^{})],$$
(20)
where (omitting the prefactor)
$$\mathrm{\Delta }_{xx}(𝝆)=(x^2y^2)/\rho ^4,\mathrm{\Delta }_{xy}(𝝆)=2xy/\rho ^4.$$
(21)
Changing to polar coordinates, $`𝝆=(\rho ,\phi )`$ and $`𝝆^{}=(\rho ^{},\phi ^{})`$, we rewrite the integral as
$$I(𝝆)=d^2\rho ^{}[1+2h(\rho ^{})]\frac{\rho ^2\mathrm{cos}[2(\phi ^{}\phi )]+\rho ^22\rho \rho ^{}\mathrm{cos}(\phi ^{}\phi )}{\rho ^2[\rho ^2+\rho ^22\rho \rho ^{}\mathrm{cos}(\phi ^{}\phi )]^2},$$
(22)
from which it is evident that $`I(𝝆)=I(\rho )`$ is independent of $`\phi `$. Hence, we may set $`\phi =0`$. This isotropy immediately implies also that the corrections to the longitudinal and transverse couplings, proportional to $`I(\rho \widehat{x})`$ and $`I(\rho \widehat{y})`$, respectively, are identical.
The calculation of $`I(𝝆)`$ via equation (22) is a bit tricky for reasons similar to those encountered in electrostatics. We exclude two small areas $`a^2`$ around the (integrable) singularities at $`\rho ^{}=0`$ and $`𝝆^{}=𝝆`$, and divide the integration into three domains: (i) $`\rho ^{}(a,\rho a)`$, $`\phi ^{}(0,2\pi )`$; (ii) $`\rho ^{}(\rho +a,\mathrm{})`$, $`\phi ^{}(0,2\pi )`$; (iii) $`\rho ^{}(\rho a,\rho +a)`$, $`\phi ^{}(a/\rho ,2\pi a/\rho )`$. The contribution from the inner domain (i) vanishes upon integration over $`\phi ^{}`$ for any $`𝝆`$. (Note that the static pair correlation $`h(\rho )`$ does not have an angular dependence.) Integration over $`\phi ^{}`$ in the outer domain (ii) gives $`2\pi [1+2h(\rho ^{})]/(\rho ^{})^3`$. Subsequent integration over $`\rho ^{}`$ in domain (ii) yields $`2\pi _\rho ^{\mathrm{}}𝑑\rho ^{}[1+2h(\rho ^{})]/(\rho ^{})^3+O(a/\rho ^3)`$. However, due to the integrable singularity at $`𝝆^{}=𝝆`$, the intermediate narrow domain (iii) has a finite contribution as well, of opposite sign: $`\pi [1+2h(\rho )]/\rho ^2+O(a/\rho ^3)`$.
Thus, in the absence of static correlations, $`h(\rho )=0`$, the leading terms in small $`a/\rho `$ from domains (ii) and (iii) cancel and $`I(\rho a)=0`$. This result has already been obtained in . A more detailed inspection reveals that the $`O(a/\rho ^3)`$ corrections from domains (ii) and (iii) cancel as well and, hence, $`I=O(a^2/\rho ^4)`$.
In the presence of static correlations, $`h0`$, we finally get
$$I(𝝆)=2\pi \left[\frac{h(\rho )}{\rho ^2}2_\rho ^{\mathrm{}}𝑑\xi \frac{h(\xi )}{\xi ^3}\right]+O(a^2/\rho ^4).$$
(23)
## References
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# The double-mode nature of the HADS star GSC 00144-03031 and the Petersen diagram of the class Based on observations partly collected at S. Pedro Mártir and Sierra Nevada Observatories
## 1 Introduction
The double–mode High Amplitude Delta Scuti (HADS) stars have been considered for a long time similar to double–mode Cepheids, i.e., evolved stars pulsating in the first two radial modes only. A remarkable exception is constituted by the AC And subclass, composed of stars pulsating in the first three radial modes (see Tab. 9 in Rodríguez 2003). Very recently, nonradial modes have been found combined with the radial ones, making these stars more similar to the low–amplitude, unevolved $`\delta `$ Sct stars which are quite common between the Zero–Age and the Terminal–Age Main Sequences (Poretti 2003). Therefore, HADS stars lost their definition of “simple double– or single–mode radial pulsators”. A natural question arises: do all the double–mode HADS stars display nonradial modes? To give a first answer to this question we investigated in detail the light curve of GSC 00144-03031, a new member of the class. Moreover, since other double–mode pulsators have been found in large–scale projects, it is timely to discuss the entire class by comparing the new, more robust observational results with the theoretical models.
The variability of GSC 00144-03031 was discovered when observing the field around HD 43587, a candidate primary target for the asteroseismic space mission COROT (COnvection, ROtation and planetary Transits; Baglin et al. 2002). The field around this star has been explored in the search for secondary targets (Poretti et al. 2005) and the STARE telescope <sup>1</sup><sup>1</sup>1http://www.hao.ucar.edu/public/research/stare/stare.html detected the variability of GSC 00144-03031 on the night of 16–17 March 2002 (JD 2452350). The STARE light curve clearly showed the effect of multiperiodicity in this large amplitude variable ($`\mathrm{\Delta }V>`$0.40 mag).
## 2 Observations
Owing to the recent growing interest in high–amplitude pulsators, new observations were planned to investigate the pulsational content of GSC 00144-03031 in detail. The first long subset was obtained at S. Pedro Mártir Observatory: 249 measurements were collected using the simultaneous $`uvby`$ Danish photometer mounted on the 1.5–m telescope, on five nights from November 14 to 23, 2003. The observations were performed with the scope of monitoring the brightest COROT targets and therefore the time sampling is not very dense (Fig. 1, upper row). HD 43913 was used as comparison star. Later in the same winter, dedicated intensive CCD monitoring was performed at Athens University Observatory: a 0.40–m telescope equipped with a SBIG ST-8 camera and a $`V`$ filter was used on eight nights from January 26, 2004 to March 2, 2004. GSC 00144-02970 was used as comparison star and GSC 00144-03140 as check star: both are located in the same field of view and the monitoring was practically continuous (Fig. 1, middle row). A third dataset was obtained at Beersel Hills Observatory (Belgium), using a 0.40–m telescope equipped with an ST10 XME CCD camera and a $`V`$ filter. The observations were performed on ten nights from February 16, 2004 to March 28, 2004 (Fig. 1, lower row). Other less conspicuous subsets were collected at SETEC Observatory (0.30–m telescope equipped with a ST–8 CCD camera and a $`V`$ filter, 168 measurements on one night), at Sierra Nevada Observatory (OSN; 116 $`uvby`$ data obtained at the 0.90–m telescope on three nights) and by MM (0.20–m telescope equipped with a Kodak 401 CCD camera and a $`V`$ filter, 483 measurements on four nights). The time baseline spans 133 d.
## 3 Frequency analysis
Observations obtained at two different longitudes (Europe and America) strongly reduced the daily aliasing effect. To perform the various steps of the frequency analysis, we used the least–squares power spectrum method (Vaniĉek 1971), allowing us to detect one by one the constituents of the light curve. After each detection, we refined the frequency values by applying the MTRAP code (Carpino et al. 1987): the refined values were introduced as known constituents in the new search. Such a procedure is particularly suitable for a multimode, high–amplitude pulsation as it keeps the relationships between the detected terms (i.e., 2$`f_1`$, $`f_1+f_2`$, …) locked. Moreover, it does not require any data prewhitening as amplitudes and phases of the known constituents are recalculated for each new trial frequency, always subtracting the exact amount of signal for any detected term.
As a first step, we analyzed each subset separately. The frequency analysis of the 249 $`uvby`$ data obtained at S. Pedro Mártir (SPM subset) clearly detected eight terms: $`f_1`$=17.22 d<sup>-1</sup>, $`f_2`$=22.31 d<sup>-1</sup>, 2$`f_1`$, 3$`f_1`$, $`f_1+f_2`$, $`f_2f_1`$, 2$`f_1+f_2`$, 2$`f_1f_2`$. The term with the smallest amplitude is the 2$`f_1f_2`$ one (0.0041 mag). The power spectrum does not show any well–defined structure after introducing these eight terms as known constituents, even though the peaks related to the $`f_1`$ and $`f_2`$ coupling terms (3$`f_1f_2`$, 3$`f_1+f_2`$,4$`f_1`$) are embedded among others, without standing out clearly. The amplitude of the noise slightly increases toward low–frequencies (0.0019 mag for $`f<`$5 d<sup>-1</sup>, 0.0016 mag for $`f<`$10 d<sup>-1</sup>), while it is quite small at high–frequencies (0.0009 mag for 40$`<f<`$50 d<sup>-1</sup>). No difference was found in the frequency analysis of different colours; the rms residuals are 0.0170, 0.0094, 0.0084 and 0.0078 in $`u,v,b`$ and $`y`$, respectively: the large scatter in $`u`$–light is due to some instrumental problems with that photomultiplier tube.
The analysis of the 1285 $`V`$ measurements carried out at Beersel Hills Observatory (BHO subset) allowed us to detect the same terms as in the SPM subset. The amplitude of the noise is 0.0008 mag in the 0$`<f<`$100 d<sup>-1</sup> interval, with a small increase at low frequencies (0.0015 mag for $`f<`$7d<sup>-1</sup>). The solution with eight terms leaves an rms residual of 0.0087 mag. Other combination terms are visible in the residual spectrum, slightly higher than the noise level. The magnitude differences between the comparison and checks star yield a standard deviation of 0.010 mag; the values on individual nights range from 0.005 to 0.014 mag.
The data obtained at Athens University Observatory (AUO subset) are the most homogeneous ones. Eleven terms have been detected: $`f_1`$, $`f_2`$, 2$`f_1`$, 3$`f_1`$, $`f_1`$+$`f_2`$, $`f_2f_1`$, 2$`f_1`$+$`f_2`$, 2$`f_1f_2`$ (i.e., those detected in the BH and SPM subsets) plus 3$`f_1f_2`$, 3$`f_1+f_2`$, 4$`f_1`$ (i.e., the terms embedded in the noise in the BH and SPM subsets). The solution with 11 terms leaves an rms residual of 0.0128 mag: it is higher than that of previous subsets, but the high number of datapoints (2421) allowed us to uncover terms with smaller amplitudes, by decreasing the amplitude of the white noise (0.0007 mag in the 0–100 d<sup>-1</sup> interval). The observed scatter is similar to that observed for the check star GSC 0144-03140; the standard deviations range from 0.010 to 0.018 mag.
The three solutions of the $`V`$ or $`y`$ light curves are very coherent, detecting the same terms with very similar amplitudes and phases (Table 1; a cosine series has been used). Therefore, once the least–squares solutions for each subset were obtained, the mean magnitudes were aligned on the standard value $`V`$=10.198. At this stage, the other subsets (SETEC, OSN and MM) were added: their samplings do not allow a reliable frequency analysis by themselves, but combined with the previous ones they supply a better frequency resolution. The alignments were made by applying the same procedure as above. No systematic difference was found between $`V`$ and $`y`$ data. Hence, we have at our disposal a global dataset composed of 4722 measurements, which constitutes a solid baseline for the careful search for small amplitude modes.
The analysis of the global dataset evidenced the large contribution of the $`f_1`$ term, i.e., 17.22 d<sup>-1</sup> (Fig. 2, upper row, left panel). This periodicity is characterized by a strong asymmetric light curve since the harmonics 2$`f_1`$ (upper row, right panel), 3$`f_1`$ (middle row, left panel) and 4$`f_1`$ (middle row, right panel) were detected when going further in the analysis. The second contribution in amplitude comes from the $`f_2`$=22.31 d<sup>-1</sup> term (upper row, right panel): this term has a sine–shaped light curve, since no harmonic was detected. The power spectra also show the contribution from the coupling terms between the two independent frequencies: the $`f_1+f_2`$ and $`f_2f_1`$ terms (middle row, left panel), the $`2f_1f_2`$ and $`2f_1+f_2`$ terms (middle row, right panel) and the $`3f_1f_2`$ and $`3f_1+f_2`$ terms (lower row, left panel). The last power spectrum (lower row, right panel) does not show any reliable residual peak. The mean noise level is around 0.52 mmag in the whole 0–80 d<sup>-1</sup> interval, ranging from 0.39 mmag in the 50–80 d<sup>-1</sup> interval to 0.67 mmag in the 0–10 d<sup>-1</sup> one.
Table 1 also summarizes the least–squares solution of the global dataset. We also inspected the residual light curves to see if there were systematic deviations in some part of them, but we found nothing. Formal errors of the two independent frequencies were calculated following Montgomery & O’Donoghue (1999). Figure 3 shows the light curve of each independent term. We stress the fact that no additional independent term with an amplitude larger than 0.6 mmag was found: therefore, GSC 00144-03031 looks like an authentic pure double–mode HADS star.
## 4 The other double–mode HADS stars
The class of double–mode HADS stars has been enriched with various new entries provided by large–scale projects. In addition to GSC 00144-03031, we discovered new double–mode pulsators when performing systematic analyses of some public databases: OGLE (Optical Gravitational Lensing Experiment; Szymański 2005 and Udalski et al., 1997), NSVS (Northern Sky Variability Survey; Woźniak et al. 2004), ASAS (All Sky Automated Survey; Pojmanski 2002 and 2003) In some cases, we planned new photometric observations. The enlarged sample allows us to define better the characteristics of the double–mode HADS stars. As shown by the analysis of GSC 00144-03031, there are double–mode radial pulsators which do not show any hint of nonradial terms, even when the amplitude threshold is very small. On the other hand, Poretti (2003) showed how the radial double–mode pulsation can be also paired to a nonradial one.
Table LABEL:hads lists the stars where the $`F`$ and 1O modes have been detected. The formal errors on the frequency values are usually better than 1$`10^4`$ d<sup>-1</sup>, which in turn implies an error on the ratio values better than 1$`10^5`$. The new variables are described below.
### 4.1 BW2\_V142 and BW1\_V207
Moskalik & Poretti (2002) reported on the double–mode nature of some stars contained in the OGLE–I database. We note that BW1\_V109$``$Macho 119.19574.1169 (Alcock et al. 2000). The first least–squares solutions of the BW2\_V142 and BW1\_V207 light curves are given in Table 3. Poretti (2000) briefly reported on BW2\_V142 just showing the light curves on the two terms (note that there is a misprint in Fig. 1 about the value of the $`f_1`$ term). GSC 00144-03031 and BW2\_V142 fill the gap between shortest periods (namely those of SX Phe and GSC 02583–00504) and the group of HADS stars which pulsate with periods around 0.10 d.
### 4.2 GSC 03949-00811
The star is catalogued as NSVS 3234596 by Woźniak et al. (2004). The frequency analysis of the NSVS data allowed us to detect the $`f_1`$=7.687632 d<sup>-1</sup>, $`f_2`$=5.890971 d<sup>-1</sup>, $`f_1+f_2`$ and $`f_1f_2`$ terms. The amplitude of the 1O–mode is larger than that of the $`F`$–mode. The scatter is quite high (0.034 mag) and it was not possible to evidence other coupling or additional terms. The least–squares solution of the NSVS light curve is reported in Tab. 3.
### 4.3 GSC 06047-00749
Three subsets are available for this star: ASAS (ASAS3 094303-1707.3; Pojmanski 2002 and 2003), NSVS (NSVS 15731136; Woźniak et al. 2004) and a new one obtained by NB (0.20–m telescope equipped with a SBIG ST7e CCD and a $`V`$ filter, 402 measurements on four nights in June, 2004). The terms $`f_1`$=13.0691 d<sup>-1</sup>, $`f_2`$=10.0829 d<sup>-1</sup> and $`f_1+f_2`$ were detected in all the subsets. The least–squares solutions supplied a smaller amplitude for the NSVS subset, since these data are in unfiltered light. The solution reported in Table 3 was obtained by combining the ASAS and the NB subsets.
### 4.4 HD 224852
The variability of HD 224852$``$NSV 14800 was discovered by Strohmeier & Ott (1967) and, independently, by ASAS (ASAS3 000116-6037.0; Pojmanski 2002 and 2003). In the framework of our project, the star was observed by NB (1671 measurements on nine nights from October to December, 2004). The $`f_1`$=8.1920 d<sup>-1</sup>, $`f_2`$=6.3356 d<sup>-1</sup>, 2$`f_1`$, $`f_1+f_2`$, $`f_1f_2`$ terms were detected in the ASAS subset. The observations carried out by NB are more dense and also the 2$`f_2`$, 2$`f_1+f_2`$ terms were detected. In spite of the small differences in the amplitudes between the two subsets, the frequency analysis of the combined dataset evidenced other coupling terms: 3$`f_1`$, 2$`f_1f_2`$, $`f_1+2f_1`$. The search for new terms was stopped when the noise at the lowest frequencies became predominant. The 1O–mode has an amplitude larger than the $`F`$–one and the light curves of both periods have a non–sinusoidal shape. Table 3 lists the least–squares solution of the combined dataset; the mean magnitude has been aligned to that of the ASAS subset. The solution with 10 terms leaves an rms residual of 0.023 mag; the amplitude of the noise is 3.3 mmag for $`f<`$5 d<sup>-1</sup>, 1.6 mmag for 10$`<f<`$50 d<sup>-1</sup>. Also in this case there is no hint of the excitation of additional modes.
### 4.5 GSC 04257-00471
The star is catalogued as NSVS 3324715$``$NSVS 3361415 by Woźniak et al. (2004). We could detect the $`f_1`$=5.7538 d<sup>-1</sup>, $`f_2`$=7.5140 d<sup>-1</sup>, 2$`f_1`$ and $`f_1f_2`$ terms (Table 3). The amplitude of the 1O–mode is quite small, but the NSVS data are in unfiltered light and therefore damping is expected with respect to $`V`$–light. We note that GSC 03949-00811 and GSC 04257-00471 fill an important gap toward long periods.
## 5 Modelling the double–mode HADS stars
The period ratios reported in Table LABEL:hads are plotted in the middle panel of Fig. 4. Most of the stars are grouped around $`\mathrm{log}P_0=1.0`$ and $`P_1/P_0`$=0.772. It appears that the period ratio decreases toward the longest periods. However, some deviating points are observed at both extrema. In particular, the decrease in the period ratio values around $`\mathrm{log}P_0=0.8`$ appears much more rapid than the one calculated by Petersen & Christensen-Dalsgaard (1996; their Fig. 4) using the star sample available at that time. Therefore, we calculated a new set of theoretical models.
### 5.1 The theoretical sequences in the Petersen diagram
Stellar models were computed with the evolutionary code CESAM (Morel 1997). The numerical precision as well as the model mesh grid (2000 mesh points approximately) were optimized to compute oscillations. The equation of state CEFF (Christensen-Dalsgaard & Daeppen 1992) was used, also including the Coulombian correction to the classical EFF (Eggleton et al. 1973). A weak electronic screening in these reactions was assumed (Clayton 1968), accordingly to the evolutionary stages considered in this work. Opacity tables were taken from the OPAL package (Iglesias & Rogers 1996), complemented at low temperatures ($`T10^4`$K) by the tables provided by Alexander & Ferguson (1994). For the atmosphere reconstruction, the Eddington $`T(\tau )`$ law (grey approximation) was considered.
The transformation from heavy element abundances with respect to Hydrogen \[M/H\] into concentration in mass $`Z`$ assumes an enrichment ratio of $`\mathrm{\Delta }Y/\mathrm{\Delta }Z=2`$ and $`Y_{\mathrm{pr}}=0.235`$ and $`Z_{\mathrm{pr}}=0`$ for Helium and heavy element primordial concentrations, respectively. Convective transport was described by the classical Mixing Length theory, with efficiency and core overshooting parameters set to $`\alpha _{ML}=l_m/\mathrm{H}_p=1.8`$ and $`d_{\mathrm{over}}=l_{\mathrm{over}}/\mathrm{H}_p=0.2`$, respectively. The latter parameter is prescribed by Schaller et al. (1992) for intermediate mass stars. $`\mathrm{H}_p`$ corresponds to the local pressure scale-height, while $`l_m`$ and $`l_{\mathrm{over}}`$ represent respectively the mixing length and the inertial penetration distance of convective elements.
Theoretical oscillation spectra were computed from the equilibrium models described above. For this purpose the oscillation code *Filou* (Suárez 2002) was used. The upper panel of Fig. 4 shows various sequences calculated from our models. The effects of changing the mass keeping \[Fe/H\]=0.00 dex are shown by the three solid lines. The other two lines show the effect of decreasing metallicity keeping M=1.80 M: from \[Fe/H\]=–0.50 dex (dashed line) to \[Fe/H\]=–1.00 dex (dotted line).
### 5.2 The evolutionary tracks
The physical parameters of HADS stars are summarized by McNamara (2000). $`uvby\beta `$ photometry of GSC 00144-03031 ($`V`$=10.198, $`by`$=+0.138, $`m_1`$=0.194, $`c_1`$=0.849, $`\beta `$=2.805; SPM and OSN observations described above) yields the following parameters: E<sub>b-y</sub>=+0.005, M<sub>V</sub>=+2.00, T<sub>eff</sub>=7590 K, $`\mathrm{log}g`$=3.98, \[Fe/H\]=–0.05 dex (TEMPLOGG method; Rogers 1995, Kupka & Bruntt 2001).
The connection between double–mode HADS stars is confirmed when considering the evolutionary tracks. The tracks have been calculated for \[Fe/H\]=0.00 dex, a value which can be applied to GSC 00144-03031, RV Ari (\[Fe/H\]=0.00 dex), BP Peg ($`0.02`$), AI Vel ($`0.15`$) and VX Hya (0.05). Figure 5 shows the positions in the $`\mathrm{log}`$T<sub>eff</sub>–M<sub>bol</sub> plane for these stars. The position of GSC 00144-03031 suggests M=1.75 M; such a position is similar to that of RV Ari. Note how the points in Fig. 5 describe a narrow instability strip for double–mode HADS stars.
## 6 Interpreting the Petersen diagram
The decreasing trends in the theoretical sequences are characterized by a long standstill between $`\mathrm{log}P_0=1.30`$ and $`\mathrm{log}P_0=0.90`$ (Fig. 4, upper panel). In this interval the $`P_1/P_0`$ values are more sensitive to the changes in metallicity than to changes in mass. On the other hand, small differences in mass produce strong differences in the sequences at the longest periods. We can compare these predictions with the Petersen diagram as determined from the stars listed in Table LABEL:hads.
### 6.1 Stars with $`P_0`$=0.10 d
The period ratios cluster around 0.772 for $`\mathrm{log}P_0=1.0`$. This point in the Petersen diagram is a well–defined one on the basis of the enlarged sample: the theoretical models have to match it. We calculated a large variety of models, but a satisfactory match is provided only for \[Fe/H\]=0.00 dex and 1.70$`<`$M$`<`$2.00 M. These constraints are the same obtained from the positions of double–mode HADS stars in the $`\mathrm{log}`$T<sub>eff</sub>$`M_{\mathrm{bol}}`$ plane (Fig. 5). We note that the changes in metallicity shift the theoretical sequences toward higher period ratio values (Fig. 4, upper panel). Therefore, the small dispersion of the points related to stars with $`P_0`$=0.10 d suggests a very similar metallicity (Fig. 4, middle panel). The period ratio value corresponding to AE UMa (Table LABEL:hads) is quite normal (0.7738), even though AE UMa is a slightly metal deficient star (\[Fe/H\]=$`0.50`$ dex). However, it is also less massive than other stars (M=1.49 M; McNamara 2000): the two effects balance each other and generate the normal period ratio value.
### 6.2 The shortest periods
The extreme values observed at the shortest periods refer to GSC 00144-03031, GSC 02583–00504, SX Phe and BW2\_V142. If we admit that stars with $`\mathrm{log}P_0=1.0`$ and $`P_1/P_0`$=0.772 constitute a homogeneous sample, GSC 00144-03031 and GSC 02583–00504 can be considered as the natural extension toward the shortest period of this sample. Indeed, the sequences with \[Fe/H\]=0.00 dex and masses in the 1.70–2.00 M range match very well also these points (Fig. 4, solid lines in the upper panel). The high period–ratio values shown by SX Phe and BW2\_V142 (crosses in the lower panel of Fig. 4) can be explained by a lower metallicity (Fig. 4, dotted line in the upper panel), which is a well–established fact for SX Phe (\[Fe/H\]=–1.40 dex, McNamara 2000), the prototype of Pop. II short–period pulsating stars. The same argument holds for BW2\_V142, as this star is located in the galactic bulge. We can infer that the spread of the period ratio values at the shortest periods can be explained by different metallicities.
### 6.3 The longest periods
V575 Lyr, V703 Sco, BrhV128, GSC 04257-00471 and GSC 03949-00811 are the stars defining the rapid decrease at $`\mathrm{log}P_0>0.90`$. This decrease is in agreement with the theoretical sequences and hence these stars are naturally linked to the ones at $`P_0`$=0.10 d. However, the smaller the mass, the more rapid the decrease is (Fig. 4, solid lines in the upper panel). Therefore, only the M=1.90 M models are able to match the points related to V575 Lyr, V703 Sco and BrhV128 (Fig. 4, open circles in the lower panel) and only the M=2.00 M models are able to match the points related to GSC 04257-00471 and GSC 03949-00811 (Fig. 4, open squares in the lower panel).
Two long–period stars show period ratios higher than 0.770, i.e., HD 224852 and VX Hya (triangles in the lower panel of Fig. 4). Such values are admitted by models with M$`>2.00`$ M only. Indeed, VX Hya is a very luminous (see Fig. 5) and massive (more than 2.30 M, McNamara 2000) star. It is for sure a Pop. I star (\[Fe/H\]=+0.05 dex) and therefore we can rule out a connection with the Pop. II stars SX Phe and BW2\_V142.
Therefore the spread in the period ratio values observed at the longest periods can be explained by differences in the masses. These differences in masses cannot be distinguished at $`\mathrm{log}P_0<0.90`$, where the sequences are almost coinciding. We note that the rapid decrease in the $`P_1/P_0`$ values for M$`<1.80`$ M could explain some unusual ratio (as for example Macho 104.20389.1202, $`\mathrm{log}P_0=0.80`$, $`P_1/P_0`$=0.751; Alcock et al. 2000), without invoking the excitation of nonradial modes (see Fig. 5 in Poretti 2003).
## 7 Conclusion
The new observations of GSC 00144-03031 confirm the existence of pure radial double–mode pulsators. Thanks to large–scale projects it has been possible to discover several new double–mode HADS stars and to discuss the properties of the class. Detailed models have been calculated and compared with the observed period ratios and the physical parameters determined on the basis of $`uvby\beta `$ photometry. Our models suggest that double–mode HADS stars constitute a very homogeneous class of variable stars. The Petersen diagram for these stars can be reconstructed using models having \[Fe/H\]=0.00 dex and M=1.70–2.00M. Different masses generate similar period ratios for $`\mathrm{log}P<`$0.90, but theoretical period ratios predict differences for $`\mathrm{log}P>`$0.90. More massive stars show higher values: this spread is actually observed in the Petersen diagram. The sequences with 1.90$`<`$M$`<`$2.00 M match most of the observed period ratios in a very satisfactory way. On the other hand, short–period Pop. II stars show slightly higher period ratios owing to their different metallicity and they clearly deviate from the Pop. I sequence.
###### Acknowledgements.
The authors wish to thank the STARE team for the attribution of observing time to the COROT project. J. Vialle checked the English form of the manuscript. This publication makes use of data partially taken from the Northern Sky Variability Survey created jointly by the Los Alamos National Laboratory and University of Michigan. The NSVS was funded by the Department of Energy, the National Aeronautics and Space Administration, and the National Science Foundation. Part of the collected data were acquired with equipment purchased thanks to a research fund financed by the Belgian National Lottery (1999). PGN, KDG and VNM acknowledge financial support from the Special Account for Research Grants No 70/3/7187 (HRAKLEITOS) and 70/3/7382 (PYTHAGORAS) of the National and Kapodistrian University of Athens, Greece. RA and JAB acknowledge financial support from grants AyA2001-1571 and ESP2001-4529-PE of the Spanish National Research plan. JCS acknowledges the financial support from the Spanish Plan of Astronomy and Astrophysics, under project $`AYA200304651`$, from the Spanish “Consejería de Innovación, Ciencia y Empresa”, from the “Junta de Andalucía” local government, and from the Spanish Plan Nacional del Espacio under project ESP2004-03855-C03-01. EP acknowledges financial support from PRIN 2003 “Continuità e discontinuità nella formazione della nostra Galassia”.
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# The high-velocity clouds and the Magellanic Clouds
## 1 Introduction
The existence of HI high-velocity clouds (HVCs) at high Galactic latitudes has intrigued astronomers from the time of their discovery in 1963. The early observations induced the belief that all HVCs had negative radial velocities, and that these clouds could be intergalactic gas falling onto the Galaxy. The recognition of violent activity in the Galactic disk and the discovery of large complexes of high positive velocity clouds and of the Magellanic Stream opened a whole gamut of possibilities for the nature and origin of the HVCs such as “Galactic fountain” clouds, material stripped from the Magellanic Clouds and extragalactic objects. The lack of adequate distance indicators makes it difficult to discriminate among the different alternatives.
The first interpretations of the HVCs were discussed by Oort (oorta (1966), oortb (1969), oortc (1970)). Wakker & van Woerden (wakker-woerden (1997)) have given a recent review of the HVC phenomenon. Pöppel (poeppel (1997)) reviewed the possible role played by the HVCs in the local interstellar medium. Although a large amount of observational data has now been accumulated, such as very sensitive HI 21 cm surveys (Morras et al. morras (2000); de Heij et al. 2002a , 2002b ; Lockman et al. lockman (2002); Putman et al. putman (2002); Putman et al. 2003a ), the detection of molecules (Richter et al. richter (2001)) and H$`\alpha `$ emission in HVCs (Tufte et al. tofte (1998); Putman et al. 2003b ) and the distance constraints to some HVCs (Wakker wakker (2001)), as well as a large number of theoretical studies that have been dedicated to this topic (eg. Quilis & Moore quilis (2001); Espresate et al. espresate (2002); Sternberg et al. sternberg (2002); Maloney & Putman maloney (2003)), the HVCs remain enigmatic.
The theories for the origin and nature of the high velocity clouds can be divided into three categories: intergalactic (eg. Bajaja et al. bajaja1 (1987); Blitz et al. blitz (1999); Braun & Burton braun (2000)), circumgalactic (eg. Kerr & Sullivan III kerr (1969); Hulsbosch & Oort hulsbosch (1973)) and Galactic (eg. Bregman bregman (1980); Verschuur verschuur (1993)). On the basis of an analysis of the sky and velocity distributions of the HVCs we will demonstrate in this paper that the majority of these clouds constitutes a circumgalactic flux or stream of clouds related to the Magellanic Stream. The idea that the HVCs may be fragments of Magellanic material precipitating toward the galactic disk was first suggested by Giovanelli (giovanelli (1981)) and Mirabel (mirabela (1981)). The main theories on the origin of the Magellanic Stream itself consider that the stream consists of gas extracted from the Magellanic Clouds. The proposed mechanisms that have been studied in detail are tidal forces of the Galaxy, friction forces of the gaseous galactic halo and a collision between the Large and Small Magellanic Clouds (Gardiner et al. gardiner2 (1994); Heller & Rohlfs heller (1994); Moore & Davis moore (1994); Gardiner & Noguchi gardiner (1996)). The tidal forces affect equally both gas and stars of the Magellanic Clouds. However there are no stars in the Stream, indicating that tidal interaction could not be the primary cause of its formation. On the other hand, the weakness of the stripping drag hypothesis is that the gas density of the Galactic halo should be very low. Hence, to account for the formation of the Magellanic Stream and the rest of the HVCs, we will here take a different point of view that includes only the collision between the Large and Small Magellanic Clouds as the mechanism that triggered the process.
## 2 A simple interpretation of the sky distribution and kinematics of the high velocity clouds
In general the observations of HVCs provide the HVC positions projected on the sky and their radial velocities. Hence, we do not know the distances and tangential velocities of the HVCs. Even with this limitation, if the HVCs are a homogeneous population with a common origin, in principle it is possible to find the space distribution and the three-dimensional velocity distribution of the HVC population from the analysis of the observed distributions of the sky positions and radial velocities of the HVCs. This is a typical case of an inverse problem in astronomy (Craig & Brown craig (1986); Merritt merritt (1996); Saha saha (1998)).
Nowadays a vast amount of observational material is available on HVCs in the form of surveys whose combination covers the whole sky, and that therefore provide complete and reliable statistical information. Wakker (wakker-privado (2002)) compiled a catalog of HVCs with 11000 entries, using the surveys of Hulsbosch & Wakker (hulsbosch-wakker (1988)), Bajaja et al. (bajaja (1985)) for $`23^{}<\delta <18^{}`$ and Morras et al. (morras (2000)) for $`\delta <23^{}`$. Each entry contains the celestial position coordinates, $`\mathrm{}`$ and $`b`$, the radial velocity $`\rho `$ and other parameters that characterize a detected high-velocity profile component (i.e. a piece of cloud intercepted by the radiotelescope beam). An HVC can be defined by a set of points ($`\mathrm{}`$, $`b`$, $`\rho `$), or cloud elements, that satisfy criteria of continuity in the ($`\mathrm{}`$, $`b`$, $`\rho `$)-space. For the study of HVC distributions, we will represent all the HVC components as separate points ($`\mathrm{}`$, $`b`$, $`\rho `$); although each point can be associated with a subset of points (defining an HVC or complex of HVCs) that are not necessarily independent. In statistical considerations one should take this into account (see Sect. 4). We will represent the distributions of HVCs in the conventional form as orthogonal projections of points in the ($`\mathrm{}`$, $`b`$, $`\rho `$)-space upon the planes $`\mathrm{}`$-$`b`$, $`\mathrm{}`$-$`\rho `$ and $`b`$-$`\rho `$, although a three-dimensional representation is more elegant and compact.
On the basis of the data of Wakker’s catalog, we display the distributions of the positions and radial velocities of the HVCs in Figs. 1-a, 1-b and 1-c.
For purposes of comparison, in Figs. 2-a, 2-b and 2-c we represent the HVC distributions corresponding to a second set of data with an improvement in spatial resolution taken from Putman et al. (putman (2002)) and Giovanelli (giovanelli2 (1980)) for $`\delta >15\mathrm{°}`$.
The enclosing curve of the longitude-velocity distribution (Figs. 1-b and 2-b) is approximately sinusoidal and that of the latitude-velocity distribution (Figs. 1-c and 2-c) follows a cosine law. In consequence, the radial velocity distribution of the HVCs can be represented by the following formula:
$$\rho =Vsin(\mathrm{})cos(b),$$
(1)
where $`V450\mathrm{km}\mathrm{s}^1`$. This distribution law admits a simple interpretation (or “inversion”), namely: it reflects a roughly uniform flow of HVCs that comes approximately from galactic co-ordinates $`\mathrm{}_0=90\mathrm{°}`$ and $`b_0=0\mathrm{°}`$ and that encloses the whole Galaxy. If $`V_{\mathrm{}}`$ and $`V_f`$ are the velocities of the Sun and of the HVC flow with respect to the Galactic center respectively, the radial velocities of the HVCs with respect to the LSR are given by
$$\rho =V_{\mathrm{}}sin(\mathrm{})cos(b)V_fcos(\mathrm{}\mathrm{}_0)cos(bb_0).$$
(2)
Substituting the values of $`\mathrm{}_0`$ and $`b_0`$ into Eq. (2), we obtain Eq. (1) with $`V=V_{\mathrm{}}V_f`$. The velocity of the HVC flow is of the order of the Solar velocity ($`220\mathrm{km}\mathrm{s}^1`$), and has the same direction, but its sense is reversed. In other words, the Sun is immersed in the flow of HVCs, moving counter stream. This flow of clouds seems to circulate largely in the halo or in the outskirts of the Galaxy. In the next section, we will find a clue that makes it possible to connect the flow of HVCs with the Magellanic Clouds and the Magellanic Stream.
## 3 Model for the origin and dynamical evolution of the high velocity clouds
In this section we will try to determinate the probable origin of the HVC flow described in Sect. 2 and the initial conditions of the phenomenon; with this we will construct a model to study the trajectories of the HVCs throughout the halo and the disk of the Galaxy. Our first objective is to show that the Magellanic Clouds have played a central role in the origin of the HVCs. For this purpose, we need to calculate the orbital plane of the Magellanic Clouds. We adopt throughout this paper a Cartesian coordinate system $`(X,Y,Z)`$ with the origin at the Galactic center, the X-axis pointing in the direction of the Sun’s Galactic rotation, the Y-axis pointing in the direction from the Galactic center to the Sun, and the Z-axis pointing toward the Galactic north pole. Since the Magellanic Clouds are subject to a central force due to the Galactic halo (we will ignore the forces of the galactic disk), their orbital plane can be characterized by the vector normal to this plane, namely $`N=r\times \dot{r}`$, where $`r=(X,Y,Z)`$ is the position and $`\dot{r}=(\dot{X},\dot{Y},\dot{Z})`$ is the velocity vector of the center of mass of the combined Clouds at the present epoch in the Galactocentric rest frame. We use the observational parameters of the LMC, which are better determined, as representative of the Magellanic Clouds as a whole. We adopt the following parameters for the LMC: $`\mathrm{}=280\mathrm{°}.5`$ and $`b=32\mathrm{°}.9`$ for the Galactic coordinates, $`50.1\text{kpc}`$ for the distance from the Sun, $`262.2\mathrm{km}\mathrm{s}^1`$ for the heliocentric systematic velocity (van der Marel et al. marel (2002)), $`\mu _\alpha =0.00486\mathrm{}\mathrm{yr}^1`$ and $`\mu _\delta =0.00034\mathrm{}\mathrm{yr}^1`$ for the mean proper motion (see Table 1 of van der Marel et al. marel (2002)). To correct for the reflex motion of the Sun and to obtain the positions and velocities in the Galactocentric frame of rest, we use the IAU values $`R_{\mathrm{}}=8.5\text{kpc}`$ and $`V_{\mathrm{}}=220\mathrm{km}\mathrm{s}^1`$ and the “basic” solar motion of $`16.5\mathrm{km}\mathrm{s}^1`$ towards the direction of $`\mathrm{}=53\mathrm{°}`$ and $`b=25\mathrm{°}`$ (Allen allen (1973)). With these adopted values, we calculate the position and velocity vectors of the LMC
$$r_{\mathrm{LMC}}=(41.30,1.29,27.43)\text{kpc}$$
(3)
and
$$\dot{r}_{\mathrm{LMC}}=(218,66,187)\mathrm{km}\mathrm{s}^1.$$
(4)
The resulting plane of the orbit of the LMC is represented in Fig. 3. The angle between the orbital plane and the plane perpendicular to the Galactic disk is only of the order of $`10\mathrm{°}`$. i.e the plane of the LMC orbit is almost perpendicular to the Galactic plane. The Galactic longitude of the nodal line is nearly $`82\mathrm{°}`$, and thus the orbital plane is almost coincident with the $`X`$-$`Z`$ plane (cf. Fig.1 of Murai & Fujimoto murai (1980)). Therefore, the three-dimensional velocity vectors representing the HVC flow described in Sect. 2 are nearly parallel to the plane of the orbit of the Magellanic Clouds. This indicates a dynamical relationship between the HVCs and the Magellanic Clouds. Our working hypothesis will be that the HVCs were ejected from the Magellanic Clouds as a result of internal processes in the Magellanic Clouds, activated likely by the collision of the LMC and SMC.
### 3.1 Basic equations of the model
The basic problem we should consider here is the motion of a test particle representing an HVC in the gravitational fields of the Galaxy and the Magellanic Clouds, and friction forces due to the gaseous disk of the Galaxy.
In our computations, we use the spherical gravitational potential of the halo $`\psi =V_c^2\mathrm{ln}r`$ for the Galaxy (Murai & Fujimoto murai (1980); Lin & Lynden-Bell lin (1982)), which gives a flat rotation curve with a constant circular velocity, $`V_c=220\mathrm{km}\mathrm{s}^1`$, so that the gravitational force of the Galaxy exerted on a particle of unit mass is $`V_c^2\frac{r}{r^2}`$. For the Magellanic Clouds, we use the gravitational potential of the LMC and ignore that of the SMC. We consider the SMC as another test particle (see Sect.3.2). We assume that the LMC has a Plummer-type potential with an effective radius $`K=3\mathrm{kpc}`$ (Murai & Fujimoto murai (1980)), giving a gravitational force per unit mass of $`\frac{GM_{\mathrm{LMC}}r^{}}{(r^2+K^2)^{3/2}}`$, where $`r^{}(=rr_{\mathrm{LMC}})`$ denotes the position of the test particle with respect to the LMC, and $`M_{\mathrm{LMC}}`$ is the total mass of the LMC estimated at $`0.87\times 10^{10}\mathrm{M}_{\mathrm{}}`$ (van der Marel et al. marel (2002)).
The friction force per unit mass that acts on an HVC can be expressed as $`F_\mathrm{d}=\epsilon \frac{nS}{m}v_\mathrm{r}v_\mathrm{r}`$, where $`v_\mathrm{r}`$ is the relative velocity between the HVC and the gaseous medium in which the cloud moves, $`n`$ is the density of the gaseous medium at the position of the HVC, S is the head cross-section of the HVC, m is the HVC mass, and $`\epsilon `$ is a dimensionless quantity $`1`$. A measure of the ratio $`\frac{m}{S}`$ is the column density $`N_\mathrm{H}`$ of the HVC being studied, an observable quantity. We adopt a mean $`N_\mathrm{H}`$ of $`7.2\times 10^{18}\mathrm{atoms}\mathrm{cm}^2`$ for the HVCs. In our simplified model, the Galactic halo is empty of gas, thus $`n`$ depends only upon the density distribution of the gaseous disk of the Galaxy. The density $`n`$ may be conveniently expressed in cylindrical coordinates as $`n(R,z)=\frac{\sigma (R)}{2h(R)}\mathrm{exp}[\frac{z}{h(R)}]`$, where $`\sigma (R)`$ and $`h(R)`$ are the surface density of interstellar HI and the scale height of the thickness of the HI gas layer at $`R(=\sqrt{x^2+y^2})`$, respectively. We denote by $`x,y`$ and $`z`$ the Cartesian components of the position vector $`r`$ of the test particle or HVC. From a fit of Wouterloot et al.’s (wouterloot (1990)) Table 1 we derive $`\sigma (R)=1.0875\times 10^{21}\mathrm{exp}[(\frac{8.24R}{1.06\times 8.24})^2]\mathrm{atoms}\mathrm{cm}^2`$ and $`h(R)=105.8\mathrm{exp}[(\frac{3.37R}{4.23\times 3.37})^2]\mathrm{pc}`$ with $`R`$ in kpc. Assuming that the gas of the Galactic disk rotates circularly with constant velocity $`V_c`$, the velocity vector of the gaseous medium at the position of the HVC is $`(\frac{V_c}{R}y,\frac{V_c}{R}x,0)`$ and hence $`v_\mathrm{r}=(\dot{x}\frac{V_c}{R}y,\dot{y}+\frac{V_c}{R}x,\dot{z})`$. Under these conditions, the equations of motion of a test particle are
$$\ddot{r}+\frac{V_c^2r}{r^2}+\frac{GM_{\mathrm{LMC}}(rr_{\mathrm{LMC}})}{(rr_{\mathrm{LMC}}^2+K^2)^{3/2}}F_\mathrm{d}(r,\dot{r})=0.$$
(5)
To solve this system of equations, we need first the solution of the equations of motion of the LMC in the gravitational field of the Galaxy, namely:
$$\ddot{r}_{\mathrm{LMC}}+\frac{V_c^2r_{\mathrm{LMC}}}{r_{\mathrm{LMC}}^2}=0.$$
(6)
### 3.2 The collision between the Large and Small Magellanic Clouds
Several authors agree that the Magellanic clouds had a close encounter (separation less than 3 kpc) between 200 and 500 Myr ago (Murai & Fujimoto murai (1980); Gardiner et al. gardiner2 (1994); Heller & Rohlfs heller (1994); Moore & Davis moore (1994); Gardiner & Noguchi gardiner (1996)). The morphological peculiarities of the Clouds such as the widely scattered distribution of HII regions around the optical bar of the LMC and the severe fragmentation of the SMC could be understood in terms of this collision.
The encounter time $`t_\mathrm{e}`$ of the Clouds is an important parameter in our model, and we want to calculate its value congruently with the other parameters adopted above. The orbit of the LMC can be computed straightforwardly by means of the numerical integration of Eq. (6) with the initial conditions given by Eqs. (3) and (4) (see Fig. 3). For simplicity we assume that the orbit of the LMC is not essentially altered by the presence of the SMC because of its lesser mass. Knowing the present position and velocity of the center of mass of the SMC, we compute the orbit of the SMC by using Eq. (5) with $`F_\mathrm{d}=0`$ and the solution of the orbit of the LMC. Thus we obtain the time of the closest approach of both orbits. However, the position and velocity of the SMC are not known with great precision. Therefore, in the equations we only introduce the radial component of the SMC velocity, based on the heliocentric radial velocity of $`148\mathrm{km}\mathrm{s}^1`$ for the SMC that was accurately measured with optical and radio-astronomical methods (Hardy et al. hardy (1989)), and leave as unknowns the proper motions and the distance, which were used for determining the tangential components of the velocity. Solving Eq. (5) with the help of Eq. (6), we obtain $`r_{\mathrm{SMC}}`$ as a function of $`\mu _\alpha `$, $`\mu _\delta `$, the distance $`d`$ and time $`t`$. By requiring that $`r_{\mathrm{SMC}}(\mu _\alpha ,\mu _\delta ,d,t)=r_{\mathrm{LMC}}(t)`$, we have three equations with three unknowns for each chosen past time; their solutions give us $`\mu _\alpha (t)`$, $`\mu _\delta (t)`$ and $`d(t)`$. Choosing the set of theoretical parameters compatible with the observational ones, we found that $`\mu _\alpha `$, $`\mu _\delta `$ and $`d`$ should be close to $`0.00432\mathrm{}\mathrm{yr}^1`$, $`0.00108\mathrm{}\mathrm{yr}^1`$ (cf. Kroupa & Bastian kroupa (1997)) and 60.3 kpc (cf. Westerlund westerlund (1990)) respectively, and that the collision between the Clouds occurred 570 $`\pm 350`$ Myr ago. The uncertainty of $`\pm 350`$ Myr in the value of $`t_\mathrm{e}`$ was derived from that of the present velocity of the LMC (Eq. 4), $`\mathrm{\Delta }\dot{r}_{\mathrm{LMC}}(\pm 23,\pm 36,\pm 35)\mathrm{km}\mathrm{s}^1`$ (see Eq. 55 of van der Marel et al. marel (2002)), with the assumption that the rest of the adopted parameters like the distances to the Large and Small Magellanic Clouds and the gravitational potential used for the halo of the Galaxy are almost exact. The orbits of the Clouds are displayed in Figs. 4-a and 4-b.
The velocity of the LMC at the time of the collision $`t_\mathrm{e}=570\text{Myr}`$ was
$$\dot{r}_{\mathrm{LMC}}(t_\mathrm{e})=(100,12,141)\mathrm{km}\mathrm{s}^1,$$
(7)
while that of the SMC was $`\dot{r}_{\mathrm{SMC}}(t_\mathrm{e})=(100,167,155)\mathrm{km}\mathrm{s}^1`$. Thus, during the encounter, the velocity of the SMC relative to the LMC was of order $`156\mathrm{km}\mathrm{s}^1`$; i.e. very close to the escape velocity from the LMC ($`v_\mathrm{e}=\sqrt{\frac{2GM_{\mathrm{LMC}}}{K}}160\mathrm{km}\mathrm{s}^1`$). To explain this, we propose that the Magellanic Clouds formed a binary system in the past and that it was broken by a sudden mass loss of the system taking place in the epoch of the collision. The collision compressed the interstellar gas of the Clouds, inducing a burst of star formation. The massive stars of this new stellar generation expelled the surrounding interstellar matter in a few millions of years at velocities greater than the escape velocity of the system, through supernova explosions and strong stellar winds. A part of the material ejected from the Clouds at that time formed the HVCs. The HVCs left likely the Magellanic Clouds as magnetized clouds of plasma (eg. Olano olano1 (1999)). Magnetic fields in HVCs have been detected by measuring Zeeman splitting of the 21-cm HI line (Kazès et. al. kazes (1991)). This magnetic self-confinement of the HVCs could explain the HVC longevity as discrete entities. This idea is coherent with the fact that the HVCs are highly ionized (Sembach et al. sembach (1999); Lehner et al. lehner (2001); Sembach et al. sembach2 (2003)). The HVCs are still generally thought of as neutral entities. However, at an early stage they were probably fully ionized. Subsequently, the evolution of their physical conditions permitted the recombination of part of the protons and electrons. On the other hand, the explanations generally given for the cloud confinement are that they are either bound by gravitationally dominant dark matter or confined by an external pressure provided by a hot gaseous Galactic halo (Espresate et al. espresate (2002); Stanimirović et al. stanimirovic (2002); Sternberg et al. sternberg (2002)).
### 3.3 Initial conditions of the high velocity clouds
We conjecture that the HVCs were ejected from the Magellanic Clouds as the consequence of very energetic interactions of a large number of stars born during the collision of the Clouds and the cloudy interstellar medium of the Clouds. The relevant initial conditions, namely the positions, velocities and times of the ejections of the HVCs, should therefore be characterized by distribution functions. For simplicity we assume that the HVCs were ejected from a relatively small region of the Clouds, at the central point of the collision, during a short period of time, i.e. at $`t_\mathrm{e}=`$-570 Myr. The position vector of the central point of the collision (see Sect. 3.2 and Figs. 4-a and 4-b) is then
$$r_{\mathrm{LMC}}(t_\mathrm{e})=(86.39,15.53,15.68)\mathrm{kpc},$$
(8)
and the initial position vector for all HVCs is $`r_0=r_{\mathrm{LMC}}(t_\mathrm{e})`$. For the initial systematic or mean velocity of the HVCs we adopt the velocity of the mass center of the LMC at the time $`t_\mathrm{e}`$, which is given by Eq. (7). The errors in Eqs. (7) and (8), i.e. in the position of the LMC in phase space at $`t_\mathrm{e}`$, are $`\mathrm{\Delta }\dot{r}_{\mathrm{LMC}}(t_\mathrm{e})(\pm 3.87,\pm 0.72,\pm 1.03)\times \mathrm{\Delta }\dot{r}_{\mathrm{LMC}}\mathrm{km}\mathrm{s}^1`$ and $`\mathrm{\Delta }r_{\mathrm{LMC}}(t_\mathrm{e})(\pm 0.79,\pm 0.20,\pm 1.21)\times \mathrm{\Delta }\dot{r}_{\mathrm{LMC}}\mathrm{kpc}`$, obtained by propagation of the errors in $`\dot{r}_{\mathrm{LMC}}`$ and $`t_\mathrm{e}`$ (Sect. 3.2).
Thus the initial total velocity of an HVC can be written as
$$\dot{r}_0=\dot{r}_{\mathrm{LMC}}(t_\mathrm{e})+v_\mathrm{p},$$
(9)
where $`v_\mathrm{p}`$ is the initial peculiar velocity of the HVC. The number of HVCs in the initial velocity volume $`dv_\mathrm{p}`$ at $`v_\mathrm{p}`$ is $`dN=f(v_\mathrm{p})H(v_\mathrm{p}v_\mathrm{e})dv_\mathrm{p}`$ at time $`t_\mathrm{e}`$. We assume that the distribution function $`f(v_\mathrm{p})`$ follows Schwarzschild’s ellipsoidal law:
$$f(v_\mathrm{p})=\frac{N}{(2\pi )^{3/2}\mathrm{\Sigma }_1\mathrm{\Sigma }_2\mathrm{\Sigma }_3}\mathrm{exp}\frac{1}{2}((\frac{\dot{x}_\mathrm{p}}{\mathrm{\Sigma }_1})^2+(\frac{\dot{y}_\mathrm{p}}{\mathrm{\Sigma }_2})^2+(\frac{\dot{z}_\mathrm{p}}{\mathrm{\Sigma }_3})^2)),$$
(10)
where $`\mathrm{\Sigma }_1`$, $`\mathrm{\Sigma }_2`$ and $`\mathrm{\Sigma }_3`$ are the three velocity dispersions in the directions of the three co-ordinate axes and $`\dot{x}_\mathrm{p}`$, $`\dot{y}_\mathrm{p}`$ and $`\dot{z}_\mathrm{p}`$ are the components of the peculiar velocity vector $`v_\mathrm{p}`$. The Heavyside function $`H`$, which is 1 for $`v_\mathrm{p}v_\mathrm{e}>0`$ and 0 otherwise, takes into account that only the clouds with velocities greater than the escape velocity $`v_\mathrm{e}`$ became HVCs. In general form, the function of initial distribution of the HVCs may be written as $`F(r_0,\dot{r}_0,t)=\delta (r_0r_{\mathrm{LMC}}(t_\mathrm{e}))\delta (tt_\mathrm{e})f(\dot{r}_0\dot{r}_{\mathrm{LMC}}(t_\mathrm{e}))H(\dot{r}_0\dot{r}_{\mathrm{LMC}}(t_\mathrm{e})v_\mathrm{e})`$, where $`\delta `$ is the standard delta function. The total number of HVCs (or the number of test particles used in the simulation) can be expressed by $`N_{\mathrm{HVC}}=_0^{\mathrm{}}F(r_0,\dot{r}_0,t)𝑑r_0𝑑\dot{r}_0𝑑t`$, which makes it possible to determine $`N`$ of Eq. (10).
## 4 Results of the model and comparison with the observations
To reproduce the main features of the position and velocity distributions of the HVCs we calculated the orbits of about 850 test particles, which represent the HVCs. Distributing these test particles in phase space according to the distribution functions defined in Sect. 3.3, we can assign an initial position and velocity to each test particle, and calculate its present position and velocity by means of the equations of motion given in Sect. 3.1. The three dispersions of the velocity ellipsoid in Eq. (10) are the free parameters of the model. Their values can be estimated by comparing the predictions of the model with the observations. We found that the case particular of a three-dimensional Gaussian distribution is compatible with the observations, namely $`\mathrm{\Sigma }_1=\mathrm{\Sigma }_2=\mathrm{\Sigma }_3=\mathrm{\Sigma }=120\mathrm{km}\mathrm{s}^1`$. Thus, $`\mathrm{\Sigma }`$ is the only free parameter of the model.
To consider the possibility of non-isotropic expulsion of material, we tried other initial ellipsoidal distributions with different values and spatial orientations of their velocity dispersions with respect to the co-ordinate axes, but none of them could produce better results than those of a Gaussian distribution. This simple distribution reproduces fairly well the kinematics and sky distributions of the Magellanic stream, the “classic” stream and the leading stream (Putman et al. 2003a ), as well as those of the rest of the HVC population. Certainly, quantitative comparisons between theoretical and observed distributions in the way developed by Saha (saha (1998)) would be desirable. However, it is beyond the scope of the present work. A proof of the robustness of our method is that the results of the simulations are rather insensitive to the uncertainties of the time of the collision, $`\mathrm{\Delta }t_\mathrm{e}`$, and the velocity, $`\mathrm{\Delta }\dot{r}_{\mathrm{LMC}}(t_\mathrm{e})`$, and position, $`\mathrm{\Delta }r_{\mathrm{LMC}}(t_\mathrm{e})`$, of the centroid of the initial Gaussian distribution.
The results of the model are shown in Figs. 5-a, 5-b, 5-c, 5-d and 5-e. Comparing Fig. 5-a with Figs. 1-a and 2-a, Fig. 5-b with Figs. 1-b and 2-b, and Fig. 5-c with Figs. 1-c and 2-c, we see that the theoretical distributions are in good agreement with the corresponding observed ones. An exception is a group of HVCs that lies in a region centered at $`\mathrm{}=180^{}`$, $`b=20^{}`$, whose velocities are not explained by the model. However, we should remember the simplifications and uncertainties of the observational parameters introduced into the model, as well as the probable existence of other sources of HVCs, apart from the Magellanic Clouds. Studies of HVCs in external galaxies (see Schulman et al. schulman (1996); Wakker & Woerden wakker-woerden (1997); Jiménez-Vicente & Battaner battaner (2000); Miller et. al. miller (2001)) will enable us to make a comparison with the HVC system associated with our own Galaxy and to decide the relevance of the other possible sources of HVCs, such as Galactic fountains (Bregman bregman (1980); Houck & Bregman houck (1990)). The part of the HVC population of Magellanic origin having LSR radial velocities between -100 and +100 $`\mathrm{km}\mathrm{s}^1`$ and lying at low Galactic latitudes (see Figs. 5-b and 5-c) should be indistinguishable from the interstellar clouds of the Galactic disk that are parts of supernova shells or that participate simply in the Galactic rotation.
There is a clear difference between the two Galactic hemispheres. In the northern hemisphere the population of HVCs is dominated by a few extended, complex structures (the A, C and M objects). When comparing the observed HVC distributions with the theoretical ones, we should remember this. An HVC sampled at regular angular intervals is characterized by a set of neighboring points in the ($`\mathrm{}`$, $`b`$, $`\rho `$) space whose number increases proportionally with the solid angle subtended by the cloud in the sky (see Sect. 2). The few northern complexes of HVCs cover a large area of the sky. Hence their densities of points in the observed distributions relative to the densities of points in the corresponding theoretical distributions should be greater than the relative densities of points corresponding to the southern HVCs. The results of our simulation agree with this picture (cf. Figs. 1-a and 5-a, and Figs 1-c and 5-c). According to our model, the explanation of these asymmetries between the northern and southern hemispheres is that in the northern Galactic hemisphere there are far fewer HVCs and they lie at much smaller distances (see Fig. 5-e). Therefore, they subtend far larger angles on the sky. The importance of the model is that it provides a probable velocity field and space distribution of the HVCs in the three dimensions (Figs. 5-d and 5-e). This information cannot be obtained from the observations so far.
The group of HVCs spreads throughout a volume of approximately cubic 300 kpc. The mean velocity of the HVC flow is similar to that of the Magellanic Clouds, since the HVCs started their trajectories with the mean or systematic velocity of the Clouds. To illustrate this effect further, we apply the model to a set of test particles with initial peculiar velocities of the same magnitude $`v_\mathrm{p}=`$cont, but with different directions. At the present time, these test particles should form a stream that moves approximately in the same direction and sense as the Magellanic Clouds (Fig. 6-a). In addition, this cloud of particles is elongated in the direction of the velocity of rotation around the Galactic center, an effect due to the differential rotation and well known in the evolution of expanding Galactic shells (e.g. Olano olano2 (1982); Palous, Franco & Tenorio-Tagle palous (1990); Olano olano3 (2001)). Notice that the degree of concentration of the particles around the Magellanic orbit increases as the initial peculiar velocity decreases (cf. Figs. 6-a and 6-b). Also, notice that the lower the peculiar velocities of the particles, the greater the number density of particles (or HVCs), according to the initial distribution function Eq. (10). Both these effects contribute to create in the sky the appearance of the Magellanic Stream (see Fig. 5-a).
From the distances to the HVCs predicted by the model and the density column data of Wakker’s (wakker-privado (2002)) list we estimated the total mass of the HVCs by means of the formula $`M_\mathrm{t}(\mathrm{M}_{\mathrm{}})=6.085\times 10^{19}\overline{d_i^2}(_{i=1}^NN_{\mathrm{H}_i})`$, where $`_{i=1}^NN_{\mathrm{H}_i}`$ is the sum of the column densities in $`\mathrm{atoms}\mathrm{cm}^2`$ of all high-velocity profile components detected with a telescope beam of $`0\mathrm{°}.5`$, and $`\overline{d_i^2}`$ is the theoretical mean squared distance of the HVCs from the Sun, expressed in $`\mathrm{kpc}^2`$. The results are $`_{i=1}^NN_{\mathrm{H}_i}=1.87\times 10^{23}\mathrm{atoms}\mathrm{cm}^2`$ and $`\sqrt{\overline{d_i^2}}=105\mathrm{kpc}`$, hence the total HI mass of the HVCs is $`1.2\times 10^9\mathrm{M}_{\mathrm{}}`$. Assuming that the HVCs have as much ionized as neutral hydrogen (Sembach et al. sembach1 (2002)), and adopting a factor 1.3 to include $`H_\mathrm{e}`$, we obtain $`M_\mathrm{t}=3.1\times 10^9\mathrm{M}_{\mathrm{}}`$. Thus the Magellanic Clouds lost about 25 per cent of their original mass in the form of HVCs. This is consistent with the idea that the Clouds constituted a binary system before the last collision, whereas the important mass loss induced by this collision transformed the Clouds into an unbound system.
To study the history of the Clouds as a binary system we should calculate the relative motion of the SMC and LMC before the collision. This can be realized by means of the equation of motion resulting from the difference of Eq. (5) and Eq. (6), the boundary conditions at the moment of the collision, and a LMC mass greater than the present one (see Sect. 3.1). However, for this purpose we merely quote the pioneering study on this topic (Murai & Fujimoto murai (1980)). According to it the LMC and the SMC formed a stable binary system during $`10^{10}\mathrm{yr}`$ (see also Yoshizawa & Noguchi yoshizawa (2003)).
Another interesting result of the model is the estimate of the total energy associated with the production of the HVCs. The total initial kinetic energy of all HVCs ($`\frac{1}{2}M_\mathrm{t}\overline{v_\mathrm{p}^2})`$ is about $`1.8\times 10^{57}\text{ergs}`$, which implies processes that generated at least a total of $`10^{58}\text{ergs}`$. This huge amount of energy should have been supplied by supernovae and winds from massive stars in a starburst originated by the collision of the Clouds. The collective action of supernovae and stellar winds of a typical starburst can drive a “superwind” that may eject $`10^710^9\mathrm{M}_{\mathrm{}}`$ from a galaxy with a mechanical energy of $`10^{57}10^{59}\text{ergs}`$ during its lifetime estimated at $`10^7\mathrm{yr}`$ (Heckman et al. heckman (1989)). There are various stellar generations in the Magellanic Clouds (Westerlund westerlund (1990)). There is nothing against the idea that the formation of the younger generations, with ages lower than 600 Myr, was initiated by the powerful starburst associated with the collision of the Clouds about 570 Myr ago.
## 5 Interaction of the high velocity clouds with the Galactic disk and generation of the Galactic warp
Most of the HVCs have not been affected by the friction force of the interstellar medium, since their orbits have not intersected with the gas layer of the Galactic disk. Here, we analyze the important group of HVCs that was strongly decelerated by the friction force of the Galactic layer and that had a strong influence on the Galactic disk.
The HVCs that penetrated the dense regions of the gaseous disk were trapped in the Galactic plane, and are participating in the Galactic rotation now (see Figs. 7-a and 7-b).
In Fig. 8, we show the positions of entry of the test particles in the gaseous layer of the Galactic disk and their present positions, resulting from the braking process. According to our model, during the last 400 Myr the Galactic disk has been profusely hit by HVCs (see Fig. 8 and Fig. 9-a).
The detailed description of the interaction of HVCs with the Galactic disk is beyond the scope of the present study. Tenorio-Tagle (tenorioa (1980), tenoriob (1981)) analyzed the physics of the collisions of HVCs with the Galactic disk. It is clear that the impact of HVCs upon the Galactic disk should be an important source for producing shells and supershells. This mechanism of creation of large-scale structures in the interstellar medium can account for many of the observational details of the large HI shells lying beyond the solar circle (Mirabel mirabelb (1982); Heiles heiles (1984)). In connection with this, it is interesting to compare the theoretical distribution of the impacts of test particles on the Galactic disk (Fig. 8), with the location of observed shells and supershells of HI (McClure-Griffiths et al. mcclure (2002)). Evidence for interactions of HVCs with the gas in the Galactic plane was found by several authors (see Tripp et al. tripp (2003) and the reviews of Tenorio-Tagle & Bodenheimer tenorio-bodenheimer (1988) and Pöppel poeppel (1997)).
In the next paragraphs of this section we will show that the flow of HVCs onto the Galactic disk over a period of 400 Myr has caused a major disturbance in the Galaxy. As a consequence of the fall of clouds, a burst of star formation involving the whole Galactic disk should have started 400 Myr ago and is still going on (Barry barry (1988); Noh & Scalo noh (1990)). There are observational details of the vertical structure of the Galactic disk, such as peculiar departures from the midplane or huge depressions delineated by young star-gas supercomplexes (Alfaro et al. alfaro (1991)), that may be explained in terms of HVC-disk collisions that induced star formation (Cabrera-Ca$`\stackrel{~}{\mathrm{n}}`$o et al. cabrera (1995)). Further predictions provided by our model are the times, velocities and angles of entry of the HVCs in the Galactic layer (see Figs. 9-a and 9-b). The angle of entry (or incidence) of an HVC is defined as the angle of the direction of motion of an HVC with respect to the vertical ($`Z`$-axis). The angle of entry is an important parameter for the hydrodynamic simulations of the collision process (Comerón & Torra comeron (1991); Franco et al. franco (1988)). The times, velocities and angles of entry correspond to the point at which the friction coefficient reaches the threshold value $`\epsilon \frac{nS}{m}=10^5\mathrm{pc}^1`$ (see Sect. 3.1). Figs 9-a and 9-b show interesting differences between the times and angles of entry on the right side of the Galactic disk (positive $`X`$-axis) and those on the left side. The HVCs arrived first at the right side of the Galactic disk, while there was a delay of $`150\mathrm{Myr}`$ for those arriving at the left side of the disk. On the right side of the disk the HVCs impinged on the disk from below with angles of entry $`80\mathrm{°}`$, i.e. almost a grazing incidence. In contrast, on the left side of disk the angles of entry are lower; and although most of the HVCs entered from below too, there were some HVCs that entered from above.
In the following we examine whether there is a causal relationship between the HVC-disk interaction and the generation of the Galactic warp. Galactic warps are very complex phenomena (see López-Corredoira et al. lopez (2002), García-Ruiz et al. garcia (2002) and references therein). We will analyze only a few aspects of the problem in the light of our model. During the last 450 Myr, the flow of HVCs has been exerting a pressure on the gaseous disk which could have caused the distortion of the disk of the Galaxy. Since about 7 per cent of the test particles of the simulation impacts within the galactocentric radius $`R=`$ 25 kpc, the mass in the form of HVCs accumulated by the disk is $`0.07\times M_\mathrm{t}\mathrm{2.2\hspace{0.17em}10}^8\mathrm{M}_{\mathrm{}}`$. Therefore, the rate of mass accretion per surface unit is $`\gamma =600\mathrm{M}_{\mathrm{}}\mathrm{kpc}^2\mathrm{Myr}^1`$, taking into account that the time during which the flow acted on the disk at a fixed point in the $`(X,Y,Z)`$ frame was nearly 200 Myr (see Fig. 9-a). Let us consider a cylinder of unit cross-section with generators parallel to the $`Z`$-axis and height $`h(R)`$ as a volume element of the gaseous disk of the Galaxy at $`R`$. The accretion process acting on the volume element generates a force in $`Z`$ per unit of mass that can be written as $`F_z=\frac{\gamma }{\sigma (R)}\overline{v}cos(\overline{i})`$, where $`\overline{v}`$, $`\overline{i}`$ are the mean velocity and angle of entry, and $`\sigma (R)`$ is the surface density of the gaseous disk at $`R`$ (see Sect. 3.1).
To describe the gas motion of the volume element under the effects of the additional force associated with the mass accretion, we use cylindrical coordinates, $`(R,\theta ,Z)`$, representing the galactocentric radius, azimuth and vertical distance from the plane $`b=0\mathrm{°}`$ of the mass center of the chosen volume element. The gas motion in the $`Z`$-direction is governed by the force $`F_z`$ and the restoring force of the stellar disk $`F=\lambda ^2Z`$. It can be expressed by
$$Z(t)=Z_0(t_\mathrm{e})+\frac{sin(\lambda (tt_\mathrm{e}))}{\lambda }_0^{tt_\mathrm{e}}F_zcos(\lambda \tau )𝑑\tau \frac{cos(\lambda (tt_\mathrm{e}))}{\lambda }_0^{tt_\mathrm{e}}F_zsin(\lambda \tau )𝑑\tau $$
(11)
, (see Eq. 5.713 of Chandrasekhar chandrasekhar (1943)). As initial conditions of the disk in $`Z`$ we put $`Z_0=\dot{Z}_0=0`$ at the time origin $`t_\mathrm{e}=570\text{Myr}`$. Ignoring the effects of the $`(X,Y)`$ component of the force associated with the mass accretion, and assuming circular rotation, we have $`\theta (t)=\theta _0(t_\mathrm{e})+\frac{V_c}{R}(tt_\mathrm{e})`$.
The flow of HVCs moves parallel to the $`X`$-axis. Therefore, we can think of the interaction of the HVCs with the disk as an accretion front perpendicular to the $`X`$-axis, propagating in the $`X`$ direction with a constant velocity $`v_\mathrm{a}`$. Then the position of the accretion front as a function of time is $`X_\mathrm{a}(t)=v_\mathrm{a}t+X_0(t_\mathrm{e})`$ (see Fig. 9-a). The volume element enters the front when $`X=Rsin(\theta (t))=X(t)_\mathrm{a}`$, a condition that allows us to determine the entry time $`t_1`$. The time at which the volume element leaves the rear zone of the accretion front, $`t_2`$, is determined by solving $`X=X(t)_\mathrm{a}+\mathrm{\Delta }X_\mathrm{a}`$ for t, where $`\mathrm{\Delta }X_\mathrm{a}`$ is the thickness of the accretion front. The times $`t_1`$ and $`t_2`$ and the duration of the process of mass accretion, $`t_\mathrm{a}=t_2t_1`$, depend on $`\theta `$ and $`R`$ of the volume element (see Fig. 10-a). At large galactocentric radii the time intervals of mass accretion $`t_\mathrm{a}`$ in regions of the third and fourth Galactic quadrants are longer than those in the first and second quadrants, i.e the outer regions in the third and fourth Galactic quadrants are the most affected by the interaction with the HVCs. For each volume element, the force $`F_z`$ acts between the corresponding times $`t_1`$ and $`t_2`$, vanishing outside this time interval. Hence, Eq. (11) becomes $`Z(\theta ,R,t)=\frac{F_z}{\lambda ^2}(cos\lambda (t_2+t_\mathrm{e}t)cos\lambda (t_1+t_\mathrm{e}t))`$. This is a general expression for the $`Z`$ deformation of the Galactic layer due to the flow of HVCs. At t=0 this equation can reproduce roughly the present configuration of the Galactic warp, if the period of vertical oscillation $`\frac{1}{2\pi \lambda }450\mathrm{Myr}`$ (see Fig. 10-b). This value seems reasonable, because at large galactocentric radii the restoring force should be small. Note that the model can explain the overall pattern of positive vertical displacements in the first and second Galactic quadrants, as well as the less pronounced displacements toward negative $`Z`$ in the third and fourth quadrants (cf. our Fig. 10-b with Figs. 57b and 60b of Burton burton (1991)).
The first explanation proposed for the Galactic warp was that the Galaxy is passing through a continuous intergalactic medium (Kahn & Woltjer kahn (1959)). In contrast, we propose an interaction of the Galactic disk with a transient circumgalactic flow of HVCs produced by a single event of mass transfer from the Magellanic Clouds. Our model shows that the flow of HVCs is able to produce the formation of the Galactic warp. Future extensions of the model should contemplate consequences of the process of mass accretion such as large deviations from the circular rotation of the gaseous disk and the sudden increment of the rate of star formation and of the input of energy in the interstellar medium. The HVCs have injected energy at a rate of $`10^{50}\mathrm{ergs}\mathrm{kpc}^2\mathrm{Myr}^1`$ into the Galactic disk during the last 400 Myr. A question we should address, among others, should be the role played by this important source of energy in the production and maintenance of the diverse phases of the interstellar medium (Kulkarni & Heiles kulkarni (1988); Elmegreen elmegreen (1991); Heiles & Troland heiles-troland (2003)). Another interesting question is whether there exists any genetic or dynamic link between the Galactic warp and a ring of stars in the plane of the Galaxy, at a Galactocentric radius of 18-20 kpc, which may completely encircle the Galaxy (Ibata et al. ibata (2003); Yanny et al. yanny (2003)).
## 6 Conclusions
We have shown that the majority of the HVCs have a common origin linked to the dynamic evolution of the Magellanic Clouds. During a long time these two small satellite galaxies have likely behaved as a binary system under the effects of the tidal perturbations of the Galaxy. The two Clouds have probably approached each other several times in the past, until the Clouds collided catastrophically about 570 Myr ago. We assumed that the collision triggered a burst of stellar formation in the Clouds. The superwind driven by the starburst blew the HVCs as magnetized clouds into the massive halo of the Galaxy. As a consequence, the Milky Way is likely surrounded by a mist of highly ionized gas at high velocity, the HVCs being the denser concentrations.
The catastrophic loss of matter was of a significant amount. Probably it permitted the SMC to reach escape velocity, disrupting the binary system. According to the results of our model the amounts of matter and mechanical energy liberated in the form of HVCs were of the order of $`3.1\times 10^9\mathrm{M}_{\mathrm{}}`$ and $`1.8\times 10^{57}\mathrm{ergs}`$, respectively. In the numerical simulation of the process, we assumed that the HVCs were ejected from the Clouds 570 Myr ago in all directions with peculiar velocities larger than the escape velocity of $`160\mathrm{km}\mathrm{s}^1`$, and following a Gaussian distribution law. Because of these initial expansion velocities, the cloud of HVCs, or metacloud, expanded from the center of collision of the Magellanic Clouds, reaching its present size of $`300\text{kpc}`$, and filling a great part of the halo volume of the Galaxy. As the initial mean velocity of the HVCs was that of the Clouds at the time of the ejections of HVCs, the centroid of the group of HVCs (metacloud) describes an orbit close to the orbit of the Magellanic Clouds (i.e. approximately in the direction $`\mathrm{}=90\mathrm{°}270\mathrm{°}`$). To an observer at the position of the Sun, this gives the impression of a flow of HVCs coming from $`\mathrm{}=90\mathrm{°}`$. As a consequence of the initial velocity distribution of the HVCs, a larger number of HVCs tended to concentrate towards the Magellanic orbit behind the present position of the Clouds, whereas other HVCs moved ahead of them, forming the Magellanic Stream.
The passage of the HVC flow through the Galactic disk had transcendent consequences for the evolution of the gaseous layer of the Galaxy. The outer gas layer of the Galaxy was considerably displaced in the $`Z`$-direction from its equilibrium position. The Galactic disk accumulated mass coming in the form of HVCs at an average rate of $`600\mathrm{M}_{\mathrm{}}\mathrm{kpc}^2\mathrm{Myr}^1`$ over a period of 200 Myr. Due to the time dependence of the mass accretion upon the azimuth of the region, the effects on the first and second Galactic quadrants at large galactocentric radii were different from those on the third and fourth ones. These are perhaps some of the clues to understand the genesis of the Galactic warp. The collisions of the HVCs with the Galactic disk might have exerted notable influences on the large-scale morphology and energy of the interstellar medium.
###### Acknowledgements.
I am particularly grateful to Dr. Virpi S. Niemela and Dr. Wolfgang G. L. Pöppel for their constant encouragement and help. Dr. Wolfgang G. L. Pöppel read the manuscript with a constructive eye. I am indebted to Dr. Bart P. Wakker for providing his catalog of HVCs via e-mail. The helpful comments of an anonymous referee led to substantial improvements in the paper. Part of this work was supported by the *Consejo Nacional de Investigaciones Científicas y Técnicas (CONICET)* project number PIP-0608/98.
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# Expressions for the Exchange Correlation Potential and Exchange–Correlation Functional of Kohn–Sham Density Functional Theory
## I Introduction
The Kohn-Sham version of density functional theory plays a major role in both quantum chemistry and condensed matter physics Dreizler and E. K. U. Gross (1990); Parr and Yang (1989); Springborg (1997); Ellis (1995); E. K. U. Gross and Dreizler (1994); Seminario and Politzer (1995); Handy (1997). Unfortunately, the exchange-correlation functional $`E_{\text{xc}}`$ is an unknown, implicit functional, and there is no systematic method to improve approximations. Recently we have derived a generalization of the Kohn–Sham approach in which the correlation energy $`E_{\text{co}}`$ is assumed to be an explicit functional of $`v`$ and $`\rho _1`$, where $`v`$ is the external potential from the interacting target-state, and $`\rho _1`$ is the spinless one-particle density matrix from the noninteracting states Finley (2005a, b). In this approach, errors from Coulomb self-interactions do not occur, nor the need to introduce functionals defined by a constraint search. Furthermore, the exchange energy $`E_\text{x}`$ is treated as in Hartree–Fock theory, as an explicit functional of $`\rho _1`$. Below, we use this approach to derive the Kohn-Sham exchange-correlation potential $`v_{\text{xc}}`$ and exchange-correlation functional $`E_{\text{xc}}`$ as explicit functionals of $`v_s`$ and $`\phi _1`$, where $`v_s`$ is the local, one-body potential from the Kohn–Sham equations, and $`\phi _1`$ is the one-particle density matrix from the Kohn–Sham noninteracting state, say $`|\phi _1`$. In other words, $`|\phi _1`$ is the ground state eigenfunction of the noninteracting Schrödinger equation with the one-body potential $`v_s`$.
## II State-specific Kohn–Sham density functional theory
In state-specific Kohn–Sham density functional theory Finley (2005a, b), we use an energy functional $`E_v[\rho _1]`$ that is assumed to be an explicit functional of the external potential $`v`$ and the spinless one-particle density matrix $`\rho _1`$, where $`\rho _1`$ comes from a closed-shell determinantal state, say $`|\rho _1`$; this energy functional is given by
$`E_v[\rho _1]={\displaystyle \frac{\stackrel{~}{\mathrm{\Psi }}_{v\rho _1}|\widehat{H}_v|\stackrel{~}{\mathrm{\Psi }}_{v\rho _1}}{\stackrel{~}{\mathrm{\Psi }}_{v\rho _1}|\stackrel{~}{\mathrm{\Psi }}_{v\rho _1}}},`$ (1)
where the trial wave function $`\stackrel{~}{\mathrm{\Psi }}_{v\rho _1}`$ generates the exact, or target, wave function, say $`\mathrm{\Psi }_n`$, under the following conditions:
$`|\stackrel{~}{\mathrm{\Psi }}_{v\varrho _1}=|\mathrm{\Psi }_n;\varrho _1n,nN,v,`$ (2)
where $`n`$ is the electron density of $`\mathrm{\Psi }_n`$, and $`\mathrm{\Psi }_n`$ is the ground state singlet eigenfunction of $`\widehat{H}_v`$. The right side notation of Eq. (2) indicates that $`\varrho _1`$ is a spin-less one particle density matrix that delivers the density $`n`$; according to the Hohenberg-Kohn theorem Hohenberg and Kohn (1964); Parr and Yang (1989); Dreizler and E. K. U. Gross (1990), $`n`$ also determines the external potential $`v`$, and $`n`$ determines the number of electrons $`N`$. Using Eqs. (1) and (2), we have
$`E_v[\varrho _1]=_n,\varrho _1n,nN,v.`$ (3)
Here, $`_n`$ is the exact electronic energy of the target state:
$`\widehat{H}_v|\mathrm{\Psi }_n=_n|\mathrm{\Psi }_n;nN,v.`$ (4)
where the Hamiltonian operator is given by
$`\widehat{H}_v=\widehat{T}+\widehat{V}_{\text{ee}}+\widehat{V}_v,`$ (5)
and we have
$`\widehat{T}`$ $`=`$ $`{\displaystyle \underset{i}{\overset{N}{}}}(\frac{1}{2}_i^2),`$ (6)
$`\widehat{V}_{\text{ee}}`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{ij}{\overset{N}{}}}r_{ij}^1,`$ (7)
$`\widehat{V}_v`$ $`=`$ $`{\displaystyle \underset{i}{\overset{N}{}}}v(i).`$ (8)
The electronic energy functional can also be written as
$`E_v[\rho _1]={\displaystyle 𝑑𝐫_1\left[\frac{1}{2}_1^2\rho _1(𝐫_1,𝐫_2)\right]_{𝐫_2=𝐫_1}}+{\displaystyle 𝑑𝐫v(𝐫)\rho _s(𝐫)}`$
$`+E_J[\rho _s]+E_\text{x}[\rho _1]+E_{\text{co}}[\rho _1,v]+{\displaystyle 𝑑𝐫v(𝐫)\stackrel{~}{\rho }_c(𝐫)},`$ (9)
where the Coulomb and exchange energies are given by the following:
$`E_J[\rho _s]`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle r_{12}^1𝑑𝐫_1𝑑𝐫_2\rho (𝐫_1)\rho (𝐫_2)},`$ (10)
$`E_\text{x}[\rho _1]`$ $`=`$ $`{\displaystyle \frac{1}{4}}{\displaystyle r_{12}^1𝑑𝐫_1𝑑𝐫_2\rho _1(𝐫_1,𝐫_2)\rho _1(𝐫_2,𝐫_1)},`$ (11)
and the correlation-energy functional is
$`E_{\text{co}}[\rho _1,v]={\displaystyle \frac{\stackrel{~}{\mathrm{\Psi }}_{v\rho _1}|\widehat{T}|\stackrel{~}{\mathrm{\Psi }}_{v\rho _1}}{\stackrel{~}{\mathrm{\Psi }}_{v\rho _1}|\stackrel{~}{\mathrm{\Psi }}_{v\rho _1}}}\rho _1|\widehat{T}|\rho _1+{\displaystyle \frac{\stackrel{~}{\mathrm{\Psi }}_{v\rho _1}|\widehat{V}_{\text{ee}}|\stackrel{~}{\mathrm{\Psi }}_{v\rho _1}}{\stackrel{~}{\mathrm{\Psi }}_{v\rho _1}|\stackrel{~}{\mathrm{\Psi }}_{v\rho _1}}}\rho _1|\widehat{V}_{\text{ee}}|\rho _1.`$ (12)
Furthermore, $`\stackrel{~}{\rho }_c`$ is the correlation density of the trial wave function, i.e, we have
$`\stackrel{~}{\rho }_c(𝐫)={\displaystyle \frac{\stackrel{~}{\mathrm{\Psi }}_{v\rho _1}|\widehat{\mathrm{\Gamma }}(𝐫)|\stackrel{~}{\mathrm{\Psi }}_{v\rho _1}}{\stackrel{~}{\mathrm{\Psi }}_{v\rho _1}|\stackrel{~}{\mathrm{\Psi }}_{v\rho _1}}}\rho _s(𝐫)=\stackrel{~}{n}\rho _s(𝐫),\stackrel{~}{\mathrm{\Psi }}_{v\rho _1}\stackrel{~}{n},\rho _1\rho _s,`$ (13)
and $`\widehat{\mathrm{\Gamma }}`$ is the density operator; $`\stackrel{~}{n}`$ is the density of $`\stackrel{~}{\mathrm{\Psi }}_{v\rho _1}`$; $`\rho _s`$ is the density of $`\rho _1`$. Since $`\stackrel{~}{\rho }_c(𝐫)`$ is a functional of $`v`$ and $`\rho _1`$, we can also write $`\stackrel{~}{\rho }_c[\rho _1,v](𝐫)`$.
A determinantal state with the density matrix $`\varrho _1`$ satisfies the following noninteracting Schrödinger equation:
$`{\displaystyle \underset{i=1}{\overset{N}{}}}\widehat{}_{\varrho _1}(𝐫_i)|\varrho _1=2\left({\displaystyle \underset{w}{}}\epsilon _w\right)|\varrho _1,`$ (14)
where the generalized, or exact, Fock operator $`\widehat{}_{\varrho _1}`$ is given by
$`\widehat{}_{\varrho _1}=\frac{1}{2}^2+v+v_J^n+\widehat{v}_\text{x}^{\varrho _1}+\widehat{v}_{\text{co}}^{\varrho _1}+\widehat{v}_{\text{ec}}^{\varrho _1}.`$ (15)
Here, the Coulomb operator is defined by
$`v_J^\rho (𝐫_1)\chi (𝐫_1)={\displaystyle 𝑑𝐫_2r_{12}^1\rho (𝐫_2)\chi (𝐫_1)},`$ (16)
and the exchange operator $`\widehat{v}_\text{x}^{\rho _1}`$, correlation operator $`\widehat{v}_{\text{co}}^{\rho _1}`$ and external-correlation operator $`\widehat{v}_{\text{ec}}^{\rho _1}`$ are defined by their kernels:
$`v_\text{x}^{\rho _1}(𝐫_1,𝐫_2)`$ $`=`$ $`{\displaystyle \frac{\delta E_\text{x}[\rho _1,v]}{\delta \rho _1(𝐫_2,𝐫_1)}}={\displaystyle \frac{1}{2}}r_{12}^1\rho _1(𝐫_1,𝐫_2),`$ (17)
$`v_{\text{co}}^{\rho _1}(𝐫_1,𝐫_2)`$ $`=`$ $`{\displaystyle \frac{\delta E_{\text{co}}[\rho _1,v]}{\delta \rho _1(𝐫_2,𝐫_1)}},`$ (18)
$`v_{\text{ec}}^{\rho _1}(𝐫_1,𝐫_2)`$ $`=`$ $`{\displaystyle \frac{\delta \left(𝑑𝐫_3v(𝐫_3)\stackrel{~}{\rho }_c(𝐫_3)\right)}{\delta \rho _1(𝐫_2,𝐫_1)}}.`$ (19)
Our energy functionals $`E_v`$ are implicit functionals of the noninteracting density $`\rho _s`$. Hence, any one-particle density-matrix that yields the interacting density minimizes our energy functional, i.e., we have
$`_n=E_v[\varrho _1]=E_v[\varrho _1^{}]=E_v[\varrho _1^{\prime \prime }]\mathrm{},`$ (20)
where
$`n(𝐫)=\varrho _1(𝐫,𝐫)=\varrho _1^{}(𝐫,𝐫)=\varrho _1^{\prime \prime }(𝐫,𝐫)\mathrm{}.`$ (21)
Assuming $`n`$ is a noninteracting $`v`$-representable density, there exist a noninteracting state, say $`|\phi _1`$, that has $`n`$ as its density:
$`n(𝐫)=\phi _1(𝐫,𝐫),`$ (22)
and this determinant—assuming it is a closed-shell determinant—is the ground-state solution of the following noninteracting Schrödinger equation:
$`{\displaystyle \underset{i=1}{\overset{N}{}}}\widehat{f}(𝐫_i)|\phi _1=2\left({\displaystyle \underset{w}{}}ϵ_w\right)|\phi _1,`$ (23)
where
$`\widehat{f}=\frac{1}{2}^2+v_s,`$ (24)
and $`v_s`$ is a local potential. Therefore, the canonical occupied orbitals from $`|\phi _1`$ satisfy the following one-particle Schrödinger equation:
$`\widehat{f}\varphi _w=\left(\frac{1}{2}^2+v+v_J^n+v_{\text{xc}}\right)\varphi _w=ϵ_w\varphi _w,\varphi _w\phi _1,`$ (25)
where, with no loss of generality, we have required $`v_s`$ to be defined by
$`v_s=v+v_J^n+v_{\text{xc}}.`$ (26)
Using the approach by Sala and Görling Sala and Görling (2001), but permitting the orbitals to be complex, it is readily demonstrated that $`v_{\text{xc}}`$ is given by
$`v_{\text{xc}}(𝐫)={\displaystyle \frac{1}{2n(𝐫)}}{\displaystyle }d𝐫_1[2w(𝐫_1,𝐫)\phi _1(𝐫,𝐫_1)\phi _1(𝐫,𝐫_1){\displaystyle }d𝐫_2\phi _1(𝐫_2,𝐫)w(𝐫_1,𝐫_2)`$ (27)
$`+\phi _1(𝐫_1,𝐫)\phi _1(𝐫,𝐫_1)v_{\text{xc}}(𝐫_1)],`$
where $`w`$ is the kernel of the nonlocal potential $`\widehat{w}_{\rho _1}`$, given by
$`\widehat{w}_{\rho _1}=\widehat{v}_\text{x}^{\rho _1}+\widehat{v}_{\text{co}}^{\rho _1}+\widehat{v}_{\text{ec}}^{\rho _1},`$ (28)
and these operators appear in the exact Fock operator $`\widehat{}_{\varrho _1}`$, given by Eq. (15). By substituting $`v_{\text{xc}}`$ repeatedly on the right side, we can obtain an expansion for $`v_{\text{xc}}`$:
$`v_{\text{xc}}(𝐫)={\displaystyle \frac{1}{2n(𝐫)}}[2w(𝐫_1,𝐫)\phi _1(𝐫,𝐫_1)\phi _1(𝐫,𝐫_1)\phi _1(𝐫_2,𝐫)w(𝐫_1,𝐫_2)`$
$`+\phi _1(𝐫_1,𝐫)\phi _1(𝐫,𝐫_1){\displaystyle \frac{1}{n(𝐫_1)}}\{w(𝐫_2,𝐫_1)\phi _1(𝐫_1,𝐫_2){\displaystyle \frac{1}{2}}\phi _1(𝐫_1,𝐫_2)\phi _1(𝐫_3,𝐫_1)w(𝐫_2,𝐫_3)\}`$
$`+\phi _1(𝐫_1,𝐫)\phi _1(𝐫,𝐫_1){\displaystyle \frac{1}{2n(𝐫_1)}}\phi _1(𝐫_2,𝐫_1)\phi _1(𝐫_1,𝐫_2){\displaystyle \frac{1}{n(𝐫_2)}}w(𝐫_3,𝐫_2)\phi _1(𝐫_2,𝐫_3)+\mathrm{}],`$ (29)
where in this relation it is understood that there are integrations over the dummy variables $`𝐫_1`$, $`𝐫_2`$ and $`𝐫_3`$. The leading term of Eq. (29) is the Slater potential Slater (1951); Harbola and Sahni (1993); Hirata et al. (2001); this term also appears within the Krieger–Li–Iafrate (KLI) approximation of the optimized potential method Fiolhais et al. (2003); Krieger et al. (1992); Li et al. (1993); Hirata et al. (2001).
The orbitals $`\varphi _w`$ satisfying Eq. (25) are the Kohn–Sham orbitals Kohn and Sham (1965); $`|\phi _1`$ is the Kohn–Sham noninteracting state. However, $`\widehat{f}`$ differs from the Kohn–Sham operator, since, in addition to depending explicitly on $`\phi _1`$, instead of $`n`$, $`\widehat{f}`$ depends explicitly on the external potential $`v`$ from the interacting Hamiltonian $`\widehat{H}_v`$. Furthermore, the external-correlation operator $`\widehat{v}_{\text{ec}}^{\rho _1}`$ does not appear in Kohn–Sham formalism. And, unlike the original Kohn–Sham approach Kohn and Sham (1965), the $`N`$-representability problem does not arise, nor the need to introduce a constraint-search definition Levy (1978, 1979, 1982); Levy and Perdew (1985) to avoid this problem.
## III The Kohn–Sham exchange–correlation potential
According to Eqs. (28), (18), and (19), $`\widehat{w}_{\rho _1}`$ is a functional of $`\rho _1`$ and $`v`$, indicating that we can, symbolically speaking, represent Eq. (29) by
$`v_{\text{xc}}(𝐫)=v_{\text{xc}}[\phi _1,v](𝐫),\phi _1nv.`$ (30)
In other words, $`v_{\text{xc}}`$ is also a functional of $`v`$ and $`\phi _1`$, where $`\phi _1`$ determines $`n`$, and $`n`$ determines $`v`$.
Note that $`\phi _1`$ and $`v`$ from Eq. (30) are not independent. Since, for a given $`v`$, the one-particle density matrix $`\phi _1`$ which determines $`v_{\text{xc}}`$ from Eq. (30), and then gives $`v_s`$ from (26), must also be the one-particle density matrix $`\phi _1`$ from the noninteracting state $`|\phi _1`$ that satisfies Eq. (23). However, for a given $`\phi _1`$ or $`v_s`$, we can also think of $`v`$ as a dependent variable to be determined. In that case, we choose $`v_s`$, construct $`\widehat{f}`$ using Eq. (24), and obtain the one-particle density matrix, say $`\phi _1`$, that satisfies Eq. (23). The external potential $`v`$ is then a simultaneous solution of Eqs. (26) and (30).
Substituting Eq. (30) into Eq. (26), gives
$`v=v_sv_J^nv_{\text{xc}}[\phi _1,v],N,v_s\phi _1nN,v,v_s,`$ (31)
where the notation on the right side indicates that $`N`$ and $`v_s`$ determine $`\phi _1`$, as indicated by Eqs. (23) and (24); also, $`\phi _1`$ determines $`n`$; $`n`$ determines $`N`$, $`v`$, and $`v_s`$.
By substituting $`v`$ on the right side of Eq. (31) repeatedly, as in Eq. (29), we can remove it, giving, symbolically speaking,
$`v=v[\phi _1,v_s],`$ (32)
Hence, this relation gives the external potential $`v`$ from the interacting system $`\mathrm{\Psi }_n`$, for $`nv`$, as a functional of the local potential $`v_s`$ and the one-particle density matrix $`\phi _1`$ from the noninteracting state $`|\phi _1`$, where $`|\phi _1`$ is an eigenfunction of the noninteracting Hamiltonian given by Eq. (23), and this noninteracting Hamiltonian is defined by the local potential $`v_s`$, as indicated by Eq. (24). Furthermore, the noninteracting state $`|\phi _1`$ shares the same density with $`\mathrm{\Psi }_n`$; $`|\phi _1`$ is the Kohn-Sham determinantal state when $`n`$ is noninteracting $`v`$-representable.
Since there is a one-to-one correspondence between local potentials $`v_s`$ and noninteracting ground states $`|\phi _1`$ Dreizler and E. K. U. Gross (1990), Eq. (32) can be written as
$`v=v[\phi _1,v_s],v_s\phi _1,`$ (33)
where the right side indicates the one-to-one correspondence.
Substituting Eq. (33) into Eq. (30) gives
$`v_{\text{xc}}=v_{\text{xc}}[\phi _1,v_s],v_s\phi _1.`$ (34)
Denoting $`v_{\text{xc}}^{\text{KS}}`$ the Kohn–Sham exchange-correlation potential Kohn and Sham (1965); Dreizler and E. K. U. Gross (1990); Parr and Yang (1989); Springborg (1997); Ellis (1995); E. K. U. Gross and Dreizler (1994); Seminario and Politzer (1995); Handy (1997), at least for non-interacting $`v`$-representable densities, we have
$`v_{\text{xc}}^{\text{KS}}[n]=v_{\text{xc}}[\phi _1,v_s],n\phi _1,v_s,v_s\phi _1.`$ (35)
In order to obtain the density functional on the left side of Eq. (35), we need the functionals $`\phi _1[n]`$ and $`v_s[n]`$. According the to Hohenberg-Kohn theorem Hohenberg and Kohn (1964); Parr and Yang (1989); Dreizler and E. K. U. Gross (1990), for noninteracting $`v`$-representable densities, these functionals exist. For other densities, however, these expressions must be generalized by, for example, using some modified constraint search definition, or, perhaps, by an approach that permits $`v_s`$ to be nonlocal.
## IV The Kohn–Sham correlation–energy functional and exchange correlation functional
Using the notation form Eq. (4), the universal functional from the Hohenberg-Kohn theorem Hohenberg and Kohn (1964); Parr and Yang (1989); Dreizler and E. K. U. Gross (1990), can be written as
$`F[n]=\mathrm{\Psi }_n|\left(\widehat{T}+\widehat{V}_{\text{ee}}\right)|\mathrm{\Psi }_n;nN,v.`$ (36)
where $`\mathrm{\Psi }_n`$ is a singlet, ground state wave function. Previously we have shown that the correlation energy from many body perturbation theory Lindgren and Morrison (1986); Harris et al. (1992); Raimes (1972) can be written as an explicit functional of $`v`$ and $`\rho _1`$ Finley (2003). In a similar manner, but using less restrictive energy denominators, the universal functionals $`F`$ can be shown to be an explicit functional of $`v`$ and $`\rho _1`$ Finley (2005c), where this functional does not implicitly depend on $`\rho _1`$, i.e., any $`|\rho _1`$ that has considerable overlap with $`|\mathrm{\Psi }_n`$ can be used as a reference state. So, we can write
$`F[n]=F[v,\rho _1],nv,n,\rho _1N.`$ (37)
For noninteracting $`v`$-representable densities $`n`$, we can use Eq. (32), giving
$`F[n]=F[\phi _1,v_s,\rho _1],n\phi _1,v_s,v_s\phi _1,\rho _1N.`$ (38)
Setting $`\rho _1=\phi _1`$, we get
$`F[n]=F[\phi _1,v_s],n\phi _1,v_s,v_s\phi _1.`$ (39)
Hence, assuming the existence of the explicit functional, given by Eq. (37), Eq. (39) gives this functional as an explicit functional of $`v_s`$ and $`\phi _1`$, where $`v_s`$ is the local, one-body potential from the Kohn–Sham equations, and $`\phi _1`$ is the spinless one-particle density matrix from the Kohn–Sham noninteracting state $`|\phi _1`$.
Using Eq. (39), the exchange-correlation functional from the Kohn-Sham formalism Kohn and Sham (1965) is
$`E_{\text{xc}}^{\text{KS}}[n]=F[\phi _1,v_s]E_J[n]\phi _1|\widehat{T}|\phi _1,n\phi _1,v_s.`$ (40)
It is well know that the Kohn–Sham exchange functional is given by Parr and Yang (1989); Levy and Perdew (1985)
$`E_\text{x}^{\text{KS}}[n]`$ $`=`$ $`E_\text{x}[\phi _1],n\phi _1.`$ (41)
Hence, the Kohn–Sham correlation-energy functional is
$`E_{\text{co}}[n]^{\text{KS}}=F[\phi _1,v_s]\phi _1|\widehat{T}|\phi _1\phi _1|\widehat{V}_{\text{ee}}|\phi _1,`$ (42)
where
$`\phi _1|\widehat{V}_{\text{ee}}|\phi _1=E_J[n]+E_\text{x}[\phi _1].`$ (43)
## Appendix A The Kohn–Sham correlation-energy functional from State–Specific Kohn–Sham Density Functional Theory
Consider the energy functional $`E_v[\rho _1,v^{}]`$ that is obtained by permitting the external potential from the trial wave function, say $`v^{}`$, to differ from the one from the Hamiltonian $`\widehat{H}_v`$; generalizing Eq. (1) in this way, we have
$`E_v[\rho _1,v^{}]={\displaystyle \frac{\stackrel{~}{\mathrm{\Psi }}_{v^{}\rho _1}|\widehat{H}_v|\stackrel{~}{\mathrm{\Psi }}_{v^{}\rho _1}}{\stackrel{~}{\mathrm{\Psi }}_{v^{}\rho _1}|\stackrel{~}{\mathrm{\Psi }}_{v^{}\rho _1}}}.`$ (44)
and Eqs. (2), (3), and (II) become
$`|\stackrel{~}{\mathrm{\Psi }}_{v^{}\varrho _1}`$ $`=`$ $`|\mathrm{\Psi }_n;v^{}=v,\varrho _1n,nN,v,`$ (45)
$`E_v[\varrho _1,v]`$ $`=`$ $`_n.`$ (46)
$`E_v[\rho _1,v^{}]=_1[\rho _1,v]+E_{\text{co}}[\rho _1,v^{}]+{\displaystyle 𝑑𝐫v(𝐫)\stackrel{~}{\rho }_c[\rho _1,v^{}](𝐫)},`$ (47)
where we have explicitly mentioned the dependence of $`\stackrel{~}{\rho }_c`$ on $`\rho _1`$ and $`v^{}`$, and the energy through the first order, $`_1`$, is given by
$`_1[\rho _1,v]=\rho _1|H_v|\rho _1={\displaystyle 𝑑𝐫_1\left[\frac{1}{2}_1^2\rho _1(𝐫_1,𝐫_2)\right]_{𝐫_2=𝐫_1}}`$ (48)
$`+{\displaystyle 𝑑𝐫_1v(𝐫_1)\rho (𝐫_1)}+{\displaystyle \frac{1}{2}}{\displaystyle 𝑑𝐫_1𝑑𝐫_2r_{12}^1\rho (𝐫_1)\rho (𝐫_2)}{\displaystyle \frac{1}{4}}{\displaystyle 𝑑𝐫_1𝑑𝐫_2r_{12}^1\rho _1(𝐫_1,𝐫_2)\rho _1(𝐫_2,𝐫_1)}.`$
Using Eq. (32), we have
$`E_v[\rho _1,v^{}[\phi _1^{},v_s^{}]]=_1[\rho _1,v]+E_{\text{co}}[\rho _1,v^{}[\phi _1^{},v_s^{}]]+{\displaystyle 𝑑𝐫v(𝐫)\stackrel{~}{\rho }_c[\rho _1,v^{}[\phi _1^{},v_s^{}]](𝐫)}.`$ (49)
Setting $`\rho _1=\phi _1^{}`$
$`E_v[\phi _1^{},v_s^{}]=_1[\phi _1^{},v]+E_{\text{co}}[\phi _1^{},v_s^{}],`$ (50)
and suppressing the primes, we have
$`E_v[\phi _1,v_s]=_1[\phi _1,v]+E_{\text{co}}[\phi _1,v_s],`$ (51)
where we used
$`\stackrel{~}{\rho }_c[\phi _1^{},v^{}[\phi _1^{},v_s^{}]]=0.`$ (52)
Comparing Eq. (51) with the energy functionals from the Kohn-Sham formalism, for noninteracting $`v`$-representable densities, we have
$`E_v^{\text{KS}}[n]`$ $`=`$ $`E_v[\phi _1,v_s],`$ (53)
$`E_{\text{co}}^{\text{KS}}[n]`$ $`=`$ $`E_{\text{co}}[\phi _1,v_s],`$ (54)
where $`n\phi _1,v_s,v_s\phi _1`$. As in $`v_{\text{xc}}^{\text{KS}}`$, in order to obtain the above density functionals, we need $`\phi _1[n]`$ and $`v_s[n]`$. Recall that it is well know that the Kohn–Sham exchange functional satisfies Eq. (41) Parr and Yang (1989); Levy and Perdew (1985).
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# Fluorescence and phosphorescence from individual C60 molecules excited by local electron tunneling
## Abstract
Using the highly localized current of electrons tunneling through a double barrier Scanning Tunneling Microscope (STM) junction, we excite luminescence from a selected C<sub>60</sub> molecule in the surface layer of fullerene nanocrystals grown on an ultrathin NaCl film on Au(111). In the observed luminescence fluorescence and phosphorescence spectra, pure electronic as well as vibronically induced transitions of an individual C<sub>60</sub> molecule are identified, leading to unambiguous chemical recognition on the single-molecular scale.
preprint: APS/123-QED
Light emission induced by electrons tunneling through the junction formed by the sample and the tip of a Scanning Tunneling Microscope (STM) has been proposed to characterize the optical properties of nanoscale objects at surfaces Gimzewski88 . Contrary to conventional non-local techniques, the local character of this method offers the unique possibility to select and probe individual atoms, molecules or clusters on surfaces.
Photon emission due to the decay of localized surface plasmons, excited by inelastic electron tunneling (IET) has been observed on metal surfaces Berndt91 ; Berndt93B , as well as on supported metallic nanoparticles Nilius00 . Luminescence spectra have been acquired from semiconductor heterostructures Alvarado91 , quantum well states of metallic films Hoffmann01 . Recently, luminescence from supported molecules has been obtained Qiu03 ; Dong04 by successfully decoupling them from the metallic substrate in order to avoid quenching of the radiative transitions Berndt93 ; Dong03 , using either a thin oxide film Qiu03 or several molecular layers Dong04 .
However, unambiguous chemical identification of single complex molecules requires the observation and identification of *several* vibrational and/or electronic-vibrational transitions, which are the spectroscopic fingerprint of the species. Here we present the first observation of energy resolved luminescence from an individually selected C<sub>60</sub> molecule excited by electrons tunneling through a double barrier STM junction. A comparison with the luminescence spectra obtained by non-local laser spectroscopy from dispersed C<sub>60</sub> molecules in rare gas and glass matrices Sassara96B ; Sassara97 ; Sassara96 ; Hung96 ; vandenHeuvel94 ; vandenHeuvel95 , and from solid C<sub>60</sub> Guss94 ; vandenHeuvel95B ; Akimoto02 enables us to demonstrate the molecular origin of the detected light and to identify the observed spectral features with pure electronic transitions and with vibronic transitions induced via Jahn-Teller (JT) and Herzberg-Teller (HT) coupling Negri92 ; Orlandi02 . The present novel observation of both, fluorescence (singlet-to-singlet transitions) *and* phosphorescence (triplet-to-singlet transitions) constitutes a solid basis for the chemical identification of an individual C<sub>60</sub> molecule.
C<sub>60</sub> nanocrystals were grown on NaCl layers deposited onto a Au(111) substrate. NaCl was evaporated from a Knudsen cell on a clean Au(111) surface at room temperature. Subsequently, the C<sub>60</sub> molecules were sublimated on the NaCl covered substrate. The experiments were performed with a homebuilt ultrahigh vacuum (UHV) STM operating at a temperature of $`50`$K, using cut PtIr tips. The photons emitted from the tunnel junction were collected by a lens placed inside the cryostat, guided through an optical system outside the UHV chamber to the spectrograph, and detected by a CCD camera. The wavelength resolution of the experiment was 8 nm, corresponding to $``$20 meV in the energy range of interest. The spectra were acquired with closed feedback loop while tunneling over a defined position on the sample, e.g. over a single molecule, with a typical acquisition time of 300 s. Bias voltages $`V`$ refer to the sample voltage with respect to the tip.
NaCl forms (100)-terminated islands on Au(111) of thickness between 1 and 3 monolayers and width up to $`1\mu `$m. Contrary to the layer-by-layer growth of C<sub>60</sub> on Au(111) leading to extended islands found by STM Altman93 , on NaCl electron microscopy studies Saito92 revealed that the C<sub>60</sub> molecules aggregate into hexagonal or truncated triangular nanocrystals with a height of several molecular layers. This situation is well illustrated in Fig. 1, where C<sub>60</sub> islands grown on both, the bare Au(111) (A) and the NaCl covered surface (B-D) are visible. The nanocrystals present a minimum height of two layers of C<sub>60</sub> molecules (island B). The nucleation of the C<sub>60</sub> nanocrystals starts at defects of the NaCl layer (protrusions or vacancies), monatomic steps of Au(111) (covered with NaCl), or edges of the second layer of NaCl. As shown in Fig. 1(b), the C<sub>60</sub> molecules form hexagonally arranged layers with an intermolecular distance of 1 nm.
Figure 2 shows STM-induced optical spectra from the bare Au(111) surface, the NaCl covered Au(111) surface and from a C<sub>60</sub> nanocrystal. Photons emitted from Au(111) originate from an IET process, involving excitation and decay of a surface plasmon localized between the tip and the surface Berndt91 ; Berndt93B ; Meguro02 . A similar spectral shape is observed over the NaCl layer, however with reduced intensity due to the dielectric NaCl spacer layer. Characteristic for this process is the energy-dependent quantum cut-off Berndt91 (not shown) and the possibility to excite the emission with both bias polarities Berndt93B , as shown in Fig. 2. The third spectrum in Fig. 2(a) was acquired over a single C<sub>60</sub> molecule in the surface of a nanocrystal ($`V=3`$V, $`I=1`$nA). Light emission from C<sub>60</sub> is observed only for bias higher than the threshold voltage of $`V`$ = $``$2.3 V. The emission onset is located at $``$680 nm and its position is independent of the voltage. For positive voltages up to $`+`$4.5 V no photon emission is detected. These observations clearly distinguish the light emission spectrum of C<sub>60</sub> from those acquired over the substrate (Au and NaCl).
The occurrence of luminescence from the C<sub>60</sub> molecule is related to the characteristics of the tunneling junction, as shown in Fig. 3. At negative bias voltage larger than $``$2.3 V, the highest-occupied molecular orbital (HOMO), which is completely filled for C<sub>60</sub> in the ground state, is higher than the Fermi level ($`E_F`$) of the tip. The electrons are extracted from the HOMO and tunnel to the tip, while the lowest-unoccupied molecular orbital (LUMO), now lower than the Fermi level of the sample, is populated by the electrons tunneling from the substrate, electrons that can radiatively decay into the partially empty HOMO (hot electron/hole injection). The fact that luminescence is not observed for tunneling from the tip to the sample for voltages up to $`+`$4.5 V may be due to the asymmetry of the HOMO-LUMO gap with respect to $`E_F`$ Lof92 and to the different properties of the two tunneling barriers (vacuum and NaCl) Wu04 .
Figure 4(a) shows the same luminescence spectrum obtained from an individual C<sub>60</sub> molecule as in Fig. 2(a), but corrected for the quantum efficiency of the detection system. In order to identify the electronic and vibronic transitions giving rise to the observed emission, we compare our results with laser-induced high-resolution photoluminescence data Sassara96B ; Sassara97 and with quantum chemical calculations Sassara97 ; Orlandi02 . It is now established that the lowest excited singlet state $`S_1`$ has mixed $`T_{1g}`$, $`T_{2g}`$ and $`G_g`$ character Sassara97 . The electric dipole transitions from this state to the ground state $`S_0`$ ($`A_g`$) are symmetry forbidden, but they occur through HT and JT electron-vibration coupling mechanisms of intensity borrowing Negri92 ; Orlandi02 . The relaxation of the selection rules due to symmetry lowering in the C<sub>60</sub> lattice gives rise to a very weak luminescence signal corresponding to the pure electronic (0-0) $`S_0S_1`$ transition, found at $``$678 nm, as indicated in Fig. 4(a). The red shift of about 40 nm with respect to C<sub>60</sub> in the gas phase is attributed to environmental effects Sassara96B ; Sassara97 ; vandenHeuvel95 . The observation of the pure electronic origin helps to determine the vibronically induced false origins as in the high-resolution photoluminescence measurements Hung96 ; Sassara97 . The fluorescence spectra are simulated using calculated oscillator strengths and experimentally determined frequencies for the HT vibronically induced $`S_0S_1`$ ($`T_{1g}`$, $`T_{2g}`$, $`G_g`$) transitions, and the experimental frequencies for the JT active modes, as presented in Fig. 4(a) and in Tab. 1 Sassara97 ; Schettino94 . The contribution of each symmetry character of $`S_1`$ varies slightly from one probed molecule to another, reflecting the known sensitivity of C<sub>60</sub> to the local environment Sassara97 ; vandenHeuvel95B ; Chergui05 . The agreement between measured and calculated spectra in Fig. 4(a) demonstrates the local character of the measurement and provides evidence for the preservation of the C<sub>60</sub> molecular properties in the van de Waals crystal, characterized by weak interactions between the molecules.
Interestingly, we also observe another type of electronic transitions, shown in Fig. 4(c). Similar spectra have been reported for laser-induced luminescence from solid C<sub>60</sub> Guss94 ; vandenHeuvel95B ; Akimoto02 , and have recently been identified as phosphorescence originating from triplet to singlet ground state transitions Chergui05 . Although symmetry and spin-forbidden, intense pure electronic (0-0) triplet to singlet transitions have been observed in C<sub>60</sub> phosphorescence spectra Sassara96 ; vandenHeuvel94 . The low-energy part of the spectrum in Fig. 4(c) arises from the $`S_0T_1`$ ($`{}_{}{}^{3}T_{2g}^{}`$) transition, characterized by the intense 0-0 origin at $``$803 nm and by the progression of a JT active mode, see Tab. 1, in agreement with phosphorescence spectra obtained from dispersed C<sub>60</sub> molecules vandenHeuvel94 ; Sassara96 ; Hung96 . The electronic origin is shifted by 27 nm to the red with respect to the estimated gas phase energy Sassara96 ; Chergui05 . The high-energy region of the spectrum presents a similar shape, i.e. an intense transition, located at $``$720 nm, and a progression of vibronic bands. The observed energy difference of $``$0.18 eV between the two most intense features in the spectrum shown in Fig. 4(c) is in agreement with electron energy loss spectroscopy results Gensterblum91 ; Chergui05 and with calculations Laszlo89 for the splitting of the lowest C<sub>60</sub> triplet states. Therefore, the high-energy part of the spectrum in Fig. 4(c) is assigned to transitions from the next higher triplet state $`S_0T_2`$ ($`{}_{}{}^{3}T_{1g}^{}`$) Chergui05 . The good agreement between the calculated and the measured spectra allows us to assign the STM-induced phosphorescence spectrum to light emission from an individual C<sub>60</sub> molecule, as in the case of fluorescence. This finding contradicts previous interpretations of similar spectra obtained by laser induced luminescence from solid C<sub>60</sub> in terms of excitonic emission Guss94 or emission delocalized over more than one C<sub>60</sub> molecule vandenHeuvel95B ; Akimoto02 .
The observation of both radiative relaxation processes, fluorescence and phosphorescence, in the STM-induced light emission may be related to (i) the sensitivity of the probed C<sub>60</sub> molecule to the local environment in the nanocrystal and/or to (ii) the actual tunneling conditions. (i) Differences in the local C<sub>60</sub> environment may induce a modification of the mixed character of the states and/or a relaxation of the selection rules, as observed in ensemble-averaged experiments Hung96 ; Sassara96 ; Sassara96B ; Sassara97 ; vandenHeuvel94 ; vandenHeuvel95 ; Guss94 ; vandenHeuvel95B ; Akimoto02 ; Chergui05 . Novel, here, is the fact that the influence of the environment is probed on an individual selected molecule. (ii) Even for equal nominal tunneling parameters, the actual tunneling conditions can vary from one measurement to the other. Tip shape and composition, current instabilities, or electric field fluctuation may influence the relaxation paths, for example by enhancing the intersystem crossing, and increasing the population of the triplet states. The identification of the physical origin of the observation of both, fluorescence and phosphorescence calls for future time-resolved luminescence studies employing both, non-local laser excitation and local STM-induced excitation of supported molecules.
To summarize, unambiguous chemical identification of individual C<sub>60</sub> molecules is obtained via their luminescence induced by tunneling electrons. Emission from three electronic states mapping more than twenty vibrational Jahn-Teller and Herzberg-Teller modes of the molecule is identified, in excellent agreement with the known energies of the electronic Hung96 ; Sassara96 ; Sassara96B ; Sassara97 ; vandenHeuvel94 ; Chergui05 and vibrational Sassara97 ; Orlandi02 ; Schettino94 levels of C<sub>60</sub>. The present observation of local fluorescence and phosphorescence demonstrates the capability of STM-induced light emission for the chemical recognition on the single-molecular scale.
The financial support of the Swiss National Science Foundation is acknowledged.
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# Lie Groupoids and Lie algebroids in physics and noncommutative geometry
## 1 Introduction
Influenced by mathematicians such as Grothendieck, Mackey, Connes, and Weinstein, the use of groupoids in pure mathematics has become respectable (though by no means widespread), at least in their respective areas of algebraic geometry, representation theory, noncommutative geometry, and symplectic geometry.<sup>1</sup><sup>1</sup>1There is a Groupoid Home Page at http://unr.edu/homepage/ramazan/groupoid/. See also http://www.cameron.edu/$``$koty/groupoids/ for an incomplete but useful list of papers involving groupoids, necessarily restricted to mathematics. Unfortunately, in physics groupoids remain virtually unknown.<sup>2</sup><sup>2</sup>2Conferences such as Groupoids in Analysis, Geometry, and Physics (Boulder, 1999, see ) and Groupoids and Stacks in Physics and Geometry (Luminy, 2004) tend te be almost exlusively attended by mathematicians. This is a pity for at least two reasons. Firstly, much of the spectacular mathematics developed in the areas just mentioned becomes inaccessible to physicists, despite its undeniable relevance to physics. This obstructs, for example, the development of a good theory for quantizing singular spaces (of the kind necessary for quantum cosmology); cf. . As a case in point, many completely natural constructions in noncommutative geometry look mysterious to physicists who are not familiar with groupoids. Secondly, in the smooth setting, Lie groupoids along with their associated infinitesimal objects called Lie algebroids provide an ideal framework for practically all aspects of both classical and quantum physics that involve symmetry in one way or the other.
Indeed, whereas in the work of Grothendieck and Connes groupoids mainly occur as generalizations of equivalence relations,<sup>3</sup><sup>3</sup>3Grothendieck (to R. Brown in a letter from 1985): “The idea of making systematic use of groupoids (notably fundamental groupoids of spaces, based on a given set of base points), however evident as it may look today, is to be seen as a significant conceptual advance, which has spread into the most manifold areas of mathematics. (…) In my own work in algebraic geometry, I have made extensive use of groupoids - the first one being the theory of the passage to quotient by a ‘pre-equivalence relation’ (which may be viewed as being no more, no less than a groupoid in the category one is working in, the category of schemes say), which at once led me to the notion (nowadays quite popular) of the nerve of a category. The last time has been in my work on the Teichmüller tower, where working with a ‘Teichmüller groupoid’ (rather than a ‘Teichmüller group’) is a must, and part of the very crux of the matter (…)” the role of groupoids as generalized symmetries has been emphasized by Weinstein : “Mathematicians tend to think of the notion of symmetry as being virtually synonymous with the theory of groups and their actions.<sup>4</sup><sup>4</sup>4Cf. Connes: “It is fashionable among mathematicians to despise groupoids and to consider that only groups have an authentic mathematical status, probably because of the pejorative suffix oid.” (…) In fact, though groups are indeed sufficient to characterize homogeneous structures, there are plenty of objects which exhibit what we clearly recognize as symmetry, but which admit few or no nontrivial automorphisms. It turns out that the symmetry, and hence much of the structure, of such objects can be characterized if we use groupoids and not just groups.
The aim of this paper is to (briefly) explain what Lie groupoids and Lie algebroids are, and (more extensively) to outline which role they play in physics (at least from the perspective of the author). Because of the close relationship between quantum theory and noncommutative geometry on the one hand, and classical mechanics and symplectic geometry on the other,<sup>5</sup><sup>5</sup>5Throughout this paper we use the term ‘symplectic geometry’ so as to include Poisson geometry. our discussion obviously relates to matters of pure mathematics as well, and here the physics perspective turns out to be quite useful in clarifying the relationship between noncommutative and symplectic geometry. This relationship is rarely studied in noncommutative geometry, which might explain the regrettable absence of the concept of a Lie algebroid from the field.<sup>6</sup><sup>6</sup>6Except for the work of the author, the sole exception known to him is .
With this goal in mind, one of our main points will be to show that the role of Lie groupoids on the quantum or noncommutative side is largely paralleled by the role Lie algebroids play on the classical or symplectic side. The highlight of this philosophy is undoubtedly the close analogy between Connes’s map $`\mathrm{\Gamma }C^{}(\mathrm{\Gamma })`$ in noncommutative geometry and Weinstein’s map $`\mathrm{\Gamma }A^{}(\mathrm{\Gamma })`$ in symplectic geometry , notably the functoriality of both . Furthermore, the transition from classical to quantum theory through deformation quantization turns out to be given precisely by the association of the $`C^{}`$-algebra $`C^{}(\mathrm{\Gamma })`$ to the Poisson manifold $`A^{}(\mathrm{\Gamma })`$ . Hence quantization is closely related to ‘integration’, in the sense of the association of a Lie groupoid to a Lie algebroid; see for an introduction to this problem, and for its solution.
We do not provide an extensive mathematical introduction to Lie groupoids and Lie algebroids, partly because we have already done so before , and partly because various excellent textbooks on this subject are now available . Instead, we start entirely on the physics side, with a crash course on canonical quantization and its reformulation by Mackey in terms of systems of imprimitivity. In its original setting Mackey’s notion of quantization was not only limited to homogeneous configuration spaces, but in addition lacked an underlying classical theory.<sup>7</sup><sup>7</sup>7More precisely, the underlying classical theory was not correctly identified . Both drawbacks are entirely removed once one adopts the perspective of Lie groupoids on the quantum side and Lie algebroids on the classical side, and we propose this as a convenient point of entry for physicists into the world of these seemingly strange and unfamiliar objects.
Once this perspective has been adopted, the entire theory of canonical quantization and its (finite-dimensional) generalizations is absorbed into a single theorem, stating that the association of $`C^{}(\mathrm{\Gamma })`$ to $`A^{}(\mathrm{\Gamma })`$ mentioned above is a ‘strict’ deformation quantization (in the sense of Rieffel ). Furthermore, in our opinion the deepest understanding of Mackey’s imprimitivity theorem comes from its derivation from the functoriality of Connes’s map $`\mathrm{\Gamma }C^{}(\mathrm{\Gamma })`$; similarly, the classical analogue of the imprimitivity theorem in symplectic geometry can be derived from the functoriality of Weinstein’s map $`\mathrm{\Gamma }A^{}(\mathrm{\Gamma })`$ already mentioned.
We finally combine the toolkit of noncommutative geometry with that of symplectic geometry in proposing a functorial approach to quantization, which is based on KK-theory on the quantum side and on symplectic groupoids on the classical side. As we see it, this approach provides the ultimate generalization of the ‘quantization commutes with reduction’ philosophy of Dirac (in physics) and Guillemin and Sternberg (in mathematics). Beside the use of the K-theory of $`C^{}`$-algebras, this generalization hinges on the use of Lie groupoids and Lie algebroids, and therefore appears to be an appropriate endpoint of this paper.
## 2 From canonical quantization to systems of imprimitivity
Quantum mechanics was born in 1925 with the work of Heisenberg, who discovered the noncommutative structure of its algebra of observables . The complementary work of Schrödinger from 1926 , on the other hand, rather started from the classical geometric structure of configuration space. Within a year, their work was unified by von Neumann, who introduced the abstract concept of a Hilbert space, in which Schrödinger’s wave functions are vectors, and Heisenberg’s observables are linear operators; see . As every physicist knows, the basic link between matrix mechanics and wave mechanics lies in the identification of Heisenberg’s infinite matrices $`p_j`$ and $`q^i`$ ($`i,j=1,2,3`$), representing the momentum and position of a particle moving in $`^3`$, with Schrödinger’s operators $`i\mathrm{}/x^j`$ and $`x^i`$ (seen as a multiplication operator) on the Hilbert space $`=L^2(^3)`$, respectively. The key to this identification lies in the canonical commutation relations
$$[p_i,q^j]=i\mathrm{}\delta _i^j.$$
(1)
Although a mathematically rigorous theory of these commutation relations (as they stand) exists , they are problematic nonetheless. Firstly, the operators involved are unbounded, and in order to represent physical observables they have to be self-adjoint; yet on their respective domains of self-adjointness the commutator on the left-hand side is undefined. Secondly, (1) relies on the possibility of choosing global coordinates on $`^3`$, which precludes at least a naive generalization to arbitrary configuration spaces.<sup>8</sup><sup>8</sup>8Mackey \[61, p. 283\]: “Simple and elegant as this model is, it appears at first sight to be quite arbitrary and ad hoc. It is difficult to understand how anyone could have guessed it and by no means obvious how to modify it to fit a model for space different from $`^r`$.
Finding an appropriate mathematical interpretation of the canonical commutation relations (1) is the subject of quantization theory; see for recent reviews. From the numerous ways to handle the situation, we here select Mackey’s approach .<sup>9</sup><sup>9</sup>9Continuing the previous quote, Mackey claims with some justification that his approach “(a) Removes much of the mystery. (b) Generalizes in a straightforward way to any model for space with a separable locally compact group of isometries. (c) Relates in an extremely intimate way to \[the theory of induced representations\].” In any case, Mackey’s approach to the canonical commutation relations, especially in its $`C^{}`$-algebraic reformulation presented below, is vastly superior to their equally $`C^{}`$-algebraic reformulation in terms of the so-called Weyl $`C^{}`$-algebra (cf. e.g. ). Indeed (see Def. IV.3.5.1), the Weyl algebra over a Hilbert space $``$ (which in the case at hand is $`^3`$) may be seen as the twisted group $`C^{}`$-algebra over $``$ as an abelian group under addition, equipped with the discrete topology. This rape of $``$ as a topological space is so ugly that it is surprising that papers on the Weyl $`C^{}`$-algebra continue to appear. Historically, Weyl’s exponentiation of the canonical commutation relations was just one of the first attempts to reformulate a problem involving unbounded operators in terms of bounded ones, and has now been superseded. The essential point is to assign momentum and position a quite different role in quantum mechanics, despite the fact that in classical mechanics $`p`$ and $`q`$ can be interchanged by a canonical transformation.<sup>10</sup><sup>10</sup>10This feature is shared by most approaches to quantization, except the one mentioned in the preceding footnote.
Firstly, the position operators $`q^j`$ are collectively replaced by a single projection-valued measure $`P`$ on $`^3`$,<sup>11</sup><sup>11</sup>11A projection-valued measure $`P`$ on a space $`\mathrm{\Omega }`$ with Borel structure (i.e. equipped with a $`\sigma `$-algebra of measurable sets defined by the topology) with values in a Hilbert space $``$ is a map $`EP_E`$ from the Borel subsets $`E\mathrm{\Omega }`$ to the projections on $`H`$ that satisfies $`P_{\mathrm{}}=0`$, $`P_\mathrm{\Omega }=1`$, $`P_EP_F=P_FP_E=P_{EF}`$ for all measurable $`E,F\mathrm{\Omega }`$, and $`P_{_{i=1}^{\mathrm{}}E_i}=_{i=1}^{\mathrm{}}P_{E_i}`$ for all countable collections of mutually disjoint $`E_i\mathrm{\Omega }`$. which is given by $`P_E=\chi _E`$ as a multiplication operator on $`L^2(^3)`$. Given this $`P`$, any multiplication operator $`f`$ defined by a measurable function $`f:^3`$ can be represented as $`f=_^3𝑑P_E(x)f(x)`$, which is defined and self-adjoint on a suitable domain.<sup>12</sup><sup>12</sup>12This domain consists of all $`\psi H`$ for which $`_^3d(\psi ,P_E(x)\psi )|f(x)|^2<\mathrm{}`$. In particular, the position operators $`q^i`$ can be reconstructed from $`P`$ by choosing $`f(x)=x^i`$.
Secondly, the momentum operators $`p_i`$ are collectively replaced by a single unitary group representation $`U(^3)`$ on $`L^2(^3)`$, defined by
$$U(y)\psi (x):=\psi (xy).$$
Each $`p_i`$ can be reconstructed from $`U`$ by means of
$$p_i\psi :=i\mathrm{}\underset{t_i0}{lim}t_i^1(U(t_i)1)\psi ,$$
where $`U(t_i)`$ is $`U`$ at $`x^i=t_i`$ and $`x^j=0`$ for $`ji`$; this operator is defined and self-adjoint on the set of all $`\psi H`$ for which the limit exists (Stone’s theorem ).
Consequently, it entails no loss of generality to work with the pair $`(P,U)`$ instead of the $`(q^j,p_i)`$. The commutation relations (1) are now replaced by
$$U(x)P_EU(x)^1=P_{xE},$$
(2)
where $`E`$ is a Borel subset of $`^3`$ and $`xE=\{x\omega \omega E\}`$. On the basis of this reformulation, Mackey proposed the following sweeping generalization of the the canonical commutation relations.<sup>13</sup><sup>13</sup>13In order to maintain the connection with the classical theory later on, we restrict ourselves to Lie groups acting smoothly on manifolds. Mackey actually formulated his results more generally in terms of separable locally compact groups acting continuously on locally compact spaces.
###### Definition 1
Suppose a Lie group $`G`$ acts smoothly on a manifold $`M`$.
1. A system of imprimitivity $`(,U,P)`$ for this action consists of a Hilbert space $``$, a unitary representation $`U`$ of $`G`$ on $``$, and a projection-valued measure $`EP_E`$ on $`M`$ with values in $``$, such that (2) holds for all $`xG`$ and all Borel sets $`EM`$.
2. A $`G`$-covariant representation $`(,U,\pi )`$ of the $`C^{}`$-algebra $`C_0(M)`$ relative to this action consists of a Hilbert space $``$, a unitary representation $`U`$ of $`G`$ on $``$, and a nondegenerate representation $`\pi `$ of $`C_0(M)`$ on $``$ satisfying
$$U(x)\pi (\phi )U(x)^1=\pi (L_x\phi )$$
(3)
for all $`xG`$ and $`\phi C_0(M)`$, where $`L_x\phi (m)=\phi (x^1m)`$.
The spectral theorem (cf. ) implies that these notions are equivalent: a projection-valued measure $`P`$ defines and is defined by a nondegenerate representation $`\pi `$ of $`C_0(M)`$ on $``$ by means of $`\pi (\phi )=_M𝑑P(m)\phi (m)`$, and (2) is then equivalent to the covariance condition (3). Hence we may interchangeably speak of systems of imprimitivity or covariant representations. As a further reformulation, it is easy to show (cf. ) that there is a bijective correspondence between $`G`$-covariant representations of $`C_0(M)`$ and nondegenerate representations of the so-called transformation group $`C^{}`$-algebra $`C^{}(G,M)G\times _\alpha C_0(M)`$ defined by the given $`G`$-action on $`M`$, which determines an automorphic action $`\alpha `$ of $`G`$ on $`C_0(M)`$ by $`\alpha _x=L_x`$.<sup>14</sup><sup>14</sup>14In one direction, this correspondence is as follows: given a $`G`$-covariant representation $`(,\pi ,U)`$, one defines a representation $`\pi _U(C^{}(G,M))`$ by extension of $`\pi _U(f)=_G𝑑x\pi (f(x,))U(x)`$, where $`fC_c^{\mathrm{}}(G\times M)C^{}(G,M)`$, and $`f(x,)`$ is seen as an element of $`C_0(M)`$.
Such a system describes the quantum mechanics of a particle moving on a configuration space $`M`$ on which $`G`$ acts by symmetry transformations; in particular, each element $`X`$ of the Lie algebra $`𝔤`$ of $`G`$ defines a generalized momentum operator
$$\widehat{X}=i\mathrm{}dU(X)$$
(4)
on $``$, which is defined and self-adjoint on the domain of vectors $`\psi `$ for which
$$dU(X)\psi :=\underset{t0}{lim}t^1(U(\mathrm{exp}(tX))1)\psi $$
exists. These operators satisfy the generalized canonical commutation relations
$$[\widehat{X},\widehat{Y}]=i\mathrm{}\widehat{[X,Y]}$$
(5)
and
$$[\widehat{X},\pi (\phi )]=\pi (\xi _X\phi ),$$
(6)
where $`\phi C_c^{\mathrm{}}(M)`$ and $`\xi _X`$ is the canonical vector field on $`M`$ defined by the $`G`$-action; of course, these should be supplemented with
$$[\pi (\phi _1),\pi (\phi _2)]=0.$$
(7)
Elementary quantum mechanics on $`^n`$ then corresponds to the special case $`M=^n`$ and $`G=^n`$ with the usual additive group structure.
## 3 The imprimitivity theorem
In the spirit of the $`C^{}`$-algebraic approach to quantum physics , the $`C^{}`$-algebra $`C^{}(G,M)`$ defined by the given $`G`$-action on $`M`$ should be seen as an algebra of observables, whose inequivalent irreducible representations define the possible superselection sectors of the system. As we have seen, these representations may equivalently be seen as systems of imprimitivity or as $`G`$-covariant representations of $`C_0(M)`$ . In any case, it is of some interest to classify these. Mackey’s imprimitivity theorem describes the simplest case where this is possible.
###### Theorem 1
Let $`H`$ be a closed subgroup of $`G`$ and let $`G`$ act on $`M=G/H`$ by left translation. Up to unitary equivalence, there is a bijective correspondence between systems of imprimitivity $`(,U,P)`$ for this action (or, equivalently, $`G`$-covariant representation of $`C_0(G/H)`$ or nondegenerate representations of the transformation group $`C^{}`$-algebra $`C^{}(G,G/H)`$) and unitary representations $`U_\chi `$ of $`H`$, as follows:
* Given $`U_\chi (H)`$ on a Hilbert space $`_\chi `$, the triple $`(^\chi ,U^\chi ,P^\chi )`$ is a system of imprimitivity, where $`^\chi =L^2(G/H,G\times _H_\chi )`$ is the Hilbert space of $`L^2`$-sections of the vector bundle $`G\times _H_\chi `$ associated to the principal $`H`$-bundle $`G`$ over $`G/H`$ by $`U_\chi `$, $`U^\chi `$ is the representation of $`G`$ induced by $`U_\chi `$, and $`P_E^\chi =\chi _E`$ acts canonically on $`^\chi `$ as a multiplication operator.
* Conversely, if $`(,U,P)`$ is a system of imprimitivity, then there exists a unitary representation $`U_\chi (H)`$ such that the triple $`(,U,P)`$ is unitarily equivalent to the triple $`(^\chi ,U^\chi ,P^\chi )`$ just described.
The correspondence $`(_\chi ,U_\chi )(^\chi ,U^\chi ,P^\chi )`$ preserves direct sums and, accordingly, irreducibility.
The simplest and at the same time most beautiful application of the imprimitivity theorem is Mackey’s recovery of the Stone–von Neumann uniqueness theorem concerning the (regular) irreducible representations of the canonical commutation relations: taking $`G=^3`$ and $`H=\{e\}`$ (so that $`M=^3`$), one finds that the associated system of imprimitivity possesses precisely one irreducible representation, since the trivial group obviously has only one such representation.<sup>15</sup><sup>15</sup>15The “uniqueness of the canonical commutation relations” has also been derived from the fact that (up to unitary equivalence) there is only one irreducible representation of any of the following objects: i) The Heisenberg Lie group with given nonzero central charge (von Neumann’s theorem ); ii) The Weyl $`C^{}`$-algebra over a finite-dimensional Hilbert space, provided one restricts oneself to the class of regular representations ; or iii) The $`C^{}`$-algebra of compact operators . Furthermore (and this was one of Mackey’s main points), one may keep $`^3`$ as a confuguration space but replace $`G=^3`$ by the Euclidean group $`G=SO(3)^3`$, so that $`H=SO(3)`$. The generalized momenta then include the angular momentum operators $`J^i`$ along with their commutation relations, and the imprimitivity theorem then asserts that the irreducible representations of (2) correspond to the usual irreducible representations $`U_j`$ of $`SO(3)`$, $`j=0,1,\mathrm{}`$.<sup>16</sup><sup>16</sup>16By the usual arguments, one may replace $`SO(3)`$ by $`SU(2)`$ in this argument, so as to obtain $`j=0,1/2,\mathrm{}`$. Mackey saw this as an explanation for the emergence of spin as a purely quantum-mechanical degree of freedom; the latter perspective of spin goes back to the pionieers of quantum theory , but is now obsolete (see Section 9 below).
Mackey’s imprimitivity theorem admits a generalization to $`G`$-actions on an arbitrary manifold $`M`$, provided the action is regular.<sup>17</sup><sup>17</sup>17In view of this simple result, $`C^{}`$-algebraists are mainly interested in nonregular actions, cf. , but for physics Proposition 1 is quite useful. In any case, an example of a nonregular action is the action of $``$ on $`𝕋`$ by irrational rotations.
###### Proposition 1
Suppose that each $`G`$-orbit in $`M`$ is (relatively) open in its closure. The irreducible representations of $`C^{}(GM)`$ are classified by pairs $`(𝒪,U_\chi )`$, where $`𝒪`$ is a $`G`$-orbit in $`M`$ and $`U_\chi `$ is an irreducible representation of the stabilizer of an arbitrary point $`m_0𝒪`$.<sup>18</sup><sup>18</sup>18The associated $`G`$-covariant representation of $`C_0(M)`$ may be realized by multiplication operators on the Hilbert space $`^\chi `$ carrying the representation $`U^\chi (G)`$ induced by $`U_\chi `$.
In view of the power of Mackey’s imprimitivity theorem, both for representation theory and quantization theory, increasingly sophisticated and insightful proofs have been published over the last five decades.<sup>19</sup><sup>19</sup>19Mackey’s own proof was rather measure-theoretic in flavour, and did not shed much light on the origin of his result. Probably the shortest proof is . All proofs relevant to noncommutative geometry are either based on or are equivalent to:
###### Theorem 2
The transformation group $`C^{}`$-algebra $`C^{}(G,G/H)`$ is Morita equivalent to $`C^{}(H)`$.
This means that there exists a so-called equivalence or imprimitivity bimodule $``$ (which in modern terms would be called a $`C^{}(G,G/H)`$-$`C^{}(H)`$ Hilbert bimodule)<sup>20</sup><sup>20</sup>20A Hilbert bimodule $`AB`$ over $`C^{}`$-algebras $`A`$ and $`B`$ consists of a Banach space $``$ that is an algebraic $`A`$-$`B`$ bimodule, and is equipped with a $`B`$-valued inner product that is compatible with the $`A`$ and $`B`$ actions. Such objects were first considered by Rieffel , who defined an ‘interior’ tensor product $`\widehat{}_B`$ of an $`A`$-$`B`$ Hilbert bimodule $``$ with a $`B`$-$`C`$ Hilbert bimodule $``$, which is an $`A`$-$`C`$ Hilbert bimodule. that allows one to set up the bijective correspondence - called for in Mackey’s imprimitivity theorem - between (nondegerenerate) representations of $`C^{}(G,G/H)`$ and those of $`C^{}(H)`$ (or equivalently, of $`H`$). Given a unitary representation $`U_\chi (H)`$ on a Hilbert space $`_\chi `$, or the associated representation $`\pi _\chi `$ of $`C^{}(H)`$ on the same space, one constructs a Hilbert space $`^\chi =\widehat{}_{\pi _\chi }_\chi `$. The action of $`C^{}(G,G/H)`$ on $``$ descends to an action $`\pi ^\chi (C^{}(G,G/H))`$ on $`^\chi `$, and extracting the associated representations of $`G`$ and of $`C_0(G/H)`$ one finds that this is precisely Mackey’s induction construction paraphrased in Theorem 1. Conversely, a given representation $`\pi ^\chi `$ of $`C^{}(G,G/H)`$ on a Hilbert space $`^\chi `$ defines $`\pi _\chi (C^{}(H))`$ on $`_\chi =\overline{}\widehat{}_{\pi ^\chi }^\chi `$, and this process is the inverse of the previous one. Replacing the usual algebraic bimodule tensor product by Rieffel’s interior tensor product $`\widehat{}_\pi `$, this entirely mimics the corresponding procedure in algebra (cf. ); in the same spirit, one infers also in general that two Morita equivalent $`C^{}`$-algebras have equivalent representation categories.
The reformulation of Theorem 1 as Theorem 2 begs the question what the deeper origin of the latter could possibly be. One answer is given by the analysis in , from which Theorem 2 emerges as merely a droplet in an ocean of imprimitivity theorems. The answer below is equally categorical in spirit, but is entirely based on the use of Lie groupoids. Namely, we will derive Theorem 2 and hence Mackey’s Theorem 1 from the functoriality of Connes’s map (8) below, which associates a $`C^{}`$-algebra to a Lie groupoid. Apart from the fact that this is very much in the spirit of noncommutative geometry, the use of Lie groupoids will enable us to formulate an analogous classical procedure in terms of Lie algebroids and Poisson manifolds. All this requires a little preparation.
## 4 Intermezzo: Lie groupoids
Recall that a groupoid is a small category (i.e. a category in which the underlying classes are sets) in which each arrow is invertible. We denote the total space (i.e. the set of arrows) of a groupoid $`\mathrm{\Gamma }`$ by $`\mathrm{\Gamma }_1`$, and the base space (i.e. the set on which the arrows act) by $`\mathrm{\Gamma }_0`$; the object inclusion map $`\mathrm{\Gamma }_0\mathrm{\Gamma }_1`$ is written $`u1_u`$. We denote the inverse $`\mathrm{\Gamma }_1\mathrm{\Gamma }_1`$ by $`xx^1`$, and the source and target maps by $`s,t:\mathrm{\Gamma }_1\mathrm{\Gamma }_0`$. Thus the composable pairs form the space $`\mathrm{\Gamma }_2:=\{(x,y)\mathrm{\Gamma }_1\times \mathrm{\Gamma }_1s(x)=t(y)\}`$, so that if $`(x,y)\mathrm{\Gamma }_2`$ then $`xy\mathrm{\Gamma }_1`$ is defined.<sup>21</sup><sup>21</sup>21Thus the axioms are: 1. $`s(xy)=s(y)`$ and $`t(xy)=t(x)`$; 2. $`(xy)z=x(yz)`$ 3. $`s(1_u)=t(1_u)=u`$ for all $`u\mathrm{\Gamma }_0`$; 4. $`x1_{s(x)}=1_{t(x)}x=x`$ for all $`x\mathrm{\Gamma }_1`$. A Lie groupoid is a groupoid for which $`\mathrm{\Gamma }_1`$ and $`\mathrm{\Gamma }_0`$ are manifolds ($`\mathrm{\Gamma }_1`$ not necessarily being Hausdorff), $`s`$ and $`t`$ are surjective submersions, and multiplication and inversion are smooth.<sup>22</sup><sup>22</sup>22It follows that object inclusion is an immersion, that inversion is a diffeomorphism, that $`\mathrm{\Gamma }_2`$ is a closed submanifold of $`\mathrm{\Gamma }_1\times \mathrm{\Gamma }_1`$, and that for each $`u\mathrm{\Gamma }_0`$ the fibers $`s^1(u)`$ and $`t^1(u)`$ are submanifolds of $`\mathrm{\Gamma }_1`$. See for recent textbooks on Lie groupoids and related matters.<sup>23</sup><sup>23</sup>23The concept of a Lie groupoid was introduced by Ehresmann.
Some examples of Lie groupoids that are useful to keep in mind are:
* A Lie group $`G`$, where $`\mathrm{\Gamma }_1=G`$ and $`\mathrm{\Gamma }_0=\{e\}`$).
* A manifold $`M`$, where $`\mathrm{\Gamma }_1=\mathrm{\Gamma }_0=M`$ with the obvious trivial groupoid structure $`s(x)=t(x)=1_x=x^1=x`$, and $`xx=x`$.
* The pair groupoid over a manifold $`M`$, where $`\mathrm{\Gamma }_1=M\times M`$ and $`\mathrm{\Gamma }_0=M`$, with $`s(x,y)=y`$, $`t(x,y)=x`$, $`(x,y)^1=(y,x)`$, $`(x,y)(y,z)=(x,z)`$, and $`1_x=(x,x)`$.
* The gauge groupoid defined by a principal $`H`$-bundle $`P\stackrel{\pi }{}M`$, where $`\mathrm{\Gamma }_1=P\times _HP`$ (which stands for $`(P\times P)/H`$ with respect to the diagonal $`H`$-action on $`P\times P`$), $`\mathrm{\Gamma }_0=M`$, $`s([p,q])=\pi (q)`$, $`t([p,q])=\pi (p)`$, $`[x,y]^1=[y,x]`$, and $`[p,q][q,r]=[p,r]`$ (here $`[p,q][q^{},r]`$ is defined whenever $`\pi (q)=\pi (q^{})`$, but to write down the product one picks $`q\pi ^1(q^{})`$).
* The action groupoid $`GM`$ defined by a smooth (left) action $`GM`$ of a Lie group $`G`$ on a manifold $`M`$, where $`\mathrm{\Gamma }_1=G\times M`$, $`\mathrm{\Gamma }_0=M`$, $`s(g,m)=g^1m`$, $`t(g,m)=m`$, $`(g,m)^1=(g^1,g^1m)`$, and $`(g,m)(h,g^1m)=(gh,m)`$.
As mentioned before, an equivalence relation on a set $`M`$ defines a groupoid, namely the obvious subgroupoid of the pair groupoid over $`M`$. However, in interesting examples this is rarely a Lie groupoid. To obtain a Lie groupoid resembling a given equivalence relation on a manifold, various refinements of the subgroupoid in question have been invented, of which the holonomy groupoid defined by a foliation is the most important example for noncommutative geometry .
For reasons to emerge from the ensuing story, we look at Lie groupoids as objects in the category of principal bibundles. To define this category, we first recall that an action of a groupoid $`\mathrm{\Gamma }`$ on a space $`M`$ is only defined if $`M`$ comes equipped with a map $`M\stackrel{\pi }{}\mathrm{\Gamma }_0`$. In that case, a left $`\mathrm{\Gamma }`$ action on $`M`$ is a map $`(x,m)xm`$ from $`\mathrm{\Gamma }_1\times _{\mathrm{\Gamma }_0}^{s,\pi }M`$ to $`M`$,<sup>24</sup><sup>24</sup>24Here we use the notation $`A\times _B^{f,g}C=\{(a,c)A\times Cf(a)=g(c)\}`$ for the fiber product of sets $`A`$ and $`C`$ with respect to maps $`f:AB`$ and $`g:CB`$. such that $`\pi (xm)=t(x)`$, $`xm=m`$ for all $`x\mathrm{\Gamma }_0`$, and $`x(ym)=(xy)m`$ whenever $`s(y)=\tau (m)`$ and $`t(y)=s(x)`$. Similarly, given a map $`M\stackrel{\rho }{}\mathrm{\Delta }_0`$, a right action of a groupoid $`\mathrm{\Delta }`$ on $`M`$ is a map $`(m,h)mh`$ from $`M\times _{\mathrm{\Delta }_0}^{\rho ,t}\mathrm{\Delta }_1`$ to $`M`$ that satisfies $`\rho (mh)=s(h)`$, $`mh=m`$ for all $`h\mathrm{\Delta }_0`$, and $`(mh)k=m(hk)`$ whenever $`\rho (m)=t(h)`$ and $`t(k)=s(h)`$. Now, if $`\mathrm{\Gamma }`$ and $`\mathrm{\Delta }`$ are groupoids, a $`\mathrm{\Gamma }`$-$`\mathrm{\Delta }`$ bibundle $`M`$, also written as $`\mathrm{\Gamma }M\mathrm{\Delta }`$, carries a left $`\mathrm{\Gamma }`$ action as well as a right $`\mathrm{\Delta }`$-action that commute.<sup>25</sup><sup>25</sup>25That is, one has $`\tau (mh)=\tau (m)`$, $`\rho (xm)=\rho (m)`$, and $`(xm)h=x(mh)`$ whenever defined. Such a bibundle is called principal when $`\pi :M\mathrm{\Gamma }_0`$ is surjective, and the $`\mathrm{\Delta }`$ action is free (in that $`mh=m`$ iff $`h\mathrm{\Delta }_0`$) and transitive along the fibers of $`\pi `$.
Suppose one has right principal bibundles $`\mathrm{\Gamma }M\mathrm{\Delta }`$ and $`\mathrm{\Delta }N\mathrm{\Theta }`$. The fiber product $`M\times _{\mathrm{\Delta }_0}N`$ carries a right $`\mathrm{\Delta }`$ action, given by $`h:(m,n)(mh,h^1n)`$ (defined as appropriate). The orbit space $`(M\times _{\mathrm{\Delta }_0}N)/\mathrm{\Delta }`$ is a $`\mathrm{\Gamma }`$-$`\mathrm{\Theta }`$ bibundle in the obvious way inherited from the original actions. Thus, regarding $`\mathrm{\Gamma }M\mathrm{\Delta }`$ as an arrow from $`\mathrm{\Gamma }`$ to $`\mathrm{\Delta }`$ and $`\mathrm{\Delta }N\mathrm{\Theta }`$ as an arrow from $`\mathrm{\Delta }`$ to $`\mathrm{\Theta }`$, one map look upon $`\mathrm{\Gamma }(M\times _\mathrm{\Delta }N)/\mathrm{\Delta }\mathrm{\Theta }`$ as an arrow from $`\mathrm{\Gamma }`$ to $`\mathrm{\Theta }`$, defining the product or composition of $`M`$ and $`N`$. However, this product is associative merely up to isomorphism, so that in order to have a category one should regard isomorphism classes of principal bibundles as arrows.
For Lie groupoids everything in these definitions has to be smooth (and $`\pi `$ a surjective submersion).
###### Definition 2
The category $`𝔊`$ of Lie groupoids and principal bibundles has Lie groupoids as objects and isomorphism classes $`[\mathrm{\Gamma }M\mathrm{\Delta }]`$ of principal bibundles as arrows. Composition of arrows is given by
$$[\mathrm{\Gamma }M\mathrm{\Delta }][\mathrm{\Delta }N\mathrm{\Theta }]=[\mathrm{\Gamma }(M\times _\mathrm{\Delta }N)/\mathrm{\Delta }\mathrm{\Theta }],$$
and the identities are given by $`1_\mathrm{\Gamma }=[\mathrm{\Gamma }\mathrm{\Gamma }\mathrm{\Gamma }]`$, seen as a bibundle in the obvious way.
Of course, it can be checked that this definition is correct in the sense that one indeed defines a category in this way. This category has the remarkable feature that (Morita) equivalence of groupoids (as defined in , a notion heavily used in noncommutative geometry) is the same as isomorphism of objects in $`𝔊`$.
## 5 From Lie groupoids to the imprimitivity theorem
A central idea in noncommutative geometry is the association
$$\mathrm{\Gamma }C^{}(\mathrm{\Gamma })$$
(8)
of a $`C^{}`$-algebra $`C^{}(\mathrm{\Gamma })`$ to a Lie groupoid $`\mathrm{\Gamma }`$ .<sup>26</sup><sup>26</sup>26See also for detailed presentations. For a Lie groupoid $`\mathrm{\Gamma }`$ Connes’s $`C^{}(\mathrm{\Gamma })`$ is the same (up to isomorphism of $`C^{}`$-algebras) as the $`C^{}`$-algebra Renault associates to a locally compact groupoid with Haar system , provided one takes the Haar system canonically defined by the smooth structure on $`\mathrm{\Gamma }`$. Here $`C^{}(\mathrm{\Gamma })`$ is a suitable completion of the function space $`C_c^{\mathrm{}}(\mathrm{\Gamma }_1)`$, equipped with a convolution-type product defined by the groupoid structue. For the above examples, this yields:
* The $`C^{}`$-algebra of a Lie group $`G`$ is the usual convolution $`C^{}`$-algebra $`C^{}(G)`$ defined by the Haar measure on $`G`$ .
* For a manifold $`M`$ one has $`C^{}(M)=C_0(M)`$.
* The pair groupoid over a connected manifold $`M`$ defines $`C^{}(M\times M)K(L^2(M))`$, i.e. the $`C^{}`$-algebra of compact operators on the $`L^2`$-space canonically defined by a manifold.
* The $`C^{}`$-algebra defined by a gauge groupoid $`P\times _HP`$ as above is isomorphic to $`K(L^2(M))C^{}(H)`$ (but any explicit isomorphism depends on the choice of a measurable section $`s:MP`$, which in general cannot be smooth).
* For an action groupoid defined by $`GM`$ one has $`C^{}(GM)C^{}(G,M)`$, the transformation group $`C^{}`$-algebra defined by the given action .
Having already defined the category $`𝔊`$ of principal bibundles for Lie groupoids, in order to make the map (8) functorial, one has to regard $`C^{}`$-algebras as objects in a suitable category $``$ as well.
###### Definition 3
The category $``$ has $`C^{}`$-algebras as objects and isomorphism classes $`[AB]`$ of Hilbert bimodules, as arrows, composed using Rieffel’s interior tensor product. The identities are given by $`1_A=AAA`$, defined in the obvious way.
A crucial feature of this construction is that the notion of isomorphism of objects in $``$ coincides with Rieffel’s (strong) Morita equivalence of $`C^{}`$-algebras.
###### Theorem 3
Connes’s map $`\mathrm{\Gamma }C^{}(\mathrm{\Gamma })`$ is functorial from the category $`𝔊`$ of Lie groupoids and principal bibundles to the category $``$ of $`C^{}`$-algebras and Hilbert bimodules.
###### Corollary 1
Connes’s map $`\mathrm{\Gamma }C^{}(\mathrm{\Gamma })`$ preserves Morita equivalence, in the sense that if $`\mathrm{\Gamma }`$ and $`\mathrm{\Delta }`$ are Morita equivalent Lie groupoids, then $`C^{}(\mathrm{\Gamma })`$ and $`C^{}(\mathrm{\Delta })`$ are Morita equivalent $`C^{}`$-algebras.
The imprimitivity bimodule $`C^{}(\mathrm{\Gamma })C^{}(\mathrm{\Delta })`$ establishing the Morita equivalence of $`C^{}(\mathrm{\Gamma })`$ and $`C^{}(\mathrm{\Delta })`$ is obtained from the principal bibundle $`\mathrm{\Gamma }E\mathrm{\Delta }`$ establishing the Morita equivalence of $`\mathrm{\Gamma }`$ and $`\mathrm{\Delta }`$ in a very simple way, amounting to the completion of $`C_c^{\mathrm{}}(\mathrm{\Gamma })C_c^{\mathrm{}}(E)C_c^{\mathrm{}}(\mathrm{\Delta })`$; see .
For example, in Mackey’s case one has $`\mathrm{\Gamma }=G(G/H)`$ and $`\mathrm{\Delta }=H`$, linked by the principal bibundle $`G(G/H)GH`$ in the obvious way;<sup>27</sup><sup>27</sup>27For example, $`(g_1,m)g_2=g_1g_2`$, defined whenever $`m=\pi (g_1g_2)`$. the associated imprimitivity bimodule for $`C^{}(G(G/H))C^{}(G,G/H)`$ and $`C^{}(H)`$ is precisely the one found by Rieffel . Thus Theorem 2, and thereby Mackey’s imprimitivity theorem, ultimately derives from the Morita equivalence
$$G(G/H)H$$
(9)
of groupoids, which is an almost trivial fact once the appropriate framework has been set up. This framework cannot be specified in terms of groups and group actions alone, despite the fact that the two groupoids relevant to Mackey’s imprimitivity theorem reduce to those.
Mackey’s analysis of the canonical commutation relations admits various other generalizations than Proposition 1, at least one of which is related to groupoids as well: instead of generalizing the action groupoid $`G(G/H)`$ to an arbitrary action groupoid $`GM`$, one may note the isomorphism of groupoids
$$G(G/H)G\times _HG,$$
(10)
where the right-hand side is the gauge groupoid of the principal $`H`$-bundle $`G`$ with respect to the natural right-action of $`H`$. This isomorphism (given by $`(xy^1,\pi (x))[x,y]`$) naturally passes to the ‘algebra of observables,’ i.e. one has
$$C^{}(G(G/H))C^{}(G\times _HG),$$
(11)
and one may see the right-hand side as a special case of $`C^{}(P\times _HP)`$ for an arbitrary principal $`H`$-bundle $`P`$.<sup>28</sup><sup>28</sup>28This generalization is closely related to Kaluza–Klein theory and the Wong equations; see . Here one has a complete analogue of Mackey’s imprimitivity theorem: the Morita equivalence
$$P\times _HPH$$
(12)
at the groupoid level<sup>29</sup><sup>29</sup>29The equivalence bibundle is $`P\times _HPPH`$, with the given right $`H`$ action on $`P`$ and the left action given by $`[x,y]y=x`$. induces a Morita equivalence
$$C^{}(P\times _HP)C^{}(H)$$
(13)
at the $`C^{}`$-algebraic level, which in turn implies that there is a bijective correspondence between (irreducible) unitary representations $`U_\chi (H)`$ and representations $`\pi ^\chi (C^{}(P\times _HP))`$.<sup>30</sup><sup>30</sup>30Given $`U_\chi (H)`$ on a Hilbert space $`_\chi `$, the representation $`\pi ^\chi `$ is naturally realized on $`L^2(P/H,P\times _H_\chi )`$, as in the homogeneous case.
In the old days, the various irreducible representations (or superselection sectors) of algebras of observables like $`C^{}(GM)`$ or $`C^{}(P\times _HP)`$ were seen as ‘inequivalent quantizations’ of a single underlying classical system. From this perspective, quantities like spin were seen as degrees of freedom peculiar to and emergent from quantum theory. Starting with geometric quantization in the mid-1960s, however, it became clear that each superselection sectors of said type is in fact the quantization of a different classical system. The language of Lie groupoids and Lie algebroids allows the most precise and conceptually clearest discussion of this situation. Mathematically, what is at stake here is the relationship between noncommutative geometry and symplectic geometry as its classical analogue.<sup>31</sup><sup>31</sup>31See also for a different approach to this relationship We now turn to this language.
## 6 Intermezzo: Lie algebroids and Poisson manifolds
Since the notion of a Lie algebroid cannot found in the noncommutative geometry literature, we provide a complete definition.<sup>32</sup><sup>32</sup>32Cf. for detailed treatments. The concept of a Lie algebroid and the relationship between Lie groupoids and Lie algebroids are originally due to Pradines.
###### Definition 4
A Lie algebroid $`A`$ over a manifold $`M`$ is a vector bundle $`A\stackrel{\pi }{}M`$ equipped with a vector bundle map $`A\stackrel{\alpha }{}TM`$ (called the anchor), as well as with a Lie bracket $`[,]`$ on the space $`C^{\mathrm{}}(M,A)`$ of smooth sections of $`A`$, satisfying the Leibniz rule
$$[\sigma _1,f\sigma _2]=f[\sigma _1,\sigma _2]+(\alpha \sigma _1f)\sigma _2$$
(14)
for all $`\sigma _1,\sigma _2C^{\mathrm{}}(M,A)`$ and $`fC^{\mathrm{}}(M)`$.
It follows that the map $`\sigma \alpha \sigma :C^{\mathrm{}}(M,A)C^{\mathrm{}}(M,TM)`$ induced by the anchor is a homomorphism of Lie algebras, where the latter is equipped with the usual commutator of vector fields.<sup>33</sup><sup>33</sup>33This homomorphism property used to be part of the definition of a Lie algebroid, but as observed by Marius Crainic it follows from the stated definition.
Lie algebroids generalize (finite-dimensional) Lie algebras as well as tangent bundles, and the (infinite-dimensional) Lie algebra $`C^{\mathrm{}}(M,A)`$ could be said to be of geometric origin in the sense that it derives from an underlying finite-dimensional geometrical object. Similar to our list of example of Lie groupoids in Section 4, one has the following basic classes of Lie algebroids.
* A Lie algebra $`𝔤`$, where $`A=𝔤`$ and $`M`$ is a point (which may be identified with the identity element of any Lie group with Lie algebra $`g`$; see below) and $`\alpha =0`$.
* A manifold $`M`$, where $`A=M`$, seen as the zero-dimensional vector bundle over $`M`$, evidently with identically vanishing Lie bracket and anchor.
* The tangent bundle over a manifold $`M`$, where $`A=TM`$ and $`\alpha =\text{id}:TMTM`$, with the Lie bracket given by the usual commutator of vector fields.
* The gauge algebroid defined by a principal $`H`$-bundle $`PM`$; here $`A=(TP)/H`$, so that $`C^{\mathrm{}}(M,A)C^{\mathrm{}}(M,TP)^H`$, which inherits the commutator from $`C^{\mathrm{}}(M,TP)`$ as the Lie bracket defining the algebroid structure, and is equipped with the projection $`\alpha :(TP)/HTM`$ induced by $`TPTM`$.
* The action algebroid $`𝔤M`$ defined by a $`𝔤`$-action on a manifold $`M`$ (i.e. a Lie algebra homomorphism $`𝔤C^{\mathrm{}}(M,TM)`$) has $`A=𝔤\times M`$ (as a trivial bundle) and $`\alpha (X,m)=\xi _X(m)T_mM`$. The Lie bracket is
$$[X,Y](m)=[X(m),Y(m)]_𝔤+\xi _YX(m)\xi _XY(m).$$
It is no accident that these examples exactly correspond to our previous list of Lie groupoids: as for groups, any Lie groupoid $`\mathrm{\Gamma }`$ has an associated Lie algebroid $`A(\mathrm{\Gamma })`$ with the same base space.<sup>34</sup><sup>34</sup>34The association $`\mathrm{\Gamma }A(\mathrm{\Gamma })`$ is functorial in an appropriate way, so that Mackenzie speaks of the Lie functor . Namely, as a vector bundle $`A(\mathrm{\Gamma })`$ is the restriction of $`\mathrm{ker}(t_{})`$ to $`\mathrm{\Gamma }_0`$, and the anchor is $`\alpha =s_{}`$. One may identify sections of $`A(\mathrm{\Gamma })`$ with left-invariant vector fields on $`\mathrm{\Gamma }`$, and under this identification the Lie bracket on $`C^{\mathrm{}}(\mathrm{\Gamma }_0,A(\mathrm{\Gamma }))`$ is by definition the commutator.
Conversely, one may ask whether a given Lie algebroid $`A`$ is integrable, in that it comes from a Lie groupoid $`\mathrm{\Gamma }`$ in the said way. That is, is $`AA(\mathrm{\Gamma })`$ for some Lie groupoid $`\mathrm{\Gamma }`$? This is not necessarily the case; see .
The modern interplay between Lie Lie groupoids and Lie algebroids on the ond hand, and symplectic geometry on the other is based on various amazing points of contact. The simplest of these is as follows.
###### Proposition 2
The dual vector bundle $`A^{}`$ of a Lie algebroid $`A`$ is canonically a Poisson manifold. The Poisson bracket on $`C^{\mathrm{}}(A^{})`$ is defined by the following special cases: $`\{f,g\}_\pm =0`$ for $`f,gC^{\mathrm{}}(M)`$; $`\{\stackrel{~}{\sigma },f\}=\alpha \sigma f`$, where $`\stackrel{~}{\sigma }C^{\mathrm{}}(A^{})`$ is defined by a section $`\sigma `$ of $`A`$ through the obvious pairing, and finally $`\{\stackrel{~}{\sigma }_1,\stackrel{~}{\sigma }_2\}=\stackrel{~}{[\sigma _1,\sigma _2]}`$.
Conversely, if a vector bundle $`EM`$ is a Poisson manifold such that the Poisson bracket of two linear functions is linear, then $`EA^{}`$ for some Lie algebroid $`A`$ over $`M`$, with the above Poisson structure.<sup>35</sup><sup>35</sup>35This establishes a categorical equivalence between linear Poisson structures on vector bundles and Lie algebroids. One can also show that in this situation the differential forms on $`A`$ form a differential graded algebra, while those on $`A^{}E`$ (or, equivalently, the so-called polyvector fields on $`A`$) are a Gerstenhaber algebra; see .
The main examples are:
* The dual $`𝔤^{}`$ of a Lie algebra $`𝔤`$ acquires its canonical Lie–Poisson structure (cf. ).
* A manifold $`M`$, seen as the dual to the zero-dimensional vector bundle $`MM`$, carries the zero Poisson structure.
* A cotangent bundle $`T^{}M`$ acquires the Poisson structure defined by its standard symlectic structure.
* The dual $`(T^{}P)/H`$ of a gauge algebroid inherits the canonical Poisson structure from $`T^{}P`$ under the isomorphism $`C^{\mathrm{}}(T^{}P)/H)C^{\mathrm{}}(T^{}P)^H`$.
* The dual $`𝔤^{}M`$ of an action algebroid acquires the so-called semidirect product Poisson structure .<sup>36</sup><sup>36</sup>36Relative to a basis of $`𝔤`$ with structure constants $`C_{ab}^c`$, this is given by $`\{f,g\}=C_{ab}^c\theta _c\frac{f}{\theta _a}\frac{g}{\theta _b}+\xi _af\frac{g}{\theta _a}\frac{f}{\theta _a}\xi _ag.`$
Combining the associations $`\mathrm{\Gamma }A(\mathrm{\Gamma })`$ and $`AA^{}`$, one has an association
$$\mathrm{\Gamma }A^{}(\mathrm{\Gamma }),$$
(15)
of a Poisson manifold to a Lie groupoid, which we call Weinstein’s map. As we shall see, this is a classical analogue of Connes’s map (8) in every possible respect.
## 7 Symplectic groupoids and the category of Poisson manifolds
Another important point of contact between Poisson manifolds and Lie algebroids that is relevant for what follows is the following construction.
###### Proposition 3
If $`P`$ is a Poisson manifold, then $`T^{}P`$ is canonically a Lie algebroid over $`P`$.
The anchor is just the usual map $`T^{}PTP`$, $`\alpha \alpha ^{\mathrm{}}`$ (e.g., $`dfX_f`$)<sup>37</sup><sup>37</sup>37The Hamiltonian vector field $`X_f`$ defined by a smooth function $`f`$ on a Poisson manifold $`P`$ is defined by $`X_fg=\{f,g\}`$. defined by the Poisson structure, whereas the Lie bracket is
$$[\alpha ,\beta ]=_\alpha ^{\mathrm{}}\beta _\beta ^{\mathrm{}}\alpha +d\pi (\alpha ,\beta ),$$
(16)
where $`\pi `$ is the Poisson tensor. Combining this with Proposition 2, one infers that $`TP`$ is a Poisson manifold whenever $`P`$ is.<sup>38</sup><sup>38</sup>38In addition, one may recover the Poisson cohomology of $`P`$ as the Lie algebroid cohomology of $`T^{}P`$ .
The following definition will play a key role for us in many ways.
###### Definition 5
A Poisson manifold $`P`$ is called integrable when the associated Lie algebroid $`T^{}P`$ is integrable (in being the Lie algebroid of some Lie groupoid).
If $`P`$ is an integrable Poisson manifold, a groupoid $`\mathrm{\Gamma }(P)`$ for which $`A(\mathrm{\Gamma }(P))T^{}P`$ (and hence $`\mathrm{\Gamma }(P)_0P`$) turns out to have the structure of a symplectic groupoid.
###### Definition 6
A symplectic groupoid is a Lie groupoid whose total space $`\mathrm{\Gamma }_1`$ is a symplectic manifold, such that the graph of $`\mathrm{\Gamma }_2\mathrm{\Gamma }\times \mathrm{\Gamma }`$ is a Lagrangian submanifold of $`\mathrm{\Gamma }\times \mathrm{\Gamma }\times \mathrm{\Gamma }^{}`$.
See also . Symplectic groupoids have many amazing properties, and in our opinion their introduction into symplectic geometry has been the biggest leap forward since the subject was founded.<sup>39</sup><sup>39</sup>39It would be tempting to say that a suitable analogue of a symplectic groupoid has not been found in noncommutative geometry so far, but in fact an analysis of the categorical significance of symplectic groupoids, Poisson manifolds, and operator algebras shows that the ‘quantum symplectic groupoid’ associated to a $`C^{}`$-algebra $`A`$ is just $`A`$ itself, whereas for a von Neumann algebra its standard form plays this role. For example:
1. There exists a unique Poisson structure on $`\mathrm{\Gamma }_0`$ such that $`t`$ is a Poisson map and $`s`$ is an anti-Poisson map.
2. $`\mathrm{\Gamma }_0`$ is a Lagrangian submanifold of $`\mathrm{\Gamma }_1`$.
3. The inversion in $`\mathrm{\Gamma }`$ is an anti-Poisson map.
4. The foliations of $`\mathrm{\Gamma }`$ defined by the levels of $`s`$ and $`t`$ are mutually symplectically orthogonal.
5. If $`\mathrm{\Gamma }`$ is s-connected,<sup>40</sup><sup>40</sup>40This means that each fiber $`s^1(u)`$ is connected, $`u\mathrm{\Gamma }_0`$. Similarly for s-simply connected. then $`s^{}C^{\mathrm{}}(\mathrm{\Gamma }_0)`$ and $`t^{}C^{\mathrm{}}(\mathrm{\Gamma }_0)`$ are each other’s Poisson commutant.
6. The symplectic leaves of $`\mathrm{\Gamma }_0`$ are the connected components of the $`\mathrm{\Gamma }_1`$-orbits.
With regard to the first point, the Poisson structure on $`\mathrm{\Gamma }(P)_0`$ induces the given one on $`P`$ under the diffeomorphism $`\mathrm{\Gamma }(P)_0P`$. For later use, we record:
###### Proposition 4
If a Poisson manifold $`P`$ is integrable, then there exists an s-connected and s-simply connected symplectic groupoid $`\mathrm{\Gamma }(P)`$ over $`P`$, which is unique up to isomorphism.
For example, suppose that $`\mathrm{\Delta }`$ is a Lie groupoid; is the Poisson manifold $`A^{}(\mathrm{\Delta })`$ it defines by (15) integrable? The answer is yes, and one may take
$$\mathrm{\Gamma }(A^{}(\mathrm{\Delta }))=T^{}\mathrm{\Delta },$$
(17)
the so-called cotangent groupoid of $`\mathrm{\Delta }`$ (see also ). This is s-connected and s-simply connected iff $`\mathrm{\Delta }`$ is.
Using the above constructions, we now define a category $`𝔓`$ of Poisson manifolds, which will play a central role in what follows. First, the objects of $`𝔓`$ are integrable Poisson manifolds; the integrability condition turns out to be necessary in order to have identities in $`𝔓`$; see below. In the spirit of general Morita theory , the arrows in $`𝔓`$ are bimodules in an appropriate sense. Bimodules for Poisson manifolds are known as dual pairs . A dual pair $`QSP`$ consists of a symplectic manifold $`S`$, Poisson manifolds $`Q`$ and $`P`$, and complete Poisson maps $`q:SQ`$ and $`p:SP^{}`$, such that $`\{q^{}f,p^{}g\}=0`$ for all $`fC^{\mathrm{}}(Q)`$ and $`gC^{\mathrm{}}(P)`$. To explain the precise class of dual pairs whose isomorphism classes form the arrows in $`𝔓`$, we need a symplectic analogue $`𝔖`$ of the category $`𝔊`$ (cf. Definition 2). In preparation, we call an action of a symplectic groupoid $`\mathrm{\Gamma }`$ on a symplectic manifold $`S`$ symplectic when the graph of the action in $`\mathrm{\Gamma }\times S\times S^{}`$ is Lagrangian .
###### Definition 7
The category $`𝔖`$ is the subcategory of the category $`𝔊`$ (of Lie groupoids and principal bibundles) whose objects are symplectic groupoids and whose arrows are isomorphism classes of principal bibundles for which the two groupoid actions are symplectic.
We call such bibundles symplectic. As we have seen (cf. Section 6), the base space of a symplectic groupoid is a Poisson manifold. Moreover, it can be shown that the base map $`S\mathrm{\Gamma }_0`$ of a symplectic action of a symplectic groupoid $`\mathrm{\Gamma }`$ on a symplectic manifold $`S`$ is a complete Poisson map such that for $`(\gamma ,y)\mathrm{\Gamma }\times _{\mathrm{\Gamma }_0}^{s,\rho }S`$ with $`\gamma =\phi _1^{t^{}f}(\rho (y))`$, one has $`\gamma y=\phi _1^{\rho ^{}f}(y)`$ (here $`\phi _t^g`$ is the Hamiltonian flow induced by a function $`g`$, and $`fC^{\mathrm{}}(\mathrm{\Gamma }_0)`$). Conversely, when $`\mathrm{\Gamma }`$ is s-connected and s-simply connected, a given complete Poisson map $`\rho :S\mathrm{\Gamma }_0`$ is the base map of a unique symplectic $`\mathrm{\Gamma }`$ action on $`S`$ with the above property . Furthermore, it is easy to show that the base maps of a symplectic bibundle form a dual pair. We call a dual pair arising from a symplectic principal bibundle in this way regular.
###### Definition 8
The objects of the category $`𝔓`$ of Poisson manifolds and dual pairs are integrable Poisson manifolds, and its arrows are isomorphism classes of regular dual pairs.
The identities in $`𝔓`$ are $`1_P=[P\mathrm{\Gamma }(P)P]`$, where $`\mathrm{\Gamma }(P)`$ is “the” s-connected and s-simply connected symplectic groupoid over $`P`$; cf. Proposition 4. As in every decent version of Morita theory, isomorphism of objects in $`𝔓`$ comes down to Morita equivalence of Poisson manifolds (in the sense of Xu ).
It is clear that $`𝔓`$ is equivalent to the full subcategory $`𝔖_c`$ of $`𝔖`$ whose objects are s-connected and s-simply connected symplectic groupoids; the advantage of working with $`𝔓`$ rather than $`𝔖_c`$ lies both in the greater intuitive appeal of Poisson manifolds and dual pairs over symplectic groupoids and symplectic principal bibundles, and also in the fact that the composition of arrows can be formulated in direct terms (i.e. avoiding arrow composition in $`𝔖`$ or $`𝔊`$) using a generalization of the familiar procedure of symplectic reduction .
For example, a strongly Hamiltonian group action $`GS`$ famously defines a dual pair
$$S/G\stackrel{\pi }{}S\stackrel{J}{}𝔤^{}$$
(where $`J`$ is the momentum map of the action) , whose product with the dual pair $`𝔤^{}0pt`$ in $`𝔓`$ equals $`S/GS//Gpt`$ (if we assume $`G`$ connected). In other words, the Marsden–Weinstein quotient $`S//G`$ may be interpreted in terms of the category $`𝔓`$ (see Section 11 below for the significance of this observation.)
## 8 The classical imprimitivity theorem
There is a complete classical analogue of Mackey’s theory of imprimitivity for (Lie) group actions . Firstly, the classical counterpart of a representation of a $`C^{}`$-algebra on a Hilbert space is a so-called realization of a Poisson manifold $`P`$ on a symplectic manifold $`S`$ ; this is a complete Poisson map $`S\stackrel{\rho }{}P`$.<sup>41</sup><sup>41</sup>41Some authors speak of a realization in case that $`\rho `$ is surjective, but not necessarily complete. The completeness of $`\rho `$ means that the Hamiltonian vector field $`X_{\rho ^{}f}`$ on $`S`$ has a complete flow for each $`fC_c^{\mathrm{}}(P)`$ (i.e. the flow is defined for all times). This condition turns out to be the classical counterpart of the requirement that $`\pi (a)^{}=\pi (a^{})`$ for representations of a $`C^{}`$-algebra. The analogy between completeness of the flow of a vector field and self-adjointness of an operator is even more powerful in the setting of unbounded operators; for example, the Laplacian on a Riemannian manifold $`M`$ is essentially self-adjoint on $`C_c^{\mathrm{}}(M)`$ when $`M`$ is geodesically complete . The appropriate symplectic notion of irreducibility is that
$$\{X_{\rho ^{}f}(x)fC^{\mathrm{}}(P)\}=T_xS$$
for all $`xS`$ (where $`X_g`$ is the Hamiltonian vector field of $`gC^{\mathrm{}}(S)`$); it is easy to show (cf. Thm. I.2.6.7 in ) that $`\rho `$ is irreducible iff $`S`$ is symplectomorphic to a covering space of a symplectic leaf of $`P`$ (and $`\rho `$ is the associated projection followed by injection). In particular, any Poisson manifold has at least one irreducible realization.<sup>42</sup><sup>42</sup>42The appropriate symplectic notion of faithfulness is simply that $`\rho `$ be surjective; it was recently shown by Crainic and Fernandes that a Poisson manifold admits a faithful realization iff it is integrable; cf. Definition 5. Along with their solution of this integrability problem , this is one of the deepest results in symplectic geometry to date.
Secondly, we provide the classical counterpart of Definition 1. It goes without saying that in the present context $`G`$ is a Lie group and $`M`$ a manifold, all actions being smooth by definition.
###### Definition 9
Given a $`G`$-action on $`M`$, a $`G`$-covariant realization of $`M`$ (seen as a Poisson manifold with zero Poisson bracket) is a complete Poisson map $`S\stackrel{\rho }{}M`$, where $`S`$ is a symplectic manifold equipped with a strongly Hamiltonian $`G`$-action,<sup>43</sup><sup>43</sup>43In the sense that the $`G`$-action has an equivariant momentum map $`J:S𝔤^{}`$ . and $`L_x(\rho ^{}f)=\rho ^{}L_x(f)`$ for all $`fC^{\mathrm{}}(M)`$.
The significance of this definition and its analogy to Definition 1 are quite obvious; instead of a representation $`\pi :C_0(M)B()`$ one now has a Lie algebra homomorphism $`\rho ^{}:C^{\mathrm{}}(M)C^{\mathrm{}}(S)`$. Its relationship to the material in the preceding section is as follows:
###### Proposition 5
When $`G`$ is connected, a $`G`$-covariant realization of $`M`$ may equivalently be defined as a realization $`S\stackrel{\sigma }{}𝔤^{}M`$ (equipped with the semidirect product Poisson structure) whose associated $`𝔤`$-action on $`S`$ is integrable (i.e. to a $`G`$-action on $`S`$).
The $`𝔤`$-action on $`S`$ in question is given by $`XX_{\sigma ^{}\stackrel{~}{X}}`$, where $`X𝔤`$ defines a linear function $`\stackrel{~}{X}:𝔤^{}`$ by evaluation (and consequently also defines a function on $`𝔤^{}\times M`$ that is constant on $`M`$, which we denote by the same symbol). Of course, given $`S\stackrel{\rho }{}M`$ as in Definition 9, one defines $`S\stackrel{\sigma }{}𝔤^{}M`$ by $`\sigma =(J,\rho )`$; the nontrivial part of the proposition lies in the completeness of $`\sigma `$, given the completeness of $`\rho `$.
One then has the following classical analogue of Mackey’s imprimitivity theorem.
###### Theorem 4
Up to symplectomorphism, there is a bijective correspondence between $`G`$-covariant realizations $`S\stackrel{\rho }{}G/H`$ of $`G/H`$ (with zero Poisson structure) and strongly Hamiltonian $`H`$-spaces $`S_\rho `$, as follows:
* Given $`S_\rho `$, the Marsden–Weinstein quotient (at zero) $`S^\rho =(T^{}G\times S_\rho )//H`$ is a $`G`$-covariant realization of $`G/H`$.<sup>44</sup><sup>44</sup>44The $`G`$-action inherited from the $`G`$-action on $`T^{}G`$ is given by pullback of left-multiplication, and the map $`S^\rho G/H`$ is inherited from the natural map $`T^{}GGG/H`$.
* Conversely, given $`S\stackrel{\rho }{}G/H`$ there exists a strongly Hamiltonian $`H`$-space $`S_\rho `$ such that $`SS^\rho `$.
This correspondence preserves irreducibility.
When $`G`$ is connected, this correspondence may be seen as being between realizations $`S\stackrel{\sigma }{}𝔤^{}(G/H)`$ whose associated $`𝔤`$-action on $`S`$ is integrable, and realizations $`S_\rho \stackrel{J_\rho }{}𝔥^{}`$ whose associated $`𝔥`$-action on $`S_\rho `$ is integrable.
The original proof of this theorem was lengthy and difficult . Fortunately, as in the quantum case, there exists a direct categorical argument, according to which at least the last part of Theorem 4 is a consequence of (9) as well. Namely, the following analogue of Theorem 3 holds:
###### Theorem 5
Weinstein’s map $`\mathrm{\Gamma }A^{}(\mathrm{\Gamma })`$ is functorial from $`𝔊_c`$ to $`𝔓`$.
Recall that $`𝔊_c`$ is the full subcategory of $`𝔊`$ whose objects are s-connected and s-simply connected Lie groupoids, and that the category $`𝔓`$ of Poisson manifolds and dual pairs has been defined in the previous section. For example, $`G\times (G/H)`$ is an object in $`𝔊_c`$ iff $`G`$ is connected and simply connected. Assume this to be the case for the moment. As already mentioned, the category $`𝔓`$ has a feature analogous to the category $``$ of $`C^{}`$-algebras, namely that two objects are isomorphic iff they are Morita equivalent Poisson manifolds in the sense of Xu . Consequently, similar to Corollary 1 one has:
###### Corollary 2
Weinstein’s map $`\mathrm{\Gamma }A^{}(\mathrm{\Gamma })`$ preserves Morita equivalence, in the sense that if $`\mathrm{\Gamma }`$ and $`\mathrm{\Delta }`$ are Morita equivalent s-connected and s-simply connected Lie groupoids, then $`A^{}(\mathrm{\Gamma })`$ and $`A^{}(\mathrm{\Delta })`$ are Morita equivalent Poisson manifolds in the sense of Xu.
Thus the Morita equivalence (9) of Lie groupoids implies the Morita equivalence
$$𝔤^{}(G/H)𝔥^{}$$
(18)
of Poisson manifolds. As for $`C^{}`$-algebras (and algebras in general), if two Poisson manifolds $`P_1,P_2`$ are Morita equivalent, then they have equivalent categories of realizations, and the equivalence bimodule implementing this Morita equivalence comes with an explicit procedure that defines a realization of $`P_2`$ given one of $`P_1`$, and vice versa. This procedure is a certain generalization of symplectic reduction (much as the corresponding Rieffel induction procedure for $`C^{}`$-algebras is a generalization of Mackey induction). In the case at hand, viz. (18), this precisely gives the prescription stated in Theorem 4, proving its last part at least for simply connected $`G`$. If $`G`$ fails to be simply connected, one passes to its universal cover $`\stackrel{~}{G}`$, and lets it act on $`G/H`$ via the projection $`\stackrel{~}{G}G`$. Hence $`G/H\stackrel{~}{G}/\widehat{H}`$ for some $`\widehat{H}\stackrel{~}{G}`$; Lie theory gives $`\stackrel{~}{𝔤}=𝔤`$ and $`\widehat{𝔥}=𝔥`$. The conclusion (18) still follows, this time as a consequence of $`\stackrel{~}{G}(\stackrel{~}{G}/\widehat{H})\widehat{H}`$ rather than of (9).
We state a rather satisfying classical analogue of Proposition 1, which is essentially a corollary to Theorem 4.
###### Proposition 6
The symplectic leaves of of the semidirect Poisson structure on $`𝔤^{}M`$ are classified by pairs $`(𝒪,𝒪^{})`$, where $`𝒪`$ is a $`G`$-orbit in $`M`$, and $`𝒪^{}`$ is a coadjoint orbit of the stabilizer of an arbitrary point in $`𝒪`$.
If we call the stabilizer in question $`H`$, the symplectic leaf $`L_{(𝒪,𝒪^{})}`$ corresponding to the pair $`(𝒪,𝒪^{})`$ is given by
$$L_{(𝒪,𝒪^{})}=\{(\theta ,q)𝔤^{}\times Q|q𝒪,(\mathrm{Co}(s(q)^1)\theta 𝔥^{})𝒪^{}\},$$
(19)
where $`s:𝒪G/HG`$ is an arbitrary section of the canonical principal $`H`$-bundle $`G`$ over $`G/H`$, and $`\mathrm{Co}`$ is the coadjoint action of $`G`$ on $`𝔤^{}`$.
Furthermore, one has a classical counterpart of (11), namely an isomorphism
$$𝔤^{}(G/H)(T^{}G)/H$$
(20)
of Poisson manifolds. This may be generalized from the principal $`H`$-bundle $`G`$ to arbitrary principal $`H`$-bundles $`P`$, provided that $`P`$ is connected and simply connected (this assumption was not necessary in the quantum case). In that case, we may apply Corollary 2 to find a Morita equivalence of Poisson manifolds
$$(T^{}P)/H𝔥^{}.$$
(21)
## 9 Deformation quantization
Largely due to the functoriality of Connes’s map (8) and its classical counterpart (15), we have observed a striking analogy between the $`C^{}`$-algebra $`C^{}(\mathrm{\Gamma })`$ and the Poisson manifold $`A^{}(\mathrm{\Gamma })`$ associated to a Lie groupoid $`\mathrm{\Gamma }`$. Beyond an analogy, the classical object $`A^{}(\mathrm{\Gamma })`$ turns out to be related to its quantum counterpart through deformation quantization in the $`C^{}`$-algebraic setting proposed by Rieffel:
###### Definition 10
A $`C^{}`$-algebraic deformation quantization of a Poisson manifold $`P`$ is a continuous field of $`C^{}`$-algebras $`(A,A_{\mathrm{}})_{\mathrm{}[0,1]}`$,<sup>45</sup><sup>45</sup>45Here $`A`$ is the $`C^{}`$-algebra of sections of the given field, which defines its continuity structure. A continuous field $`(A,A_x)_{xX}`$ of $`C^{}`$-algebras comes with surjective morphisms $`\pi _x:AA_x`$. where $`A_0=C_0(P)`$, with a Poisson algebra $`\stackrel{~}{A}_0`$ densely contained in $`C_0(P)`$ and a cross-section $`Q:\stackrel{~}{A}_0A`$ of $`\pi _0`$, such that, in terms of $`Q_{\mathrm{}}=\pi _{\mathrm{}}Q`$, for all $`f,g\stackrel{~}{A}_0`$ one has
$$\underset{\mathrm{}0}{lim}\frac{i}{\mathrm{}}[Q_{\mathrm{}}(f),Q_{\mathrm{}}(g)]Q_{\mathrm{}}(\{f,g\})_{\mathrm{}}=0.$$
(22)
This has turned out to be an fruitful definition of quantization (cf. ). In many interesting examples the fiber algebras are non-isomorphic even away from $`\mathrm{}=0`$ (cf. and Footnote 53 below), but in the case at hand the situation is simpler.<sup>46</sup><sup>46</sup>46Technically, the field in Theorem 6 is said to be trivial away from $`\mathrm{}=0`$, in the sense that $`A_{\mathrm{}}=B`$ for all $`\mathrm{}(0,1]`$ and one has a short exact sequence $`0CBAA_00`$ (where $`CB=C_0((0,1],B)`$ is the cone of $`B`$).
###### Theorem 6
<sup>47</sup><sup>47</sup>47See also for a version of this result in the setting of formal deformation quantization (i.e. star products), and also cf. . For any Lie groupoid $`\mathrm{\Gamma }`$, the field $`A_0=C_0(A^{}(\mathrm{\Gamma }))`$, $`A_{\mathrm{}}=C^{}(\mathrm{\Gamma })`$ for $`\mathrm{}0`$, and $`A=C^{}(\mathrm{\Gamma }^T)`$, the $`C^{}`$-algebra of the tangent groupoid $`\mathrm{\Gamma }^T`$ of $`\mathrm{\Gamma }`$,<sup>48</sup><sup>48</sup>48Following Connes’s definition of the special case of the pair groupoid $`\mathrm{\Gamma }=M\times M`$ around 1980 (see ), the tangent groupoid (or adiabatic groupoid) of an arbitrary Lie groupoid was independently defined in . See also . defines a $`C^{}`$-algebraic deformation quantization of $`A^{}(\mathrm{\Gamma })`$.<sup>49</sup><sup>49</sup>49The same statement holds for the corresponding reduced groupoid $`C^{}`$-algebras.
We refer to the literature cited for the specification of $`\stackrel{~}{A}_0`$, as well as for the proof of (22). The proof of the remainder of the theorem actually covers a much more general situation, as follows .<sup>50</sup><sup>50</sup>50This setting was originally suggested by Skandalis.
###### Definition 11
A field of Lie groupoids is a triple $`(𝖦,X,p)`$, with $`𝖦`$ a Lie groupoid, $`X`$ a manifold, and $`p:𝖦X`$ a surjective submersion such that $`p=p_0r=p_0s`$, where $`p_0=p𝖦_0`$.
It follows that each $`𝖦_x=p^1(x)`$ is a Lie subgroupoid of $`𝖦`$ over $`𝖦_0p^1(x)`$, so that $`𝖦=_{xX}𝖦_x`$ as a groupoid. One may then form the convolution $`C^{}`$-algebras $`C^{}(𝖦)`$ and $`C^{}(𝖦_x)`$. Each $`aC_c(𝖦)`$ (or $`C_c^{\mathrm{}}(𝖦)`$) defines $`a_x=a𝖦_x`$ as an element of $`C_c(𝖦_x)`$ (etc.). These maps $`C_c(𝖦)C_c(𝖦_x)`$ are continuous in the appropriate norms, and extend to maps $`\pi _x:C^{}(𝖦)C^{}(𝖦_x)`$. Hence one obtains a field of $`C^{}`$-algebras
$$(A=C^{}(𝖦),A_x=C^{}(𝖦_x))_{xX}$$
(23)
over $`X`$, where $`aC^{}(𝖦)`$ defines the section $`x\pi _x(a)`$.<sup>51</sup><sup>51</sup>51A similar statement applies to the corresponding reduced $`C^{}`$-algebras. The question now arises when this field is continuous.
###### Lemma 1
The field (23) is continuous at all points where $`𝖦_x`$ is amenable .<sup>52</sup><sup>52</sup>52And similarly for the case of reduced $`C^{}`$-algebras.
For example, the tangent groupoid $`\mathrm{\Gamma }^T`$ of a given Lie groupoid $`\mathrm{\Gamma }`$ forms a field of Lie groupoids over $`[0,1]`$, with $`\mathrm{\Gamma }_0^T=A(\mathrm{\Gamma })`$ (seen as a Lie groupoid instead of a Lie algebroid in the way every vector bundle $`E\stackrel{\pi }{}M`$ defines a Lie groupoid over its base space, namely by $`s=t=\pi `$ and fiberwise addition) and $`\mathrm{\Gamma }_{\mathrm{}}^T=\mathrm{\Gamma }`$ for $`\mathrm{}(0,1]`$. This eventually implies Theorem 6 (except for (22)); the same strategy also leads to far-reaching generalizations thereof.<sup>53</sup><sup>53</sup>53Lemma 1 applies much more generally to fields of locally compact groupoids. In the context of $`C^{}`$-algebraic deformation quantization, there are two typical situations. In the smooth (Lie) case studied in this paper, all $`𝖦_{\mathrm{}}`$ are the same for $`\mathrm{}0`$ but possibly not amenable, whereas $`𝖦_0`$ is amenable. The former property then yields continuity at $`\mathrm{}=0`$ by the lemma, whereas the latter gives continuity on $`(0,1]`$. In the context of Definition 10, the reason why $`G_0`$ is amenable is that $`A_0`$ must be commutative, which implies that $`G_0`$ is a bundle of abelian groups. But such groupoids are always amenable . In the étale case all $`𝖦_{\mathrm{}}`$ are typically different from each other, but they are all amenable. See for a description of noncommutative tori and the noncommutative four-spheres of Connes and Landi (and of many other examples) as deformation quantizations along these lines.
In physics, Theorem 6 describes the quantization of particles with both internal and spatial degrees of freedom in a very wide setting. In noncommutative geometry, certain constructions of Connes in index theory turn out to be special cases of Theorem 6.<sup>54</sup><sup>54</sup>54One instance is the map $`p!:K^{}(F^{})K_{}(C^{}(V,F))`$ on p. 127 of , which plays a key role in the definition of the analytic assembly map for foliated manifolds. This is the K-theory map induced by the continuous field of Theorem 6, where $`\mathrm{\Gamma }`$ is the holonomy groupoid of the foliation. The index groupoid for a vector bundle map $`L:EF`$ defined in \[12, §II.6\] is another example. Here one has a Lie groupoid $`\mathrm{\Gamma }=\mathrm{Ind}_L=F_LE`$ over $`F`$, whose Lie algebroid is $`F\times _BE`$. This is a vector bundle over $`B`$, and in the above formalism it should be regarded as a groupoid over $`F`$ under addition in each fiber. Hence $`A_0=C^{}(F\times _BE)C_0(F\times E^{})`$. The corresponding K-theory map occurs in Connes’s construction of the Gysin map $`f_!:K^{}(X)K^{}(Y)`$ induced by a smooth map $`f:XY`$ between manifolds. As to the ideology of noncommutative geometry, the theorem shows that the two fundamental classes of noncommutative manifolds, namely the ones defined by a singular quotient and the ones defined by deformation , overlap. For in case that the equivalence relation defining the quotient in question can be codified by a Lie groupoid $`\mathrm{\Gamma }`$, the noncommutative space $`C^{}(\mathrm{\Gamma })`$ associated with the quotient space is at the same time a deformation of the dual of its Lie algebroid.
Furthermore, Connes’s philosophy in dealing with singular quotients, and especially his description of the Baum–Connes conjecture in Ch. II of , actually suggests a procedure for the quantization of such spaces. We explain this in a simple example . Suppose a Lie group $`G`$ acts on a manifold $`M`$; it acts on $`T^{}M`$ by pull-back, and we happen to be interested in quantizing the quotient $`(T^{}M)/G`$. In case that the $`G`$-action is free and proper the situation is completely understood: the quotient is a Poisson manifold of the type $`A^{}(\mathrm{\Gamma })`$ for $`\mathrm{\Gamma }=M\times _GM`$, to which Theorem 6 applies (see also for a detailed study of this case). However, if the $`G`$-action is not free (but still assumed to be proper), the quotient $`(T^{}M)/G`$ may fail to be a manifold, let alone a Poisson manifold. According to Connes, one should replace the space $`(T^{}M)/G`$ by the groupoid $`T^{}MG`$, and regard the associated noncommutative space $`C^{}(T^{}MG)`$ as a classical space. If the $`G`$-action is free, one has a Morita equivalence of Lie groupoids
$$T^{}MG(T^{}M)/G$$
(24)
which by Corollary 1 implies a Morita equivalence
$$C^{}(T^{}MG)C^{}((T^{}M)/G)$$
(25)
of $`C^{}`$-algebras.<sup>55</sup><sup>55</sup>55See for the original, non-groupoid proof of (25). In general, we propose to quantize the singular space $`(T^{}M)/G`$ by deforming $`C^{}(T^{}MG)`$, which may be done by the field of Lie groupoids defined by the tangent groupoid $`\mathrm{\Gamma }^T`$ of $`\mathrm{\Gamma }=(M\times M)G`$. This field has fibers $`\mathrm{\Gamma }_0^T=TMG`$ (where $`TM`$ is seen as a Lie groupoid, as explained above), and $`\mathrm{\Gamma }_{\mathrm{}}^T=(M\times M)G`$. By Lemma 1 (which applies because $`TMG`$ is amenable; see Lemma 2 in ), this field of groupoids leads to a continuous field of $`C^{}`$-algebras with $`A=C^{}(\mathrm{\Gamma }^T)`$, etc., in the familiar way. The fibers of the latter field are simply $`A_0=C_0(T^{}M)G`$ and $`A_{\mathrm{}}=K(L^2(M))G`$ for all $`\mathrm{}(0,1]`$. To what extent this reflects physical desiderata remains to be seen.
## 10 Functorial quantization
The final application of groupoids to physics and noncommutative geometry we wish to describe in this paper is a functorial approach to quantization. In our opinion this forms the natural outcome of the categorical approach to Mackey’s imprimitivity theorem described above. Beyond the desire to complete Mackey’s program, why should one wish to turn quantization into a functor? Historically, quantum mechanics started with Heisenberg’s paper Über die quantentheoretische Umdeutung kinematischer und mechanischer Beziehungen<sup>56</sup><sup>56</sup>56On the quantum-theoretical reinterpretation of kinematical and mechanical relations.. One might argue that the proper mathematical reading of Heisenberg’s idea of Umdeutung (reinterpretation) is that the transition from classical to quantum mechanics should be given by a functor. Indeed, attempts to make quantization functorial date back at least to van Hove’s famous paper from 1951 (see also ), the general conclusion being that functorial quantization is impossible (see and refs. therein). However, all no-go theorems in this direction start from wrong and naive categories, both on the classical and on the quantum side.
Instead, though we have to warn the reader that we are presenting a program rather than a theorem here, it seems possible to interpret quantization as a functor $`𝒬`$ from either the category $`𝔖`$ (cf. Definition 2), or, more straightforwardly, from the category $`𝔓`$ (see Definition 8; recall that $`𝔓`$ is equivalent to a full subcategory of $`𝔖`$) to the category $`𝔎𝔎`$ defined by Kasparov’s bivariant K-theory (see ).<sup>57</sup><sup>57</sup>57The objects of $`𝔎𝔎`$ are separable $`C^{}`$-algebras, and the arrows are $`\mathrm{Hom}_{𝔎𝔎}(A,B)=KK(A,B)`$, composed with Kasparov’s product $`KK(A,B)\times KK(B,C)KK(A,C)`$. This was first proposed in . Beyond the defining property of making quantization functorial, this program would:
* Unify deformation quantization and geometric quantization into a single operation (the former becoming the object side of the quantization functor and the latter the arrow side);
* Imply the functoriality of shriek maps in K-theory , in particular providing a natural home for Connes-style proofs and generalizations of index theorems ;
* Imply the “quantization commutes with reduction” conjecture of Guillemin and Sternberg ;
* Provide unlimited generalizations of this conjecture, e.g., to noncompact Lie groups and Lie groupoids (see for the former).
It should be clear that the use of groupoids is essential in this program, since the classical category $`𝔖`$ of symplectic groupoids and principal symplectic bibundles either forms the domain of the quantization functor $`𝒬`$, or, in case one more naturally starts from $`𝔓`$, plays an essential role in the definition of the latter category.
Let us indeed construe quantization as a functor $`𝒬:𝔓KK`$. This means that quantization sends (isomorphism classes of) dual pairs into (homotopy classes of) Kasparov bimodules. More precisely, if Poisson manifolds $`P_1`$ and $`P_2`$ are quantized by (separable) $`C^{}`$-algebras $`𝒬(P_1)`$ and $`𝒬(P_2)`$, respectively, a dual pair $`P_1MP_2`$ should be quantized by an element
$$𝒬(P_1MP_2)KK(𝒬(P_1),𝒬(P_2)),$$
(26)
where $`KK(,)`$ is the usual Kasparov group . Roughly speaking, the construction of $`𝒬(P)`$ should be done by some $`C^{}`$-algebraic version of deformation quantization, whereas that of $`𝒬(P_1MP_2)`$ should come from a far-reaching generalization of geometric quantization first proposed, in special cases, by Raoul Bott;<sup>58</sup><sup>58</sup>58This was done in seminars and conversations; no paper by Bott containing his proposal seems to exist. (V. Guillemin and R. Sjamaar, private communications.) see . This proposal turns out to be closely related to Connes’s construction of shriek maps .
To explain the construction of (26), we assume that the symplectic manifold $`(M,\omega )`$ is prequantizable. Cf. for details of the following approach to geometric quantization. One picks an almost complex structure $`𝒥`$ on $`M`$ that is compatible with $`\omega `$ (in that $`\omega (,𝒥)`$ is positive definite and symmetric). This $`𝒥`$ canonically induces a $`\mathrm{Spin}^c`$ structure on $`TM`$, which should subsequently be twisted by a prequantization line bundle $`L`$ line bundle over $`M`$ to obtain a $`\mathrm{Spin}^c`$ structure $`(P,)`$ on $`M`$.<sup>59</sup><sup>59</sup>59We here define a $`\mathrm{Spin}^c`$ structure on $`M`$ as an equivalence class of principal $`\mathrm{Spin}^c(n)`$-bundle $`P`$ over $`M`$ with an isomorphism $`P\times _\pi ^nTM`$ of vector bundles. Here $`n=dim(M)`$ and the bundle on the left-hand side is the bundle associated to $`P`$ by the defining representation of $`SO(n)`$. Connes’s construction of shriek maps lacks the twisting with the prequantization line bundle. Denote the (complex) spin representation of $`\mathrm{Spin}^c(n)`$ on the finite-dimensional Hilbert space $`S_n`$ by $`\mathrm{\Delta }_n`$. One may then form the associated spinor bundle $`𝒮_n=P\times _{\mathrm{\Delta }_n}S_n`$, with Dirac operator $`D\text{/}:C^{\mathrm{}}(M,𝒮_n)C^{\mathrm{}}(M,𝒮_n)`$. For even $`n`$ (the case that applies here, as $`M`$ is symplectic) the spin representation decomposes into two irreducibles $`\mathrm{\Delta }_n=\mathrm{\Delta }_n^+\mathrm{\Delta }_n^{}`$ on $`S_n=S_n^+S_n^{}`$, so that also the vector bundle $`𝒮_n`$ decomposes accordingly as $`𝒮_n=𝒮_n^+𝒮_n^{}`$. Being odd with respect to this decomposition, the Dirac operator then splits accordingly as $`D\text{/}^\pm =C^{\mathrm{}}(M,𝒮^\pm )C^{\mathrm{}}(M,𝒮^{})`$.
Given a dual pair $`P_1MP_2`$, the fundamental idea is to use the map $`MP_2`$ to turn the appropriate completion of $`C_c^{\mathrm{}}(M,𝒮_n)`$ to a graded Hilbert $`C^{}(𝒬(P_2))`$ module $``$, and subsequently, to use the map $`P_1M`$ to construct an action of $`C^{}(𝒬(P_1))`$ on $``$, producing a $`C^{}(𝒬(P_1))`$-$`C^{}(𝒬(P_2))`$ graded Hilbert bimodule. The final step is to employ the Dirac operator $`D\text{/}`$ to enrich this bimodule into a Kasparov cycle, whose homotopy class defines the element (26) we are after.
This procedure has so far been carried through in a few cases only, namely those in which Theorem 6 states how the Poisson manifolds $`P_i`$ are to be quantized, and in which simultaneously techniques from the literature on the Baum–Connes conjecture are available to construct (26) according to the procedure just sketched. The simplest case is $`P_1=P_2=pt`$ (i.e. a point) and $`M`$ an arbitrary compact prequantizable symplectic manifold.<sup>60</sup><sup>60</sup>60Let us note that the associated dual pair $`ptMpt`$ does not define an element of our category $`𝔓`$, but this nuisance does not stop us from proceeding. Most people would agree that $`𝒬(pt)=`$, and under the isomorphism $`KK(,)`$ the Kasparov cycle defined by $`D\text{/}`$ is just the Fredholm index of $`D\text{/}_+`$ . This number, then, is Bott’s quantization of $`(M,\omega )`$. Consequently, we have
$$𝒬(ptMpt)=\mathrm{Index}(D\text{/}_+).$$
(27)
## 11 Quantization commutes with reduction
The above definition of quantization gains in substance when one passes to a dual pair $`M/GM𝔤^{}`$ defined by a strongly Hamiltonian group action $`GM`$ in the usual way . For simplicity, we will actually use the dual pair $`ptM𝔤^{}`$.<sup>61</sup><sup>61</sup>61This dual pair does not define an element of $`𝔓`$, but this does not affect any of our arguments. Theorem 6 tells us that $`𝒬(𝔤^{})=C^{}(G)`$, where $`G`$ is any Lie group with Lie algebra $`G`$; we take the connected and simply connected one.<sup>62</sup><sup>62</sup>62Here the use of the category $`𝔖`$ as the domain of the quantization functor $`𝒬`$ is more satisfactory. The classical data is then formed by the $`G`$-action on $`M`$ itself (in the guise of the associated symplectic action of the symplectic groupoid $`T^{}G`$), instead of the associated momentum map $`M𝔤^{}`$. This refinement is, of course, essential when $`G`$ is discrete. Hence the quantization of the dual pair $`ptM𝔤^{}`$ should be an element of the Kasparov group $`KK(,C^{}(G))`$.
This element can be defined when the $`G`$-action is proper and cocompact (i.e. $`M/G`$ is compact), and lifts to an action on the principal bundle $`P`$ defining the $`\mathrm{Spin}^c`$ structure. Namely, in that case one regards $`D\text{/}`$ as an operator on the graded Hilbert space $`L^2(M,𝒮_n)`$ of $`L^2`$-sections of $`𝒮_n`$, which at the same time carries a natural representation $`\pi `$ of $`C_0(M)`$ by multiplication operators, as well as a natural unitary representation $`U(G)`$. Provided that in addition the Dirac operator $`D\text{/}`$ is almost $`G`$-invariant in the sense that $`[U(x),D\text{/}]`$ is bounded for each $`xG`$, these data specify an element $`[L^2(M,𝒮_n),\pi (C_0(M)),U(G),D\text{/}]`$ of the equivariant analytic K-homology group $`K_0^G(M)=KK^G(C_0(M),)`$ . Here we suppress the grading of the Hilbert space in question in our notation. Let
$$\mathrm{Index}_G:K_0^G(M)K_0(C^{}(G))$$
be the analytic assembly map as defined by Baum, Connes, and Higson , seen however as a map taking values in $`K_0(C^{}(G))`$ instead of $`K_0(C_r^{}(G))`$ (cf. for this point). For simplicity we write
$$\mathrm{Index}_G(D\text{/}_+)\mathrm{Index}_G([L^2(M,𝒮_n),\pi (C_0(M)),U(G),D\text{/}]).$$
(28)
We then define the quantization of the dual pair $`ptM𝔤^{}`$ as
$$𝒬(ptM𝔤^{})=\alpha _{C^{}(G)}(\mathrm{Index}_G(D\text{/}_+)),$$
(29)
where $`\alpha _A:K_0(A)KK(,A)`$ is the natural isomorphism one has for any separable $`C^{}`$-algebra $`A`$ . As required, (29) defines an element of
$$KK(𝒬(pt),𝒬(𝔤^{}))=KK(,C^{}(G)).$$
For a much simpler example, whose significance will become clear shortly, consider the dual pair $`𝔤_{}^{}0pt`$, where $`0`$ (seen as a coadjoint orbit of $`G`$) is the zero element of the vector space $`𝔤^{}`$, equipped with minus the Lie–Poisson structure. Its quantization should be an element of the Kasparov representation ring $`KK(C^{}(G),)`$, which we simply take to be the graded Hilbert space $`=0`$ carrying the trivial representation of $`G`$, with $`F=0`$. We denote this element by $`[,0,0]`$, so that
$$𝒬(𝔤^{}0pt)=[,0,0].$$
(30)
Let
$$\tau _{}:KK(,C^{}(G))KK(,)$$
be the map functorially induced by the morphism $`\tau :C^{}(G)`$ given by the trivial representation of $`G`$.<sup>63</sup><sup>63</sup>63For $`fC_c(G)`$ one simple has $`\tau (f)=_G𝑑xf(x)`$. This is the reason why we use $`C^{}(G)`$ rather than $`C_r^{}(G)`$, as is customary in the Baum–Connes conjecture: for $`\tau `$ is not continuous on $`C_r^{}(G)`$ (unless $`G`$ is amenable). Functoriality of the Kasparov product
$$KK(,C^{}(G))\times KK(C^{}(G)),)KK(,)\stackrel{}{}$$
then yields
$$y\times [,0,0]=\tau _{}(y)$$
(31)
for any $`yKK(,C^{}(G))`$. In particular, (29) and (30) give
$$𝒬(ptM𝔤^{})\times 𝒬(𝔤_{}^{}0pt)=\tau _{}(\mathrm{Index}_G(D\text{/}_+^M));$$
(32)
to avoid confusion later on, we have added a suffix $`M`$ to the pertinent Dirac operator.
On the classical side, in the category $`𝔓`$ we compute
$$(ptM𝔤^{})(𝔤_{}^{}0pt)=ptM//Gpt,$$
(33)
where $`M//G`$ is the Marsden–Weinstein quotient. Assuming that $`M//G`$ is prequantizable (this is a theorem in the compact case ), we have already seen from (27) that
$$𝒬(ptM//Gpt)=\mathrm{Index}(D\text{/}_+^{M//G}),$$
(34)
where we have denoted the appropriate Dirac operator on $`M//G`$ by $`D\text{/}^{M//G}`$.
Functoriality of quantization would imply
$$𝒬(ptM𝔤^{})\times 𝒬(𝔤_{}^{}0pt)=𝒬((ptM𝔤^{})(𝔤_{}^{}0pt)).$$
(35)
Using (32) and (33), this amounts to
$$\tau _{}(\mathrm{Index}_G(D\text{/}_+^M))=\mathrm{Index}(D\text{/}_+^{M//G}).$$
(36)
For $`G`$ and $`M`$ compact, this is precisely the so-called Guillemin–Sternberg conjecture that “quantization commutes with reduction” in its modern form .<sup>64</sup><sup>64</sup>64This conjecture is, in fact, a theorem , but the name “conjecture” is still generally used. To see this, note that for $`M`$ compact the Dirac operator $`D\text{/}_+`$ is Fredholm, whereas for $`G`$ compact one has $`K_0(C^{}(G))R(G)`$, the representation ring of $`G`$. Consequently, $`\mathrm{Index}_G(D\text{/}_+)`$ defines an element of $`R(G)`$, and the map $`\tau _{}:R(G)R(e)`$ is just $`[V][W]dim(V_0)dim(W_0)`$, where $`V_0V`$ is the space of $`G`$-invariant vectors, etc.
For $`G`$ countable (acting properly and cocompactly on $`M`$, as stated before), (36) boils down to the naturality of the Baum–Connes assembly map for countable discrete groups . Combining this fact with the validity of (36) for compact $`G`$ and $`M`$, it can be shown that (36) holds for any strongly Hamiltonian proper cocompact action of $`G`$ on a possibly noncompact symplectic manifold, provided that $`G`$ contains a discrete normal subgroup $`\mathrm{\Gamma }`$ with $`G/\mathrm{\Gamma }`$ compact .
Let us close this paper in the right groupoid spirit by pointing out that all arguments in this section should be carried out for Lie groupoids instead of Lie groups. For example, the pertinent symplectic reduction procedure (generalizing Marsden–Weinstein reduction) was first studied in , and can be reinterpreted in terms of the product in the category $`𝔓`$ just as in the group case. A very interesting special case comes from foliation theory, as follows (cf. ). Let $`(V_i,F_i)`$, $`i=1,2`$, be foliations with associated holonomy groupoids $`G(V_i,F_i)`$ (assumed to be Hausdorff for simplicity). A smooth generalized map $`f`$ between the leaf spaces $`V_1/F_1`$ and $`V_2/F_2`$ is defined as a principal bibundle $`M_f`$ between the Lie groupoids $`G(V_1,F_1)`$ and $`G(V_2,F_2)`$. Classically, such a bibundle defines a dual pair $`T^{}F_1T^{}M_fT^{}F_2`$ . Here $`TF_iTV_i`$ is the tangent bundle to the foliation $`(V_i,F_i)`$, whose dual bundle $`T^{}F_i`$ has a canonical Poisson structure.<sup>65</sup><sup>65</sup>65The best way to see this is to interpret $`TF_i`$ as the Lie algebroid of $`G(V_i,F_i)`$. Quantum mechanically, $`f`$ defines an element
$$f_!KK(C^{}(G(V_1,F_1)),C^{}(G(V_2,F_2))).$$
In the functorial approach to quantization, $`f_!`$ is interpreted as the quantization of the dual pair $`T^{}F_1T^{}M_fT^{}F_2`$. The functoriality of quantization among dual pairs of the same type should then follow from the computations in on the quantum side and on the classical side. The construction and functoriality of shriek maps in is a special case of this, in which the $`V_i`$ are both trivially foliated.
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# Anisotropic Hubbard model on a triangular lattice — spin dynamics in HoMnO₃
## Abstract
The recent neutron-scattering data for spin-wave dispersion in $`\mathrm{HoMnO}_3`$ are well described by an anisotropic Hubbard model on a triangular lattice with a planar (XY) spin anisotropy. Best fit indicates that magnetic excitations in $`\mathrm{HoMnO}_3`$ correspond to the strong-coupling limit $`U/t>15`$, with planar exchange energy $`J=4t^2/U2.5`$meV and planar anisotropy $`\mathrm{\Delta }U0.35`$meV.
There has been renewed interest in correlated electron systems on triangular lattices, as evidenced by recent studies of antiferromagnetism, superconductivity and metal-insulator transition in the organic systems $`\kappa (\mathrm{BEDT}\mathrm{TTF})_2\mathrm{X}`$,review1 ; review2 the discovery of superconductivity in $`\mathrm{Na}_\mathrm{x}\mathrm{CoO}_2.\mathrm{yH}_2\mathrm{O}`$,watersup the observation of low-temperature insulating phases in some $`\sqrt{3}`$-adlayer structures such as K on Si,weitering and quasi two-dimensional $`120^0`$ spin ordering and spin-wave excitations in $`\mathrm{RbFe}(\mathrm{MoO}_4)_2`$ (Refs. 5,6) and the multiferroic materials $`\mathrm{YMnO}_3`$ and $`\mathrm{HoMnO}_3`$.sato ; holmium
Recent neutron-scattering studies of the multiferroic material $`\mathrm{HoMnO}_3`$ have revealed a non-collinear $`120^0`$ antiferromagnetic (AF) ordering below $`T_\mathrm{N}72`$ K of the $`S=2`$ Mn<sup>3+</sup> spins arranged in offset layers of two-dimensional (2D) triangular lattice.holmium Measurements of the spin wave dispersion were found to be well described by a nearest-neighbour Heisenberg AF with exchange energy $`J=2.44`$ meV and a planar anisotropy $`D=0.38`$ meV at 20 K. No discernible dispersion was observed in the out-of-plane direction, indicating primarily 2D spin dynamics.
Recently spin-wave excitations in the $`120^0`$ AF state of the Hubbard model on a triangular lattice were studied within the random phase approximation (RPA) in the full $`U`$ range.tri The spin wave energy in the large $`U`$ limit was shown to asymptotically approach the corresponding result for the Quantum Heisenberg antiferromagnet (QHAF), thus providing a continuous interpolation between weak and strong coupling limits. However, competing interactions and frustration were found to significantly modify the dispersion at finite $`U`$, resulting in vanishing spin stiffness at $`U6`$ and a magnetic instability at $`U7`$ corresponding to vanishing spin-wave energy at wave vector $`𝐪_M=(\pi /3,\pi /\sqrt{3})`$. The sharp fall-off of $`\omega _M`$ near $`U7`$ provides a sensitive indicator of finite-$`U`$ effects in the AF state. Indeed, recent high-resolution neutron-scattering studies of the spin-wave dispersion in the square-lattice S=1/2 AF $`\mathrm{La}_2\mathrm{CuO}_4`$ have revealed noticeable spin-wave dispersion along the MBZ edge,spinwave associated with finite-$`U`$ double-occupancy effects.spectrum
In this brief report we extend the spin-wave analysis to include planar spin anisotropy, and show that the neutron-scattering data for $`\mathrm{HoMnO}_3`$ are also well described by a Hubbard model on a triangular lattice, thus providing a microscopic description of the most essential features of spin dynamics in the $`120^0`$ AF state of $`\mathrm{HoMnO}_3`$, including the spin-wave dispersion and energy scale. We examine the behaviour of spin-wave anisotropy gap with spin anisotropy and also suggest a sensitive measure of finite-$`U`$, double occupancy effects.
However, the Mn spin planar anisotropy is treated here only at a phenomenological level, equivalent to the effective anisotropy $`DS_{iz}^2`$ included in recent investigations using spin models.holmium ; sato In a detailed study of single-ion anisotropy and crystal-field effects in the layered rare-earth cuprates $`\mathrm{R}_2\mathrm{CuO}_4`$ (R=Nd,Pr,Sm), the magnetic behaviour (including spin-reorientation transitions) has been attributed to coupling of Cu with the rare-earth magnetic subsystem which exhibits a large single-ion anisotropy resulting in preferential ordering of rare-earth moments along specific lattice directions.sachi It has been suggested that the anisotropy of Mn spins, its observed temperature dependence, and the reorientation transitions in $`\mathrm{HoMnO}_3`$ also originate from a similar anisotropic exchange coupling with the rare-earth Holmium,holmium resulting in magnetic behaviour as seen in layered rare-earth cuprates, where frustrated interlayer coupling allows for weaker, higher-order interactions to control the magnetic structure. Especially relevant for non-collinear ordering, the anisotropic exchange (Dzyaloshinski-Moriya) interaction $`𝐃.(𝐒_i\times 𝐒_j)`$ originating from spin-orbit coupling has been suggested as responsible for the clamping of ferroelectric and antiferromagnetic order parameters in $`\mathrm{YMnO}_3`$.hanamura
Hund’s rule coupling responsible for the $`S=2`$ spin state of Mn<sup>+++</sup> ions and crystal-field splitting have also not been realistically incorporated here. However, these realistic details do not qualitatively affect the spin-rotation symmetry and spin dynamics, as discussed below. Hund’s rule coupling in the generalized Hubbard model considered here is maximal as inter-orbital Coulomb interaction for parallel spins is dropped completely. Including an inter-orbital density-density interaction $`V_0`$ and an intra-atomic exchange interaction $`F_0`$ favouring parallel-spin alignment (Hund’s rule coupling), the more realistic orbital Hubbard modelroth ; held
$`H`$ $`=`$ $`t{\displaystyle \underset{i,\delta ,\gamma ,\sigma }{}}a_{i\gamma \sigma }^{}a_{i+\delta ,\gamma \sigma }+U{\displaystyle \underset{i\gamma }{}}n_{i\gamma }n_{i\gamma }`$ (1)
$`+`$ $`{\displaystyle \underset{i,\gamma <\gamma ^{},\sigma ,\sigma ^{}}{}}(V_0\delta _{\sigma \sigma ^{}}F_0)n_{i\gamma \sigma }n_{i\gamma ^{}\sigma ^{}}`$
$`+`$ $`F_0{\displaystyle \underset{i,\gamma <\gamma ^{},\sigma \sigma ^{}}{}}a_{i\gamma \sigma }^{}a_{i\gamma ^{}\sigma }a_{i\gamma ^{}\sigma ^{}}^{}a_{i\gamma \sigma ^{}}`$
remains spin-rotationally invariant under a global rotation of the fermion spin $`𝐒_{i\gamma }=\mathrm{\Psi }_{i\gamma }^{}\frac{𝝈}{2}\mathrm{\Psi }_{i\gamma }`$ (where $`\mathrm{\Psi }_{i\gamma }^{}(a_{i\gamma }^{}a_{i\gamma }^{})`$ is the fermion field operator), even if orbitals $`\gamma `$ are identified with the Mn orbitals $`(t_{2g},e_g)`$ resulting from crystal-field splitting of the atomic 3d orbitals. As the intra-atomic exchange interaction $`F_0`$ is much larger than the spin excitation energy scale ($`\frac{t^2}{U+F_0}`$, within a strong-coupling expansion), all fermion spins on a site are effectively coupled, yielding a composite quantum spin $`S`$ and a corresponding multiplication by factor $`2S`$ to obtain the effective spin-wave energy scale. Therefore, orbital multiplicity does not change the Goldstone-mode structure, and the spin-dynamics energy scale in the orbital Hubbard model is essentially determined by $`t`$ and $`U_{\mathrm{eff}}=U+F_0`$.
As a simplest extension to phenomenologically include spin-space anisotropy, while retaining only the relevant energy scales $`t`$ and $`U_{\mathrm{eff}}`$, we consider the generalized $`𝒩`$-orbital Hubbard modelquantum
$`H`$ $`=`$ $`t{\displaystyle \underset{i,\delta ,\gamma ,\sigma }{}}a_{i\gamma \sigma }^{}a_{i+\delta ,\gamma ,\sigma }+{\displaystyle \frac{U_1}{𝒩}}{\displaystyle \underset{i,\gamma ,\gamma ^{}}{}}a_{i\gamma }^{}a_{i\gamma }a_{i\gamma ^{}}^{}a_{i\gamma ^{}}`$ (2)
$`+`$ $`{\displaystyle \frac{U_2}{𝒩}}{\displaystyle \underset{i,\gamma ,\gamma ^{}}{}}a_{i\gamma }^{}a_{i\gamma ^{}}a_{i\gamma ^{}}^{}a_{i\gamma }`$
on a triangular lattice with nearest-neighbour (NN) hopping between sites $`i`$ and $`i+\delta `$. Here $`\gamma ,\gamma ^{}`$ refer to the (fictitious) degenerate $`𝒩`$ orbitals per site. The factor $`\frac{1}{𝒩}`$ is included to render the energy density finite in the $`𝒩\mathrm{}`$ limit. The two correlation terms involve density-density and exchange interactions with respect to the orbital indices. The Hartree-Fock (HF) approximation and Random Phase Approximation (RPA) are of O(1) whereas quantum fluctuation effects appear at higher order within the inverse-degeneracy expansion and thus $`1/𝒩`$, in analogy with $`1/S`$ for quantum spin systems, plays the role of $`\mathrm{}`$.
The key feature of spin-rotation symmetry of the generalized Hubbard model is highlighted by writing the two interaction terms as
$$H_{\mathrm{int}}=\frac{U_2}{𝒩}\underset{i}{}𝐒_i.𝐒_i+\frac{U_2U_1}{𝒩}\underset{i}{}S_{iz}^2$$
(3)
in terms of the total spin operator
$$𝐒_i=\underset{\gamma }{}\psi _{i\gamma }^{}\frac{𝝈}{2}\psi _{i\gamma }\underset{\gamma }{}\frac{𝝈_{i\gamma }}{2}$$
(4)
where $`\psi _{i\gamma }^{}(a_{i\gamma }^{}a_{i\gamma }^{})`$. An Ising (uniaxial) anisotropy is obtained for $`U_1>U_2`$, a planar (XY) anisotropy for $`U_2>U_1`$, and full spin-rotation symmetry for $`U_1=U_2`$.
As appropriate for $`\mathrm{HoMnO}_3`$, we consider the case $`U_2>U_1`$ corresponding to preferential ordering of spins in the $`xy`$ plane in spin space and an anisotropy gap for out-of-plane excitations.holmium Magnetic excitations were analyzed in terms of a Heisenberg model with a similar anisotropy term in Ref. . At the HF level, the interaction term for orbital $`\gamma `$ then reduces to
$$H_{\mathrm{int}}^\gamma =\underset{i}{}𝝈_{i\gamma }.𝚫_i$$
(5)
where the self-consistently determined mean field $`𝚫_i=U_2𝐒_{i\gamma ^{}}_{\mathrm{HF}}`$ lies in the $`xy`$ plane in spin space. The HF sublattice magnetization depends only on $`U_2`$ and is determined from the self-consistency condition $`𝐒_\alpha _{\mathrm{HF}}=_{𝐤,l}𝐤,l|\frac{𝝈}{2}|𝐤,l_\alpha `$ in terms of the HF states $`|𝐤,l`$ on sublattice $`\alpha `$, exactly as for the isotropic case.tri In the strong coupling limit the sublattice magnetization $`𝐒_\alpha 1/2`$ at the HF level, and is reduced by about 50% (in the isotropic case) due to quantum fluctuations, as found in different calculations cited in Ref. .
We consider the $`120^0`$ ordered AF state on the triangular lattice, and examine transverse spin fluctuations about the broken-symmetry state. At the RPA level, the magnon propagator reduces to a sum of all bubble diagrams where the interaction vertices involving $`S_{ix}^2`$ and $`S_{iy}^2`$ appear with interaction $`U_2`$, whereas those involving $`S_{iz}^2`$ with interaction $`U_1`$. Introducing a planar spin rotation, so that spins are oriented along the $`x^{}`$ direction for all three sublattices, we obtain for the transverse spin-fluctuation propagator
$$[\chi (𝐪,\omega )]_{\alpha \beta }^{\mu \nu }=\frac{[\chi ^0(𝐪,\omega )]}{\mathrm{𝟏}2[U][\chi ^0(𝐪,\omega )]}$$
(6)
in the $`23`$ spin-sublattice basis of the two transverse spin directions $`\mu ,\nu =y^{},z^{}`$ and the three sublattices $`\alpha ,\beta =A,B,C`$. The sublattice-diagonal interaction matrix
$$[U]=\left[\begin{array}{cc}U_2\mathrm{𝟏}& \mathrm{𝟎}\\ \mathrm{𝟎}& U_1\mathrm{𝟏}\end{array}\right]$$
(7)
in the $`y^{},z^{}`$ basis, and the bare particle-hole propagator
$`[\chi ^0(𝐪,\omega )]_{\alpha \beta }^{\mu \nu }`$
$`=`$ $`{\displaystyle \frac{1}{4}}{\displaystyle \underset{𝐤,l,m}{}}\left[{\displaystyle \frac{\sigma _\mu _\alpha ^+\sigma _\nu _\beta ^+}{E_{𝐤𝐪,m}^+E_{𝐤,l}^{}+\omega }}+{\displaystyle \frac{\sigma _\mu _\alpha ^+\sigma _\nu _\beta ^+}{E_{𝐤,m}^+E_{𝐤𝐪,l}^{}\omega }}\right]`$
involves integrating out the fermions in the broken-symmetry state. In the particle-hole matrix elements
$$\sigma _\mu _\alpha ^+𝐤𝐪,m|\sigma _\mu |𝐤,l_\alpha $$
(9)
of the rotated spins, the spin orientation angles $`\varphi _\alpha `$ in the fermion states $`|𝐤,l`$ are transformed out. The numerical evaluation of $`[\chi ^0(𝐪,\omega )]`$ in terms of the HF-level AF-state energies and amplitudes is exactly as for the isotropic case studied earlier.tri
The spin-wave energies $`\omega _𝐪`$ are then obtained from the poles $`1\lambda _𝐪(\omega _𝐪)=0`$ of Eq. (6) in terms of the eigenvalues $`\lambda _𝐪(\omega )`$ of the matrix $`2[U][\chi ^0(𝐪,\omega )]`$. As $`\omega _𝐪`$ corresponds to spin 1/2, it is scaled by the factor $`2S`$ for arbitrary spin $`S`$.magimp As expected for planar anisotropy, there is only one Goldstone mode corresponding to planar rotation of spins, and the two out-of-plane modes become massive with an anisotropy gap $`\omega _{\mathrm{gap}}`$.
Spin-wave dispersion $`\omega _𝐪`$ for the three modes is shown in Fig. 1 along two symmetry directions in the magnetic Brillouin zone (MBZ), along with the neutron-scattering data for $`\mathrm{HoMnO}_3`$ at 20 K from Ref. . The anisotropic Hubbard model provides a remarkably good description of the spin dynamics. We find that best fits with the magnetic excitations in $`\mathrm{HoMnO}_3`$ are only obtained in the strong-coupling limit $`U/t>15`$, with a planar exchange energy $`J=4t^2/U2.5`$meV and anisotropy $`U_2U_1\mathrm{\Delta }U0.35`$meV, the individual values of $`t`$ and $`U`$ not being resolvable within experimental resolution. To estimate the order of magnitude of the ratio $`U/t`$, if we nominally take $`U=1`$eV, we obtain $`t=25`$meV and $`U/t=40`$. Spin-wave dispersion calculated in the intermediate coupling regime cannot be fitted with the neutron scattering data, indicating no evidence of finite-$`U`$ double occupancy effects, as discussed below.
For the square-lattice $`S=1/2`$ AF $`\mathrm{La}_2\mathrm{CuO}_4`$, finite-$`U`$, double-occupancy effects associated with higher-order $`(t^4/U^3`$) spin couplings are manifested in noticeable spin-wave dispersion along the MBZ boundary.spinwave ; spectrum Similarly, for the isotropic triangular-lattice AF, the ratio $`\omega _1/\omega _{2,3}`$ of the non-degenerate and degenerate spin-wave energies at wave vector $`𝐪_M=(\pi /3,\pi /\sqrt{3})`$ is a sensitive measure of the $`U/t`$ ratio,tri which asymptotically approaches $`2/\sqrt{10}=0.632`$ in the strong-coupling $`(U/t\mathrm{})`$ limit, decreases monotonically with $`U/t`$, and eventually vanishes at $`U/t7`$.tri Variation of $`\omega _1/\omega _{2,3}`$ with $`U/t`$ is shown in Fig. 2 for different values of anisotropy $`\mathrm{\Delta }U/t`$. Neutron-scattering data for $`\mathrm{HoMnO}_3`$ shows that $`\omega _1/\omega _{2,3}11\mathrm{m}\mathrm{e}\mathrm{V}/16\mathrm{m}\mathrm{e}\mathrm{V}0.7`$, which yields a lower bound $`(15)`$ on the ratio $`U/t`$ from Fig. 2.
The spin-wave anisotropy gap $`\omega _{\mathrm{gap}}/t`$, corresponding to out-of-plane fluctuations with $`q=0`$, varies as $`(\mathrm{\Delta }U/U)^{1/2}`$ with spin anisotropy (Fig. 3), which translates to $`\omega _{\mathrm{gap}}\sqrt{J\mathrm{\Delta }U}`$, the geometric mean of the two energy scales.
In conclusion, the anisotropic Hubbard model provides a simple extension to phenomenologically include spin-space anisotropy and study spin excitations in the full range from weak to strong coupling. Comparison of spin-wave dispersion with RPA calculations allows for a quantitative determination of the effective anisotropy, in addition to highlighting any finite-$`U`$, double occupancy effects as seen in $`\mathrm{La}_2\mathrm{CuO}_4`$. The recent neutron-scattering data for spin-wave dispersion in $`\mathrm{HoMnO}_3`$ are found to be well described by an anisotropic Hubbard model on a triangular lattice with a planar (XY) spin anisotropy. We find that the twin constraints $`\omega _{\mathrm{gap}}/\omega _{1,2}1/3`$ and $`\omega _1/\omega _{2,3}2/3`$ in the neutron scattering data cannot be satisfied in the intermediate coupling regime, and best fit indicates that magnetic excitations in $`\mathrm{HoMnO}_3`$ correspond to the strong-coupling limit $`U/t>15`$, with $`J2.5`$meV and $`\mathrm{\Delta }U0.35`$meV. The ratio $`\omega _1/\omega _{2,3}`$ of the non-degenerate and degenerate spin-wave energies at wave vector $`𝐪_M=(\pi /3,\pi /\sqrt{3})`$ is suggested as a sensitive measure of finite-$`U`$, double occupancy effects in a triangular-lattice antiferromagnet. Finally, in view of the formal resemblance with the orbital Hubbard model (1), our RPA calculation provides a step forward towards investigating magnetic excitations using realistic models including Hund’s rule coupling, crystal-field splitting etc.
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# Finite-temperature Mott transitions in multi-orbital Hubbard model
## I Introduction
The Mott transition has been one of the most attractive topics in strongly correlated electron systems.Mott ; Gebhard There are a number of prototype materials in transition metal oxides. A well-known example is $`\mathrm{V}_2\mathrm{O}_3`$,V2O3 ; ImadaRev which exhibits the first-order metal-insulator transition at a certain critical temperature. The phase diagram obtained by systematic experiments is consistent with theoretical studies of the single orbital Hubbard model by means of dynamical mean field theory (DMFT).GeorgesRev ; KotliarPT This implies that the DMFT treatment of the Hubbard model, which properly takes into account local electron correlations, captures the essence of the Mott transition. However, in order to give more quantitative discussions about the Mott transitions, it is indispensable to incorporate the effect of the orbital degeneracy, which often gives rise to rich phase diagrams, as observed for $`\mathrm{La}_{1\mathrm{x}}\mathrm{Sr}_\mathrm{x}\mathrm{MnO}_3`$,TokuraScience $`R\mathrm{TiO}_3`$Katsufuji97 ; Mochizuki03 , etc. More recently, the orbital-selective Mott transition has attracted considerable attention, Anisimov02 ; Sigrist04 ; Liebsch ; KogaOS ; Ferrero05 ; Medici05 ; Arita05 in connection with the materials such as $`\mathrm{Ca}_{2\mathrm{x}}\mathrm{Sr}_\mathrm{x}\mathrm{RuO}_4`$Nakatsuji and $`\mathrm{La}_{\mathrm{n}+1}\mathrm{Ni}_\mathrm{n}\mathrm{O}_{3\mathrm{n}+1}`$LaNiO ; Kobayashi96 .
These interesting experimental findings have stimulated a number of theoretical works on the Mott transitions in the two-orbital Hubbard model by means of DMFT.Kotliar96 ; Rozenberg97 ; Bunemann:Gutzwiller ; Hasegawa98 ; Held98 ; Han98 ; Momoi98 ; Klejnberg98 ; Imai01 ; Koga:ED\_DMFT ; OnoDP ; Ono03 ; Pruschke04 ; Liebsch ; KogaOS ; Ferrero05 ; Medici05 ; Arita05 Although the two-orbital model has been studied in detail, a systematic study on the systems having more orbitals is still lacking at finite temperatures. For example, the finite-temperature Mott transitions for the multi-orbital Hubbard models has not been investigated quantitatively by DMFT. One of the difficulties lies in the practical calculation within the DMFT framework. Quantum Monte Carlo simulations, which can be a powerful numerical method to treat local correlations in DMFT, encounter sign problems. Also, Wilson’s renormalization group method gets more difficult to apply as the number of orbitals increases. A much more simplified approach, the two-site DMFT,Potthoff01 is not efficient enough to treat finite-temperature properties. It is thus desirable to systematically investigate electron correlations in the multi-orbital Hubbard model at finite temperatures.
Motivated by the above hot topics, we consider the Mott transitions in the multi-orbital Hubbard model at zero as well as finite temperatures. For this purpose, we make use of a self-energy functional approach (SFA) proposed by Potthoff recently.SFA ; Potthoff04 This method, which is based on the Luttinger-Ward functional approach,Luttinger gives a powerful tool to discuss electron correlations. A remarkable point is that this approach provides an efficient way to deal with finite-temperature properties of the multi-orbital Hubbard model. The main purpose of the present paper is to determine the finite-temperature phase diagram of the Hubbard model with up to four-fold bands, and quantitatively discuss how the orbital degeneracy affects the Mott transitions.
The paper is organized as follows. In the next section, we introduce the model Hamiltonian and briefly summarize the formulation of SFA. In Sec. III we present the detailed analysis based on SFA by exploiting the two-orbital model as a prototype system. Then in Sec. IV we determine the phase diagram, and discuss the nature of the Mott transitions in multi-orbital systems in detail. A brief summary is given in Sec. V.
## II Model and Method
We consider a correlated electron system having $`M`$ degenerate orbitals, which is described by the following multi-orbital Hubbard Hamiltonian,
$``$ $`=`$ $`_0+^{},`$ (1)
$`_0`$ $`=`$ $`{\displaystyle \underset{<i,j>}{}}{\displaystyle \underset{\alpha =1}{\overset{M}{}}}{\displaystyle \underset{\sigma }{}}t_\alpha c_{i\alpha \sigma }^{}c_{j\alpha \sigma },`$ (2)
$`^{}`$ $`=`$ $`U{\displaystyle \underset{i}{}}{\displaystyle \underset{\alpha =1}{\overset{M}{}}}n_{i\alpha }n_{i\alpha }`$ (3)
$`+`$ $`U^{}{\displaystyle \underset{i}{}}{\displaystyle \underset{\alpha <\alpha ^{}=1}{\overset{M}{}}}{\displaystyle \underset{\sigma \sigma ^{}}{}}n_{i\alpha \sigma }n_{i\alpha ^{}\sigma ^{}},`$
where $`c_{i\alpha \sigma }^{}(c_{i\alpha \sigma })`$ creates (annihilates) an electron with spin $`\sigma (=,)`$ and orbital $`\alpha (=1,2,3\mathrm{},M)`$ at the $`i`$ th lattice site, $`t_\alpha `$ denotes the hopping integral for orbital $`\alpha `$, $`U(U^{})`$ is the intra-orbital (inter-orbital) Coulomb interaction. In the following, we mainly consider the case of $`U=U^{}`$, and give brief discussions for more general cases including the Hund coupling in the end of the paper.
In order to address the Mott transitions in the multi-orbital Hubbard model, we use SFA.SFA ; Potthoff04 We briefly summarize the essence of SFA. For a correlated electron system, the grand potential is generally expressed as
$`\mathrm{\Omega }[𝚺]`$ $`=`$ $`F[𝚺]+\mathrm{Tr}\mathrm{ln}[(𝐆_0^1𝚺)^1],`$ (4)
where $`F[𝚺]`$ is the Legendre transformation of the Luttinger-Ward potential $`\mathrm{\Phi }[𝐆]`$, and $`𝐆`$ ($`𝐆_0`$) and $`𝚺`$ are the full (bare) Green’s function and the self-energy, respectively. The condition imposed on the functional (4),
$`{\displaystyle \frac{\mathrm{\Omega }[𝚺]}{𝚺}}=0,`$ (5)
gives the Dyson equation $`𝐆^1=𝐆_0^1𝚺`$. An important point is that the functional form of $`F[𝚺]`$ does not depend on the detail of the Hamiltonian $`_0`$ as far as the interaction term $`^{}(𝐔)`$ keeps its shape unchanged. This fact allows us to introduce a proper reference system having the same interaction term, which we denote as $`_{\mathrm{ref}}=_0(𝐭^{})+^{}(𝐔)`$. Then the grand potential is written as,
$`\mathrm{\Omega }[𝚺(𝐭^{})]`$ $`=`$ $`\mathrm{\Omega }(𝐭^{})`$ (6)
$`+`$ $`\mathrm{Tr}\mathrm{ln}\left[(𝐆_0(𝐭)^1𝚺(𝐭^{}))^1\right]`$
$``$ $`\mathrm{Tr}\mathrm{ln}\left[(𝐆_0(𝐭^{})^1𝚺(𝐭^{}))^1\right],`$
where $`\mathrm{\Omega }(𝐭^{})`$ and $`𝚺(𝐭^{})`$ are the grand potential and the self-energy for the reference system, whose bare Green’s function is denoted as $`𝐆_0(𝐭^{})^1=\omega +\mu 𝐭^{}`$ ($`\mu `$: chemical potential). The variational condition (5) is rewritten as
$$\frac{\mathrm{\Omega }[𝚺(𝐭^{})]}{𝐭^{}}=0.$$
(7)
By choosing the parameters $`𝐭^{}`$ for the reference system so as to satisfy the condition (7), we can find the best system within the reference Hamiltonian, which can approximately describe the original correlated system.
It should be noticed here that SFA provides us with an efficient and tractable way to deal with finite-temperature properties of the multi-orbital system, where standard DMFT combined with numerical techniques faces difficulties in a practical computation when the number of orbitals increases. Another notable point in this approach is that the critical behavior can be discussed more precisely than the DMFT analysis,GeorgesRev when one chooses the same type of the effective impurity model. In fact, by comparing the results obtained by the two-site DMFT,Bulla00 ; Potthoff01 the critical point for a single orbital Hubbard model obtained by SFA with the two-site model SFA is in good agreement with that obtained by DMFT with the aid of numerical techniques.Caffarel94 ; Moeller95 ; Bulla99
To discuss the multi-orbital Hubbard model (3), we exploit the Anderson impurity model as a reference system in SFA.SFA ; Potthoff04 The Hamiltonian for the reference system is explicitly given as
$`_{\mathrm{ref}}`$ $`=`$ $`{\displaystyle \underset{i}{}}_{\mathrm{ref}}^{(i)},`$ (8)
$`_{\mathrm{ref}}^{(i)}`$ $`=`$ $`{\displaystyle \underset{\alpha \sigma }{}}\epsilon _{0\alpha }^{(i)}c_{i\alpha \sigma }^{}c_{i\alpha \sigma }+{\displaystyle \underset{k=1}{\overset{N_b}{}}}{\displaystyle \underset{\alpha \sigma }{}}\epsilon _{k\alpha }^{(i)}a_{k\alpha \sigma }^{(i)}a_{k\alpha \sigma }^{(i)}`$ (9)
$`+`$ $`{\displaystyle \underset{k=1}{\overset{N_b}{}}}{\displaystyle \underset{\alpha \sigma }{}}V_{k\alpha }^{(i)}(c_{i\alpha \sigma }^{}a_{k\alpha \sigma }^{(i)}+H.c.)`$
$`+`$ $`U{\displaystyle \underset{\alpha }{}}n_{i\alpha }n_{i\alpha }+U^{}{\displaystyle \underset{\alpha <\alpha ^{}}{}}{\displaystyle \underset{\sigma \sigma ^{}}{}}n_{i\alpha \sigma }n_{i\alpha ^{}\sigma ^{}},`$
where $`a_{k\alpha \sigma }^{(i)}(a_{k\alpha \sigma }^{(i)})`$ creates (annihilates) an electron with $`\sigma `$ spin and $`\alpha `$ orbital at the $`k(=1,2,\mathrm{}N_b)`$ th site, which is connected to the $`i`$ th site in the original lattice. Since we consider the Mott transitions without symmetry breaking, we set $`\epsilon _{k\alpha }`$ and $`V_{k\alpha }`$ site-independent. Note that in the limit of $`N_b\mathrm{}`$, Eq. (7) is equivalent to the self-consistent equation in DMFT.SFA Since the Green’s function and the self-energy are diagonal with respect to the site, spin, and orbital indices, the grand potential per site reads
$`\mathrm{\Omega }/L`$ $`=`$ $`\mathrm{\Omega }_{\mathrm{imp}}`$ (10)
$``$ $`2{\displaystyle \underset{\alpha }{}}{\displaystyle 𝑑\omega f(\omega )R_\alpha \left[\omega +\mu \mathrm{\Sigma }_\alpha (\omega )\right]}`$
$`+`$ $`2{\displaystyle \underset{\alpha }{}}{\displaystyle \underset{k=0}{\overset{N_b}{}}}{\displaystyle 𝑑\omega f(\omega )\theta \left[G_{k\alpha }^{}(\omega )\right]},`$
$`G_{0\alpha }^{}(\omega )`$ $`=`$ $`[\omega +\mu \epsilon _{0\alpha }\mathrm{\Delta }_\alpha (\omega )\mathrm{\Sigma }_\alpha (\omega )]^1,`$ (11)
$`G_{k\alpha }^{}(\omega )`$ $`=`$ $`(\omega +\mu \epsilon _{k\alpha })^1,(k=1,2,\mathrm{},N_b),`$ (12)
$`\mathrm{\Delta }_\alpha (\omega )`$ $`=`$ $`{\displaystyle \underset{k=1}{\overset{N_b}{}}}V_{k\alpha }^2G_{k\alpha }^{}(\omega ),`$ (13)
where $`R_\alpha (\omega )=_{\mathrm{}}^\omega \rho _\alpha (z)𝑑z`$, $`f(\omega )=1/(1+e^{\beta \omega })`$ and $`\theta (\omega )`$ is a step function. The grand potential and the self-energy for an impurity system $`_{\mathrm{ref}}^{(i)}`$ are denoted as $`\mathrm{\Omega }_{\mathrm{imp}}`$ and $`\mathrm{\Sigma }_\alpha (\omega )`$.
In the following, we focus on the half-filled Hubbard model ($`\mu =(M1/2)U`$) to discuss how the orbital degeneracy $`(M=1,2,3,4)`$ affects the Mott transitions. We take the simplest reference system $`(N_b=1)`$, where we fix the parameters $`(\epsilon _{0\alpha },\epsilon _{1\alpha },V_{1\alpha })=(0,\mu ,V)`$ for each orbital $`\alpha `$. This simplified treatment is somewhat analogous to the two-site DMFT at zero temperature. However, the present approach enables us to treat finite-temperature quantities systematically, which is a notable advantage beyond the two-site DMFT. The condition (7) is now reduced to $`\mathrm{\Omega }/V=0`$. We note that the hopping integral $`V`$ between a given site on the original lattice and a fictitious site corresponds, roughly speaking, to the renormalized bandwidth of quasi-particles. In fact, Fermi-liquid properties at zero temperature are determined by the value of $`V`$. For instance, small $`V`$ implies the formation of heavy quasi-particles, and the system enters the insulating phase at $`V=0`$. By calculating the effective hybridization $`V`$, we thus discuss the stability of the metallic state in the multi-orbital systems.
## III Detail of calculations: two-orbital system
Let us start with the two-orbital Hubbard model, by which we illustrate the basic procedure of SFA to determine the phase diagram for more general multi-orbital cases. We adopt here a semielliptic density of states $`\rho (\epsilon )=(\pi /W)\sqrt{1(2\epsilon /W)^2}`$ with the bandwidth $`W=4`$.
We first look at the ground-state properties by examining the stationary points in the grand potential.
In Fig. 1, we show the grand potential at zero temperature for several values of the Coulomb interaction $`U`$. When $`U`$ is small, the grand potential has a minimum at finite $`V`$. Thus the effective bandwidth is finite, stabilizing the metallic ground state. It is seen that the stationary point moves towards the origin continuously with increasing $`U`$. For large $`U`$, the minimum is located at the origin, indicating that the Mott insulating state is stabilized by strong electron correlations. In the vicinity of the critical point $`U_c`$, the grand potential can be expanded in powers of $`V`$,
$$\mathrm{\Omega }(V)=\mathrm{\Omega }(0)+AV^2+𝒪(V^4).$$
(14)
Therefore, the critical point separating the metallic and the insulating phases is characterized by the condition $`A=0`$.SFA By solving this analytically, we obtain the self-consistent equation for the critical point $`U_c`$ as,
$`U_c`$ $`=`$ $`{\displaystyle \frac{100}{39}}{\displaystyle _{\mathrm{}}^0}\left[z\xi (z,U_c)\right]\rho (z)𝑑z`$ (15)
$``$ $`{\displaystyle \frac{U_c}{39}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}{\displaystyle \frac{317U_c+83\xi (z,U_c)}{[17U_c+8\xi (z,U_c)]\xi (z,U_c)}}\rho (z)𝑑z`$
with $`\xi (z,U)=\sqrt{U^2+z^2}`$. The detail of the derivation is shown in Appendix A. The critical point $`U_c9.217`$ thus obtained is more accurate than the results $`(U_c=10)`$ estimated by the two-site DMFT method.Koga02
We now consider the competition between the metallic and the Mott insulating phases at finite temperatures, following the way for the single band case.SFA The grand potential is shown in Fig. 2 at finite temperatures.
It is seen that two minima appear at low temperatures. One of the minima is located at larger $`V`$, which corresponds to the metallic solution, since it is continuously connected to the metallic one at zero temperature. The other is adiabatically connected to $`V=0`$ at zero temperature, so that this solution characterizes the Mott insulating state. Such a double-well structure at low temperatures causes the first-order transition accompanied by hysteresis. In fact, as increasing temperature, the stationary point for the metallic state disappears around $`T_{c2}=0.05`$, where the Mott transition occurs to the insulating phase. On the other hand, as decreasing temperature, the Mott insulating phase realized at high temperatures is stable except for zero temperature, since the corresponding local minimum always exists at finite temperatures. The first-order transition temperature $`T_c=0.045`$ for $`U=7`$ is determined by the crossing point of the two minima in the grand potential, as shown in Fig. 2. This hysteresis induced by the first-order transition is also observed when the interaction $`U`$ is varied. We show the effective hybridization $`V`$ at $`T=0.04`$ in Fig. 3.
Starting from the metallic state, the increase of the Coulomb interaction decreases the effective hybridization $`V`$, and finally triggers the Mott transition to the insulating phase at $`U_{c2}=7.2`$, where we observe the discontinuity in $`V`$. On the contrary, the Mott transition occurs at $`U_{c1}=6.7`$ when the interaction decreases. In this parameter regime, the level crossing in the minima of the grand potential appears at $`U_c=7.1`$ (inset of Fig. 3), which defines the first-order transition point at $`T=0.04`$.
We have seen here that the two-orbital Hubbard model exhibits qualitatively similar properties to the single-band case GeorgesRev ; SFA as far as the nature of the Mott transitions is concerned. In the following section, we give more quantitative discussions about the Mott transitions by using the phase diagrams obtained for the multi-orbital Hubbard model with $`M=14`$.
## IV Phase diagrams of multi-orbital systems
We now determine the phase diagrams of the multi-orbital Hubbard model. We first calculate the renormalization factor $`Z=(1\mathrm{\Sigma }(\omega )/\omega )^1`$, which is inversely proportional to the effective mass, to characterize the metallic ground state at zero temperature. The results are shown in Fig. 4.
The introduction of the Coulomb interaction decreases the renormalization factor, making the mass of quasi-particles heavier. In the small $`U`$ region, there is not a big difference in the behavior of $`Z`$ among different $`M`$ cases. If we look at $`Z`$ in more detail (inset of Fig. 4), the mass gets slightly heavier (smaller $`Z`$) as $`M`$ increases. This comes from the fact that the electron correlations are somewhat enhanced for multi-orbital cases by the inter-orbital Coulomb interaction. On the other hand, further increase of the Coulomb interaction leads to quite different behavior: as the number of orbitals $`M`$ increases, the mass is less renormalized, which makes the metallic state more stable up to large $`U`$.Lu94 ; Fresard:SlaveBoson ; Florens02 This implies that the inter-orbital interaction enhances orbital fluctuations for large $`U`$, which in turn stabilize the metallic state, as pointed out in the two-orbital case.Koga02 In fact, the critical points $`U_{c2}(T=0)=9.2173,12.6044`$ and $`15.9958`$ for $`M=2,3`$ and 4 (see Appendix), are much larger than $`U_c=5.84`$ for the single-orbital model.SFA Therefore, enhanced orbital fluctuations play a key role to stabilize the metallic state at zero temperature.
We now move to the finite-temperature properties. The phase diagrams obtained in the way mentioned in the previous section are shown in Fig. 5.
There are three phase boundaries for each system, $`U_{c1}`$, $`U_{c2}`$ and $`U_c`$. The double-well structure in the grand potential gives rise to two kinds of transitions with the discontinuity in the physical quantities. On the phase boundary $`U_{c1}(U_{c2})`$, the Mott transition occurs from the insulating (metallic) state to the metallic (insulating) state as $`U`$ decreases (increases). The area surrounded by $`U_{c1}`$ and $`U_{c2}`$ is referred to as the coexistence phase.GeorgesRev ; SFA The critical point $`U_c`$ is determined so that the two minima of the grand potential take the same value. At zero temperature, the critical points $`U_c`$ and $`U_{c2}`$ merge to give the continuous Mott transition, because the double-well structure disappears at $`T=0`$. Therefore, the introduction of the temperature drastically changes the phase boundary as discussed in the single-orbital case. GeorgesRev ; SFA It is found that as temperature increases, the phase boundaries $`U_{c1}`$, $`U_{c2}`$, and $`U_c`$ merge at the critical temperature $`T_c`$, where the second-order phase transition occurs.
The above characteristic properties in the Mott transitions are qualitatively the same as the single-orbital case. We now discuss the effects due to orbital fluctuations quantitatively. First note that the coexistence region bounded by the phase boundaries $`U_{c1}`$ and $`U_{c2}`$ is enlarged when the number of orbital $`M`$ increases. Furthermore, we see that the $`M`$-dependence of the critical points $`U_{c1}`$, $`U_c(T_c)`$ and the critical temperature $`T_c`$ is different from that of $`U_{c2}`$. To clarify this point, the critical values are plotted as a function of $`M`$ in Fig. 6.
It is found that $`U_{c2}`$ is proportional to the orbital degeneracy $`M`$, while $`U_{c1}`$ and $`U_c(T_c)`$ the square root of $`M`$. These results are consistent with the previous results. Lu94 ; Fresard:SlaveBoson ; Hasegawa98 ; Florens02 ; Ono03 However, we wish to note that the finite-temperature properties such as the phase diagrams have not been clarified quantitatively so far. In particular, the results on the $`M`$-dependence of the critical temperature $`T_c`$ obtained here are beyond the qualitative discussions by Florens et al.,Florens02 who obtained only the upper bound for $`T_cM`$.
To conclude this section, we would like to briefly discuss the effects of the Hund exchange coupling. For this purpose, we add the following term to the Hamiltonian (3),
$`_J`$ $`=`$ $`J{\displaystyle \underset{i}{}}{\displaystyle \underset{\alpha <\alpha ^{}}{}}{\displaystyle \underset{\sigma }{}}n_{i\alpha \sigma }n_{i\alpha ^{}\sigma }`$ (16)
$`J{\displaystyle \underset{i}{}}{\displaystyle \underset{\alpha <\alpha ^{}}{}}(c_{i\alpha }^{}c_{i\alpha }c_{i\alpha ^{}}^{}c_{i\alpha ^{}}+H.c.)`$
$`J{\displaystyle \underset{i}{}}{\displaystyle \underset{\alpha <\alpha ^{}}{}}(c_{i\alpha }^{}c_{i\alpha }^{}c_{i\alpha ^{}}c_{i\alpha ^{}}+H.c.)`$
with $`J>0`$, where we impose the condition $`U=U^{}+2J`$ due to the symmetry requirement. Since the calculation gets somewhat difficult in the presence of $`J`$, we show the results obtained for the representative case of $`M=2`$. The obtained phase diagram is shown in Fig. 7 and compared with the $`J`$ =0 case. First, we notice that upon introducing $`J`$, the phase boundaries are immediately shifted to the weak-interaction regime, and therefore the metallic state gets unstable for large $`U`$. Correspondingly, the coexistence region surrounded by the first order transitions shrinks as $`J`$ increases. This remarkable tendency indeed reflects the fact that the metallic state is stabilized by enhanced orbital fluctuations in the case of $`U=U^{}`$: the Hund coupling suppresses such orbital fluctuations, and stabilizes the Mott-insulating phase. Another point to be mentioned is that the Mott transition becomes of first-order even at zero temperature in the presence of $`J`$,Bunemann:Gutzwiller ; Ono03 ; Pruschke04 since the subtle balance realized at $`T=0`$ in the case of $`U=U^{}`$ is not kept anymore for finite $`J`$. Since there is another claim that the second order transition is possible for certain choices of parameters at $`T=0`$,Pruschke04 which may depend on the strength of Hund coupling, we need further detailed discussions to clarify this problem. Although we have presented only the results for the $`M=2`$ case, we expect that the effects of $`J`$ on the phase diagram should be essentially the same as shown here: the introduction of $`J`$ suppresses orbital fluctuations, thus favors the Mott insulating phase even in the small $`U`$ regime. It also shrinks the coexistence phase dramatically.
## V Summary
We have investigated the Mott transitions in the multi-orbital Hubbard model at zero and finite temperatures by means of the self-energy functional approach. By choosing the Anderson impurity model as a reference system, we have discussed how the orbital degeneracy affects the nature of the Mott transitions. We have obtained the finite-temperature phase diagram for the system having $`M`$($`4`$) degenerate bands. Although the phase diagrams show qualitatively similar features irrespective of the orbital degeneracy, there are some remarkable effects due to degenerate orbitals. In particular, it has been shown that enhanced orbital fluctuations make the metallic phase more stable even at finite temperatures. Therefore, if such fluctuations are suppressed, the metallic state is expected to be unstable, which has been shown to be the case by considering the system with the Hund coupling. Also, it has quantitatively been clarified how distinctly the $`M`$-dependence appears in the critical points $`U_{c1}`$ and $`U_{c2}`$: the critical point $`U_{c1}`$ depends on the square root of the orbital degeneracy $`M`$, while the critical point $`U_{c2}`$ is proportional to $`M`$. This analysis concludes that the critical temperature $`T_c`$ is proportional to the square root of $`M`$, which may be important to understand the Mott transitions in the real materials with orbital degeneracy.
The self-energy functional approach used in the paper allows quantitatively reliable discussions for multi-orbital systems at finite temperatures. Since this formalism provides a tractable way to incorporate spin and charge ordered states, which have been neglected in this paper, it may be used for more detailed study of the finite-temperature properties in multi-orbital Mott systems.
###### Acknowledgements.
Numerical computations were carried out at the Supercomputer Center, the Institute for Solid State Physics, University of Tokyo. This work was supported by a Grant-in-Aid for Scientific Research from the Ministry of Education, Culture, Sports, Science, and Technology, Japan.
## Appendix A
We analytically determine the critical value of $`U`$ for the Mott transition in the multi-orbital Hubbard model at zero temperature. The grand potential Eq. (10) is rewritten as
$`\mathrm{\Omega }/L=\mathrm{\Omega }_{\mathrm{imp}}`$ $`+`$ $`2M{\displaystyle \underset{i}{}}{\displaystyle 𝑑z\rho (z)\omega _i(z)\theta \left[\omega _i(z)\right]}`$ (17)
$``$ $`2M{\displaystyle \underset{i}{}}\omega _i\theta (\omega _i^{}),`$
where $`\mathrm{\Omega }_{\mathrm{imp}}`$ is the grand potential of the reference system and $`\theta (x)`$ is the step function. $`\omega _i^{},\omega _i(z)`$ are the poles of the impurity Green’s function of the reference system, $`G^{}(\omega )`$, and the approximated Green function of the original system, $`G(\omega ;z)=1/(\omega +\mu z\mathrm{\Sigma }(\omega ))`$, respectively.
First we consider the two-orbital system. In the atomic limit $`V=0`$, the Green’s function and the self-energy at the impurity site for the reference system are given by
$`G_{\mathrm{atom}}^{}(\omega )`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left[{\displaystyle \frac{1}{\omega U/2}}+{\displaystyle \frac{1}{\omega +U/2}}\right]`$ (18)
$`\mathrm{\Sigma }_{\mathrm{atom}}(\omega )`$ $`=`$ $`\mu +{\displaystyle \frac{U^2}{4\omega }},`$ (19)
where the poles of $`G_{\mathrm{atom}}^{}`$ are $`\pm U/2`$. Then the Green’s function of the original lattice model is given by
$`G_{\mathrm{atom}}^1(\omega ;z)`$ $`=`$ $`\omega z+{\displaystyle \frac{U^2}{4\omega }},`$ (20)
which has the poles at
$$\omega _\pm (z)=\frac{1}{2}\left(z\pm \sqrt{z^2+U^2}\right).$$
(21)
As discussed in the text, we can expand the quantities in $`V`$ around the atomic limit. The grand potential is $`\mathrm{\Omega }^{}=2U24V^2/U+𝒪(V^4)`$. We also analyze the Green’s functions $`G^{}(\omega )`$ and $`G(\omega ;z)`$ around the atomic limit and obtain their poles up to the second order in $`V`$. Here we need only the poles in the negative energy region, because the poles in the positive energy region do not affect $`\mathrm{\Omega }`$, see Eq. (17). The Green’s functions $`G^{}(\omega )`$ and $`G(\omega ;z)`$ have eight poles in the $`M=2`$ case, and thus four poles in the negative region. The negative poles of $`G^{}(\omega )`$ are
$`\omega _1^{}`$ $`=`$ $`10{\displaystyle \frac{V^2}{U}}+𝒪(V^4),`$
$`\omega _2^{}`$ $`=`$ $`{\displaystyle \frac{U}{2}}(30+4\sqrt{3}){\displaystyle \frac{V^2}{U}}+𝒪(V^4),`$
$`\omega _3^{}`$ $`=`$ $`{\displaystyle \frac{U}{2}}(304\sqrt{3}){\displaystyle \frac{V^2}{U}}+𝒪(V^4),`$
$`\omega _4^{}`$ $`=`$ $`2U26{\displaystyle \frac{V^2}{U}}+𝒪(V^4),`$
and the negative poles of $`G(\omega ;z)`$ are
$`\omega _1(z)`$ $`=`$ $`100zV^2/U^2+𝒪(V^4),`$
$`\omega _2(z)`$ $`=`$ $`\omega _{}(z)+B(z/U)V^2/U+𝒪(V^4),`$
$`\omega _3(z)`$ $`=`$ $`U/230V^2/U+𝒪(V^4),`$
$`\omega _4(z)`$ $`=`$ $`2U{\displaystyle \frac{168z+390U}{15U+8z}}V^2/U+𝒪(V^4),`$
where
$`B(x)={\displaystyle \frac{6525+8714x^23200x^4}{(22564x^2)\sqrt{1+x^2}}}{\displaystyle \frac{50x(23764x^2)}{22564x^2}}.`$
The grand potential is written down as
$`\mathrm{\Omega }/L`$ $`=`$ $`4{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑z\rho (z)\xi (z,U)`$
$`+`$ $`{\displaystyle \frac{4V^2}{U^2}}\{39U+100{\displaystyle _{\mathrm{}}^0}dz\rho (z)(z\xi (z,U))`$
$`+`$ $`U{\displaystyle _{\mathrm{}}^{\mathrm{}}}dz\rho (z){\displaystyle \frac{317U+83\xi (z,U)}{[17U+8\xi (z,U)]\xi (z,U)}}\}+𝒪(V^4),`$
where $`\xi (z,U)=\sqrt{U^2+z^2}`$. The condition $`^2\mathrm{\Omega }/V^2|_{V=0}=0`$ is satisfied at the Mott transition point $`U_c`$.SFA Given that the free density of state is symmetric $`\rho (z)=\rho (z)`$, we derive the self-consistent equation for $`U_c`$, Eq. (15).
Similarly, we can analyze the cases of $`M=3`$ and $`M=4`$. The self-consistent equation for $`U_c`$ at $`M=3`$ reads
$`U_c`$ $`=`$ $`{\displaystyle \frac{49}{20}}{\displaystyle _{\mathrm{}}^0}[z\xi (z,U_c))]\rho (z)dz`$
$`{\displaystyle \frac{U_c}{80}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}{\displaystyle \frac{545U_c+80\xi (z,U_c)}{[17U_c+8\xi (z,U_c)]\xi (z,U_c)}}\rho (z)𝑑z,`$
and at $`M=4`$ is
$`U_c`$ $`=`$ $`{\displaystyle \frac{324}{135}}{\displaystyle _{\mathrm{}}^0}\left[z\xi (z,U_c)\right]\rho (z)𝑑z`$
$`{\displaystyle \frac{U_c}{135}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}{\displaystyle \frac{837U_c+63\xi (z,U_c)}{[17U_c+8\xi (z,U_c)]\xi (z,U_c)}}\rho (z)𝑑z.`$
The solution of these equations gives the accurate values of the critical point: $`U_c=12.6044`$ and $`15.9958`$ for $`M=3`$ and $`4`$, respectively.
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# Measurement of Double Charmonium Production in 𝑒⁺𝑒⁻ Annihilations at √𝑠=10.6 GeV
BABAR-PUB-05/021
SLAC-PUB-11287
thanks: Deceased
The BABAR Collaboration
## Abstract
We study $`e^+e^{}J/\psi c\overline{c}`$ by measuring the invariant mass distribution recoiling against fully reconstructed $`J/\psi `$ decays, using 124$`\text{ fb}^1`$ of data collected with a center-of-mass energy of 10.6$`\mathrm{GeV}`$ with the BABAR detector. We observe signals for $`\eta _c(1S)`$, $`\chi _{c0}`$, and $`\eta _c(2S)`$ in the recoil mass distribution, thus confirming previous measurements. We measure $`\sigma (e^+e^{}J/\psi +c\overline{c})(c\overline{c}>2\mathrm{charged})`$ to be $`17.6\pm 2.8(\mathrm{stat})_{2.1}^{+1.5}(\mathrm{syst})`$ fb, $`10.3\pm 2.5(\mathrm{stat})_{1.8}^{+1.4}(\mathrm{syst})`$ fb, and $`16.4\pm 3.7(\mathrm{stat})_{3.0}^{+2.4}(\mathrm{syst})`$ fb with $`c\overline{c}=\eta _c(1S)`$, $`\chi _{c0}`$, and $`\eta _c(2S)`$, respectively.
preprint: BABAR-PUB-05/021preprint: SLAC-PUB-11287
Prompt $`J/\psi `$ and $`\psi (2S)`$ production in $`e^+e^{}`$ annihilations around $`\sqrt{s}=10.6`$$`\mathrm{GeV}`$ has been observed by both the BABAR ct:BaBar-jsi and Belle ct:Belle-psi experiments. These interactions provide an opportunity to study both perturbative and non-perturbative effects in QCD and to search for new charmonium states ct:search-charmonium-1 ; ct:singleGamma-Chao .
Belle ct:Belle-double\_ccbar reported the observation of $`\eta _c(1S)`$, $`\chi _{c0}`$, and $`\eta _c(2S)`$ in the mass distribution of the system recoiling against a reconstructed $`J/\psi `$ in $`e^+e^{}`$ annihilations. The production cross sections measured by Belle are about one order of magnitude higher than those predicted by non-relativistic QCD (NRQCD) calculations ct:singleGamma-Chao ; ct:singleGamma-Braaten-Lee ; ct:singleGamma2-Chao for $`e^+e^{}\gamma ^{}J/\psi c\overline{c}`$ reactions, where $`c\overline{c}`$ is a charmonium state with even C-parity. There have been attempts ct:Zhang-JpsiEtac ; ct:Bondar-JpsiEtac ; ct:Bodwin-twoPhoton ; ct:Bodwin-JpsiJpsi ; ct:Brodsky-Glueball to reconcile the large discrepancy between the observed cross section and predictions, and the validity of NRQCD approximations has been questioned ct:Bondar-JpsiEtac ; ct:Brambilla-QWG2004 . It has also been suggested that at least part of the double charmonium production might be due to two virtual-photon interactions ct:Bodwin-twoPhoton , $`i.e.`$, $`e^+e^{}\gamma ^{}\gamma ^{}J/\psi c\overline{c}`$, where odd C-parity states could be produced. Belle updated its observation and explored the origin of the $`J/\psi c\overline{c}`$ events ct:Belle-double\_ccbar-update .
In this paper we present a measurement of the cross sections for $`e^+e^{}J/\psi \eta _c(1S)`$, $`e^+e^{}J/\psi \chi _{c0}`$, and $`e^+e^{}J/\psi \eta _c(2S)`$, and set limits on the yields for other known charmonium states produced in association with a $`J/\psi `$. We calculate the mass ($`M_{\text{rec}}`$) of the system recoiling against a fully reconstructed $`J/\psi `$ via:
$$M_{\text{rec}}^2=(\sqrt{s}E_{J/\psi }^{})^2p_{J/\psi }^2,$$
(1)
where $`\sqrt{s}`$ is the $`e^+e^{}`$ annihilation energy in the center-of-mass (CM) system, and $`E_{J/\psi }^{}`$ and $`p_{J/\psi }^{}`$ are the energy and momentum of the $`J/\psi `$ candidate in the CM system.
In this paper, we analyze 112$`\text{ fb}^1`$ of data collected at the peak of the $`\mathrm{{\rm Y}}(4S)`$ resonance and 12$`\text{ fb}^1`$ at $`\sqrt{s}=10.54\mathrm{GeV}`$, just below the $`\mathrm{{\rm Y}}(4S)`$, with the BABAR detector ct:BaBar-detector operating at the asymmetric energy PEP-II $`e^+e^{}`$ storage ring. The BABAR detector includes a five-layer, double-sided silicon vertex tracker (SVT) and a 40-layer drift chamber (DCH) in a 1.5-T solenoidal magnetic field, which detects charged particles and measures their momenta and specific ionizations ($`\mathrm{d}E/\mathrm{d}x`$). Photons and electrons are detected with a CsI(Tl)-crystal electromagnetic calorimeter (EMC). An internally reflecting ring-imaging Cherenkov (DIRC) is used for particle identification. Penetrating muons are identified by an array of resistive-plate chambers (RPC) embedded in the steel of the flux return (IFR).
We select events with at least five well reconstructed charged tracks in the DCH, within the fiducial volume 0.41$`<\theta <`$2.54, where $`\theta `$ is the polar angle. Electron candidates have a pattern of specific ionization ($`\mathrm{d}E/\mathrm{d}x`$) in the DCH, a Cherenkov cone angle, an EMC shower energy divided by momentum, and a number of EMC crystals that are consistent with an electron hypothesis. A muon candidate is selected on the basis of energy deposited in the EMC, the number and distribution of hits in the IFR, and the match between the IFR hits and the extrapolation of the DCH track into the IFR. A more detailed explanation of particle identification is given elsewhere ct:BaBar-jsi .
A pair of oppositely charged lepton candidates originating from a common vertex is selected as a $`J/\psi `$ candidate if its mass ($`m(\mathrm{}^+\mathrm{}^{})`$) falls within \[-50,30\]$`\mathrm{MeV}/c^2`$ (for $`e^+e^{}`$) or \[-30,30\]$`\mathrm{MeV}/c^2`$ (for $`\mu ^+\mu ^{}`$), of the nominal $`J/\psi `$ mass of 3.097$`\mathrm{MeV}/c^2`$ ct:PDG2004 . In the calculation of $`m(e^+e^{})`$, electron candidates are combined with nearby photon candidates in order to recover some of the energy lost through bremsstrahlung radiation. These mass intervals are referred to as the $`J/\psi `$ mass windows. In order to improve the $`p_{J/\psi }^{}`$ resolution, we perform a kinematic fit where the $`J/\psi `$ candidate is constrained to have the nominal $`J/\psi `$ mass.
There are two main background sources in this analysis: events with genuine $`J/\psi `$ mesons and combinatorial background. The region 60$`\mathrm{MeV}/c^2`$ $`<|M(\mathrm{}^+\mathrm{}^{})M(J/\psi )|<`$200$`\mathrm{MeV}/c^2`$, defined as the $`J/\psi `$ mass sidebands, where $`M(J/\psi )`$ is the nominal $`J/\psi `$ mass, is used to estimate the combinatorial background due to random tracks. This background is largely rejected by particle identification, and by a requirement on the lepton helicity angle in the $`J/\psi `$ decay, $`|\mathrm{cos}\theta _l|<0.9`$, as shown in Fig. 1(a) and 1(c).
The largest backgrounds are due to real $`J/\psi `$ mesons from QED processes such as $`J/\psi `$ or $`\psi (2S)`$ mesons produced via initial state radiation (ISR). $`J/\psi `$ mesons from $`B`$ meson decay have $`p^{}<2\mathrm{GeV}/c`$ and do not constitute a background for recoil masses below 6.6$`\mathrm{GeV}/c^2`$. Most QED backgrounds have low multiplicity, and may have electrons or photons escaping detection along the beam line. These backgrounds are suppressed by the requirement of at least five charged tracks and the following requirement: for each event we calculate the energy deposited in the EMC plus the energy that can be attributed to an undetected electron or photon,
$$E_{\text{QED}}=E_{\text{EMC}}+p_{\text{miss}},$$
(2)
where $`E_{\text{EMC}}`$ is the total energy deposited in the EMC, and $`p_{\text{miss}}`$ is the missing momentum in the lab frame in the event. We require $`E_{\text{QED}}E_{\text{beams}}<1.0`$$`\mathrm{GeV}`$ as shown in Fig. 1(b) and 1(d), where $`E_{\text{beams}}`$ is the sum of the $`e^+e^{}`$ beam energies calculated in the lab frame.
We reject the $`J/\psi `$ background from $`\psi (2S)`$ events by vetoing events if the invariant mass of the $`J/\psi `$ candidates combined with any pair of oppositely charged tracks with pion mass hypothesis is within 15$`\mathrm{MeV}/c^2`$ of the $`\psi (2S)`$ mass.
The recoil mass distribution for events in the $`J/\psi `$ mass window is shown as points with error bars in Fig. 2. The ISR $`\psi (2S)`$ background is estimated using a Monte Carlo sample of ISR $`\psi (2S)`$ events. The $`\psi (2S)`$ feeddown background from continuum production is estimated using continuum $`\psi (2S)`$ events selected in the data.
The spectrum in Fig. 2 is fit to the sum of signal functions representing the $`\eta _c(1S)`$, $`\chi _{c0}`$, and $`\eta _c(2S)`$ lineshapes, plus a second-order polynomial background function. The signal line shapes are obtained by convoluting the Breit-Wigner line shape of each resonance with a fixed-width Gaussian representing the recoil mass resolution function. The widths of the Gaussians are determined from a Monte Carlo simulation of the momentum of the reconstructed $`J/\psi `$; the $`J/\psi `$ momentum resolution is different for the $`J/\psi e^+e^{}`$ and $`J/\psi \mu ^+\mu ^{}`$ samples, but independent of the recoiling system. This shape in turn is convolved with a long radiative tail that is calculated to $`𝒪(\alpha ^2)`$ ct:generator-kk2f for ISR photons that carry off an energy greater than 10 MeV. The free parameters in the data fit are the coefficients for the background parameterization, the event yields for each resonance, the masses of the resonances, and the $`\eta _c(2S)`$ total width. The total widths for the $`\eta _c(1S)`$ and the $`\chi _{c0}`$ are fixed to their world average values ct:PDG2004 of 17.3$`\mathrm{MeV}/c^2`$ and 10.1$`\mathrm{MeV}/c^2`$, respectively. The fit is performed simultaneously to the recoil mass spectra in the $`J/\psi e^+e^{}`$ and $`J/\psi \mu ^+\mu ^{}`$ samples, and the total event yield for each resonance is given by the sum of the yields in each mode.
The fit result is given in Table 1 and is shown as the solid curve in Fig. 2. Other known charmonium states may also be produced in association with the $`J/\psi `$ via two virtual-photon interactions. We therefore attempt to include in our primary fit each one of the other known charmonium resonances in turn to determine their event yields, which are presented in Table 1. We find no evidence for $`J/\psi `$, $`\chi _{c1}`$, $`\chi _{c2}`$, or $`\psi (2S)`$ in the mass spectrum of the system recoiling against a $`J/\psi `$.
The topological branching fraction is unknown for the $`\eta _c(1S)`$, $`\chi _{c0}`$, and $`\eta _c(2S)`$, so we report the product of the branching fraction for final states with more than two charged tracks ($`_{>2}(c\overline{c}>2\mathrm{charged})`$) times the double charmonium production cross section. In order to include the effect of ISR, the yields reported in Table 1 are calculated with a line shape based on a model of the $`\sqrt{s}`$ dependence of double charmonium production model. To allow a direct comparison of experimental results, we follow the same method used by Belle ct:Belle-double\_ccbar-update to remove this model dependence by determining cross section values that correspond to the non-tail fraction of the fit shape ($`f_{rad}=0.61`$ct:generator-kk2f where no ISR photon with an energy greater than 10$`\mathrm{MeV}`$ is radiated. We use
$$\sigma (e^+e^{}J/\psi c\overline{c})_{>2}=\frac{N_{c\overline{c}}f_{rad}}{(J/\psi \mathrm{}^+\mathrm{}^{})\epsilon },$$
(3)
where $`N_{c\overline{c}}`$ is the event yield, $``$ is the integrated luminosity, $`(J/\psi \mathrm{}^+\mathrm{}^{})`$ is the $`J/\psi `$ branching fraction, and $`\epsilon `$ is the detection efficiency. The value of $`\epsilon `$ is determined using a Monte Carlo simulation with the assumption that exclusive $`J/\psi \eta _c(nS)`$ production is $`P`$ wave and that exclusive $`J/\psi \chi _{c0}`$ production is $`S`$ wave, as expected for a single virtual-photon process. The efficiency is determined to be (28.8$`\pm `$0.7)% for the $`\eta _c(1S)`$, (31.5$`\pm `$0.7)% for the $`\chi _{c0}`$, and (28.9$`\pm `$0.8)% for the $`\eta _c(2S)`$.
The systematic error is estimated taking into account contributions from the event selection and the fitting procedure, the particle identification efficiency, and the recoil-mass scale uncertainty. The contributions from uncertainties in integrated luminosity and $`J/\psi `$ branching fraction are negligible. The contributions from individual sources (listed in Table 2) are added in quadrature, except for the systematic errors due to the mass-scale uncertainty, which are added linearly, to determine the total systematic errors.
We obtain $`\sigma (e^+e^{}J/\psi c\overline{c})(c\overline{c}>2\mathrm{charged})`$ to be $`17.6\pm 2.8_{2.1}^{+1.5}`$ fb for $`J/\psi \eta _c(1S)`$, $`10.3\pm 2.5_{1.8}^{+1.4}`$ fb for $`J/\psi \chi _{c0}`$, and $`16.4\pm 3.7_{3.0}^{+2.4}`$ fb for $`J/\psi \eta _c(2S)`$. Throughout this paper, the first error is statistical and the second systematic. Our values of the cross sections are consistent with Belle’s measurements ct:Belle-double\_ccbar-update for all three resonances. The cross sections measured by both experiments are much larger than those predicted by many NRQCD calculations.
From the fit to the recoil mass spectrum we determine the $`\eta _c(2S)`$ mass to be $`3645.0\pm 5.5_{7.8}^{+4.9}`$$`\mathrm{MeV}/c^2`$, and the total width to be $`22\pm 14`$$`\mathrm{MeV}/c^2`$. The systematic errors are mainly due to the uncertainty on the $`J/\psi `$ momentum measurement. We use ISR $`J/\psi `$ and ISR $`\psi (2S)`$ data samples to determine the momentum shifts away from the expectations for ISR events. Assuming a constant momentum shift, we obtain the recoil mass uncertainty for $`J/\psi c\overline{c}`$ processes due to the $`J/\psi `$ momentum uncertainty. The mass difference ($`\mathrm{\Delta }M`$) between the $`\eta _c(2S)`$ and $`\eta _c(1S)`$ does not significantly depend on the absolute momentum scale and common systematic errors mostly cancel. We measure $`\mathrm{\Delta }M=660.2\pm 6.8_{7.6}^{+7.1}`$$`\mathrm{MeV}/c^2`$, which is in good agreement with the mass difference previously reported by this experiment ct:BaBar-yy-etac2S and by other experiments ct:CLEO-yy-etac2S ; ct:Belle-double\_ccbar-update .
In summary, we have measured the cross section for double charmonium production $`\sigma (e^+e^{}J/\psi c\overline{c})(c\overline{c}>2\mathrm{charged})`$ for $`J/\psi \eta _c(1S)`$, $`J/\psi \chi _{c0}`$, and $`J/\psi \eta _c(2S)`$. We confirm the unexpectedly large cross sections previously reported by the Belle experiment for these processes. No evidence is found for $`e^+e^{}J/\psi J/\psi `$, $`J/\psi \chi _{c1}`$, $`J/\psi \chi _{c2}`$, or $`J/\psi \psi (2S)`$. We also measure the mass difference between the $`\eta _c(2S)`$ and the $`\eta _c(1S)`$.
We are grateful for the extraordinary contributions of our PEP-II colleagues in achieving the excellent luminosity and machine conditions that have made this work possible. The success of this project also relies critically on the expertise and dedication of the computing organizations that support BABAR. The collaborating institutions wish to thank SLAC for its support and the kind hospitality extended to them. This work is supported by the US Department of Energy and National Science Foundation, the Natural Sciences and Engineering Research Council (Canada), Institute of High Energy Physics (China), the Commissariat à l’Energie Atomique and Institut National de Physique Nucléaire et de Physique des Particules (France), the Bundesministerium für Bildung und Forschung and Deutsche Forschungsgemeinschaft (Germany), the Istituto Nazionale di Fisica Nucleare (Italy), the Foundation for Fundamental Research on Matter (The Netherlands), the Research Council of Norway, the Ministry of Science and Technology of the Russian Federation, and the Particle Physics and Astronomy Research Council (United Kingdom). Individuals have received support from CONACyT (Mexico), the A. P. Sloan Foundation, the Research Corporation, and the Alexander von Humboldt Foundation.
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# Sampling errors of correlograms with and without sample mean removal for higher-order complex white noise with arbitrary mean
## 1 Introduction
The correlogram, an estimate of the autocorrelation function (ACF) of a time series, is one of the cornerstone analyses in the signal processing toolbox and therefore an understanding of its statistical errors for various random signals is of fundamental importance. Since the first published work the sampling properties of various ACF estimators for a wide variety of different processes have been investigated . Yet it seems that one fundamental process has not draw full attention: general independent and identically distributed (IID) processes. In particular, IID processes introduce two novel aspects to the statistical errors of ACF estimators: a non-zero mean and non-analytic signals.
In order to investigate these novel effects and in particular the non-zero mean in particular, we compare the standard correlogram estimator both with and without sample mean subtracted from the data samples. These two estimators are introduced in section 2.1 and section 2.2 respectively. We conclude by discussing how to treat the mean of a process when estimating correlation.
Our main finding is that an unknown mean introduces non-zero lag dependent covariance irrespective of whether the sample mean is removed or not. This is surprising seeing that white noise signals are not non-zero dependent.
## 2 Lag windowed correlogram estimators
In dealing with an uncertain mean when estimating the autocovariance of a give sequence of data there are two possible procedures: either one estimates the mean from the given data or the mean can be guessed in some way not based on the given data. These two procedures taken together with the subsequent correlogram computation can be seen as two different types of autocovariance estimators. That is, there are autocovariance estimators that include some sort of mean estimation and there correlograms that do not. In what follows we consider one version of each, namely, we consider the standard windowed correlogram and the standard windowed correlogram with sample mean removal.
### 2.1 Estimator without sample mean removal
The classical correlogram does not involve any explicit mean estimation. Allowing for lag windowing, we define it for positive lags as
$$\widehat{R}_{zz}^{(\mu )}[l]:=𝒲[l]\underset{k=1}{\overset{Nl}{}}z^{}[k]z[k+l],l=0,1,2\mathrm{}N1,$$
(1)
where $`z[1],z[2],\mathrm{}z[N]`$ are possibly complex data samples, $`N`$ is the number of data samples, $`𝒲[l]`$ is an arbitrary real valued lag weighting function, and $`z^{}`$ denotes the complex conjugate of $`z`$. We use square brackets $`[]`$ to refer to a function of a discrete valued variable.
There are two standard choices for $`𝒲[l]`$. If
$$𝒲[l]=𝒲_U[l]:=\frac{1}{N|l|}$$
(2)
the estimators is called the traditional unbiased correlogram, while
$$𝒲[l]=𝒲_P[l]:=\frac{1}{N}$$
(3)
is known as the asymptotically unbiased correlogram. The latter is exactly the Fourier transform of the periodogram. The negative lags of the correlogram are determined from the positive lags in (1) through
$$\widehat{R}_{zz}^{(\mu )}[l]:=\left(\widehat{R}_{zz}^{(\mu )}[l]\right)^{},l=1,2,\mathrm{},N+1.$$
(4)
The correlogram $`\widehat{R}_{zz}^{(\mu )}`$ can be interpreted as an estimator of either the autocovariance sequence (ACVS) or the autocorrelation sequence (ACS) of the discrete-time, complex-valued random variable sequence $`\{Z[n];n=\mathrm{}2,1,0,1,2\mathrm{}\}`$. By the ACS of $`\{Z[l]\}`$, we mean explicitly the function
$$R_{ZZ}[l]:=\mathrm{E}\left\{Z^{}[n]Z[n+l]\right\}$$
where $`\mathrm{E}\{\}`$ represents the expectation operator and $`Z^{}`$ denotes the complex conjugate of $`Z`$; and by the ACVS of $`\{Z[n]\}`$, we mean
$$C_{ZZ}[l]:=\mathrm{Cov}\{(Z[n]\mu ),(Z[n+l]\mu )\}=\mathrm{E}\left\{(Z[n]\mu )^{}(Z[n+l]\mu )\right\}$$
see for details. Thus, the ACVS differs from ACS in that the mean $`\mu `$ has been removed from the data sequence. In other words, ACVS is equivalent to its ACS if the mean of the data sequence is zero. For continuously sampled processes, the ACS and ACVS are known as the autocorrelation function (ACF) and autocovariance function (ACVF) respectively.
Note that $`\widehat{R}_{zz}^{(\mu )}`$, for a nonzero mean $`\mu `$ signal, can also be interpreted as an autocovariance estimate in which, based on other, separate information, an assumed mean $`\mu _g`$ was subtracted from the data (or that $`\mu _g`$ was assumed 0 and nothing was done to the data) which actually had the mean $`\mu _t`$. In other words, $`\mu =\mu _t\mu _g`$ could be seen as the error in the assumed mean. If the mean $`\mu `$ of the sequence $`\{z[n]\}`$ is known exactly the signal can be converted into a zero mean sequence $`\{z^{}[n]\}:=\{z[n]\mu \}`$ by subtracting the mean. We distinguish this case of the estimator $`\widehat{R}_{zz}^{(\mu )}[l]`$ by replacing the $`(\mu )`$ index with $`(0)`$, viz $`\widehat{R}_{z^{}z^{}}^{(0)}[l]`$. Thus we can summarize formally the relationship between the autocorrelation and autocovariance estimators mentioned above as
$$\widehat{R}_{z^{}z^{}}^{(0)}[l]=\widehat{C}_{zz}[l]$$
which says that the estimator of the form (1) is an ACVS estimator if the process has a zero mean. The sampling properties of this special case estimator, $`\widehat{R}_{zz}^{(0)}[l]`$, is well documented and goes back to . The sampling properties of $`\widehat{R}_{zz}^{(\mu )}`$ with nonzero $`\mu `$ however, is what will concern us subsequently.
### 2.2 Estimator with sample mean removal
If we extend the standard correlogram introduced in the previous section to include the subtraction of the sample mean from data samples we get the ACVS estimator
$$\widehat{C}_{zz}^{(\overline{\mu })}[l]:=𝒲[l]\underset{k=1}{\overset{Nl}{}}(z[k]\overline{z})^{}(z[k+l]\overline{z}),l=0,1,2\mathrm{}N1,$$
(5)
where
$$\overline{z}:=\frac{1}{N}\underset{k=1}{\overset{N}{}}z[k]$$
is the sample mean. The negative lags can be estimated through a formula analogous to (4).
There are alternatives to simply subtracting the sample mean as in (5) when estimating the ACVS if the mean is not known, see e.g. . Here we will only consider $`\widehat{C}_{zz}^{(\overline{\mu })}`$ as it is still commonly used and of fundamental importance.
## 3 General complex higher-order white noise
The correlograms will now be applied to higher-order complex white noise with arbitrary mean. We seek to derived the sampling properties of each correlogram up to second-order properties. As it turns out, we only need to consider moments of the process up to fourth order. Thus, for our purposes it suffices to define a test sequence $`ϵ[]`$ with the following properties
$`E\{ϵ[n]\}`$ $`=\mu `$ (6)
$`\mathrm{Cov}\{ϵ^{}[n],ϵ^{}[n+l]\}=\mathrm{Cov}\{ϵ^{}[n]^{},ϵ^{}[n+l]^{}\}`$ $`=\sigma ^2\delta _{0,l}=C_{ϵϵ}[l]`$ (7)
$`\mathrm{Cov}\{ϵ^{}[n]^{},ϵ^{}[n+l]\}`$ $`=m_2\delta _{0,l}=s^2\mathrm{exp}(i\theta _2)\delta _{0,l}`$ (8)
$`E\left\{ϵ^{}[n]^{}ϵ^{}[n+l]ϵ^{}[n+l^{}]\right\}`$ $`=\kappa _3\delta _{0,l}\delta _{0,l^{}}`$ (9)
$`\mathrm{Cov}\{ϵ^{}[n]^{}ϵ^{}[n+l],ϵ^{}[n^{}]^{}ϵ^{}[n^{}+l^{}]\}`$ $`=\kappa _4\delta _{n,n^{}}\delta _{0,l}\delta _{0,l^{}}+\sigma ^4\delta _{n,n^{}}\delta _{l,l^{}}+s^4\delta _{n+l,n^{}}\delta _{l,l^{}}`$ (10)
where $`ϵ^{}:=ϵ\mu `$ is the centralised version of the process, $`\delta _{l,m}`$ is the Kronecker delta, $`\mu `$ is the mean, $`\sigma ^2`$ is the central variance, $`m_2`$ is the second central moment, $`\kappa _3`$ and $`\kappa _4`$ are the third and fourth order cumulants respectively and $`n`$, $`n^{}`$, $`l`$ and $`l^{}`$ are all arbitrary integers. The quantity $`m_2`$ and $`s^2`$ is what we will call the *quadratic variance* and the *quadratic variance amplitude* respectively. These names reflect the fact that they are not hermitian in contrast with the ordinary variance $`\sigma ^2`$.
The first property implies that the process has an arbitrary mean. The second is that the autocovariance is zero except for the zero lag. The third property is the nonhermitian quadratic autocovariance of the process which usually is either zero or equal to the autocovariance. The fourth property is a third order lagged cumulant of the process. The last property basically implies that there is no covariance between different autocovariance lags. Fifth orders and above are not specified and so $`ϵ[]`$ could have higher order correlation. Thus $`ϵ[]`$ is more general than IID processes yet is equivalent in fourth and lower order properties. See for definitions of higher order lagged cumulants of random processes.
For reference, we have the following relations in special cases: for a purely real process (hence nonanalytic),
$$\mathrm{}\{ϵ[n]\}=0s^2=\sigma ^2$$
where $`\mathrm{}\{\}`$ is the real part operator, while for an analytical process \[3, sec. 13.2.3\]
$$ϵ[]\mathrm{is}\mathrm{analytic}s=0$$
and for zero mean complex circular Gaussian white noise
$$k,ϵ[k]𝒩\{0,\sigma \}\kappa _4=\kappa _3=s=\mu =0,$$
i.e. only $`\sigma ^2`$is non-zero, while for nontrivial Poissonian real white noise
$$k,ϵ[k]\mathrm{Poi}\{\lambda \}\kappa _4=\kappa _3=\sigma ^2=\mu =\lambda ,s=0.$$
## 4 Sampling properties of the estimators to second order
We now present some of the sampling properties, namely the bias, the variance, the covariance, and the mean-square error (MSE) for the case of the general noise process $`ϵ[]`$. These quantities were derived employing the usual techniques of estimation theory and assuming that the process has the properties given in the previous section. The details of the derivations are given in a companion paper .
### 4.1 Sampling properties of correlogram without sample mean removal
The estimator $`\widehat{R}^{(\mu )}`$ was defined in (1). With respect to the $`ϵ[]`$ process, its sampling properties to second order are as follows. The bias for all lags is found to be
$$\mathrm{Bias}\left\{\widehat{R}_{ϵϵ}^{(\mu )}[l]\right\}=E\left\{\widehat{R}_{ϵϵ}^{(\mu )}[l]\right\}R_{ϵϵ}[l]=((N|l|)𝒲[l]1)(\sigma ^2\delta _{0l}+|\mu |^2)$$
(11)
This shows that the so called unbiased estimator, $`𝒲_U[l]`$, is in fact always unbiased even when the mean is not zero. All other $`𝒲[l]`$ will result biased estimators if the mean is not zero.
The covariance/variance of $`\widehat{R}^{(\mu )}`$ for any two lags $`l`$ and $`l^{}`$ of the same sign is
$`\mathrm{Cov}\{\widehat{R}_{ϵϵ}^{(\mu )}[l],\widehat{R}_{ϵϵ}^{(\mu )}[l^{}]\}`$ $`=𝒲[l]𝒲[l^{}](\delta _{0l}\delta _{0l^{}}(\kappa _4+s^4)N+2\delta _{0l}\mathrm{}\{\mu \kappa _3^{}\}(N|l^{}|)+`$
$`+2\delta _{0l^{}}\mathrm{}\{\mu \kappa _3^{}\}(N|l|)+\delta _{ll^{}}\sigma ^4(N|l|)+2|\mu |^2\sigma ^2\left(N\mathrm{max}(|l|,|l^{}|)\right)`$
$`+2\mathrm{}\{\mu ^2m_2^{}\}(N\mathrm{min}(|l|+|l^{}|,N))),ll^{}0`$ (12)
and when the lags have different signs the covariance is
$`\mathrm{Cov}\{\widehat{R}_{ϵϵ}^{(\mu )}[l],\widehat{R}_{ϵϵ}^{(\mu )}[l^{}]\}`$ $`=𝒲[l]𝒲[l^{}](\delta _{0l}\delta _{0l^{}}(\kappa _4+\sigma ^4)N+2\delta _{0l}\mathrm{}\{\mu \kappa _3^{}\}(N|l^{}|)+`$
$`+2\delta _{0l^{}}\mathrm{}\{\mu \kappa _3^{}\}(N|l|)+\delta _{l,l^{}}s^4(N|l|)+2|\mu |^2\sigma ^2\left(N\mathrm{max}(|l|,|l^{}|)\right)`$
$`+2\mathrm{}\{\mu ^2m_2\}(N\mathrm{min}(|l|+|l^{}|,N))),ll^{}0`$ (13)
These expressions can broken down into all combinations of covariances and variances for zero lags and nonzero lags. The terms without $`\delta `$ factors are the covariance of the nonzero lags of the estimator. They are zero if the mean is zero but otherwise they are non-zero and have a piece-wise linear dependence on the lags before applying the weight functions $`𝒲[l]𝒲[l^{}]`$. The $`\delta _{ll^{}}`$ term plus the terms without $`\delta `$ for $`l=l^{}0`$ constitute the variance of the nonzero lags. This is nonzero and linear in lag before weighting even when the mean is zero. The terms with single $`\delta _{0l}`$ or $`\delta _{0l^{}}`$ factors plus the terms without $`\delta `$ evaluated at either $`l=0`$ or $`l^{}=0`$, are the covariance between zero lag and non-zero lags. These are nonzero if the odd order moments $`\mu `$ and $`\mu _3`$ are nonzero. Finally, the terms with the $`\delta _{0l}\delta _{0l^{}}`$ factor plus all other terms evaluated at $`l=l^{}=0`$ is the variance of zero lag. It depends on all the moments. An example of the covariance structure is shown in figure 1a), 1b), and 1c).
The bias and the variance given above can be combined to give the mean-square error for the zero lag and the nonzero lag respectively
$`\mathrm{MSE}\left\{\widehat{R}_{ϵϵ}^{(\mu )}[0]\right\}`$ $`=(\kappa _4+s^4+4\mathrm{}\{\mu \kappa _3^{}\}+2\mathrm{}\{\mu ^2m_2^{}\}+|\mu |^4N+\sigma ^2(\sigma ^2+2|\mu |^2)(N+1))N𝒲^2[0]+`$
$`2N(\sigma ^2+|\mu |^2)^2𝒲[0]+(\sigma ^2+|\mu |^2)^2`$ (14)
$`\mathrm{MSE}\left\{\widehat{R}_{ϵϵ}^{(\mu )}[l0]\right\}`$ $`=\left(\sigma ^2(\sigma ^2+2|\mu |^2)\left(N|l|\right)+2\mathrm{}\{\mu ^2m_2^{}\}(N\mathrm{min}(2|l|,N))+|\mu |^4(N|l|)^2\right)𝒲^2[l]+`$
$`2|\mu |^4(N|l|)𝒲[l]+|\mu |^4`$ (15)
These expressions are exact for all sample sizes $`N`$. Asymptotically, that is as $`N\mathrm{}`$, the MSE behave, assuming $`\mu 0`$, as
$`\mathrm{MSE}\left\{\widehat{R}_{ϵϵ}^{(\mu )}[0]\right\}`$ $`=\left(|\mu |^2+\sigma ^2\right)^2\left(N^2𝒲^2[0]2N𝒲[0]+1\right),N1,\mu 0`$ (16)
$`\mathrm{MSE}\left\{\widehat{R}_{ϵϵ}^{(\mu )}[l0]\right\}`$ $`=|\mu |^4\left((N|l|)^2𝒲^2[l]2(N|l|)𝒲[l]+1\right),N1,\mu 0`$ (17)
where we have kept only the leading terms in $`N`$ and $`l`$ for each power of $`𝒲[l]`$.
From the asymptotic expression, we see that for the unbiased lag weights $`𝒲_U[l]`$ the leading terms in the MSE are zero. The asymptotic MSE for the nonzero lags in this case is
$$\mathrm{MSE}\left\{\widehat{R}_{ϵϵ}^{(\mu )}[l0]\right\}=\frac{\sigma ^2(\sigma ^2+2|\mu |^2)+2\mathrm{}\{\mu ^2m_2^{}\}\left(N\mathrm{min}(2|l|,N)\right)}{N|l|},𝒲𝒲_U,N1.$$
(18)
While for the periodogram weighting $`𝒲_P[l]`$, the asymptotic MSE is $`|\mu |^4l^2/N^2`$ which does not tend to zero for large lags and so it is not a consistent estimator.
This is in contrast to the well known case of zero mean. If $`\mu =0`$ then the asymptotic MSE of the nonzero lags is instead $`\sigma ^4(N|l|)𝒲^2[l]`$ for all $`N`$ as expected. The MSE of the unbiased estimator in this case tends asymptotically to $`\sigma ^4/(N|l|)`$ and so it does not converge for large lags, while the MSE for the periodogram estimator tends to $`\sigma ^4(1|l|/N)/N`$ which tends to zero for large lags. It is for this reason that the periodogram weighting function $`𝒲_P[l]`$ is preferred instead of the unbiased weighting $`𝒲_U[l]`$.
For the special case of zero mean see also .
### 4.2 Sampling properties of correlogram with sample mean removal
The sampling properties of the correlogram estimator$`\widehat{C}^{(\overline{\mu })}`$, defined in (5), for the noise process $`ϵ[]`$ were found to be as follows.
The bias is
$$\mathrm{Bias}\left\{\widehat{C}_{ϵϵ}^{(\overline{\mu })}[l]\right\}=\sigma ^2\left((N𝒲[0]1)\delta _{0l}\frac{(Nl)𝒲[l]}{N}\right).$$
(19)
From this expression we find, remarkably, that all nonzero lags are biased irrespective of the choice of weights. For instance, the weights $`𝒲_U[l]`$ known as unbiased in relation to the zero-mean correlogram $`\widehat{R}^{(0)}`$, lead to a bias of $`\sigma ^2/N`$ for all lags. It is however possible to get an unbiased estimate for the zero lag if one chooses $`𝒲[0]=1/(N1)`$. For this choice, the zero lag of $`\widehat{C}^{(\overline{\mu })}`$ is equivalent to the usual sample variance.
The covariance/variance of $`\widehat{C}^{(\overline{\mu })}`$ between any two lags $`l`$ and $`l^{}`$ of the same sign is
$`\mathrm{Cov}\{\widehat{C}_{ϵϵ}^{(\overline{\mu })}[l],\widehat{C}_{ϵϵ}^{(\overline{\mu })}[l^{}]\}`$ $`=𝒲[l]𝒲[l^{}](\delta _{0,l}\delta _{0,l^{}}(\kappa _4+s^4)N+`$
$`\delta _{0,l}\kappa _4\left(1{\displaystyle \frac{|l^{}|}{N}}\right)\delta _{0,l^{}}\kappa _4\left(1{\displaystyle \frac{|l|}{N}}\right)+\delta _{l,l^{}}\sigma ^4(N|l|)+`$
$`+{\displaystyle \frac{\kappa _4}{N}}\left(1{\displaystyle \frac{2\mathrm{max}(|l|,|l^{}|)+2\mathrm{min}(|l|+|l^{}|,N)3(|l|+|l^{}|)}{N}}3{\displaystyle \frac{ll^{}}{N^2}}\right)+`$
$`\sigma ^4\left(1{\displaystyle \frac{2\mathrm{max}(|l|,|l^{}|)|l||l^{}|}{N}}{\displaystyle \frac{ll^{}}{N^2}}\right)+`$
$`s^4(1{\displaystyle \frac{2\mathrm{min}(|l|+|l^{}|,N)|l||l^{}|}{N}}{\displaystyle \frac{ll^{}}{N^2}})),ll^{}0`$ (20)
and for lags of different signs
$`\mathrm{Cov}\{\widehat{C}_{ϵϵ}^{(\overline{\mu })}[l],\widehat{C}_{ϵϵ}^{(\overline{\mu })}[l^{}]\}`$ $`=𝒲[l]𝒲^{}[l^{}](\delta _{0,l}\delta _{0,l^{}}(\kappa _4+\sigma ^4)N+`$
$`\delta _{0,l}\kappa _4\left(1{\displaystyle \frac{|l^{}|}{N}}\right)\delta _{0,l^{}}\kappa _4\left(1{\displaystyle \frac{|l|}{N}}\right)+\delta _{|l|,|l^{}|}s^4(N|l|)+`$
$`+{\displaystyle \frac{\kappa _4}{N}}\left(1{\displaystyle \frac{2\mathrm{max}(|l|,|l^{}|)+2\mathrm{min}(|l|+|l^{}|,N)3(|l|+|l^{}|)}{N}}3{\displaystyle \frac{|ll^{}|}{N^2}}\right)+`$
$`s^4\left(1{\displaystyle \frac{2\mathrm{max}(|l|,|l^{}|)|l||l^{}|}{N}}{\displaystyle \frac{|ll^{}|}{N^2}}\right)+`$
$`\sigma ^4(1{\displaystyle \frac{2\mathrm{min}(|l|+|l^{}|,N)|l||l^{}|}{N}}{\displaystyle \frac{|ll^{}|}{N^2}})),ll^{}0.`$ (21)
Again, this expression unifies all combinations of variances and covariances between zero and nonzero lags; see the discussion in the previous subsection. The main differences with that of the estimator without mean removal is that in this case there is no dependence on the odd order moments $`\mu `$ and $`\kappa _3`$, and further, that the lag dependence before applying the weighting functions factor is quadratic in lag. An example of the covariance matrix is shown in figure 1d), 1e), and 1f).
The MSE of the $`\widehat{C}^{(\overline{\mu })}`$ estimator can be determined for the expressions for the bias and covariance given above. The result is
$`\mathrm{MSE}\left\{\widehat{C}_{ϵϵ}^{(\overline{\mu })}[0]\right\}`$ $`=\left(\sigma ^4(N1)N+(\kappa _4+s^4)Ns^42\kappa _4+{\displaystyle \frac{\kappa _4}{N}}\right)𝒲^2[0]2(N1)\sigma ^4𝒲[0]+\sigma ^4`$ (22)
$`\mathrm{MSE}\left\{\widehat{C}_{ϵϵ}^{(\overline{\mu })}[l0]\right\}`$ $`=(\sigma ^4(N|l|)(12{\displaystyle \frac{|l|}{N^2}})+{\displaystyle \frac{\kappa _4}{N}}(1{\displaystyle \frac{2|l|+2\mathrm{min}(2|l|,N)6|l|}{N}}3{\displaystyle \frac{l^2}{N^2}})+`$
$`s^4(1{\displaystyle \frac{2\mathrm{min}(2|l|,N)2|l|}{N}}{\displaystyle \frac{l^2}{N^2}}))𝒲^2[l]`$ (23)
The MSE in the asymptotic case is
$`\mathrm{MSE}\left\{\widehat{C}_{ϵϵ}^{(\overline{\mu })}[0]\right\}`$ $`=\sigma ^4\left(N^2𝒲^2[0]2N𝒲[0]+1\right),N1`$ (24)
$`\mathrm{MSE}\left\{\widehat{C}_{ϵϵ}^{(\overline{\mu })}[l0]\right\}`$ $`=\sigma ^4(N|l|)𝒲^2[l],N1`$ (25)
where we have kept only the leading terms in $`N`$ or $`l`$ for each coefficient of $`𝒲[l]`$. For the weighting sequence $`𝒲_U[l]`$, the MSE tends to zero as $`\sigma ^4/(N|l|)`$; while for $`𝒲_P[l]`$, the MSE tends to zero as $`\sigma ^4(1|l|/N)/N`$. So both weight sequences lead to consistent estimators. However, the asymptotic MSE in both cases is different for different lags.
For the special case of real valued processes see .
## 5 Comparison of the correlograms with and without sample mean removal
Now that we have the second order sampling of the correlograms with and without sample mean removal we can compare them. Examples of the covariance of the estimators are shown in figure 1 for various kinds of white noise both with and without zero means; and a comparison of the MSEs is shown in figure 2 for a zero mean noise process.
Inspection of the sampling properties derived here reveals some novel features. One such feature is that the estimators nominally have nonzero covariance, that is, the estimates at different lags are not independent. This is down to the fact that the same data samples are reused in the evaluation of different lags of the correlogram leading to a relationship between lag estimates. More surprising is the fact that the variance of both the estimators are nominally lag dependent. This is despite the fact that the test sequence $`ϵ`$ is white and hence has no nonzero lag dependence, see (7). Also the covariance is lag dependent.
### 5.1 $`\widehat{R}^{(0)}`$and $`\widehat{C}^{(\overline{\mu })}`$ as autocovariance estimators when mean is known
In the preceding sections we have assumed that the mean of the data sequence we are trying to estimate the autocovariance of is unknown. What can be said if the mean is known? It would seem natural that if we knew the mean we would subtract it from the data samples and use the estimator $`\widehat{R}^{(0)}`$. Thus we would not use the estimator $`\widehat{C}^{(\overline{\mu })}`$ as it involves estimation of the mean which we already know. Surprisingly however, inspection of the sampling properties presented here show that this choice of estimators is not immediately obvious.
The MSE of the estimator $`\widehat{C}^{(\overline{\mu })}`$ is, in fact, smaller than that of the estimator $`\widehat{R}^{(0)}`$ for general white noise, assuming the same lag window is used in both estimators. The difference in the MSE is
$$2\sigma ^4\frac{l}{N^2}(Nl)𝒲^2[l]$$
(26)
and is plotted in Figure 2. This difference tends asymptotically to zero as $`N`$ is increased, but for finite sample sizes the $`\widehat{C}^{(\overline{\mu })}`$ is always better than the $`\widehat{R}^{(0)}`$ in the MSE sense. This is counter-intuitive as it seems to violate the principle that the more one knows about something, the better one can estimate it.
The solution to this conundrum is that although the error in the covariance estimate is smaller, the error in the estimate of the mean is on the other hand larger. Specifically the difference in the MSE in the mean estimate is $`\sigma ^4/N`$. If take, e.g., the weights $`𝒲_P[l]`$ and use the inequality $`l<N`$ then the MSE for lag $`l`$ must be smaller than $`2\sigma ^4l/N^3`$ and so the total MSE of the estimator is of the order $`\sigma ^4/N`$. This is comparable to the MSE in the sample mean. Thus the total error, autocovariance and mean estimation, is not better for the $`\widehat{C}^{(\overline{\mu })}`$ compared with $`\widehat{R}^{(0)}`$.
Furthermore $`\widehat{C}^{(\overline{\mu })}`$ has a nonzero covariance between its lag estimates which $`\widehat{R}^{(0)}`$does not have. This can be understood from the following observation: that the sum over all non-negative lags of the estimator $`\widehat{R}^{(\mu )}`$ with periodogram weighting is equal to sample size times the square of the sample mean, i.e.
$$\underset{l=0}{\overset{N1}{}}\widehat{R}_{zz}^{(P)}[l]=\frac{1}{N}\underset{l=0}{\overset{N1}{}}\underset{k=1}{\overset{Nl}{}}z^{}[k]z[k+l]=N\left|\underset{k=1}{\overset{N1}{}}z[k]\right|^2=N\left|\overline{z}\right|^2$$
for any process $`\{z[k]\}`$. See for further discussions. This shows that $`\widehat{R}_{zz}^{(P)}`$ and the sample mean $`\overline{z}`$ are related. But since $`\widehat{C}^{(\overline{\mu })}`$ is based on $`\widehat{R}_{zz}^{(P)}`$ and $`\overline{z}`$ enters into every lag estimate, this suggests that there is an interdependence between lags and explain the nonzero covariance of the lags.
### 5.2 When to use the sample mean if the mean is assumed small
Let us now look at the situation when we believe that the mean is small. The question is which estimator is better: the one without or the one with sample mean removal. If we use the former we run the risk that the error in the estimated mean is large than the true mean. On the other hand, the mean may not be exactly zero so the latter estimate will also be off. There is a trade off here and we wish to find a criterion for when the sample mean should be removed.
Inspection of the results presented earlier suggest that as a rule of thumb the two estimators are roughly equal when
$$|\mu |=\frac{\sigma }{\sqrt{2N}}$$
(27)
and when $`|\mu |<\sigma /\sqrt{2N}`$ , $`\widehat{R}^{(\mu )}`$ is preferable to $`\widehat{C}^{(\overline{\mu })}`$ and when $`|\mu |>\sigma /\sqrt{2N}`$, $`\widehat{C}^{(\overline{\mu })}`$ is preferable to $`\widehat{R}^{(\mu )}`$. This condition can be understood in an intuitive way as follows. The sample mean is an estimate of the population mean with a relative error of $`(\sigma /\sqrt{N})/|\mu |`$. It makes sense to use this estimate in the correlogram instead of the unknown mean only if this relative error is smaller than 1, i.e. when $`|\mu |>\sigma /\sqrt{2N}`$. This is because the absolute error if we do not remove anything is of course only $`|\mu |`$. Thus we arrive at the condition as the found above (27).
Naturally, if we do not know the mean we will not be able asses the equality (27) exactly, but one could instead get it approximately by using the sample mean and the sample variance of the given data samples.
## 6 Conclusion
We have presented expressions for the bias, variance, covariance, mean-square error of the classical correlogram estimator with and without sample mean estimation for general complex white noise with arbitrary mean. A summary of the sampling properties of the estimators for $`ϵ[]`$, i.e., complex higher order white noise with arbitrary mean is as follows. The second-order sampling properties of the correlogram without sample mean removal, $`\widehat{R}^{(\mu )}`$, are in summary:
* $`ϵ[]`$ moment dependence: cumulants up to fourth order $`\{\mu ,\sigma ^2,s^2,\kappa _3,\kappa _4\}`$
* Bias: unbiased for $`1/(Nl)`$ weighting, even when $`\mu 0`$
* Covariance: piecewise linear lag dependent covariance before weighting
* Mean square error: is asymptotically proportional to $`|\mu |^4`$ when $`\mu 0`$ while if $`\mu =0`$ the MSE is asymptotically equal to $`\sigma ^4(N|l|)𝒲^2[l]`$
For the correlogram with sample mean removal, $`\widehat{C}^{(\overline{\mu })}`$, a summary of second-order sampling properties is:
* $`ϵ[]`$ moment dependence: only even order cumulants $`\{\sigma ^2,s^2,\kappa _4\}`$
* Bias: nonzero lags biased for all weighting functions, zero lag unbiased for $`𝒲[0]=1/(N1)`$
* Covariance: piecewise quadratic lag dependent covariance before weighting
* Mean square error: asymptotically equal to $`\sigma ^4(N|l|)𝒲^2[l]`$
In terms of the properties of the process, we have found that a nonzero mean can lead to covariance between lags of the correlogram, and that complex valued processes lift the degeneracy between positive and positive lags, (that is that the sampling properties of positive lags and negative lags are different), and distinguishes between analytic and non-analytic processes.
For both estimators, the both the variance and the covariance are in general lag dependent. We conclude therefore that neither $`\widehat{R}^{(\mu )}`$ nor $`\widehat{C}^{(\overline{\mu })}`$ with the standard lag weights $`𝒲_U[l]`$ and $`𝒲_P[l]`$ are optimal for estimating the autocorrelation/autocovariance of noise processes with an unknown mean.
Finally, we found that the sampling properties suggest that one should always remove the mean if it is known apriori. If the mean is not known, the sample mean should be removed only when $`|\mu |>\sigma /\sqrt{2N}`$, i.e. $`\widehat{R}_{ϵϵ}^{(\mu )}`$ is preferable to $`\widehat{C}_{ϵϵ}^{(\overline{\mu })}`$ if $`|\mu |<\sigma /\sqrt{2N}`$.
## Acknowledgments
This work was sponsored by PPARC ref: PPA/G/S/1999/00466 and PPA/G/S/2000/00058.
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# Enhancement of spin stiffness with dilution in a ferromagnetic Kondo lattice model
## I Introduction
The recent discovery of ferromagnetism in diluted magnetic semiconductors (DMS) such as $`\mathrm{Ga}_{1\mathrm{x}}\mathrm{Mn}_\mathrm{x}\mathrm{As}`$,Ohno2 ; Matsu ; Pot ; Edmonds with transition temperature $`T_c110`$ K for Mn concentration $`x5\%`$,Matsu and $`150`$ K in films for $`x`$ in the range $`6.78.5\%`$,Pot ; Edmonds has led to intensive efforts to increase $`T_c`$ in view of potential technological applications. Tremendous interest has also been generated in the novel ferromagnetism exhibited by these systems in which magnetic interaction between localized spins is mediated by doped carriers. The long-range oscillatory nature of the carrier-mediated spin couplings results in a variety of interesting features of the ferromagnetic state, such as significant sensitivity of spin stiffness and transition temperature $`T_c`$ on carrier concentration, competing antiferromagnetic interaction and noncollinear ordering, spin-glass behaviour, spin clustering and disorder-induced localization etc., as recently reviewed.bhatt ; Timm
DMS such as $`\mathrm{Ga}_{1\mathrm{x}}\mathrm{Mn}_\mathrm{x}\mathrm{As}`$ are mixed spin-fermion systems in which the $`S=5/2`$ Mn<sup>++</sup> impurities replace Ga<sup>+++</sup>, thereby contributing a hole to the semiconductor valence band. However, large compensation due to As antisite defects reduces the hole density to nearly $`10\%`$ of Mn concentration, which plays a key role in the stabilization of long-range ferromagnetic order, and also provides a complimentary limit to Kondo systems. The interplay between itinerant carriers in a partially filled band and the localized moments is conventionally studied within a diluted ferromagnetic Kondo lattice model (FKLM), wherein $`J𝐒_I.𝝈_I`$ represents the exchange interaction between the magnetic impurity spin $`𝐒_I`$ and the electron spin $`𝝈_I`$.
In order to obtain deeper understanding of the role of dilution in terms of the perturbation to carrier spin bands due to spatially varying impurity field as well as impurity-field-induced Zeeman splitting, which control the long-range oscillatory nature of the impurity-spin couplings through the carrier spin polarization $`\chi ^0(𝐪,\omega )`$, we will consider a diluted FKLM with ordered impurity arrangements in this paper. Use of k-space representation allows for much larger lattices, thus permitting a more refined study of the competing spin interactions with respect to carrier concentration, impurity separation, interaction strength, and wave vector. More importantly, it provides for a quantitative comparison with different theoretical pictures, which we briefly review below.
Long range ferromagnetic interaction between the impurity spins is mediated, in the mean-field (Zener model) picture, Dietl1 ; Taka1 ; Jung1 ; Dietl2 ; Jung2 ; Dietl3 by a uniform carrier spin polarization, which is caused, in turn, by a uniform site-averaged impurity magnetic field within the virtual crystal approximation (VCA). In the weak-field limit ($`xJS<<ϵ_\mathrm{F}`$), the carrier spin polarization is proportional to the Pauli susceptibility $`\chi _\mathrm{P}`$, and the transition temperature ($`T_cxJ^2\chi _\mathrm{P}`$) is therefore proportional to the Mn concentration $`x`$, $`J^2`$, the carrier effective mass $`m^{}`$, and $`N(ϵ_\mathrm{F})p^{1/3}`$, where $`p`$ is the hole concentration.
An alternative mechanism for the coupling between impurity spins involves the carrier-mediated Ruderman-Kittel-Kasuya-Yosida (RKKY) interaction.Matsu The local magnetic field $`𝐁_J=J𝐒_J`$ of an impurity spin at site $`J`$ polarizes the electrons locally, and the mobile band electrons spread this magnetic polarization in a characteristic manner: $`𝐦_I=\chi _{IJ}^0𝐁_J`$, where $`\chi _{IJ}^0`$ represents the carrier spin polarization. Another impurity spin $`\stackrel{}{S}_I`$ couples to this local electronic magnetization, resulting in an effective RKKY interaction $`J^2\chi _{IJ}^0𝐒_I.𝐒_J`$. Within a mean-field (MF) treatment of the resulting Heisenberg model,Dietl1 a similar behaviour for the transition temperature ($`T_cxJ^2\chi _\mathrm{P}`$) is obtained.
Within the generalized RKKY theory involving the non-linear magnetic response,dms the carrier-induced spin couplings $`J^2\chi _{IJ}^0(J)`$ go through a maximum with respect to both carrier doping concentration $`p`$ and the fermion-spin interaction strength $`J`$. This optimization behaviour can be qualitatively understood in terms of a competition between the increasing magnitude of the carrier spin polarization $`\chi _{IJ}^0(J)`$ and the increasing rapidity of its oscillation, which limits the growth of the spin couplings. Similar behaviour was observed in the diluted Hubbard modelSingh2 ; Pandey for the effective spin coupling $`U^2\chi _{IJ}^0(U)`$.
Long-wavelength magnon excitations provide a composite measure of the carrier-induced spin couplings in the ferromagnetic state, with vansihing spin stiffness signalling instability due to competing antiferromagnetic (AF) spin interactions. Magnon excitations as function of electron density $`n`$ in the conduction band and the spin-fermion coupling $`J`$ have been studied within the concentrated FKLM having a magnetic impurity at every lattice site, in the context of heavy fermion materials,Sig ferromagnetic metals Gd, Tb, Dy, doped EuXDonath and manganites.Furu ; Wang ; Yunoki ; Vogt Magnon dispersion has also been obtained in the context of DMS, where the spatially varying impurity field has been treated exactly numerically,Berciu2 within the VCA where a uniform impurity-induced Zeeman splitting of the carrier spin bands is assumed,Konig2 and within the coherent potential approximation (CPA).Bouzerar ; Nol2 Magnon spectrum and transition temperature have also been obtained recently for $`\mathrm{Ga}_{1\mathrm{x}}\mathrm{Mn}_\mathrm{x}\mathrm{As}`$ and $`\mathrm{Ga}_{1\mathrm{x}}\mathrm{Mn}_\mathrm{x}\mathrm{N}`$ by means of an effective Heisenberg model, whose exchange parameters are obtained from first-principle calculations.hilbert\_nolting
A sublattice-basis representation for ordered impurity arrangements is introduced in section II for two-dimensional (2d) and three dimensional (3d) cases, and the formation of sub-bands due to spatially varying impurity field is discussed. The two energy scales in the coupled spin-fermion problem are also introduced, and temperature-dependence of impurity magnetization is obtained self-consistently. Spin-wave excitations are obtained using the Holstein-Primakoff transformation in the large $`S`$ limit (section III). A detailed study of the behaviour of spin stiffness with respect to dimensionality, carrier concentration, interaction strength as well as impurity separation (dilution) is presented in section IV. Ferromagnetic transition temperature obtained within the renormalized spin-fluctuation theory is compared with RKKY and VCA results in section V.
## II Diluted Kondo lattice model
We consider a diluted ferromagnetic Kondo lattice Hamiltonian
$`H=t{\displaystyle \underset{i,\delta ,\sigma }{}}a_{i,\sigma }^{}a_{i+\delta ,\sigma }{\displaystyle \frac{J}{2}}{\displaystyle \underset{I}{}}𝐒_I.𝝈_I`$ (1)
on square and cubic lattices with positive hopping $`t`$ between nearest-neighbour (NN) sites $`i,i+\delta `$, which yields a host valence band. The second term represents the ferromagnetic exchange interaction between localized impurity spins at site $`I`$ and the fermion spin density $`𝝈_I=\mathrm{\Psi }_I^{}[𝝈]\mathrm{\Psi }_I`$, where the fermionic field operator $`\mathrm{\Psi }_I=\left(\begin{array}{c}a_I\\ a_I\end{array}\right)`$. The host on-site energy is set to zero.
### II.1 Sublattice-basis representation
We consider ordered (superlattice) arrangements of magnetic impurities on square and cubic host lattices with impurity spacings $`2a`$ and $`3a`$. Translational symmetry in the sublattice basis convenientally allows Fourier transformation to momentum space. For concreteness, we consider a square host lattice in the following, with magnetic impurities placed at every other host site, corresponding to superlattice spacing $`2a`$ and impurity concentration $`x=25\%`$. There are four sublattices, numbered $`\alpha =1,2,3,4`$, corresponding to the four sites in the unit cell, as shown in Fig. 1. We choose length and energy units such that the lattice spacing $`a=1`$ and the hopping term $`t=1`$.
In the MF approximation, the interaction term
$`H_{MF}^{\mathrm{int}}={\displaystyle \frac{J}{2}}{\displaystyle \underset{I}{}}𝐒_I.𝝈_I+𝝈_I.𝐒_I`$ (2)
represents a magnetic coupling of impurity and fermion spins with the self-consistently determined fermion and impurity fields $`\frac{J}{2}𝝈_I`$ and $`\frac{J}{2}𝐒_I𝐁_I`$, respectively. Assuming, without loss of generality, a uniform impurity field $`𝐁_I`$ = $`B\widehat{z}`$, Fourier transformation within the sublattice basis yields the following MF Hamiltonian for spin-$`\sigma `$ fermion
$$H_{\mathrm{MF}}^\sigma (𝐤)=\underset{𝐤}{}\mathrm{\Psi }_{𝐤\sigma }^{}\left[\begin{array}{cccc}\hfill \sigma B& ϵ_𝐤^x& 0& ϵ_𝐤^y\\ \hfill ϵ_𝐤^x& 0& ϵ_𝐤^y& 0\\ \hfill 0& ϵ_𝐤^y& 0& ϵ_𝐤^x\\ \hfill ϵ_𝐤^y& 0& ϵ_𝐤^x& 0\end{array}\right]\mathrm{\Psi }_{𝐤\sigma }$$
(3)
where $`ϵ_𝐤^{x(y)}=2t\mathrm{cos}k_{x(y)}`$ is the hopping energy along the $`x(y)`$ direction and $`\sigma =`$+1 and $`1`$ for spin up and down, respectively. Here the field operator $`\mathrm{\Psi }_𝐤=(a_𝐤^1a_𝐤^2a_𝐤^3a_𝐤^4)`$ defines the sublattice basis, where $`a_𝐤^\alpha `$ refer to the fermion operator for sublattice index $`\alpha `$. The MF Hamiltonian is numerically diagonalized to obtain the four eigenvalues $`E_{𝐤\mu }^\sigma `$, corresponding to the four sub-band indices $`\mu =1,2,3,4`$, and the four-component eigenvectors $`\varphi _{𝐤\mu }^{\sigma \alpha }`$ yield the fermion amplitudes on the sublattice $`\alpha `$.
### II.2 Quasiparticle density of states
The quasiparticle density of states (Fig. 2) shows the formation of four and eight sub-bands for the square and cubic lattices, respectively, corresponding to the number of sublattices. The reflection symmetry about the origin of the energy axis reflects the magnetic coupling. The sub-band corresponding to impurity sublattice splits off and forms an impurity band at energy $`\pm B`$, which narrows with increasing $`B`$ due to decreasing effective hopping $`t_{\mathrm{eff}}t^2/B`$ between impurity sites. In the undoped case, both bands are completely filled, and carrier doping introduces holes at top of the spin-$``$ impurity band. With further hole doping, Fermi energy moves into the spin-$``$ band, particularly for small $`B`$.
### II.3 Impurity magnetization
With only one energy scale, the NN Heisenberg ferromagnet yields a transition temperature which differs from the MF value only by an O(1) number. However, with long-range ferromagnetic order originating from the carrier-induced, oscillating impurity-spin couplings, the transition temperature for DMS is very different from the MF value, which is determined from the bare spin-fermion exchange interaction $`J`$, and physically represents a moment-melting temperature rather than a spin-disordering temperature. In order to highlight the order of magnitude difference arising from competing interactions and spin softening, we have also obtained $`T_c^{\mathrm{MF}}`$ in the MF approximation, which is interesting in its own right, and studied recently for the $`x=1`$ case.Das2
In the MF approximation, the coupled spin-fermion problem is self-consistently solved from the finite-temperature expectation values
$`\sigma _I^z`$ $`=`$ $`{\displaystyle \underset{𝐤,\mu }{}}[(\varphi _{𝐤\mu }^{\alpha =1})^2(\varphi _{𝐤\mu }^{\alpha =1})^2]f_{\mathrm{FD}}(E_{𝐤\mu }^\sigma )`$ (4)
$`S_I^z`$ $`=`$ $`{\displaystyle \underset{m=S}{\overset{+S}{}}}me^{\beta mJ\sigma _I^z/2}/{\displaystyle \underset{m=S}{\overset{+S}{}}}e^{\beta mJ\sigma _I^z/2}`$ (5)
of fermion and impurity spins, where $`f_{\mathrm{FD}}`$ is the Fermi-Dirac distribution function and $`\beta 1/k_\mathrm{B}T`$.
Consider the $`x=1`$ case, for simplicity. There are two distinct temperature scales associated with the fermion field $`\frac{J}{2}\sigma _I^z\frac{J}{2}p`$ and the impurity field $`\frac{J}{2}S_I^z\frac{J}{2}S`$, resulting in a characteristic concave $`S_I^z`$ vs. $`T`$ behaviour in the intermediate temperature regime $`Jpk_\mathrm{B}TJS`$, wherein the nearly free impurity spins yield a typical paramagnetic $`(1/T)`$ response to impurity magnetization. In the diluted case, minority-spin holes preferentially sit on impurity sites, resulting in an enhanced fermion magnetization $`\sigma _I^zp/x`$, which blurs the distinction between the two temperature scales. Consequently, the concave behaviour in the temperature dependence of impurity magnetization \[Fig. 3\] becomes prominent only at low doping concentration $`p`$ and large $`J`$. With increasing temperature and decreasing impurity field, the rapid decrease in $`\sigma _I^z`$ due to redistribution of holes from impurity to host sites results in a first-order transition, in contrast to the continuous transition obtained for the $`x=1`$ case.
## III Spin-wave excitations
Transverse spin fluctuations are gapless, low-energy excitations in the broken-symmetry state of magnetic systems possessing continuous spin-rotational symmetry. Therefore, at low temperature they play an important role in diverse macroscopic properties such as existence of long-range order, magnitude and temperature dependence of the order parameter, magnetic transition temperature, spin correlations etc. In the following we consider $`T=0`$ and obtain the spin-wave excitations at the non-interacting level. Finite impurity concentration, interaction strength $`J`$, and dynamical effects are treated exactly in this approach.
Applying the approximate Holstein-Primakoff transformation in the large $`S`$ limit
$`S_I^+`$ $`=`$ $`b_I\sqrt{2S}`$
$`S_I^{}`$ $`=`$ $`b_I^{}\sqrt{2S}`$
$`S_I^z`$ $`=`$ $`Sb_I^{}b_I`$ (6)
from the spin-lowering ($`S_I^{}`$) and spin-raising ($`S_I^+`$) operators to boson (magnon) creation and annihilation operators $`b_I^{}`$ and $`b_I`$, which neglects magnon interaction terms of order $`1/2S`$, the Kondo-lattice Hamiltonian reduces to
$`H`$ $`=`$ $`t{\displaystyle \underset{i,\delta ,\sigma }{}}a_{i,\sigma }^{}a_{i+\delta ,\sigma }`$ (7)
$``$ $`{\displaystyle \frac{J}{2}}{\displaystyle \underset{I}{}}\left[{\displaystyle \frac{\sqrt{2S}}{2}}\left(b_I\sigma _I^{}+b_I^{}\sigma _I^+\right)+\left(Sb_I^{}b_I\right)\sigma _I^z\right],`$
where $`\sigma _I^\pm \sigma _I^x\pm i\sigma _I^y`$.
In the MF approximation, $`b_I^{}=b_I=b_I^{}b_I=0`$, and the Hamiltonian decouples into a fermion part
$$_{\mathrm{fermion}}^0=t\underset{i,\delta ,\sigma }{}a_{i,\sigma }^{}a_{i+\delta ,\sigma }\frac{JS}{2}\underset{I}{}\sigma _I^z$$
(8)
and a boson part
$$_{\mathrm{boson}}^0=\frac{J}{2}\underset{I}{}\sigma _I^zb_I^{}b_I.$$
(9)
We now obtain the time-ordered, transverse spin-fluctuation (magnon) propagator
$`\chi _{IJ}^+(tt^{})=i\mathrm{\Psi }_\mathrm{G}T[b_I(t)b_J^{}(t^{})]\mathrm{\Psi }_\mathrm{G}`$ (10)
within the random phase approximation (RPA), by summing over all bubble diagrams
where the particle-hole bubble in $`𝐪,\omega `$ space
$`\chi ^0(𝐪,\omega )`$ $`=`$ $`i{\displaystyle \frac{d\omega ^{}}{2\pi }G^{}(𝐤,\omega ^{})G^{}(𝐤𝐪,\omega ^{}\omega )}`$ (11)
$`=`$ $`{\displaystyle \underset{E_{𝐤\nu }^{}<E_F}{\overset{E_{𝐤𝐪\mu }^{}>E_F}{}}}{\displaystyle \frac{|\varphi _{𝐤\nu }^{\alpha =1}|^2|\varphi _{𝐤𝐪\mu }^{\alpha =1}|^2}{E_{𝐤𝐪\mu }^{}E_{𝐤\nu }^{}+\omega }}`$
$`+`$ $`{\displaystyle \underset{E_{𝐤\nu }^{}>E_F}{\overset{E_{𝐤𝐪\mu }^{}<E_F}{}}}{\displaystyle \frac{|\varphi _{k\nu }^{\alpha =1}|^2|\varphi _{𝐤𝐪\mu }^{\alpha =1}|^2}{E_{𝐤\nu }^{}E_{𝐤𝐪\mu }^{}\omega }}`$
involves integrating out the fermions in the broken-symmetry state. It is the particle-hole bubble $`\chi ^0(𝐪,\omega )`$ which mediates the carrier-induced impurity spin couplings in the ferromagnetic state, and the oscillatory, long-range nature of the spin couplings is effectively controlled by the Fermi wavevector ($`k_Fp^{1/d})`$ and the impurity-field-induced Zeeman splitting of the carrier spin bands.
The magnon propagator can then be expressed as
$`\chi ^+(𝐪,\omega )`$ $`=`$ $`{\displaystyle \frac{g^0(\omega )}{1+\frac{J^2S}{2}g^0(\omega )\chi ^0(𝐪,\omega )}}`$ (12)
$`=`$ $`{\displaystyle \frac{1}{\omega (_0_𝐪(\omega ))}}`$
where the site-diagonal zeroth-order magnon propagator
$`g^0(\omega )={\displaystyle \frac{1}{\omega _{\mathrm{boson}}^0}}={\displaystyle \frac{1}{\omega _0}}`$ (13)
involves the magnon on-site energy $`_0\frac{J}{2}\sigma _I^z`$, corresponding to the energy cost of a spin deviation, and $`_𝐪(\omega )=\frac{J^2S}{2}\chi ^0(𝐪,\omega )`$ is the delocalization energy due to magnon hopping associated with spin couplings $`𝒥_{IJ}=J^2\chi _{IJ}^0(\omega )`$. The energy cost $`_0`$ of creating a local spin deviation is exactly offset by this delocalization-induced energy gain for $`q,\omega =0`$, consistent with the Goldstone mode. Poles in the magnon propagator yield the magnon-mode energies by solving
$`\omega _𝐪=_0_𝐪(\omega _𝐪).`$ (14)
The spin-wave dispersion along the main symmetry directions is shown in Fig. 4 for the 2d and 3d cases.
## IV Spin stiffness
The spin stiffness $`D=\omega _𝐪/q^2`$ for long-wavelength modes provides a composite measure of the impurity spin couplings and stability of the ferromagnetic state. Variation of spin stiffness with doping concentration $`p`$ and interaction strength $`J`$ is shown in Figs. 5 and 6 for the 2d and 3d cases. The spin stiffness exhibits optimization behaviour with respect to both $`p`$ and $`J`$, which can be understood within the generalized RKKY theory in terms of a competition between increasing magnitude of carrier-spin polarization and increasing rapidity of its oscillation.dms An interesting feature of Figs. 5(a) and 6 is that due to impurity-host decoupling for higher $`J`$ values, the instability of the ferromagnetic state shifts towards ”half-filling” ($`px`$) where AF ordering is favoured. For a fixed $`p`$, a minimum $`J`$ is required to stabilize the ferromagnetic state, as seen in Fig. 5(b), and the $`1/J`$-type behaviour for large $`J`$ is associated with narrowing of the impurity band due to decreasing effective impurity hopping $`t_{\mathrm{eff}}(t^2/J)`$.
Fig. 7 highlights the remarkable effect of dilution on spin stiffness — dilution actually enhances the spin stiffness at low hole doping, demonstrating an optimization with respect to impurity concentration as well. The above result highlights the key requirements of impurity spin dilution as well as low hole doping concentration for a robust ferromagnetic state, both of which are met in the heavily compensated material $`\mathrm{Ga}_{1\mathrm{x}}\mathrm{Mn}_\mathrm{x}\mathrm{As}`$, with $`x5\%`$. It is the reduced particle-hole energy gap in the fermion polarization bubble $`\chi ^0(𝐪,\omega )`$ for the diluted case which relatively enhances the oscillating carrier-spin polarization, resulting in greater spin stiffness for low doping and higher sensitivity to competing interactions with increasing doping. Indeed, the limit of vanishing energy gap (RKKY theory) yields an upper bound to the spin stiffness and transition temperature (see Fig. 9).
Fig. 8 shows the enhancement of spin stiffness with dilution at low doping for the 3d case. Also shown is a comparison of the spin stiffness with the VCA result for a uniform impurity field $`B=xJS/2`$ with identical lattice-average. The VCA result provides a fair approximation in the low doping regime, which can be understood in terms of an averaged impurity field seen by doped carriers with wavelength longer than the impurity spacing.
## V Transition Temperature
Determination of the magnon spectrum in the carrier-induced ferromagnetic state allows for an estimation of the Curie temperature $`T_c`$ in three dimensions. As the ferromagnetic state is characterized by small spin stiffness due to competing interactions, the dominant contribution to reduction in magnetization is from the thermal excitation of long-wavelength magnon modes. Therefore, $`T_c`$ can be estimated in terms of an equivalent Heisenberg model with matching spin stiffness. For a spin-$`S`$ Heisenberg ferromagnet with nearest-neighbour interaction $`\stackrel{~}{J}`$ on a simple cubic lattice (coordination number $`z`$=6), the magnon energy
$`\omega _𝐪=z\stackrel{~}{J}S(1\gamma _𝐪),`$ (15)
where $`\gamma _𝐪`$=(cos2$`q_x+`$cos2$`q_y+`$cos2$`q_z`$)/3, corresponding to the magnetic lattice spacing 2. Indeed, the magnon energy is maximum for $`q_x=q_y=q_z=\pi /2`$, as also seen in Fig. 4 (at R). Considering the small $`q`$ limit, we obtain the spin stiffness $`D=\omega _𝐪/q^2=4\stackrel{~}{J}S=10\stackrel{~}{J}`$ for $`S=5/2`$.
Within the renormalized spin-fluctuation theory, the transition temperature
$`T_c^{\mathrm{SF}}={\displaystyle \frac{z\stackrel{~}{J}S(S+1)}{3}}f_{\mathrm{SF}}^1={\displaystyle \frac{1}{2}}D(S+1)f_{\mathrm{SF}}^1`$ (16)
where the spin-fluctuation factor $`f_{\mathrm{SF}}=_𝐪(1\gamma _𝐪)^1`$ 1.5 for the cubic lattice. As $`T_c^{\mathrm{SF}}T_c^{\mathrm{MF}}`$, the impurity mean field is essentially unchanged at low temperature; therefore we have assumed the zero-temperature value for the spin stiffness. In the RKKY-VCA approach, the lattice-averaged impurity field is assumed to decrease with impurity magnetization and vanishes as $`TT_c`$.Konig2 ; Bouzerar Considering a bandwidth of the order of that of $`\mathrm{GaAs}`$ ($`W=12t10`$ eV), we obtain $`T_c^{\mathrm{SF}}=(35/36)D`$ in eV. The calculated $`T_c^{\mathrm{SF}}`$ (Fig. 9) is in qualitative agreement with experimental results for Ga<sub>1-x</sub>Mn<sub>x</sub>As.
Figure 9 also provides a comparison with the RKKY resultDietl1
$`T_c^{\mathrm{RKKY}}=xJ^2{\displaystyle \frac{S(S+1)}{12k_\mathrm{B}}}N(E_F),`$ (17)
obtained within the mean-field approximation from the carrier-induced RKKY spin couplings $`𝒥_{IJ}=J^2\chi _{IJ}^0`$. In the limit of low doping concentration ($`p1`$), the host density of states
$`N(E_F)={\displaystyle \frac{1}{2\pi ^2}}(3\pi ^2p)^{1/3},`$ (18)
yielding the characteristic $`p^{1/3}`$ behaviour of $`T_c^{\mathrm{RKKY}}`$. It is seen in Fig. 9 that the RKKY result drastically over-estimates $`T_c`$, mainly due to neglect of the impurity-field-induced energy gap in the fermion polarization bubble, and also fails to capture the optimization behaviour with carrier doping, impurity spin dilution, or interaction strength. The MF transition temperature $`T_c^{\mathrm{MF}}`$, which physically represents the moment-melting temperature, is an order-of-magnitude still higher.
## VI Conclusions
Magnon excitations in the carrier-induced ferromagnetic state of the diluted FKLM were studied for several ordered impurity arrangements. Use of momentum-space representation allowed for a refined study with respect to variation of doping concentration, interaction strength, impurity spin dilution, and wave vector. The spin stiffness was found to exhibit a characteristic optimization behaviour with respect to hole doping concentration, interaction strength, as well as impurity spin dilution, which can be qualitatively understood in terms of a competition between increasing magnitude of the carrier-spin polarization and increasing rapidity of its oscillation.
The exact treatment of impurity spin dilution and interaction strength allowed for a detailed comparison with results of existing theories such as MF, RKKY, and VCA. We found that the RKKY result drastically over-estimates $`T_c`$ due to neglect of the impurity-field-induced Zeeman splitting of the carrier spin bands, and also does not capture the optimization behaviour with carrier doping, impurity spin dilution, or interaction strength. The VCA approach involving a uniform lattice-averaged impurity field was found to provide a fair approximation for spin stiffness in the low doping regime, but fails to describe the competing interactions at higher doping when the hole wavelength becomes comparable to impurity spacing.
The enhancement of spin stiffness with dilution at low doping concentration due to the reduced impurity-field-induced energy gap in the fermion polarization bubble highlights the key requirements of impurity spin dilution as well as low hole doping concentration for a robust ferromagnetic state of diluted magnetic semiconductors, both of which are realized in the heavily compensated material $`\mathrm{Ga}_{1\mathrm{x}}\mathrm{Mn}_\mathrm{x}\mathrm{As}`$, with $`x5\%`$.
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# Correlations in interacting systems with a network topology
## I Introduction.
The generic features of real-world networks (the Internet, the WWW, biological, social and economical networks and many others) are a complex organization of their connections and the small-world phenomenon s01 ; ab01a ; dm01c ; n03 ; w99 ; dm03 ; pvbook04 ; ajb99 . In the networks with the small-world effect, a mean intervertex distance grows with a total number of vertices, $`N`$, slower than any power of $`N`$, e.g., grows as $`\mathrm{ln}N`$. Concerning economical and social networks, the small-world property is often considered as an evidence for growing close interrelations and globalization.
Many of the real-world networks are formed by interacting objects and demonstrate complicated dynamics. The dependence of the correlations between interacting objects in systems with a network topology on time and distance provides a useful information about the network dynamics. In the present paper we discuss general properties of the correlations between a pair of interacting objects on a complex network. Recall that in interacting systems on lattices, with a few exceptions, pair correlations decrease exponentially with distance apart from a critical point, where the decrease is power-law. From a naive point of view, one might expect that the small-world property of a network would enhance pair correlations between distant objects in comparison to lattices. However, it is not the case. We demonstrate that correlations between $`\mathrm{}`$-th nearest neighbors, on average, decay with $`\mathrm{}`$ as $`1/(\mathrm{}z_{\mathrm{}})`$ or faster. In networks, where the mean intervertex separation $`\overline{\mathrm{}}(N)\mathrm{ln}N`$, this corresponds to the exponential decay of correlations with $`\mathrm{}`$ (even at the critical point). However, in networks with a divergent mean number of the second nearest neighbors \[where $`\overline{\mathrm{}}(N)`$ grows slower than $`\mathrm{ln}N`$\], we observe a dramatic weakening of pair correlations between the second and more distant nearest neighbors. In these networks, only the nearest neighbors are strongly correlated, while the correlations between more distant vertices are suppressed and approach zero in the infinite network limit.
In Section II we consider pair correlations in an interacting system defined on the top of a complex network in the framework of a phenomenological approach. This approach allows us to understand general properties of pair correlations irrespective of details of an interacting system and the nature of interactions. In Section III we support these results by calculations of the static pair correlation function of the Ising model on a complex network (more precisely, the configuration model of a network bbk72 ).
## II Phenomenological approach
Let a quantity $`X_i(t)`$ describe a dynamic process on a network, where indices $`i`$ label vertices, and $`t`$ is time. $`x_i(t)X_i(t)X_i`$ describes fluctuations around the average value $`X_i`$. If the system under consideration is in an equilibrium state, then pair correlations between two arbitrary vertices $`i`$ and $`j`$ may be characterized by the following correlation function:
$`G_{ij}(t_1,t_2)`$ $`=`$ $`t_0^1{\displaystyle \underset{0}{\overset{t_0}{}}}x_i(t+t_1)x_j(t+t_2)𝑑t`$ (1)
$``$ $`x_i(t_1)x_j(t_2).`$
In the present section the brackets $`\mathrm{}`$ mean the average over the observation time $`t_0`$ which must be much larger than the maximum relaxation time in the system under consideration. In the equilibrium state we have $`G_{ij}(t_1,t_2)=G_{ij}(t_1t_2)`$ in the limit $`t_0\mathrm{}`$. Assuming the Hamiltonian dynamics, one can introduce a generalized field $`H_i(t)`$ conjugated with $`X_i(t)`$. The function
$`\chi _{ij}(t_1t_2)`$ $`=`$ $`t_0^1{\displaystyle \underset{0}{\overset{t_0}{}}}x_i(t+t_1)/H_j(t+t_2)dt`$ (2)
$``$ $`x_i(t_1)/H_j(t_2)`$
is a generalized non-local susceptibility which characterizes the averaged response of $`x_i(t_1)`$ at time $`t_1`$ on a field applied at a vertex $`j`$ at time $`t_2`$. There is a simple relationship between $`G_{ij}(t)`$ and $`\chi _{ij}(t)`$:
$$G_{ij}(t)\chi _{ij}(t),$$
(3)
where coefficient of the proportionality is unimportant for our purpose.
In general case, the total susceptibility $`\chi (t)`$ of a system (per vertex) is finite in the limit $`N\mathrm{}`$ where $`N`$ is the number of vertices in a network:
$$\chi (t_1t_2)=\frac{1}{N}\underset{i=1}{\overset{N}{}}\frac{x_i(t_1)}{H_j(t_2)}|_{H_1=H_2=\mathrm{}=H}=O(1).$$
(4)
Here, a statement $`A=O(1)`$ means that a quantity $`A`$ is finite in the limit $`N\mathrm{}`$. The susceptibility $`\chi (t)`$ may diverge only at a critical point. We rewrite Eq. (4) as follows:
$$\chi (t)=\frac{1}{N}\underset{i}{}\left(\chi _{ii}(t)+\underset{j:\mathrm{}_{ij}=1}{}\chi _{ij}(t)+\underset{j:\mathrm{}_{ij}=2}{}\chi _{ij}(t)+\mathrm{}\right),$$
(5)
where the first sum in the parentheses is a sum over the nearest neighbors of avertex $`i`$ being at the distance $`\mathrm{}_{ij}=1`$. The second sum is over thesecond nearest neighbors, $`\mathrm{}_{ij}=2`$, and so on. One can introduce the average value $`\chi (t,\mathrm{})`$, that is the average value of the susceptibility $`\chi _{ij}(t)`$ at $`\mathrm{}_{ij}=\mathrm{}`$ over all possible network configurations. In the large network limit, this quantity coincides with the following expression calculated for a single network:
$$\chi (t,\mathrm{})=\frac{_{i,j:\mathrm{}_{ij}=\mathrm{}}\chi _{ij}(t)}{_{i,j:\mathrm{}_{ij}=\mathrm{}}1}=\frac{_{i,j:\mathrm{}_{ij}=\mathrm{}}\chi _{ij}(t)}{Nz_{\mathrm{}}},$$
(6)
where the sums are over all pairs of vertices at the intervertex distance $`\mathrm{}`$, $`z_{\mathrm{}}=N^1_{i,j:\mathrm{}_{ij}=\mathrm{}}1`$ is the mean number of $`\mathrm{}`$-th nearest neighbors. Consequently,
$$\chi (t)=\underset{\mathrm{}}{}z_{\mathrm{}}\chi (t,\mathrm{}),$$
(7)
where $`z_0=1`$ and $`\chi (t,0)`$ is the average local susceptibility. Let us first assume that all $`\chi (t,\mathrm{})0`$. In order to satisfy the condition $`\chi (t)=O(1)`$, the terms of the sum in Eq. (7) must decay with $`\mathrm{}`$ faster than $`1/\mathrm{}`$. This leads to the following restriction:
$$\chi (t,\mathrm{})<O\left(\frac{1}{\mathrm{}z_{\mathrm{}}}\right).$$
(8)
It is important that here $`z_{\mathrm{}}`$ is a function of the network size $`N`$, so formula (8) shows how the nonlocal susceptibility (and spacial correlations) vary with $`N`$.
If the nonlocal susceptibilities $`\chi _{ij}(t)`$ are of varying sign, formula (8) is not applicable. In this case, for the root-mean-square with averaging over all pairs of vertices $`i,j`$ with a given $`\mathrm{}_{ij}=\mathrm{}`$, we arrive at
$$\left[\overline{\chi _{ij}^2(t,\mathrm{}_{ij}=\mathrm{})}\right]^{1/2}<O(z_{\mathrm{}}^{1/2}).$$
(9)
In accordance to Eq. (3), the $`\mathrm{}`$ dependence of the average correlation function is the same as for the average susceptibility, i.e., it is described by formulas (8) or (9).
In sparse networks, the mean number of the nearest neighbors (the mean degree) $`z_1\overline{k}=kP(k)`$ is finite. In general, in networks, where the mean intervertex distance $`\overline{\mathrm{}}(N)N`$, the mean numbers of $`\mathrm{}`$-th nearest neighbors grow exponentially with $`\mathrm{}`$. In particular, in uncorrelated networks (without degree–degree correlations between nearest neighbor vertices), $`z_{\mathrm{}}=z_1(z_2/z_1)^\mathrm{}1`$, where $`z_2=k(k1)P(k)`$ nsw01 . In this case, even at the critical point of a cooperative model, the nonlocal susceptibility and pair correlations decrease exponentially rapidly with $`\mathrm{}`$.
The effect is even more dramatic, if the degree distribution is fat-tailed, and its second moment diverges as $`N\mathrm{}`$. In scale-free networks, this takes place if the $`\gamma `$ exponent of the degree distribution $`P(k)k^\gamma `$ is equal or below $`3`$. In this case, already the mean number $`z_2`$ of the second nearest neighbors of a vertex diverges as $`N\mathrm{}`$. Consequently, according to Eq. (8) or Eq. (9), in the limit of a large network, pair correlations between the second (and further) nearest neighbors vanish, and only pair correlations between the nearest neighbors are observable.
## III Correlations in the Ising model
The simplicity of the Ising model and the fact that many microscopic models may be mapped to this model make the Ising model to be very attractive. Recent investigations have revealed that the critical behavior of the Ising and Potts models on complex networks strongly differs from the standard mean-field behavior on a regular lattice ahs01 ; dgm02 ; lvvz02 ; b02 ; ceah02 ; lgkk04 ; em04 . In the present section we analyze pair correlations in the equilibrium Ising model on an uncorrelated random complex network.
### III.1 The model
We consider the ferromagnetic Ising model:
$$=J\underset{ij}{}S_iS_j\underset{i}{}H_iS_i,$$
(10)
where $`S_i=\pm 1`$, and $`H_i`$ is a local magnetic field at a vertex $`i`$. The sum is over all nearest neighboring vertices. The static pair-correlation function
$$G_{ij}S_iS_jS_iS_j$$
(11)
is related to the non-local magnetic susceptibility $`\chi _{ij}`$:
$$\chi _{ij}=M_i/h_j=\beta G_{ij},$$
(12)
where $`\beta =1/T`$ and $`T`$ is temperature, $`M_iS_i`$. In this section $`\mathrm{}`$ means the statistical average with the Hamiltonian $``$.
As a substrate, we use the standard model of an uncorrelated random network—the configuration model bbk72 . This is the maximally random graph with a given degree distribution or, as it is called in graph theory, a labelled random graph with a given degree sequence. It is important that the uncorrelated networks have a locally tree-like structure. That is, in the large network limit, there are no loops in a finite neighborhood of a vertex. These networks may be considered as the random Bethe lattices, which, by definition, have no boundary. One should note that in contrast, the Cayley tree contains boundary vertices which are dead ends bbook82 .
### III.2 How to solve the Ising model on a tree-like graph
Let us outline an effective approach to solution of cooperative models on tree-like networks, which we use for obtaining $`\chi _{ij}`$. This method was implemented for the Ising and Potts models on regular bbook82 and random dgm02 ; dgm04 Bethe lattices. Consider an arbitrary tree-like graph. Consider a spin $`S_i`$ on a vertex $`i`$ of this graph with $`k_i`$ adjacent spins $`S_j`$, $`j=1,2,\mathrm{},k_i`$. As any vertex in this graph, vertex $`i`$ may be treated as a root of a tree. In turn, vertices $`j`$ may be treated as roots of subtrees growing from the vertex $`i`$. In order to characterize the subtree with the root spin $`S_j`$ we introduce a quantity
$`g_{ij}(S_i)`$ $`=`$ $`{\displaystyle \underset{\{S_l\}=\pm 1}{}}\mathrm{exp}[\beta J{\displaystyle \underset{nm}{}}S_nS_m`$ (13)
$`+\beta JS_iS_j+\beta {\displaystyle \underset{n}{}}H_nS_n].`$
Here the indices $`n`$ and $`m`$ run only over spins that belong to the sub-tree, including $`S_j`$. The statistical sum in Eq. (13) is taken only over the spins on the subtree.
We introduce a quantity
$$x_{ij}g_{ij}(1)/g_{ij}(+1).$$
(14)
In such a way, for each vertex $`i`$ we can introduce $`k_i`$ parameters $`x_{ij}`$, $`j=1,2,\mathrm{}k_i`$. It should be noted that $`x_{ij}x_{ji}`$ because $`x_{ij}`$ and $`x_{ji}`$ characterize different sub-trees. So, each edge is characterized by a pair: $`x_{ij}`$ and $`x_{ji}`$. In sum, the Ising model on an arbitrary tree-like graph is described by $`2L=_ik_i`$ parameters $`x_{ij}`$. Here, $`L`$ is the total number of edges of the graph.
Using Eqs. (13) and (14), the parameter $`x_{ij}`$ may be related to parameters $`x_{jl}`$ which characterize the edges (sub-trees) outgoing from the vertex $`j`$ as follows bbook82 :
$$x_{ij}=y(H_{j,}\underset{l=1}{\overset{k_j1}{}}x_{jl}),$$
(15)
where $`l`$ is the index of the first nearest neighbors of the vertex $`j`$. The vertex $`i`$ is excluded from the product over $`l`$. In other words, a vertex $`l`$ is a second neighbor of the vertex $`i`$ and the first one of the vertex $`j`$. The function $`y(H,x)`$ in Eq. (15) depends on a cooperative model. For the Ising model,
$$y(H,x)=\frac{e^{(J+H)\beta }+e^{(JH)\beta }x}{e^{(J+H)\beta }+e^{(JH)\beta }x}.$$
(16)
$`2L`$ independent parameters $`x_{ij}`$ should be found by solution of the set of $`2L`$ equations (15) at a given temperature $`T`$ and local magnetic fields $`H_i`$. For an arbitrary tree-like graph, these equations may be solved numerically, e.g., by use of the population dynamic algorithm, seeRef. em04 where this method has been applied to the Potts model on a tree-like graph.
Observable thermodynamic quantities of the Ising model may be written as functions of the parameters $`x_{ij}`$. For example, a magnetic moment $`M_i`$ is given by the following equation:
$$M_i=(e^{2\beta H_i}\underset{j=1}{\overset{k_i}{}}x_{ij})/\left(e^{2\beta H_i}+\underset{j=1}{\overset{k_i}{}}x_{ij}\right).$$
(17)
Finally, for a random network, the resulting observables should be averaged over various configurations with appropriate statistical weights.
### III.3 Derivation of the pair correlation function
Let us find a non-local susceptibility $`\chi _{ij}`$ for two arbitrary vertices $`i`$ and $`j`$ at the distance $`l_{ij}=\mathrm{}`$ from each other. On a tree-like graph there is the only shortest way connecting these vertices. It starts from $`i`$, then goes through vertices $`i_1`$, $`i_2,\mathrm{},i_\mathrm{}1`$ and ends at $`j`$. Using Eqs. (17) and (15) we get
$$\chi _{ij}=\frac{M_i}{x_{ii_1}}\frac{x_{ii_1}}{x_{i_1i_2}}\frac{x_{i_1i_2}}{x_{i_2i_3}}\mathrm{}\frac{x_{i_{l2}i_\mathrm{}1}}{x_{i_\mathrm{}1j}}\frac{x_{i_\mathrm{}1j}}{H_j}.$$
(18)
If we know $`x_{ij}`$ we can get $`\chi _{ij}`$.
First, for the purpose of comparison, we find a non-local susceptibility $`\chi _{ij}`$ of a regular Bethe lattice with a coordination number $`k`$.
In a uniform magnetic field $`H_1=H_2=\mathrm{}.=H`$ all vertices and all edges in a regular Bethe lattice are equivalent. Therefore, the parameters $`x_{ij}`$ are equal, i.e. $`x_{ij}=x`$, and Eq. (15) is reduced to
$$x=y(H,x^{k1}).$$
(19)
This equation determines $`x`$ as a function of $`T`$ and $`H`$. From Eqs. (18) and (19) we get
$$\chi (\mathrm{})=\frac{1}{k}\frac{M}{x}\left[\frac{1}{(k1)}\frac{y(H,x^{k1})}{x}\right]^\mathrm{}1\frac{y(H,x^{k1})}{H},$$
(20)
where in accordance with Eq. (17) we have $`M=(e^{2\beta H}x^k)/(e^{2\beta H}+x^k)`$. At zero magnetic field $`H=0`$, Eq. (20) takes a form:
$$\chi (\mathrm{})=\frac{4\beta x^k}{(1+x^k)^2}\left[\frac{2x^{k2}\mathrm{sinh}(2J\beta )}{(e^{J\beta }+e^{J\beta }x^{k1})^2}\right]^{\mathrm{}}.$$
(21)
This equation is valid at all temperatures.
In the paramagnetic state, i.e. at temperatures $`T`$ above the critical temperature $`T_c=2J/\mathrm{ln}[k/(k2)]`$ of the ferromagnetic phase transition, the self-consistent equation (19) has the only solution $`x=1`$, and we get
$$\chi (\mathrm{})=\beta G(l)=\beta \mathrm{tanh}^{\mathrm{}}(J\beta )$$
(22)
in agreement with the result obtained in Ref. f75 for a Cayley tree in the framework of another method. Furthermore, using the fact that $`x_{ij}=1`$ at $`TT_c`$, one can show that in the paramagnetic phase, $`\chi (\mathrm{})`$ of the Ising model on an arbitrary tree-like graph is the same as that for an arbitrary Cayley tree, e.g., for a spin chain (see below). Recall that there is no phase transition in spin models on a Cayley tree due to the presence of boundary spins. The number of these spins is of the order of the total number of spins on a given Cayley tree.
In the ferromagnetic phase, i.e. at $`T<T_c`$, or at $`H0`$ the parameter $`x`$ is smaller than 1. Note that in accordance with Eqs. (21) and (22) the non-local susceptibility $`\chi (\mathrm{})`$ of a regular Bethe lattice is finite at all temperatures and magnetic fields.
Now let us consider an uncorrelated random network. We calculate a value of $`\chi _{ij}(\mathrm{}_{ij}=\mathrm{})`$ averaged over the ensemble of uncorrelated random graphs with a given degree distribution function $`P(k)`$:
$`\chi (\mathrm{})\overline{\chi _{ij}(\mathrm{}_{ij}=\mathrm{})}={\displaystyle \underset{k_i}{}}{\displaystyle \underset{k_1}{}}\mathrm{}{\displaystyle \underset{k_{l1}}{}}{\displaystyle \underset{k_j}{}}\chi _{ij}`$
$`\times \left[{\displaystyle \frac{P(k_j)k_j}{z_1}}\left\{{\displaystyle \underset{n=1}{\overset{\mathrm{}1}{}}}{\displaystyle \frac{P(k_n)k_n(k_n1)}{z_2}}\right\}{\displaystyle \frac{P(k_i)k_i}{z_1}}\right].`$ (23)
Here, the quantity in the square brackets is the probability that a vertex $`i`$ with degree $`k_i`$ is connected with a vertex $`j`$ having degree $`k_j`$ by a path that goes through vertices with degrees $`k_1`$, $`k_2`$, …$`k_\mathrm{}1`$.
In the paramagnetic phase at $`TT_c=2J/\mathrm{ln}[k^2/(k^22k)]`$ and $`H=0`$ we have $`x_{ij}=1`$, and Eq. (23) leads to the Eq. (22). Therefore, in zero magnetic field both on regular and random Bethe lattices the non-local susceptibility and $`\chi _{ij}`$ have the same temperature dependence determined by Eq. (22).
At a magnetic field and in the ferromagnetic phase an approximate expression for the function $`\chi (\mathrm{})`$ may be obtained in the framework of the approach proposed in dgm02 . We introduce $`x_{ij}=\mathrm{exp}(h_{ij})`$. Here, the quantities $`h_{ij}`$ are positive and independent random parameters. We use the following ansatz,
$$\underset{j=1}{\overset{k}{}}h_{ij}kh+𝒪(k^{1/2}),$$
(24)
where $`h\overline{h_{ij}}`$ is the average value of the parameter. Applying the ansatz (24) to Eq. (15) yields
$$h=\frac{1}{\overline{k}}\underset{k}{}P(k)k\mathrm{ln}y(H,e^{(k1)h})$$
(25)
which determines $`h`$ as a function of $`T`$ and $`H`$. In fact, the parameter $`h`$ plays the role of the order parameter. At $`H=0`$, we get $`h=0`$ above $`T_c`$ and non-zero below $`T_c`$.
With this ansatz,
$$\chi (\mathrm{})=z_1^2z_2^{(\mathrm{}1)}AB^\mathrm{}1C,$$
(26)
where
$`A={\displaystyle \underset{q}{}}qP(q){\displaystyle \frac{2e^{qh}}{(e^{H\beta }+e^{H\beta qh})^2}},`$ (27)
$`B={\displaystyle \underset{q}{}}q(q1)P(q){\displaystyle \frac{2e^{(q2)h}\mathrm{sinh}(2J\beta )}{(e^{(J+H)\beta }+e^{(J+H)\beta (q1)h})^2}},`$ (28)
$`C={\displaystyle \underset{q}{}}qP(q){\displaystyle \frac{4e^{(q1)h}\mathrm{sinh}(2J\beta )}{(e^{(J+H)\beta }+e^{(J+H)\beta (q1)h})^2}}.`$ (29)
The quantities $`A,B`$ and $`C`$ are finite (independent of $`N`$) for an arbitrary degree distribution $`P(k)`$. The finiteness of the sums in Eqs. (27)–(29) is ensured by the exponential multiplier $`e^{qh}`$.
Result (26) is in agreement with the conclusions of the preceding section \[compare relations (26) and (8)\]. Note that the conclusion that $`\chi (\mathrm{}2)G(\mathrm{}2)0`$ as $`N\mathrm{}`$ if the second moment of the degree distribution diverges,does not depend on the fact that we used the approximation (24). Indeed, substituting Eq. (18) into Eq. (23), one can prove that all the sums over degrees $`k_n`$ converge both in the ordered state and at $`H0`$. Therefore, we get $`\chi (\mathrm{})z_2^{(\mathrm{}1)}`$ in agreement with Eq. (26).
## IV Discussion and conclusions
In the framework of a phenomenological approach and by using the Ising model we have studied time and static pair correlations $`x_i(t_1)x_j(t_2)`$ between interacting objects defined on a random uncorrelated complex network. We have demonstrated that if a complex network has a divergent number of the second neighbors, then the pair correlations between second and more distant neighbors are strongly suppressed in large networks. In fact, the correlations tend to zero in the limit $`N\mathrm{}`$. One should note that a particular form of the characteristic $`z_2(N)`$ strongly depends on a network model. It is important that the size dependence of the pair correlation function is determined by the factor $`1/z_2^\mathrm{}1(N)`$.
The effect of weakening of the pair correlations is well known to people leaving in large towns and forming a large network of acquaintances. Close relationships are only being kept between close relatives inside family, close friends or between people working in the same office. These people are our nearest neighbors. Our “second” and more distant nearest neighbors—people living on the next floor or working in a neighboring office —usually weakly influence our personal and social life. Thisis contrary to a small village where strong correlations are present between almost all inhabitants.
Our analysis may be generalized to correlated networks. Many natural networks exhibit correlations between degrees of adjacent vertices, see, for example, Ref. corr . Note that in large uncorrelated networks, the divergence of the mean number $`z_{\mathrm{}}`$ of the $`\mathrm{}`$-th nearest neighbor can occur only simultaneously at all $`\mathrm{}2`$. In contrast, in large correlated networks, it is possible in principle that, say, $`z_2,z_3`$ is finite and only $`z_\mathrm{}4`$ diverges. In this example, pair correlations are observable between the first, second and third nearest neighbors and vanish starting from the fourth nearest neighbors.
We restricted ourselves to pair correlations. One should emphasize that pair correlations in cooperative systems on networks, unlike lattices, never show critical behavior. Even at the critical point of an interacting system on a network, we observe an exponentially rapid decay of pair correlations. In contrast, volume correlations in interacting systems on networks demonstrate a critical feature. Let us explain this point. The distribution of the full response of a system to a small local field is $`P(\epsilon )=_i\delta (_j\chi _{ij}\epsilon )`$, where $`\delta (\epsilon )`$ is the delta-function. This distribution is a rapidly decreasing function both below and above a phase transition on a network. On the other hand, at the critical point, $`P(\epsilon )`$ is power-law, that is, “critical” nsw01 ; bcd05 .
One can conclude that a network structure of an interacting system strongly influences the dependence of pair correlations on the distance between interacting objects. So, investigations of the correlations may be a useful method to understand the network topology. If an interacting system is defined on a network with a divergent number $`z_n`$ of $`n`$-th nearest neighbors of a vertex, then dynamic and static pair correlation are strongly suppressed at distances $`n`$ and greater.
This work was partially supported by projects POCTI: FAT/46241/2002, MAT/46176/2002, FIS/61665/2004, and BIA-BCM/62662/2004. S.N.D. and J.F.F.M. were also supported by project DYSONET-NEST/012911. Authors thank A.N. Samukhin for useful discussions.
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# 𝑀-Theory Solutions in Multiple Signatures from 𝐸₁₁
## 1 Introduction
It has been conjectured that $`M`$-theory possesses a large Kac-Moody symmetry which is non-linearly realised and thought to contain $`E_{11}`$. Indeed, it was shown that the bosonic sector of eleven dimensional supergravity could be described by a non-linear realisation , that is, if one assumes that the non-linear realisation was extended such that it was a Kac-Moody algebra then this algebra must contain $`E_{11}`$. Subsequently, two variants of this conjecture have been made . The first differs from the original suggestion of in that it only adopted the sub-algebra $`E_{10}`$ as a symmetry. However, it also differed from the way the original suggestion of , subsequently evolved in , where spacetime was incorporated by extending the translation generators to be part of an $`E_{11}`$ representation. Indeed, in spacetime was supposed to be contained in effect in the Kac-Moody algebra. The proposal of also adopted $`E_{11}`$ as the symmetry but took a similar approach to spacetime as that of . For a discussion of the merits of these approaches to spacetime see reference .
One of the hopes for such a large symmetry algebra is that it might provide insights into the solutions of the theory. Indeed, relationships have been uncovered between the half BPS solution of eleven dimensional supergravity and $`E_{11}`$ and more generally between all the oxidised theories and the suspected $`G^{+++}`$ symmetries of their extensions . In it was pointed out that any solution of the non-linearly realised $`E_{11}`$ theory could be expressed as an $`E_{11}`$ group element and it was found that all the known half BPS solutions were generated by a very simple group element which was specified by the choice of a single root of $`E_{11}`$,
$$g=\mathrm{exp}(\frac{1}{(\beta ,\beta )}\mathrm{ln}N\beta H)\mathrm{exp}((1N)E_\beta )$$
(1)
The root, $`\beta `$, is $`\alpha _1+2\alpha _2+3(\alpha _3+\mathrm{}\alpha _8)+2\alpha _9+\alpha _{10}+\alpha _{11}`$ for the $`M2`$-brane and $`\alpha _1+2\alpha _2+3\alpha _3+4\alpha _4+5\alpha _5+6\alpha _6+5\alpha _7+4\alpha _8+3\alpha _9+2\alpha _{10}+2\alpha _{11}`$ for the $`M5`$-brane. In fact only used a subset of the possible roots of $`E_{11}`$ and in particular those which when decomposed to the $`A_{10}`$ sub-algebra of $`E_{11}`$ corresponded to lowest weights of $`A_{10}`$. Indeed one can explicitly check that the $`M2`$ and $`M5`$ branes are electric solutions for the lowest weights, see appendix A.
In this paper we will consider the other weights of $`E_{11}`$. We will find that the group element does indeed lead to solutions for all weights. However, some of these are to $`M`$-theories which are not in the $`(1,10)`$ signature. Since the different weights are related by Weyl transformations these will play an important role in this paper. The Weyl transformations of $`E_{11}`$ were found to correspond in the dimensionally reduced theory to the U-duality transformations . Weyl transformations of some solutions in eleven dimensions were discussed in .
The non-linearly realised theory requires that the local sub-algebra is specified. This was initially taken to be the Cartan invariant involution, but this lead to a Euclidean theory and to find the usual theory in $`(1,10)`$ signature required a simple Wick rotation. However, it was realised that one could directly find the theory in $`(1,10)`$ signature by adopting a different local sub-algebra by inserting a single minus sign in the generating set of the local sub-algebra as compared to that for the Cartan invariant involution. This was intially suggested in (see equations 2.17-19) and a more systematic discussion was given in reference and subsequently in . It was then realised that by inserting other minus signs that one could find $`M`$-theories in other signatures. Furthermore, it was shown by Keurentjes , that the Weyl reflections of $`E_{11}`$ do not commute with the choice of local sub-algebra, with the consequence that, by carrying out Weyl transformations, the usual signature of $`M`$-theory, $`(1,10)`$, occurs within $`E_{11}`$ alongside the signatures $`(2,9)`$, $`(5,6)`$ and their inverses. These were the theories that were called the $`M`$ and $`M^{}`$-theories, respectively, in reference . The relation between $`E_{11}`$ and $`E_{10}`$ was discussed in and the relationship between the different 10-dimensional theories and Weyl reflections was elaborated. Recently, this line of thought has been extended to all $`G^{++}`$ theories .
In section two we evaluate the bosonic equations of motion of eleven dimensional supergravity in an arbitrary signature for a brane ansatz that includes the possibility of several times in the world-volume of the brane and also in the transverse space. We show how given one solution to a theory in a given signature we can generate classes of solutions in theories of different signature by inverting the signature in the transverse space of the solution, and in certain cases by inversion of the brane world-volume signature. These considerations are extended to find all possible signatures for which there exist solutions related to the $`M2`$ and $`M5`$-branes.
In section three we first review the work of Keurentjes on the effect of $`E_{11}`$ Weyl transformations on the signature of the theory. We then examine the effect of the Weyl transformation on the well known $`M2`$, $`M5`$ and $`pp`$-wave solutions in the $`(1,10)`$ theory. Commencing with the group element (1) and the lowest weights of $`A_{10}`$, different weights are obtained by Weyl reflections. Since it is possible to use a series of Weyl reflections that alter the signature but not the root associated to the weight we find that each weight, including the lowest, has an ambiguity concerning which signature it exists in. Indeed, we find all possible such transformed solutions and the corresponding signatures of the theory in which each solution should exist. We compare these with the results of section two and find that the Weyl transformations do indeed lead to solutions in the theories of the required signature, exactly reproducing all the known solutions of the $`M`$ and $`M^{}`$-theories . However, the Weyl transformations also lead to examples which are not solutions in the theory of the given signature. In these cases, which are as numerous as the solutions, the example is a solution to a theory with an alternative sign in front of the kinetic $`F^2`$ term in the action.
Spacelike brane solutions, or $`S`$-branes, are branes which map out a spacelike world-volume .<sup>1</sup><sup>1</sup>1The convention for naming spacelike branes is that an S$`p`$-brane has a $`(p+1)`$ Euclidean world-volume. Consequently an $`S`$-brane only exists for a moment in time. We show that these are naturally contained in the group element (1). In section 3.2, we propose that a choice of local sub-algebra which invokes a time coordinate in the transverse space of an $`M`$-theory p-brane solution, as oppose to inducing the usual longitudinal time coordinate, reproduces the known $`S2`$ and $`S5`$-branes of $`M`$-theory . We also find encoded in the group element (1) their relations in $`M`$ and $`M^{}`$-theories.
## 2 General Signature Formulation of the Einstein Equations
In this paper we will be working beyond the usual signatures of supergravity, it will be useful to express the Einstein equations in a form that is applicable to different signatures. The Einstein equations and the gauge equations for a single brane solution can be derived by varying the truncated form of a gravity action, where the truncation is the restriction to the kinetic term for one of the theory’s n-form field strengths, the dilaton term and the Ricci scalar, e.g.
$$A=\frac{1}{16\pi G_D}d^Dx\sqrt{g}(R\frac{1}{2}^\mu \varphi _\mu \varphi \frac{1}{2.n!}e^{a_i\varphi }F_{\mu _1\mathrm{}\mu _n}F^{\mu _1\mathrm{}\mu _n})$$
(2)
The usual Chern-Simons term appearing in the eleven dimensional supergravity action has been omitted here since for the class of extremal branes we will consider it plays no role in the dynamics, and will not affect our discussions. The equations of motion determined by varying with respect to the metric, $`g_{\mu \nu }`$, are the Einstein equations and the equation that comes from varying the gauge field is the gauge equation. There is also an equation of motion coming from the variation of the dilaton field. For the generic truncated action these are
$`R_{}^{\mu }{}_{\nu }{}^{}={\displaystyle \frac{1}{2}}^\mu \varphi _\nu \varphi +{\displaystyle \frac{1}{2n!}}e^{a_i\varphi }(nF^{\mu \lambda _2\mathrm{}\lambda _n}F_{\nu \lambda _2\mathrm{}\lambda _n}{\displaystyle \frac{n1}{D2}}\delta _{}^{\mu }{}_{\nu }{}^{}F_{\lambda _1\mathrm{}\lambda _n}F^{\lambda _1\mathrm{}\lambda _n})`$
$`_\mu (\sqrt{g}e^{a_i\varphi }F^{\mu \lambda _2\mathrm{}\lambda _n})=0`$ (3)
$`{\displaystyle \frac{1}{\sqrt{g}}}_\mu (\sqrt{g}^\mu \varphi ){\displaystyle \frac{a_i}{2.n!}}e^{a_i\varphi }F_{\mu \lambda _2\mathrm{}\lambda _n}F^{\mu \lambda _2\mathrm{}\lambda _n}=0`$
We compute the curvature components in the spin-connection formalism as described in , but we commence with a line element for a more general than usual brane solution,
$$ds^2=A^2(\underset{i=1}{\overset{i=q}{}}dt_i^2+\underset{j=1}{\overset{j=p}{}}dx_j^2)+B^2(\underset{a=1}{\overset{a=c}{}}du_a^2+\underset{b=1}{\overset{b=d}{}}dy_b^2)$$
(4)
The coordinates are split into two groups, those that are longitudinal to the brane, $`t_i`$ and $`x_i`$, we indicate with indices $`\{i,j,k,\mathrm{}\}`$, and those that are transverse, $`u_a`$ and $`y_a`$ with $`\{a,b,c,\mathrm{}\}`$. The given line element is the world-volume of a brane with signature $`(q,p)`$ on the brane and signature $`(c,d)`$ in the bulk; the corresponding global signature is $`(q+c,p+d)`$. We adopt the notation $`[(q,p),(c,d)]`$ to express a single signature for our ansatz in terms of its longitudinal, $`(q,p)`$ and transverse components $`(c,d)`$. For a global signature $`(t,s)`$ then $`q+c=t`$ and $`p+d=s`$. The case when $`c=0`$ and $`q=1`$, corresponds to the usual single brane solutions. In eleven-dimensional supergravity there is no dilaton, if we set the dilaton coupling to zero, the coefficients $`A`$ and $`B`$, functions of the transverse coordinates $`(u_a,y_a)`$, take the form in the extremal case
$$A=N_{(c,d)}^{(\frac{1}{p+1})},B=N_{(c,d)}^{(\frac{1}{Dp3})}$$
(5)
Where,
$$N_{(c,d)}=1+\frac{1}{Dp3}\sqrt{\frac{\mathrm{\Delta }}{2(D2)}}\frac{𝐐}{r^{(Dp3)}}$$
(6)
Where $`𝐐`$ is the conserved charge associated with the $`p`$-brane solution, $`\mathrm{\Delta }=(p+1)(d2)+\frac{1}{2}a_i^2(D2)`$ and $`r`$ is the radial distance in the transverse coordinates such that $`r^2=u_au_a+y_by_b`$. That is, $`N_{(c,d)}`$ are independent of the longitudinal coordinates, and are harmonic functions in the transverse coordinates, $`u_a`$ and $`y_b`$, so that $`^\mu _\mu N_{(c,d)}=0`$.
The full non-zero curvature components are,
$`R_{}^{t_i}{}_{t_i}{}^{}=B^2\{`$ $`_{u_a}_{u_a}\mathrm{ln}A\widehat{\delta }_{}^{u_a}{}_{u_a}{}^{}_{y_a}_{y_a}\mathrm{ln}A\widehat{\delta }_{}^{y_a}{}_{y_a}{}^{}`$
$`+_{u_a}\mathrm{ln}A_{u_a}\mathrm{\Psi }\widehat{\delta }_{}^{u_a}{}_{u_a}{}^{}_{y_a}\mathrm{ln}A_{y_a}\mathrm{\Psi }\widehat{\delta }_{}^{y_a}{}_{y_a}{}^{}\}\widehat{\delta }^{t_i}_{t_i}`$
$`R_{}^{x_i}{}_{x_i}{}^{}=B^2\{`$ $`_{u_a}_{u_a}\mathrm{ln}A\widehat{\delta }_{}^{u_a}{}_{u_a}{}^{}_{y_a}_{y_a}\mathrm{ln}A\widehat{\delta }_{}^{y_a}{}_{y_a}{}^{}`$
$`+_{u_a}\mathrm{ln}A_{u_a}\mathrm{\Psi }\widehat{\delta }_{}^{u_a}{}_{u_a}{}^{}_{y_a}\mathrm{ln}A_{y_a}\mathrm{\Psi }\widehat{\delta }_{}^{y_a}{}_{y_a}{}^{}\}\widehat{\delta }^{x_i}_{x_i}`$
$`R_{}^{u_a}{}_{u_a}{}^{}=B^2\{`$ $`_{u_a}_{u_a}\mathrm{ln}B\widehat{\delta }_{}^{u_a}{}_{u_a}{}^{}_{y_a}_{y_a}\mathrm{ln}B\widehat{\delta }_{}^{y_a}{}_{y_a}{}^{}`$
$`+_{u_a}\mathrm{ln}B_{u_a}\mathrm{ln}B(\widehat{\delta }_{}^{y_a}{}_{y_a}{}^{}+\widehat{\delta }_{}^{u_a}{}_{u_a}{}^{}2)`$ (7)
$`+_{u_a}\mathrm{ln}A_{u_a}\mathrm{ln}A(\widehat{\delta }_{}^{t_i}{}_{t_i}{}^{}+\widehat{\delta }_{}^{x_i}{}_{x_i}{}^{})_{y_a}\mathrm{ln}B_{y_a}\mathrm{\Psi }\widehat{\delta }_{}^{y_a}{}_{y_a}{}^{}`$
$`+_{u_a}_{u_a}\mathrm{\Psi }+_{u_a}\mathrm{ln}B_{u_a}\mathrm{\Psi }(\widehat{\delta }_{}^{u_a}{}_{u_a}{}^{}2)\}\widehat{\delta }^{u_a}_{u_a}`$
$`R_{}^{y_a}{}_{y_a}{}^{}=B^2\{`$ $`_{u_a}_{u_a}\mathrm{ln}B\widehat{\delta }_{}^{u_a}{}_{u_a}{}^{}_{y_a}_{y_a}\mathrm{ln}B\widehat{\delta }_{}^{y_a}{}_{y_a}{}^{}`$
$`_{y_a}\mathrm{ln}B_{y_a}\mathrm{ln}B(\widehat{\delta }_{}^{y_a}{}_{y_a}{}^{}+\widehat{\delta }_{}^{u_a}{}_{u_a}{}^{}2)`$
$`_{y_a}\mathrm{ln}A_{y_a}\mathrm{ln}A(\widehat{\delta }_{}^{t_i}{}_{t_i}{}^{}+\widehat{\delta }_{}^{x_i}{}_{x_i}{}^{})+_{u_a}\mathrm{ln}B_{u_a}\mathrm{\Psi }\widehat{\delta }_{}^{u_a}{}_{u_a}{}^{}`$
$`_{y_a}_{y_a}\mathrm{\Psi }_{y_a}\mathrm{ln}B_{y_a}\mathrm{\Psi }(\widehat{\delta }_{}^{y_a}{}_{y_a}{}^{}2)\}\widehat{\delta }^{y_a}_{y_a}`$
Where $`\widehat{\delta }_{}^{x_i}{}_{x_i}{}^{}`$ counts the number of $`x_i`$ coordinates in the line element (e.g. for the $`M2`$-brane, $`\widehat{\delta }_{}^{t_i}{}_{t_i}{}^{}=1`$, $`\widehat{\delta }_{}^{x_i}{}_{x_i}{}^{}=2`$, $`\widehat{\delta }_{}^{u_a}{}_{u_a}{}^{}=0`$ and $`\widehat{\delta }_{}^{y_a}{}_{y_a}{}^{}=8`$); repeated lowered indices a, b, i and j are not summed - sums are taken care of via the counting symbols $`\widehat{\delta }`$; and, for the ansatz (4) with extremal coefficients A and B (5),
$`\mathrm{\Psi }`$ $`(p+1)\mathrm{ln}A+(Dp3)\mathrm{ln}B`$
$`=`$ $`(p+1)\mathrm{ln}N_p^{2(\frac{1}{p+1})}+(Dp3)\mathrm{ln}N_p^{2(\frac{1}{Dp3})}`$ (8)
$`=`$ $`0`$
The curvature terms reduce to
$`R_{}^{t_i}{}_{t_i}{}^{}=B^2\{`$ $`_{u_a}_{u_a}\mathrm{ln}A\widehat{\delta }_{}^{u_a}{}_{u_a}{}^{}_{y_a}_{y_a}\mathrm{ln}A\widehat{\delta }_{}^{y_a}{}_{y_a}{}^{}\}\widehat{\delta }^{t_i}_{t_i}`$
$`R_{}^{x_i}{}_{x_i}{}^{}=B^2\{`$ $`_{u_a}_{u_a}\mathrm{ln}A\widehat{\delta }_{}^{u_a}{}_{u_a}{}^{}_{y_a}_{y_a}\mathrm{ln}A\widehat{\delta }_{}^{y_a}{}_{y_a}{}^{}\}\widehat{\delta }^{x_i}_{x_i}`$
$`R_{}^{u_a}{}_{u_a}{}^{}=B^2\{`$ $`_{u_a}_{u_a}\mathrm{ln}B\widehat{\delta }_{}^{u_a}{}_{u_a}{}^{}_{y_a}_{y_a}\mathrm{ln}B\widehat{\delta }_{}^{y_a}{}_{y_a}{}^{}`$
$`+_{u_a}\mathrm{ln}B_{u_a}\mathrm{ln}B(\widehat{\delta }_{}^{y_a}{}_{y_a}{}^{}+\widehat{\delta }_{}^{u_a}{}_{u_a}{}^{}2)`$ (9)
$`+_{u_a}\mathrm{ln}A_{u_a}\mathrm{ln}A(\widehat{\delta }_{}^{t_i}{}_{t_i}{}^{}+\widehat{\delta }_{}^{x_i}{}_{x_i}{}^{})\}\widehat{\delta }^{u_a}_{u_a}`$
$`R_{}^{y_a}{}_{y_a}{}^{}=B^2\{`$ $`_{u_a}_{u_a}\mathrm{ln}B\widehat{\delta }_{}^{u_a}{}_{u_a}{}^{}_{y_a}_{y_a}\mathrm{ln}B\widehat{\delta }_{}^{y_a}{}_{y_a}{}^{}`$
$`_{y_a}\mathrm{ln}B_{y_a}\mathrm{ln}B(\widehat{\delta }_{}^{y_a}{}_{y_a}{}^{}+\widehat{\delta }_{}^{u_a}{}_{u_a}{}^{}2)`$
$`_{y_a}\mathrm{ln}A_{y_a}\mathrm{ln}A(\widehat{\delta }_{}^{t_i}{}_{t_i}{}^{}+\widehat{\delta }_{}^{x_i}{}_{x_i}{}^{})\}\widehat{\delta }^{y_a}_{y_a}`$
The $`M2`$, $`M5`$, and $`pp`$-wave solutions of $`M`$-theory are encoded in a group element of the non-linear realisation of $`E_{11}`$ . For reference and comparison with later solutions, we demonstrate that these cases satisfy the Einstein equations in appendix A.
### 2.1 New Solutions from Signature Change
A change of signature has the potential to alter both the Einstein equations and the field content of a theory. In preparation for the next section, we pose a question: given a solution to the Einstein equations in one signature, are there any other signatures which would also carry a related version of that solution in the new signature? We note that the harmonic function, $`N_{(c,d)}`$, is a function of the transverse coordinates and may be transformed by a signature change. By a ’related solution’ we specifically mean that if the Einstein equations for a given solution were reexpressed in terms of the functions carrying the new signature then they would remain balanced and we would find a new solution.
Signature change can be brought about in two equivalent ways, the first is as a mapping of a coordinate, or a subset of the coordinates, $`x^\mu ix^\mu ,x_\mu ix_\mu `$, leaving the metric unaltered. For example, a Lorentzian signature can be made Euclidean by making the change on the temporal coordinate, $`t^1ix^1`$, having the effects,
$`ds^2`$ $`=g_{\mu \nu }x^\mu x^\nu `$
$`=g_{t_1t_1}dt_1^2+g_{x_2x_2}dx_2^2+\mathrm{}+g_{x_Dx_D}dx_D^2`$
$`=f_1(N)dt_1^2+f_2(N)dx_2^2+\mathrm{}+f_D(N)dx_D^2`$
$``$ (10)
$`ds^2`$ $`=f_1(N)dx_1^2+f_2(N)dx_2^2+\mathrm{}+f_D(N)dx_D^2`$
$`=g_{x_1x_1}dx_1^2+g_{x_2x_2}dx_2^2+\mathrm{}+g_{x_Dx_D}dx_D^2`$
Where $`\pm f_i(N)`$, some function of $`N`$, is the metric component in each case. An electric field $`A_\mu `$ transforms as
$$A_\mu dx^\mu =A_{t_1}dt^1iA_{x_1}dx^1$$
(11)
Equivalently, signature change at the quadratic level of the line element can be thought of as a transformation of the metric components, or subset of the metric components, where appropriate, $`g_{\mu \nu }g_{\mu \nu }`$ as opposed to the coordinates. Correspondingly we can view the example above (10), as the transformation $`g_{t_1t_1}g_{x_1x_1}`$. Both methods realise the change of line element but only when applied independently.
In equation (9) we have written out the curvature coefficients for our ansatz, (4). We observe that the expressions for $`R_{}^{t_i}{}_{t_i}{}^{}`$ and $`R_{}^{x_i}{}_{x_i}{}^{}`$ satisfying our explicit ansatz of (4) are interchanged under the interchange of longitudinal temporal and spatial coordinates, given by $`t^iix^i`$ and $`x^jit^j`$. These transformations corresponds to the notational swap $`t_ix_i`$ in the equations (9) and the set of curvature terms as a whole is unaffected. In terms of signature this corresponds to a signature inversion on only the longitudinal coordinates. Alternatively, the swap $`u_a`$ for $`y_a`$ and vice-versa, given by $`u^aiy^a`$ and $`y^biu^b`$ and corresponding to a signature inversion on only the transverse coordinates, interchanges the expressions for $`R_{}^{u_a}{}_{u_a}{}^{}`$ and $`R_{}^{y_a}{}_{y_a}{}^{}`$ and introduces a minus into all the curvature components, $`R_{}^{\mu }{}_{\nu }{}^{}R_{}^{\mu }{}_{\nu }{}^{}`$. One could also achieve these signature changes by transforming the metric components: a signature inversion on all of the longitudinal coordinates $`g^{t_it_i}g^{x_ix_i}`$ and $`g^{x_ix_i}g^{t_it_i}`$ leaves the set of curvature terms unaltered, whereas $`g^{u_au_a}g^{y_ay_a}`$ and $`g^{y_ay_a}g^{u_au_a}`$ introduces a minus sign for all the curvature terms.
To find a new solution under longitudinal and transverse signature inversions the signs induced in the curvature components must match the sign changes in the remaining terms of the Einstein equations, those derived from the field strength and the dilaton. For the eleven dimensional case, which we consider in this paper, there is no dilaton, $`\varphi 0`$, $`a_i0`$, so we shall disregard it in the following discussion.
The field strength terms in each of the Einstein equations for the usual single brane solutions, with only one temporal coordinate longitudinal to the brane and none transverse, given in appendix A, are proportional to
$$g^{tt^{}}g^{y_1y_1^{}}\mathrm{}g^{y_py_p^{}}g^{y_ay_a^{}}(F_{t^{}y_1^{}\mathrm{}y_p^{}y_a^{}})(F_{ty_1\mathrm{}y_py_a})$$
(12)
With the generalisation to our ansatz (4) to include multiple time coordinates longitudinal and transverse to the brane, the equivalent proportional term is
$$g^{t_1t_1^{}}\mathrm{}g^{t_qt_q^{}}g^{y_1y_1^{}}\mathrm{}g^{y_py_p^{}}g^{\mu \mu ^{}}(F_{t_1^{}\mathrm{}t_q^{}y_1^{}\mathrm{}y_p^{}\mu ^{}})(F_{t_1\mathrm{}t_qy_1\mathrm{}y_p\mu })$$
(13)
Where the radial coordinate $`\mu `$ may now be a spatial or temporal transverse coordinate. An inversion of the longitudinal coordinates only, causes a sign change $`(1)^{p+q}`$ in this term. The effect of inverting only the transverse coordinates $`g^{\mu \mu ^{}}g^{\mu \mu ^{}}`$ introduces a minus sign as there is only one occurrence of the metric component with transverse coordinates in (13).
A new solution is found under a signature inversion when the sign changes induced in the Riemann curvature components match the sign changes in the term (13). For example, since we have observed that a signature inversion on only the transverse coordinates introduces a minus sign in both the curvature components and the remaining terms in the Einstein equations (13), then a solution with signature components $`[(q,p),(c,d)]`$, will always find a new solution under an inversion of the transverse signature, taking the signature to $`[(q,p),(d,c)]`$. Additionally, if $`p+q`$ is even we may invert just the longitudinal coordinates and find another new solution $`[(q,p),(c,d)][(p,q),(c,d)]`$, and furthermore in this case we may invert the full signature, an inversion of both longitudinal and transverse coordinates together, $`[(q,p),(c,d)][(p,q),(c,d)]`$ and find yet another new solution. However if $`p+q`$ is odd an inversion of the longitudinal signature introduces an unbalanced minus sign and no new solution is found.<sup>2</sup><sup>2</sup>2However, we shall observe later that the $`F^2`$ term in the action may change its sign and in such cases the ’lost’ solution of a $`F^2`$ theory is a new solution of a $`+F^2`$ theory, but for the time being we continue to consider the usual $`F^2`$ theory For example, a solution in $`(1,D1)`$ with longitudinal and transverse signature components $`[(1,p),(0,Dp1)]`$ is the familiar $`p`$-brane solution. We find that for even $`p+q`$, we have the following set of signature components that carry a related solution, $`[(1,p),(Dp1,0)]`$, $`[(p,1),(0,Dp1)]`$ and $`[(p,1),(Dp1,0)]`$ which give spacetime signatures $`(Dp,p)`$, $`(p,Dp)`$ and $`(1,D1)`$ respectively. For the case of odd $`(p+q)`$ we have only one alternative signature, coming from an inversion of only the transverse coordinates, which gives a new solution, namely $`(Dp,p)`$ with longitudinal and transverse components $`[(1,p),(Dp1,0)]`$.
Following the preceding considerations we are in a position to write down a set of signatures in which the $`M2`$ and $`M5`$ branes remain solutions. The $`M2`$ brane has $`p+q=3`$ and hence has related solutions in $`(1,10)`$ and $`(9,2)`$, with the following parameters,
| Signature | Longitudinal | Transverse |
| --- | --- | --- |
| | Signature | Signature |
| $`(1,10)`$ | $`(1,2)`$ | $`(0,8)`$ |
| $`(9,2)`$ | $`(1,2)`$ | $`(8,0)`$ |
Equivalently, the $`M5`$ brane has $`p+q=6`$ and has related solutions in $`(1,10)`$, $`(6,5)`$, $`(5,6)`$ and $`(10,1)`$,
| Signature | Longitudinal | Transverse |
| --- | --- | --- |
| | Signature | Signature |
| $`(1,10)`$ | $`(1,5)`$ | $`(0,5)`$ |
| $`(6,5)`$ | $`(1,5)`$ | $`(5,0)`$ |
| $`(5,6)`$ | $`(5,1)`$ | $`(0,5)`$ |
| $`(10,1)`$ | $`(5,1)`$ | $`(5,0)`$ |
In addition to inverting components of the signature we may also transform individual temporal coordinates into spacelike coordinates and vice versa. It is observed by following the computations in appendix A and the term in equation (13), that a given solution will give a new solution by converting an even number of longitudinal temporal coordinates into longitudinal spatial coordinates, while leaving the transverse coordinates unaltered. Such a transformation introduces a sign change $`(1)^{2m}=1`$ into term (13), where $`m`$ is an integer such that $`q\pm 2m,p2m0`$, and only alters the counting symbols $`\widehat{\delta }_{}^{t_i}{}_{t_i}{}^{}`$ and $`\widehat{\delta }_{}^{x_i}{}_{x_i}{}^{}`$ in the curvature components so that a new solution is found. That is, given a solution in a signature with components $`[(q,p),(c,d)]`$ then functions carrying the signature $`[(q\pm 2m,p2m),(c,d)]`$ will give a new solution. Applying this to each of the signatures containing solutions related to the $`M2`$ brane, we find the additional solutions,
| Signature | Longitudinal | Transverse |
| --- | --- | --- |
| | Signature | Signature |
| $`(3,8)`$ | $`(3,0)`$ | $`(0,8)`$ |
| $`(11,0)`$ | $`(3,0)`$ | $`(8,0)`$ |
And similarly, additional signatures for the $`M5`$ solutions are,
| Signature | Longitudinal | Transverse |
| --- | --- | --- |
| | Signature | Signature |
| $`(3,8)`$ | $`(3,3)`$ | $`(0,5)`$ |
| $`(8,3)`$ | $`(3,3)`$ | $`(5,0)`$ |
We can carry this argument to the transverse coordinates, but we are no longer restricted to transforming even multiples of temporal coordinates into spatial coordinates, indeed any integer is possible giving a range of new solutions in signatures $`[(q,p),(c\pm m,dm)]`$ where m is an integer such that $`c\pm m,dm0`$. For the solutions related to the $`M2`$-brane we find further solutions,
| Signature | Longitudinal | Transverse |
| --- | --- | --- |
| | Signature | Signature |
| $`(2,9)`$ | $`(1,2)`$ | $`(1,7)`$ |
| $`(3,8)`$ | $`(1,2)`$ | $`(2,6)`$ |
| $`(4,7)`$ | $`(1,2)`$ | $`(3,5)`$ |
| $`(5,6)`$ | $`(1,2)`$ | $`(4,4)`$ |
| $`(6,5)`$ | $`(1,2)`$ | $`(5,3)`$ |
| $`(7,4)`$ | $`(1,2)`$ | $`(6,2)`$ |
| $`(8,3)`$ | $`(1,2)`$ | $`(7,1)`$ |
| $`(4,7)`$ | $`(3,0)`$ | $`(1,7)`$ |
| $`(5,6)`$ | $`(3,0)`$ | $`(2,6)`$ |
| $`(6,5)`$ | $`(3,0)`$ | $`(3,5)`$ |
| $`(7,4)`$ | $`(3,0)`$ | $`(4,4)`$ |
| $`(8,3)`$ | $`(3,0)`$ | $`(5,3)`$ |
| $`(9,2)`$ | $`(3,0)`$ | $`(6,2)`$ |
| $`(10,1)`$ | $`(3,0)`$ | $`(7,1)`$ |
And for the $`M5`$ solution we find the further related solutions,
| Signature | Longitudinal | Transverse |
| --- | --- | --- |
| | Signature | Signature |
| $`(2,9)`$ | $`(1,5)`$ | $`(1,4)`$ |
| $`(3,8)`$ | $`(1,5)`$ | $`(2,3)`$ |
| $`(4,7)`$ | $`(1,5)`$ | $`(3,2)`$ |
| $`(5,6)`$ | $`(1,5)`$ | $`(4,1)`$ |
| $`(4,7)`$ | $`(3,3)`$ | $`(1,4)`$ |
| $`(5,6)`$ | $`(3,3)`$ | $`(2,3)`$ |
| $`(6,5)`$ | $`(3,3)`$ | $`(3,2)`$ |
| $`(7,4)`$ | $`(3,3)`$ | $`(4,1)`$ |
| $`(5,6)`$ | $`(5,1)`$ | $`(1,4)`$ |
| $`(7,4)`$ | $`(5,1)`$ | $`(2,3)`$ |
| $`(8,3)`$ | $`(5,1)`$ | $`(3,2)`$ |
| $`(9,2)`$ | $`(5,1)`$ | $`(4,1)`$ |
This discussion is exhaustive, we have found all signatures that give solutions in $`F^2`$ theories related to the $`M2`$ and $`M5`$-brane solutions of $`M`$-theory. We note that a universal shorthand for assessing whether or not a given signature contains a solution is to count the number of temporal longitudinal coordinates and if this is odd we have a solution to $`F^2`$ theories.
For the purposes of this paper it will be useful to consider not only the truncated action of the form given in equation (2), but also theories constructed from an action under the double Wick rotation transforming $`A_{\mu _1\mathrm{}\mu _n}iA_{\mu _1\mathrm{}\mu _n}`$. Such an action, upto Chern-Simons terms, looks identical to that of equation (2) except that the sign of the kinetic $`F^2`$ term has changed from ”$``$” to ”$`+`$”. Such a theory can be imagined as originating with an imaginary brane charge and will be relevant to our later consideration of spacelike branes. Let us now apply the reasoning of this section to such a $`+F^2`$ action. Adjusting the considerations of this section to $`+F^2`$ theories introduces an extra minus sign in front of the field strength terms in the Einstein equations (3) and as a consequence into the term proportional to the non-curvature components given in equation (13). In this case we will have a solution if the number of longitudinal time coordinates is even, meaning that exactly all the signatures not listed in this section, from the set of all signatures in 11-dimensions, will admit solutions to a $`+F^2`$ theory.
Let us introduce a new term $`\kappa `$ to keep track of solutions in both $`+F^2`$ and $`F^2`$ theories. For reasons that will become clear we use $`f(\alpha _{11})`$ to indicate the sign in front of the $`F^2`$ term, if $`f(\alpha _{11})=0`$ we have a $`F^2`$ theory and if $`f(\alpha _{11})=1`$ we have a $`+F^2`$ theory. We define $`\kappa `$ to be
$$\kappa (1)^{(\widehat{\delta }_{}^{t_i}{}_{t_i}{}^{}+f(\alpha _{11}))}$$
(14)
If $`\kappa =1`$ we have a brane solution, otherwise we do not; this criterion will be used to check for solutions throughout this paper. We note that this implies that there are no extremal $`S`$-branes ($`\widehat{\delta }_{}^{t_i}{}_{t_i}{}^{}=0`$) in $`F^2`$ theories ($`f(\alpha _{11})=0`$) indicating the known result that $`S`$-branes in $`M`$-theory have an associated imaginary charge $`𝐐`$ in equation (6). More simply, the transformation $`𝐐i𝐐`$ induces the transformation $`F^2+F^2`$.
## 3 The Local Sub-Algebra and Weyl Reflections
In the non-linear realisation of $`E_{11}`$ as a coset symmetry $`E_{11}/H_{11}`$, the local denominator sub-algebra, $`H_{11}`$, was chosen to be Cartan involution invariant, yielding the maximal compact sub-algebra. It was understood that a Wick rotation would then give a Lorentz invariant sub-algebra. We recall that the Cartan involution, $`\mathrm{\Omega }`$, takes the generators of the positive roots, $`E_i=K_{}^{i}{}_{i+1}{}^{}`$, to the negative of the generators of the negative roots, $`F_i=K_{}^{i+1}{}_{i}{}^{}`$ and vice-versa,
$$\mathrm{\Omega }(E_i)=F_i\text{and}\mathrm{\Omega }(F_i)=E_i$$
(15)
such that the set of generators, $`E_iF_i`$, is invariant under the Cartan involution, $`\mathrm{\Omega }(E_iF_i)=F_i(E_i)=E_iF_i`$, and form a basis for the local denominator sub-algebra, $`H_{11}`$. It was noted that the Cartan involution can be generalised to what has been called the temporal involution, $`\widehat{\mathrm{\Omega }}`$, whose action is:
$$\widehat{\mathrm{\Omega }}(K_{}^{i}{}_{i+1}{}^{})=ϵ_iK_{}^{i+1}{}_{i}{}^{}$$
(16)
Where $`ϵ_i=\pm 1`$. The generalisation redefines $`H_{11}`$ to be the sub-algebra left invariant under the temporal involution, as opposed to the Cartan involution; we denote the local sub-algebra invariant under the temporal involution as $`\widehat{H}_{11}`$. This redefinition allows $`\widehat{H}_{11}`$ to include non-compact generators, and in this way to differentiate between temporal and spatial coordinates, imposing a signature on the sub-algebra and a Lorentzian invariance. The information about which coordinates are timelike is carried by the new variable $`ϵ_i`$. For example, we may impose a $`(1,10)`$ signature where $`x^1`$ is the temporal coordinate, by taking $`ϵ_1=1`$ and $`ϵ_i=1`$ for the remaining spatial coordinates, giving:
$$\widehat{\mathrm{\Omega }}(K_{}^{1}{}_{2}{}^{})=K_{}^{2}{}_{1}{}^{},\widehat{\mathrm{\Omega }}(K_{}^{i}{}_{i+1}{}^{})=K_{}^{i+1}{}_{i}{}^{}\text{For}2i10$$
(17)
We find a basis for the local denominator algebra consisting of both compact and non-compact generators that is invariant under the temporal involution:
$`\widehat{\mathrm{\Omega }}(E_1+F_1)`$ $`=ϵ_1F_1ϵ_1E_1=E_1+F_1`$ (18)
$`\widehat{\mathrm{\Omega }}(E_iF_i)`$ $`=ϵ_iF_i+ϵ_iE_i=E_iF_i\text{For}2i10`$
It was recently pointed out by Keurentjes that the temporal involution does not commute with the Weyl reflections, $`S_i`$, where,
$$S_i\beta \beta 2\frac{(\alpha _i,\beta )}{(\alpha _i,\alpha _i)}\alpha _i$$
(19)
Under the action of the Weyl group, the choice of local sub-algebra is not preserved and we obtain a new set of $`ϵ_i`$’s, corresponding to a different temporal involution $`\widehat{\mathrm{\Omega }}^{}`$. This idea is described in detail elsewhere and may be summarised as
$`S_i\widehat{\mathrm{\Omega }}(K_{}^{j}{}_{j+1}{}^{})=S_i(ϵ_jK_{}^{j}{}_{j+1}{}^{})=ϵ_j\rho _jK_{}^{j+1}{}_{j}{}^{}\widehat{\rho }_j\mathrm{\Omega }^{}(K_{}^{j}{}_{j+1}{}^{})`$ (20)
Where $`S_iK_{}^{j}{}_{j+1}{}^{}=\rho _jK_{}^{j+1}{}_{j}{}^{}`$ and $`\rho _j=\pm 1`$ arises because $`K_{}^{j}{}_{j+1}{}^{}`$ are representations of the Weyl group up to a sign. A consequence of the new temporal involution is that a new set of compact and non-compact generators form the basis of the local denominator sub-algebra $`\widehat{H}_{11}^{}`$, potentially corresponding to a new set of temporal and spatial coordinates and signature.
The Weyl reflections corresponding to nodes on the the gravity line of $`E_{11}`$ preserve the signature of spacetime while the reflection in the exceptional root $`\alpha _{11}`$ can change it. Keurentjes makes use of a $`_2`$ valued function on the root lattice, $`f`$, which encodes the values of the $`ϵ_i`$’s. We may regard the function $`f`$ as a member of the weight space, it may be written $`f=_in_i\lambda _i`$ where $`\lambda _i`$ is in the weight space. Its action on the simple roots is $`(f,\alpha _i)=n_i`$ where $`n_i`$ take the values $`0`$ or $`1`$, a value $`f(\alpha _i)=1`$ corresponds to a Chevalley generator $`K_{}^{i}{}_{i+1}{}^{}`$ which has $`ϵ_i=1`$. Put more simply, one of $`x^i`$ or $`x^{i+1}`$ is timelike and the other spacelike. Alternatively $`f(\alpha _i)=0`$ implies that the mixed coordinates, $`x^i`$ or $`x^{i+1}`$, are of the same type, either both timelike or both spacelike. It is worth highlighting that the root associated with the group element does not determine $`f`$, to obtain a putative solution we must specify both a position in the root lattice (a root) and a vector ($`f`$) encoding the signature; hitherto the group element has been used to find solutions in signature $`(1,10)`$.
Keurentjes offers a useful shorthand notation for following the effect of Weyl reflections upon the function $`f`$, which we will describe here before making some use of it. First the values of the function are written out on the Dynkin diagram, with the value $`f(\alpha _i)`$ written in the position of $`\alpha _i`$ on the diagram. As an example a $`(1,10)`$ signature might have the diagram:
$`0`$
$`0001100`$ $`000`$
$`sssstsss`$ $`sss`$
We have indicated beneath the diagram with a series of $`s`$ (spatial) and $`t`$’s (temporal) the nature of the coordinates obtained by commencing with a spacelike $`x^1`$ on the far left of the gravity line. But we may also consider the case where $`x^1`$ is a temporal coordinate, giving a mostly timelike set of coordinates. In general each signature diagram is ambiguous, representing both $`(t,s,\pm )`$ and $`(s,t,\pm )`$, but as we shall see these signatures do not always contain related solutions, so some care must be taken to specify the nature of one of the coordinates so that a signature diagram is not ambiguous.
The ten values of $`f`$ on the gravity line gives the nature of all eleven dimensions of spacetime. The value of $`f(\alpha _{11})`$ plays no part in this, but it is argued in that it determines the sign in front of the kinetic term in the action derived from the non-linear realisation. Specifically, in this paper, we impose the choice that $`f(\alpha _{11})=0`$ corresponds to the usual minus sign in front of $`F^2`$, while the alternative implies a plus sign in front of the $`F^2`$ term. It will be useful to follow the convention and label our signatures as $`(t,s,\pm )`$, where we will always write the number of timelike coordinates first and where $`{}_{}{}^{}+_{}^{}`$ implies that the sign of $`F^2`$ is negative, and $`{}_{}{}^{}_{}^{}`$ that our action has a positive $`F^2`$ term.
Let us find the effect of a Weyl reflection, $`S_i`$, on the function $`f`$, we have,
$`S_i(f)`$ $`={\displaystyle \underset{j}{}}n_j\lambda _j{\displaystyle \underset{j}{}}n_j(\alpha _i,\lambda _j)\alpha _i`$
$`{\displaystyle \underset{j}{}}m_j\lambda _j`$ (21)
Where $`m_j`$ are the components of $`S_i(f)`$. Taking the inner product with $`\alpha _k`$ gives the relation between $`n_k`$ and $`m_k`$
$`m_k=n_kn_iA_{ki}`$ (22)
Thus we find Keurentjes’ diagrammatic prescription for following signature change: to apply $`S_i`$ to $`f`$ we simply add the value of $`f(\alpha _i)`$ to all the nodes it is connected to, and subsequently reduce modulo two. The reduction modulo two comes from the size of the fundamental lattice which has edge length $`\left|\alpha _i\right|+\left|+\alpha _i\right|=2\left|\alpha _i\right|`$, so the sub-algebra repeats with this unit and we need only consider a version of it upto $`\pm n2\left|\alpha _i\right|,n`$. For example in the signature diagram above, $`S_i`$ where $`i=\{1\mathrm{}3,6\mathrm{}11\}`$ have no effect upon the signature, whereas $`S_4`$ and $`S_5`$ bring about the following two signature diagrams respectively.
$`0`$
$`0011000`$ $`000`$
$`ssstssss`$ $`sss`$
$`0`$
$`0000110`$ $`000`$
$`ssssstss`$ $`sss`$
### 3.1 Invariant Roots and Signature Orbits
Single brane solutions have been found from the decomposition of $`E_{11}`$ with respect to its gravity line, encoded within the group element (1). The prescription for finding the brane solution used so far made use of the lowest weight generator in a representation. It was not clear how to interpret the other weights in a representation, in particular the highest weight. For the antisymmetric representations that we shall consider herein, other generators are reached from the lowest weight by a series of Weyl reflections, which in the light of the previous discussion implies a potential signature change. If we commence with the $`R^{123}`$ and $`R^{123456}`$ generators in $`(1,10,+)`$, where $`x^1`$ is the time coordinate, and raise them to their highest weights with the Weyl reflections $`S_0^1=(S_1)(S_2S_1)\mathrm{}(S_{10}S_9\mathrm{}S_1)`$<sup>3</sup><sup>3</sup>3$`S_0=(S_1\mathrm{}S_{10})\mathrm{}(S_1S_2)(S_1)`$ is the series of Weyl reflections that takes the highest weight to the lowest weight for all representations of $`A_{10}`$, so that $`S_0R^{91011}=R^{123}`$ and $`S_0R^{67891011}=R^{123456}`$., we find that the effect on the signature diagram is,
$`0`$ $`1`$
$`1000000`$ $`000`$ $`\stackrel{S_0^1}{}0000000`$ $`001`$ (23)
$`tsssssss`$ $`sss`$ $`tttttttt`$ $`tts`$
Up to insisting that $`x^1`$ is preserved as a time coordinate, we obtain the signature $`(10,1,)`$, where the singled out spatial coordinate is in the longitudinal sector, specifically the spatial coordinate here is $`x^{11}`$, and the gauge fields in each case are $`A_{91011}`$ and $`A_{67891011}`$. From the observations of section 2.1, it is known that the $`M2`$ brane has a related solution in $`(10,1,)`$ with signature components $`[(2,1),(8,0)]`$, alternatively the $`M5`$ brane does not have a related solution $`[(5,1),(5,0)]`$ in $`(10,1,)`$.
Some Weyl reflections may change the signature diagram, and even the signature, without changing the root, so there is an ambiguity about which signature the root and its associated solution exist in. We now consider the example of the exceptional root $`\alpha _{11}`$, associated with the generator $`R^{91011}`$, the highest weight of the ‘$`M2`$ representation’ and find what we shall refer to as its signature orbit.
#### 3.1.1 Membrane Solution Signatures
Let us first consider the trivial Weyl reflections of the gravity line on the root $`\alpha _{11}`$. It is noted that $`\alpha _{11}`$ is invariant under the reflections $`\{S_1,\mathrm{}S_7\}`$ and $`\{S_9,S_{10}\}`$ and we may apply any number of these reflections, without altering the root, although we may trivially change the signature diagram but not the signature. Furthermore, we can observe from the signature diagram of the highest weight, shown on the right of (23), that as only $`f(\alpha _{10})=1`$ along the gravity line, only a series of reflections composed of $`\{S_9,S_{10}\}`$ may have an effect on the signature diagram without effecting the root. Explicitly, the only possible different signature diagrams that may be reached without changing the root are,
$`1`$ $`1`$
$`0000000`$ $`001`$ $`\stackrel{S_{10}}{}0000000`$ $`011`$
$`tttttttt`$ $`tts`$ $`tttttttt`$ $`tst`$
$`1`$ $`1`$ (24)
$`0000000`$ $`001`$ $`\stackrel{S_9S_{10}}{}0000000`$ $`110`$
$`tttttttt`$ $`tts`$ $`tttttttt`$ $`stt`$
Here $`x^1`$ has been held as a temporal coordinate. The interpretation is that the trivial Weyl reflection in the roots of the gravity line shift the singled-out coordinate between $`\{x^9,x^{10},x^{11}\}`$, the longitudinal brane coordinates. While it is always true that the gravity line Weyl reflections do not alter the signature, it is not generally true that they do not alter the value of $`f(\alpha _{11})`$. For example, consider the root, $`\alpha _8+\alpha _9+\alpha _{10}+\alpha _{11}`$, associated to the generator $`R^{8910}`$, for which only a series of reflections composed of $`\{S_8,S_9\}`$ may alter the signature diagram without changing the root. In this example the reflection $`S_8`$ may change the value of $`f(\alpha _{11})`$ in addition to shifting the singled-out coordinate amongst the longitudinal coordinates. In general, the gravity line, or trivial, Weyl reflections preserve a signature $`(t,s)`$ but do not necessarily preserve whether we are working with a $`F^2`$ or a $`+F^2`$ theory.
We now turn our attention to the non-trivial signature changes that may be applied to a root, $`\beta `$, without altering it and we outline here a prescription for finding alternative signatures without altering the root. In order to consider the effects of the reflection $`S_{11}`$ on the signature we first transform our root to a new root that is invariant under $`S_{11}`$, we call the series of Weyl reflections applied to acheive this $`U`$. Furthermore, we restrict ourselves to using only trivial Weyl reflections in this transformation, $`U`$, so that we only effect a non-trivial signature change after we have transformed to an $`S_{11}`$ invariant root. These restrictions identify a unique $`S_{11}`$ invariant root for a given non-zero coefficient of $`\alpha _{11}`$ in the simple root expansion of the root, $`\beta `$, or the level of $`\beta `$. At this stage $`S_{11}`$ may be applied without changing the root, but with the potential of altering the signature. Our original root in the new signature may be reobtained by applying $`U^1`$. These steps allow an algorithmic exploration of the related signatures for a specific root.
For $`\alpha _{11}`$, at level one, the $`S_{11}`$ invariant root is $`\alpha _7+2\alpha _8+\alpha _9+\alpha _{11}`$. It is obtained from $`\alpha _{11}`$ by reflections $`S_8S_7S_9S_8U`$, so that $`\alpha _{11}=U^1S_{11}U\alpha _{11}`$, and a new class of signature diagrams is obtained that is not trivially related to the first class. This process is repeated for every trivially related signature diagram and in this way all possible Weyl reflections preserving $`\alpha _{11}`$ are applied and the associated set of signature diagrams including $`(10,1,)`$ is obtained, we call this set the signature orbit of $`\alpha _{11}`$. An equivalent approach would be to apply all possible trivial reflections to $`S_{11}U\alpha _{11}`$, before transforming back to $`\alpha _{11}`$ with $`U^1`$. This procedure is simply completed by a computer program, with the results shown in table 1, where we have only listed the cases where we have taken $`x^1`$ to be a temporal coordinate. The equivalent signature orbit where $`x^1`$ is taken to be spacelike is found by inverting all signatures, while keeping $`f(\alpha _{11})`$ constant. The signature orbits for the lowest weight, associated with generator $`R^{123}`$, are found by applying $`(S_1\mathrm{}S_{10})(S_1\mathrm{}S_9)(S_1\mathrm{}S_8)`$ and the results are shown in table 2.
We have found the signature orbits of the highest and lowest weights in the membrane representation, but we have not checked whether each prescribed signature offers a solution to the Einstein and gauge equations. Using our observations of section 2.1 it is, in fact, a quick exercise to check all signatures and see if they offer a solution. We have evaluated $`\kappa `$, defined in equation (14) for each putative solution given in the tables, and wherever we find $`\kappa =1`$ we have a solution of the Einstein equations. If we count the number of trivially related signatures for these cases we notice that exactly half of the total signature orbit are solutions, that is $`3+7+105+21=136`$ solutions associated to the generator $`R^{91011}`$. For spacelike $`x^1`$ we find $`21+105+7+3=136`$ solutions too. For the lowest weight generator $`R^{123}`$ considered in table 2, we again find $`1+15+70+35+13+2=136`$ solutions for timelike $`x^1`$, and similarly $`2+13+70+35+15+1=136`$ solutions for spacelike $`x^1`$.
#### 3.1.2 Fivebrane Solution Signatures
We now turn our attention to the representation that gives the $`M5`$ solution. Its highest weight generator is $`R^{67891011}`$ which is associated with the root $`\beta \alpha _6+2\alpha _7+3\alpha _8+2\alpha _9+\alpha _{10}+2\alpha _{11}`$. Only the Weyl reflections $`\{S_5,S_{11}\}`$ alter the root. The level two $`S_{11}`$-invariant root is obtained from $`\beta `$ by acting upon it with the series of reflections given by $`US_{10}S_9S_8S_7S_6S_5`$ and the lowest weight representation, with generator $`R^{123456}`$, is obtained under the reflection $`(S_1\mathrm{}S_{10})\mathrm{}(S_1\mathrm{}S_5)`$. The signature orbits are shown in tables 3 and 4 respectively. Again the cases where $`\kappa =1`$ give solutions, but we note that there is a difference to the $`M2`$ case when we consider the solutions for spacelike $`x^1`$. As before the spacelike $`x^1`$ case is found by a global signature inversion, however for the $`M5`$ case this does not bring about a change of sign in $`\kappa `$. The solutions for spacelike and timelike $`x^1`$ are no longer complementary but identical. For the highest weight representation we find $`3+3+12+60+40+3+12+3=136`$ solutions, the same number of solutions as for the $`M2`$ case, but for the lowest weight we find $`5+1+3+15+40+40+15+3+5+1=128`$ solutions.
#### 3.1.3 $`pp`$-Wave Solution Signatures
The $`pp`$-wave is treated in the same manner. Its highest weight is associated to the root $`\beta =\alpha _1+\mathrm{}\alpha _{10}`$ and its lowest weight is associated to the root $`S_0\beta =\beta `$. In both cases only the Weyl reflection $`\{S_1,S_{10},S_{11}\}`$ alter the root. However, the negative roots have generators of the form $`K_{}^{a}{}_{b}{}^{}`$, where $`a>b`$, which are projected out of the general group element of $`E_{11}`$, see equation (2.24) in , so the lowest weight representation has generator $`K_{}^{1}{}_{2}{}^{}`$, and associated root $`\alpha _1`$. We note that $`\alpha _1`$ is only altered by the Weyl reflections $`\{S_1,S_2\}`$, and is related to the highest weight by $`\alpha _1=S_2S_3\mathrm{}S_{10}\beta `$. There are a number of possible level 0 $`S_{11}`$-invariant roots, we make use of $`US_8S_7S_6S_5S_4S_3S_2S_1`$ to transform $`\beta `$ into $`\alpha _9+\alpha _{10}`$ and then effect the signature-changing Weyl reflection, $`S_{11}`$, before transforming back to $`\beta `$. The signature orbits containing the $`M`$-theory $`pp`$-wave coming from the root associated to the highest weight and the lowest weight with a positive root generator are listed in tables 5 and 6 respectively.
Our analysis of solutions to the Einstein equations from $`\kappa `$ is not appropriate for the $`pp`$-wave, instead we have a $`pp`$-wave solution if the ansatz given in appendix A is satisfied . From the tables we obtain $`1+2+7=10`$ solutions for the $`pp`$-wave for each weight of the representation and, in addition, we note that there is no associated $`M^{}`$-theory $`pp`$-wave, in our signature orbits.
### 3.2 $`S`$-Branes from a Choice of Local Sub-Algebra
Spacelike branes or $`S`$-branes were discovered as a constituent of string theory by Gutperle and Strominger , who argued that they were a timelike kink in the tachyon field on the world volume of an unstable D-brane, or D-brane anti-D-brane pair. There is a wealth of literature on the rolling tachyon whose association with $`S`$-branes was first highlighted by Sen. Supergravity $`S`$-branes were found by Chen, Gal’tsov and Gutperle in arbitrary dimension, D, in and in D=10 by Kruczenski, Myers and Peet in . These solutions were shown to be equivalent under a coordinate transformation by Bhattacharya and Roy in . General $`S`$-brane solutions in eleven dimensions were also found in reference , where intersection rules are also considered. We concentrate here on simply identifying the spacelike branes of $`M`$-theory and the related solutions in other theories that may be constructed from the brane spectrum of $`E_{11}`$.
The group element (1) has been used to find brane solutions in exotic signatures by Weyl reflecting the known electric brane solutions of $`M`$-theory. The group element itself does not know which signature its associated solution exists in, indeed signature information comes from the choice of local sub-algebra. In our solutions we have singled out electric field strengths, those which always have a temporal coordinate, and used these as a starting point for the signature orbits of the previous section. It was observed in section 3.1.1 that the Weyl reflections that did not alter the root kept the temporal coordinate on the brane world-volume for the $`M`$-theory solutions, presenting an obstacle to finding $`S`$-branes from the electric solutions by Weyl reflecting the group element (1). Let’s look at this in more detail. If one rotates the coordinates, using a Weyl reflection, to obtain a new root and associated gauge field, the effect of $`S_i`$ on the expansion of $`\beta H`$ in 1 is to interchange $`K_{}^{i}{}_{i}{}^{}`$ and $`K_{}^{i+1}{}_{i+1}{}^{}`$, where $`i=1\mathrm{}10`$, and in terms of the coordinate indices on the gauge potential the indices $`x^i`$ and $`x^{i+1}`$ are swapped. For example consider the lowest weight of the $`M2`$ representation with gauge field $`A_{123}`$, whose indices are transformed in the following manner,
$`A_{123}`$ $`\stackrel{S_1}{}A_{123}`$
$`A_{123}`$ $`\stackrel{S_2}{}A_{123}`$ (25)
$`A_{123}`$ $`\stackrel{S_3}{}A_{124}`$
If we pick $`x^3`$ to be the timelike coordinate we might think that to remove it from the gauge field would require a Weyl reflection in $`S_3`$. The effect of $`S_3`$ on the signature diagram is to change the timelike coordinate from $`x^3`$ to $`x^4`$,
$`0`$ $`0`$
$`0110000`$ $`000`$ $`\stackrel{S_3}{}0011000`$ $`000`$
$`sstsssss`$ $`sss`$ $`ssstssss`$ $`sss`$
Consequently the new gauge field $`A_{124}`$ remains electric and similar considerations for each possible choice of timelike coordinate show that an electric gauge field remains electric under Weyl reflections. Importantly the $`S`$-brane solutions of $`M`$-theory are not related to the electric solutions by Weyl reflections, and are not found in the signature orbits of the usual electric solutions.
However, given any real form of an algebra that leaves a Lorentzian form, e.g. $`t^2+x_1^2+\mathrm{}x_{10}^2`$, invariant we may consider the complex extension of the algebra such that a Euclidean form is left invariant. A specific example of how the generators transform is $`K_{}^{1}{}_{2}{}^{}iK_{}^{1}{}_{2}{}^{}`$ where $`x^1`$ is the temporal coordinate. To reintroduce a Lorentzian symmetry we apply the inverse transformation. For example to make $`x^{10}`$ temporal, we transform $`K_{}^{10}{}_{11}{}^{}iK_{}^{10}{}_{11}{}^{}`$, and obtain a complexified version of the original set of generators preserving a real Lorentzian form, $`t^2+x_1^2+\mathrm{}+x_9^2x_{10}^2`$. The result is that the generators, $`A_{123},A_{123456}`$ and $`K_{}^{1}{}_{2}{}^{}`$ used to find the electric solutions of appendix A become $`iA_{123},iA_{123456}`$ and $`iK_{}^{1}{}_{2}{}^{}`$, as would be expected by the Wick rotations as in equation 11, and all their indices are now spacelike.
Equivalently, there is a different choice of local sub-algebra with a different set of generators, all real, that preserve the same Lorentzian form, that give an identical theory but with a different sign in front of $`F^2`$ in the action. For example the $`S2`$ and $`S5`$-branes are solutions in signature $`(1,10,)`$ with a real set of generators. Our Weyl reflections lead us to pick out the real form of the sub-algebra, and the sign of $`F^2`$. The $`pp`$-wave has a null field strength, hence we make use of the complex generators to find its spacelike solution. These spacelike solutions are verified in appendix B.
We may commence with the $`S2`$ and $`S5`$-brane solutions of $`M`$-theory given in appendix B, encoded in the group element and find their signature orbits. This solution is identical to commencing using a local subgroup, $`H_{11}`$, whose temporal coordinate with respect to our ansatz (4) is not part of the brane world-volume. The metric for the $`M`$-theory spacelike solution takes the same form as our ansatz (4), explicitly,
$$ds^2=A^2(\underset{j=1}{\overset{j=p}{}}dx_j^2)+B^2(du^2+\underset{b=1}{\overset{b=d}{}}dy_b^2)$$
(26)
Where $`A`$ and $`B`$, are as defined in (5), using the harmonic function,
$$N_{(1,Dp2)}=1+\frac{1}{Dp3}\sqrt{\frac{\mathrm{\Delta }}{2(D2)}}\frac{𝐐}{\widehat{r}^{(Dp3)}}$$
(27)
Where $`\widehat{r}^2=(u^1)^2+(y^1)^2+\mathrm{}(y^{Dp2})^2`$. For the $`S2`$ and $`S5`$ branes we have the associated field strengths
$`F_{x_1x_2x_3\widehat{r}}=`$ $`4_{[\widehat{r}}A_{x_1x_2x_3]}=_{\widehat{r}}A_{x_1x_2x_3}=_{\widehat{r}}N_{(1,7)}^1`$ (28)
$`F_{x_1x_2x_3x_4x_5x_6\widehat{r}}=`$ $`7_{[\widehat{r}}A_{x_1x_2x_3x_4x_5x_6]}=_{\widehat{r}}A_{x_1x_2x_3x_4x_5x_6}=_{\widehat{r}}N_{(1,4)}^1`$
Weyl reflections of these solutions then give new orbits of possible solutions and completes the range of signature configurations related by $`E_{11}`$, in that we find solutions of $`M`$-theory where the temporal coordinate may be any of $`\{x^1\mathrm{}x^{11}\}`$ for each weight. We list these signature orbits for the case of the highest and lowest weights of the $`S2`$-brane in the tables 7, 8 and similarly for the $`S5`$-brane in tables 9, 10.
The $`pp`$-wave solution distinguishes three sets of coordinates, namely a longitudinal coordinate, a transverse coordinate and the nine remaining coordinates in 11-dimensions, $`\mathrm{\Omega }_9`$. Consequently there are two alternative choices of local sub-algebra that may be made, the first introduces a transverse time coordinate and the second introduces a time coordinate into $`\mathrm{\Omega }_9`$, which we then relabel $`\mathrm{\Omega }_{(1,8)}`$. The former case is similar to the original $`pp`$-wave solution under an interchange of $`KK`$, having a line element:
$$ds^2=(1+K)dt_{1}^{}{}_{}{}^{2}+(1K)dx_{1}^{}{}_{}{}^{2}+2Kdt_1dx_1+d\mathrm{\Omega }_9^2$$
(29)
Consequently this second choice of local sub-algebra leads to the same solution as the usual sub-algebra, but the $`pp`$-wave in this case is completely out of phase with the original $`pp`$-wave. This solution is still dependent on a static harmonic function.
There is a time-dependent solution, which we label the $`Spp`$-wave, coming from the choice of local sub-algebra that introduces a time coordinate into $`\mathrm{\Omega }_9`$, which we indicate by $`\mathrm{\Omega }_{(1,8)}`$. The solution is given in appendix B. We list the signature orbit of the highest weight case in table 11, as in the case of the $`pp`$-wave the lowest weight signature orbit is identical.
Let us count the solutions in the $`S2`$-brane signature orbit from the M, $`M`$ and $`M^{}`$-theories<sup>4</sup><sup>4</sup>4That is, only those with signature $`(1,10,\pm )`$,$`(10,1,\pm )`$,$`(2,9,\pm )`$,$`(9,2,\pm )`$,$`(5,6,\pm )`$ and $`(6,5,\pm )`$. From the highest weight signature orbit in table 7, we find $`1+63+35+21=120`$ solutions, where $`x^1`$ is timelike and $`7+105+21+3=136`$ where $`x^1`$ is spacelike. Recollect that we found $`136`$ solutions related to the equivalent highest weight $`M2`$ signature orbit, giving a total of $`256`$ solutions with timelike $`x^1`$ and $`272`$ with spacelike $`x^1`$; in all we have $`528`$ solutions from the root associated to the $`R^{91011}`$ generator. Similarly for the lowest weight associated to the $`S2`$-brane we find $`3+50+31+6+5+25=120`$ solutions with timelike $`x^1`$ and $`5+62+25+10+3+31=136`$ with spacelike $`x^1`$, giving a total of $`528`$ solutions from the root associated to the $`R^{123}`$ generator.
The $`S5`$-brane signature orbit coming from the highest weight given in table 9 has $`4+1+9+36+36+24+9=119`$ solutions of the three $`M`$-theories for both timelike and spacelike $`x^1`$. If we include the $`136`$ solutions from the highest weight of the $`M5`$ representation for both timelike and spacelike $`x^1`$, we find a total of $`510`$ solutions associated to the $`R^{67891011}`$ generator. Perhaps most interesting are the results from the lowest weight generator $`R^{123456}`$; from the $`S5`$ signature orbit in table 8 we find $`2+30+60+20+10+6=128`$ solutions to the three $`M`$-theories for both timelike and spacelike $`x^1`$. Recalling that we earlier counted $`128`$ solutions from the lowest weight signature orbit containing the $`M5`$ brane for each choice of $`x^1`$ giving a total of $`512`$ solutions associated to the $`R^{123456}`$ generator. We note that there is a difference between both the total number and type of solutions associated the generator $`R^{a_1\mathrm{}a_6}`$ at different weights, in contrast to the $`R^{a_1a_2a_3}`$ generator to which a consistent set of solutions is associated at each weight.
## 4 Discussion
It has previously been noted that the adjoint representation of $`E_{11}`$ may be decomposed to representations $`A_{10}`$ whose lowest weight has an associated group element (1) giving half BPS solutions to the usual $`(1,10)`$ theory. In this paper we considered representations of $`E_{11}`$ which were not the lowest weights under $`A_{10}`$. These can be related to the lowest weight by a series of $`E_{11}`$ Weyl reflections. The group element does indeed give solutions, but these are to theories which are not only in $`(1,10)`$, but also $`(2,9)`$, $`(5,6)`$ and their inverses, the transformation of signature being induced by Weyl reflections. We found that there exists an ambiguity in the signature for each solution due to there being a series of Weyl reflections that may preserve the root uniquely associated to the weight but may change the signature. We also carried out a general search for all the alternative signatures in which we found the equivalent of the $`M2`$, $`M5`$ and $`pp`$-wave solutions. The results of this search is probably not completely understood as some, and in most cases half, of the cases are not solutions, but would be under a change of sign of the $`F^2`$ term in the action for the theory, sign changes of $`F^2`$ being specified by the Weyl reflections .
In particular the solutions we have found in signatures $`(2,9)`$ and $`(5,6)`$, derived from solutions of $`M`$-theory using the Weyl reflections of $`E_{11}`$, reproduce the solutions of $`M`$ and $`M^{}`$ theories . The set of solutions to these theories bear an intimate relation to the string of weights of each representation, and in shifting between weights the number of solutions to each theory is, in general, not preserved. We considered the highest and lowest weight associated to the generators $`R^{a_1a_2a_3}`$ and $`R^{a_1\mathrm{}a_6}`$ in sections 3.1.1 and 3.1.2 and counted the number of solutions in each case. It was noted that while the representation of the generator $`R^{a_1a_2a_3}`$ provided 136 solutions from both its highest and lowest weight, the representation of the generator $`R^{a_1\mathrm{}a_6}`$ provided 136 solutions from its highest weight, but only 128 solutions from its lowest weight.
The solution counting differentiated between trivially related versions of the same solution, for example the $`M2`$-brane in $`(1,10)`$ with signature $`[(1,2),(0,8)]`$ contributes three related solutions to the count coming from the three $`(=\left(\genfrac{}{}{0pt}{}{3}{1}\right)\times \left(\genfrac{}{}{0pt}{}{8}{0}\right))`$ different ways to associate a single timelike coordinate with the three coordinates $`x^\mu `$ of the brane world-volume. Of course in the signature $`(1,10)`$ there are eleven ways to choose the local sub-algebra, so one might expect to find eleven solutions in $`(1,10)`$ associated to each weight of the generator $`R^{a_1a_2a_3}`$ instead of three. The remaining eight solutions are the trivially related $`S2`$-brane solutions of signature $`[(0,3),(1,7)]`$.
The arguments reviewed herein lead to the conclusion that readers who are predisposed to the $`E_{11}`$ conjecture are also encouraged to take up the $`M`$ and $`M^{}`$-theories. Furthermore there appears to be no clear way to single out any of these three $`M`$-theories, as $`E_{11}`$ treats them all on an equal footing through the choice of local sub-algebra and the Weyl reflections.
In ordinary quantum field theory in Minkowski space one does require solutions in Euclidean space, namely instantons, which occur when the particle enters a region forbidden by energy considerations. One could imagine a similar scenario for branes and, in particular, one brane in the potential of another. We note that for every brane solution in the signature orbits given in this paper there is an exact pairing between putative solutions in $`[(p,q),(c,d)]`$ and $`[(q,p),(c,d)]`$, that is, under an signature inversion of the brane coordinates. For the fivebrane case either both or none of these signatures will carry solutions, while for the membrane case only one of the two signatures will carry a solution. Specifically the M5-brane, $`[(1,5),(0,5)]=(1,10)`$, is paired with a solution in $`[(5,1),(0,5)]=(5,6)`$ which is an $`M^{}`$-theory solution; and the M2-brane, $`[(1,2),(0,8)]=(1,10)`$, has a solution in $`[(2,1),(0,8)]=(2,9)`$, with a ’$`+F^2`$’ term, which is an $`M`$-theory solution. From this viewpoint, we are always considering the theory in a $`(1,10)`$ signature but certain calulations would require theories in other signatures, namely $`(2,9)`$ and $`(5,6)`$.
## Acknowledgements
PPC thanks Joe Gillard and Vid Stojevic for fruitful conversations.
This research was supported by the PPARC grants PPA/S/S/2003/03644, PPA/G/O/2000/00451 and PPARC senior fellowship PPA/Y/S/2002/001/44.
## Appendix A Electric Brane Solutions from $`E_{11}`$
There is a substantial literature deriving single brane solutions in generic supergravity theories , for a review see . The $`M2`$, $`M5`$ and $`pp`$-wave solutions of $`M`$-theory have been derived from $`E_{11}`$ in .
For $`E_{11}`$ we have no dilaton and find the single brane solutions to be determined from a truncated action of the form
$$A=\frac{1}{16\pi G_{11}}d^{11}x\sqrt{g}(R\frac{1}{2.4!}F_{a_1\mathrm{}a_n}F^{a_1\mathrm{}a_n})$$
(30)
$`F_{a_1\mathrm{}a_n}`$ is a general $`n`$-form field strength formed in the non-linear realistion from the gauge fields. From we have two field strengths, which are dual to each other, which we consider as arising from the a 3-form and a 6-form gauge field
$`F_{t_1x_1x_2r}=`$ $`4_{[r}A_{t_1x_1x_2]}=_rA_{t_1x_1x_2}=_rN_{(0,8)}^1`$ (31)
$`F_{t_1x_1x_2x_3x_4x_5r}=`$ $`7_{[r}A_{t_1x_1x_2x_3x_4x_5]}=_rA_{t_1x_1x_2x_3x_4x_5}=_rN_{(0,5)}^1`$
The appropriate Einstein and gauge equations may be found by setting $`\varphi =0`$, $`\widehat{\delta }_{}^{t_i}{}_{t_i}{}^{}=1`$, $`\widehat{\delta }_{}^{x_i}{}_{x_i}{}^{}=p`$, $`\widehat{\delta }_{}^{u_a}{}_{u_a}{}^{}=0`$ and $`\widehat{\delta }_{}^{y_a}{}_{y_a}{}^{}=Dp1`$ in equations (3) and (9). The line element of an electric solution is derived from the solution generating group element and coincides with line element for single brane BPS solutions specified by equations (5) . The non-linear realisation decomposes $`E_{11}`$ with respect to its longest gravity line $`A_{10}`$ obtaining a gravitational theory in 11 dimensions and an infinite array of irreducible representations. These representations are classified by level, the coefficient of the exceptional root $`\alpha _{11}`$ associated to the representation; all the usual solutions of eleven dimensional supergravity appear below level three. In the group element above, $`\beta `$ is the root associated to the lowest weight of each representation. The following solutions are found,
### A.1 The $`M2`$-Brane
The line element of the $`M2`$-brane solution is
$$ds^2=N_{(0,8)}^{\frac{2}{3}}(dt_{1}^{}{}_{}{}^{2}+dx_1^2+dx_2^2)+N_{(0,8)}^{\frac{1}{3}}(dy_1^2+\mathrm{}dy_8^2)$$
(32)
Giving a metric
$$g_{\mu \nu }=\begin{array}{ccccccccc}& & t_1& x_1& x_2& y_1& \mathrm{}& y_8& \\ t_1& (\mathrm{}& N_{(0,8)}^{\frac{2}{3}}& 0& 0& 0& 0& 0& )\mathrm{}\\ x_1& 0& N_{(0,8)}^{\frac{2}{3}}& 0& 0& 0& 0\\ x_2& 0& 0& N_{(0,8)}^{\frac{2}{3}}& 0& 0& 0\\ y_1& 0& 0& 0& N_{(0,8)}^{\frac{1}{3}}& 0& 0\\ \mathrm{}& 0& 0& 0& 0& \mathrm{}& 0\\ y_8& 0& 0& 0& 0& 0& N_{(0,8)}^{\frac{1}{3}}\end{array}$$
So $`\sqrt{g}=N_{(0,8)}^{\frac{1}{3}}`$ and we see that the gauge equation is satisfied by $`F_{t_1x_1x_2y_a}=_{y_a}N_{(0,8)}^1`$ in the following manner,
$`_{y_a}(\sqrt{g}F^{t_1x_1x_2y_a})=`$ $`_{y_a}(N_{(0,8)}^{\frac{1}{3}}g^{t_1t_1^{}}g^{x_1x_1^{}}g^{x_2x_2^{}}g^{y_ay_a^{}}F_{t_1^{}x_1^{}x_2^{}y_a^{}})`$
$`=`$ $`_{y_1}(N_{(0,8)}^2F_{t_1x_1x_2y_1})+\mathrm{}_{y_8}(N_{(0,8)}^2F_{t_1x_1x_2y_8})`$
$`=`$ $`_{y_a}(N_{(0,8)}^2_{y_a}N_{(0,8)}^1)\widehat{\delta }_{}^{y_a}{}_{y_a}{}^{}`$
$`=`$ $`_{y_a}_{y_a}N_{(0,8)}\widehat{\delta }_{}^{y_a}{}_{y_a}{}^{}`$
$`=`$ $`0`$
In the last line we have used the fact that $`N_{(0,8)}`$ is an harmonic function in $`y_a`$, as can be checked from its definition in equation (6). We see that the Einstein equations are satisfied by checking that the right-hand-side of each equation equals the curvature term for the electric solution. A term that appears frequently is $`F^{t_1x_1x_2y_a}F_{t_1x_1x_2y_a}`$ which we evaluate at the outset
$`F^{t_1x_1x_2y_a}F_{t_1x_1x_2y_a}`$ $`=F^{t_1x_1x_2y_1}F_{t_1x_1x_2y_1}+\mathrm{}F^{t_1x_1x_2y_8}F_{t_1x_1x_2y_8}`$ (33)
$`=N_{(0,8)}^{\frac{1}{3}}N_{(0,8)}^2(F_{t_1x_1x_2y_1})^2\mathrm{}N_{(0,8)}^{\frac{1}{3}}N_{(0,8)}^2(F_{t_1x_1x_2y_8})^2`$
$`=8N_{(0,8)}^{\frac{1}{3}}N_{(0,8)}^2(_{y_a}N_{(0,8)})^2`$
We proceed to check the Einstein equations,
$`({}_{t_i}{}^{t_i}){\displaystyle \frac{1}{2.4!}}(4`$ $`F^{t_1\mu _1\mu _2\mu _3}F_{t_1\mu _1\mu _2\mu _3}{\displaystyle \frac{1}{3}}F^{\mu _1\mu _2\mu _3\mu _4}F_{\mu _1\mu _2\mu _3\mu _4})`$
$`={\displaystyle \frac{1}{2.4!}}(4.3!{\displaystyle \frac{4!}{3}})F^{t_1x_1x_2y_a}F_{t_1x_1x_2y_a}`$ (34)
$`={\displaystyle \frac{8}{3}}N_{(0,8)}^{\frac{1}{3}}N_{(0,8)}^2(_{y_a}N_{(0,8)})^2`$
$`=N_{(0,8)}^{\frac{1}{3}}\{_{y_a}_{y_a}\mathrm{ln}N_{(0,8)}^{\frac{1}{3}}\}\widehat{\delta }_{}^{y_a}{}_{y_a}{}^{}`$
$`R_{}^{t_i}{}_{t_i}{}^{}`$
$`({}_{x_i}{}^{x_i}){\displaystyle \frac{1}{2.4!}}(4`$ $`F^{x_i\mu _1\mu _2\mu _3}F_{x_i\mu _1\mu _2\mu _3}{\displaystyle \frac{1}{3}}F^{\mu _1\mu _2\mu _3\mu _4}F_{\mu _1\mu _2\mu _3\mu _4}\widehat{\delta }_{}^{x_i}{}_{x_i}{}^{})`$
$`={\displaystyle \frac{1}{2.4!}}(4.3!{\displaystyle \frac{4!}{3}})F^{t_1x_1x_2y_a}F_{t_1x_1x_2y_a}\widehat{\delta }_{}^{x_i}{}_{x_i}{}^{}`$ (35)
$`=N_{(0,8)}^{\frac{1}{3}}\{_{y_a}_{y_a}\mathrm{ln}N_{(0,8)}^{\frac{1}{3}}\}\widehat{\delta }_{}^{y_a}{}_{y_a}{}^{}\widehat{\delta }_{}^{x_i}{}_{x_i}{}^{}`$
$`R_{}^{x_i}{}_{x_i}{}^{}`$
$`({}_{y_a}{}^{y_a}){\displaystyle \frac{1}{2.4!}}(4`$ $`F^{y_a\mu _1\mu _2\mu _3}F_{y_a\mu _1\mu _2\mu _3}{\displaystyle \frac{1}{3}}F^{\mu _1\mu _2\mu _3\mu _4}F_{\mu _1\mu _2\mu _3\mu _4})`$
$`={\displaystyle \frac{1}{2.4!}}(4.3!{\displaystyle \frac{4!}{3}}\widehat{\delta }_{}^{y_a}{}_{y_a}{}^{})F^{t_1x_1x_2y_a}F_{t_1x_1x_2y_a}`$ (36)
$`={\displaystyle \frac{5}{6}}N_{(0,8)}^{\frac{1}{3}}N_{(0,8)}^2(_{y_a}N_{(0,8)})^2\widehat{\delta }_{}^{y_a}{}_{y_a}{}^{}`$
$`=N_2^{\frac{1}{3}}\{8_{y_a}_{y_a}\mathrm{ln}N_{(0,8)}^{\frac{1}{6}}3(_{y_a}\mathrm{ln}N_{(0,8)}^{\frac{1}{3}})^2`$
$`6(_{y_a}\mathrm{ln}N_{(0,8)}^{\frac{1}{6}})^2\}\widehat{\delta }^{y_a}_{y_a}`$
$`R_{}^{y_a}{}_{y_a}{}^{}`$
### A.2 The $`M5`$-Brane
The line element of the $`M5`$-brane solution is
$$ds^2=N_{(0,5)}^{\frac{1}{3}}(dt_{1}^{}{}_{}{}^{2}+dx_1^2+\mathrm{}dx_5^2)+N_{(0,5)}^{\frac{2}{3}}(dy_1^2+\mathrm{}dy_5^2)$$
(37)
Giving a metric
$$g_{\mu \nu }=\begin{array}{cccccccccc}& & t_1& x_1& \mathrm{}& x_5& y_1& \mathrm{}& y_5& \\ t_1& (\mathrm{}& N_{(0,5)}^{\frac{1}{3}}& 0& 0& 0& 0& 0& 0& )\mathrm{}\\ x_1& 0& N_{(0,5)}^{\frac{1}{3}}& 0& 0& 0& 0& 0\\ \mathrm{}& 0& 0& \mathrm{}& 0& 0& 0& 0\\ x_5& 0& 0& 0& N_{(0,5)}^{\frac{1}{3}}& 0& 0& 0\\ y_1& 0& 0& 0& 0& N_{(0,5)}^{\frac{2}{3}}& 0& 0\\ \mathrm{}& 0& 0& 0& 0& 0& \mathrm{}& 0\\ y_5& 0& 0& 0& 0& 0& 0& N_{(0,5)}^{\frac{2}{3}}\end{array}$$
So $`\sqrt{g}=N_{(0,5)}^{\frac{2}{3}}`$ and we see that the gauge equation is satisfied by $`F_{t_1x_1\mathrm{}x_5y_a}=_{y_a}N_{(0,5)}^1`$ in the following manner,
$`_{y_a}(\sqrt{g}F^{t_1x_1\mathrm{}x_5y_a})=`$ $`_a(N_{(0,5)}^2_aN_{(0,5)}^1)\widehat{\delta }_{}^{y_a}{}_{y_a}{}^{}`$
$`=`$ $`_{y_a}_{y_a}N_{(0,5)}\widehat{\delta }_{}^{y_a}{}_{y_a}{}^{}`$
$`=`$ $`0`$
As $`N_{(0,5)}`$ is an harmonic function in $`y_a`$. The Einstein equations are satisfied in the same way as the $`M2`$-brane solution, but it will be useful to express the equations in terms of the harmonic function $`N_{(0,5)}`$ for reference.
$`({}_{t_i}{}^{t_i}){\displaystyle \frac{1}{2.7!}}(7.6!{\displaystyle \frac{2.7!}{3}})F^{t_1x_1\mathrm{}x_5y_a}F_{t_1x_1\mathrm{}x_5y_a}={\displaystyle \frac{1}{6}}N_{(0,5)}^{\frac{2}{3}}N_{(0,5)}^2(_{y_a}N_{(0,5)})^2\widehat{\delta }_{}^{y_a}{}_{y_a}{}^{}`$
$`({}_{x_i}{}^{x_i}){\displaystyle \frac{1}{2.7!}}(7.6!{\displaystyle \frac{2.7!}{3}})5F^{t_1x_1\mathrm{}x_5y_a}F_{t_1x_1\mathrm{}x_5y_a}={\displaystyle \frac{5}{6}}N_{(0,5)}^{\frac{2}{3}}N_{(0,5)}^2(_{y_a}N_{(0,5)})^2\widehat{\delta }_{}^{y_a}{}_{y_a}{}^{}`$
$`({}_{y_a}{}^{y_a}){\displaystyle \frac{1}{2.7!}}(7.6!5{\displaystyle \frac{2.7!}{3}})F^{t_1x_1\mathrm{}x_5y_a}F_{t_1x_1\mathrm{}x_5y_a}={\displaystyle \frac{7}{6}}N_{(0,5)}^{\frac{2}{3}}N_{(0,5)}^2(_{y_a}N_{(0,5)})^2\widehat{\delta }_{}^{y_a}{}_{y_a}{}^{}`$
### A.3 The $`pp`$-Wave
The $`pp`$-wave solution arises from considering the lowest weight generator associated to a positive root, namely $`K_{}^{1}{}_{2}{}^{}`$, in the weight chain whose highest weight has root $`\beta =\alpha _1+\alpha _2+\mathrm{}\alpha _{10}`$. $`K_{}^{1}{}_{2}{}^{}`$ is the generator of the root $`\alpha _1`$ and we find an associated line element,
$$ds^2=(1K)dt_{1}^{}{}_{}{}^{2}+(1+K)dx_{2}^{}{}_{}{}^{2}2Kdt_1dx_2+d\mathrm{\Omega }_9^2$$
(38)
Where we have made the substitution $`N_{pp}=1+K`$ and we note that $`N_{pp}=1+\frac{𝐐}{7r^7}`$, and $`r^2=y_1^2+\mathrm{}y_9^2`$. We note that $`K=K(y_1,\mathrm{}y_9)`$, which is less general than the solution in , but fits with the generic harmonic functions we have used for all brane solutions in this paper.
## Appendix B Spacelike Brane Solutions from $`E_{11}`$
In this appendix we demonstrate that the $`S`$-brane solutions discussed in section 3.2 satisfy the Einstein equations (3) in signature $`(1,10,)`$ for our ansatz (4). As discussed in section 3.2 our field strength may be constructed out of a complexified version of the generators that give rise to the usual electric solutions in $`(1,10,+)`$, such that these solutions are derived from a truncated action with a $`+F^2`$ term. Equivalently we may use a real form of the sub-algebra generators in signature $`(1,10,)`$ to construct our putative solutions. We follow the same approach as in appendix A and have two field strengths derived from a 3-form and a 6-form gauge field both of which have purely spatial indices, given in equation (28). The appropriate Einstein and gauge equations may be found by setting $`\varphi =0`$, $`\widehat{\delta }_{}^{t_i}{}_{t_i}{}^{}=0`$, $`\widehat{\delta }_{}^{x_i}{}_{x_i}{}^{}=p+1`$, $`\widehat{\delta }_{}^{u_a}{}_{u_a}{}^{}=1`$ and $`\widehat{\delta }_{}^{y_a}{}_{y_a}{}^{}=Dp2`$ in equations (3) and (9). The line element of a spacelike solution is derived from the solution generating group element in the same way as the electric case but using a choice of local sub-algebra that invokes a non-compact timelike generator in the transverse coordinates. The following solutions are associated with the lowest weights,
### B.1 The $`S2`$-Brane
We now demonstrate that there exists an $`S2`$-brane solution in signature $`(1,10,)`$ for our ansatz (4). The line element of the $`S2`$-brane solution is
$$ds^2=N_{(1,7)}^{\frac{2}{3}}(dx_1^2+\mathrm{}dx_3^2)+N_{(1,7)}^{\frac{1}{3}}(du_1^2+dy_1^2+\mathrm{}dy_7^2)$$
(39)
Giving a metric
$$g_{\mu \nu }=\begin{array}{cccccccccc}& & x_1& x_2& x_3& u_1& y_1& \mathrm{}& y_7& \\ x_1& (\mathrm{}& N_{(1,7)}^{\frac{2}{3}}& 0& 0& 0& 0& 0& 0& )\mathrm{}\\ x_2& 0& N_{(1,7)}^{\frac{2}{3}}& 0& 0& 0& 0& 0\\ x_3& 0& 0& N_{(1,7)}^{\frac{2}{3}}& 0& 0& 0& 0\\ u_1& 0& 0& 0& N_{(1,7)}^{\frac{1}{3}}& 0& 0& 0\\ y_1& 0& 0& 0& 0& N_{(1,7)}^{\frac{1}{3}}& 0& 0\\ \mathrm{}& 0& 0& 0& 0& 0& \mathrm{}& 0\\ y_7& 0& 0& 0& 0& 0& 0& N_{(1,7)}^{\frac{1}{3}}\end{array}$$
So $`\sqrt{g}=N_{(1,7)}^{\frac{1}{3}}`$ and we see that the gauge equation is satisfied by $`F_{x_1x_2x_3\widehat{r}}=_{\widehat{r}}N_{(1,7)}^1`$ in the following manner,
$`_{\widehat{r}}(\sqrt{g}F^{x_1x_2x_3\widehat{r}})=`$ $`_{\widehat{r}}(N_{(1,7)}^{\frac{1}{3}}g^{x_1x_1^{}}g^{x_2x_2^{}}g^{x_3x_3^{}}g^{\widehat{r}\widehat{r}^{}}F_{x_1^{}x_2^{}x_3^{}\widehat{r}^{}})`$
$`=`$ $`_{u_1}(N_{(1,7)}^2_{u_1}N_{(1,7)}^1)+_{y_a}(N_{(1,7)}^2_{y_a}N_{(1,7)}^1)\widehat{\delta }_{}^{y_a}{}_{y_a}{}^{}`$
$`=`$ $`_{u_1}_{u_1}N_{(1,7)}_{y_a}_{y_a}N_{(1,7)}\widehat{\delta }_{}^{y_a}{}_{y_a}{}^{}`$ (40)
$`=`$ $`0`$
In the last line we have used the fact that $`N_{(1,7)}`$ is an harmonic function in $`y_a`$, as can be checked from its definition in equation (6),
$`_{u_1}_{u_1}N_{(1,7)}+_{y_a}_{y_a}N_{(1,7)}\widehat{\delta }_{}^{y_a}{}_{y_a}{}^{}`$ $`={\displaystyle \frac{6k}{\widehat{r}^8}}{\displaystyle \frac{8.6ku_1^2}{\widehat{r}^{10}}}{\displaystyle \frac{6k}{\widehat{r}^8}}\widehat{\delta }_{}^{y_a}{}_{y_a}{}^{}+{\displaystyle \frac{8.6ky_a^2}{\widehat{r}^{10}}}\widehat{\delta }_{}^{y_a}{}_{y_a}{}^{}`$
$`={\displaystyle \frac{8.6k}{\widehat{r}^8}}+{\displaystyle \frac{8.6k\widehat{r}^2}{\widehat{r}^{10}}}`$ (41)
$`=0`$
Where $`k=\pm \frac{𝐐}{6}`$ for the $`S2`$-brane, as defined in equation (6).
We now check that the Einstein equations are satisfied by verifying that the right-hand-side of each equation equals the curvature term for the spacelike solution. A term that appears frequently is $`F^{x_1x_2x_3\widehat{r}}F_{x_1x_2x_3\widehat{r}}`$ which we evaluate at the outset
$`F^{x_1x_2x_3\widehat{r}}F_{x_1x_2x_3\widehat{r}}`$ $`=F^{x_1x_2x_3u_1}F_{x_1x_2x_3u_1}+F^{t_1x_1x_2y_1}F_{t_1x_1x_2y_1}+\mathrm{}`$
$`F^{t_1x_1x_2y_7}F_{t_1x_1x_2y_7}`$ (42)
$`=N_{(1,7)}^{\frac{1}{3}}N_{(1,7)}^2(F_{x_1x_2x_3u_1})^2+7N_{(1,7)}^{\frac{1}{3}}N_{(1,7)}^2(F_{x_1x_2x_3y_a})^2`$
We proceed to check the Einstein equations,
$`({}_{x_i}{}^{x_i}){\displaystyle \frac{1}{2.4!}}(4`$ $`F^{x_i\mu _1\mu _2\mu _3}F_{x_i\mu _1\mu _2\mu _3}{\displaystyle \frac{1}{3}}F^{\mu _1\mu _2\mu _3\mu _4}F_{\mu _1\mu _2\mu _3\mu _4}\widehat{\delta }_{}^{x_i}{}_{x_i}{}^{})`$
$`={\displaystyle \frac{1}{2.4!}}(4.3!{\displaystyle \frac{4!}{3}})F^{x_1x_2x_3\widehat{r}}F_{x_1x_2x_3\widehat{r}}\widehat{\delta }_{}^{x_i}{}_{x_i}{}^{}`$ (43)
$`=N_{(1,7)}^{\frac{1}{3}}\{_{u_1}_{u_1}\mathrm{ln}N_{(1,7)}^{\frac{1}{3}}_{y_a}_{y_a}\mathrm{ln}N_{(1,7)}^{\frac{1}{3}}\widehat{\delta }_{}^{y_a}{}_{y_a}{}^{}\}\widehat{\delta }_{}^{x_i}{}_{x_i}{}^{}`$
$`R_{}^{x_i}{}_{x_i}{}^{}`$
$`({}_{u_a}{}^{u_a}){\displaystyle \frac{1}{2.4!}}(4`$ $`F^{u_1\mu _1\mu _2\mu _3}F_{u_1\mu _1\mu _2\mu _3}{\displaystyle \frac{1}{3}}F^{\mu _1\mu _2\mu _3\mu _4}F_{\mu _1\mu _2\mu _3\mu _4})`$
$`={\displaystyle \frac{1}{2.4!}}(4.3!F^{x_1x_2x_3u_1}F_{x_1x_2x_3u_1}{\displaystyle \frac{4!}{3}}F^{x_1x_2x_3\widehat{r}}F_{x_1x_2x_3\widehat{r}})`$ (44)
$`=N_{(1,7)}^{\frac{1}{3}}N_{(1,7)}^2\{{\displaystyle \frac{1}{3}}(_{u_1}N_{(1,7)})^2+{\displaystyle \frac{7}{6}}(_{y_a}N_{(1,7)})^2\}`$
$`R_{}^{u_a}{}_{u_a}{}^{}`$
$`({}_{y_a}{}^{y_a}){\displaystyle \frac{1}{2.4!}}(4`$ $`F^{y_a\mu _1\mu _2\mu _3}F_{y_a\mu _1\mu _2\mu _3}{\displaystyle \frac{1}{3}}F^{\mu _1\mu _2\mu _3\mu _4}F_{\mu _1\mu _2\mu _3\mu _4})`$
$`={\displaystyle \frac{1}{2.4!}}(4.3!F^{x_1x_2x_3y_a}F_{x_1x_2x_3y_a}{\displaystyle \frac{4!}{3}}\widehat{\delta }_{}^{y_a}{}_{y_a}{}^{}F^{x_1x_2x_3\widehat{r}}F_{x_1x_2x_3\widehat{r}})`$ (45)
$`=N_{(1,7)}^{\frac{1}{3}}N_{(1,7)}^2\{{\displaystyle \frac{1}{6}}(_{u_1}N_{(1,7)})^2+{\displaystyle \frac{2}{3}}(_{y_a}N_{(1,7)})^2\}\widehat{\delta }_{}^{y_a}{}_{y_a}{}^{}`$
$`R_{}^{y_a}{}_{y_a}{}^{}`$
### B.2 The $`S5`$-Brane
We now demonstrate that there exists an $`S5`$-brane solution in signature $`(1,10,)`$ for our ansatz (4). The line element of the $`S5`$-brane solution is
$$ds^2=N_{(1,4)}^{\frac{1}{3}}(dx_1^2+\mathrm{}dx_6^2)+N_{(1,4)}^{\frac{2}{3}}(du_1^2+dy_1^2+\mathrm{}dy_4^2)$$
(46)
Giving a metric
$$g_{\mu \nu }=\begin{array}{cccccccccc}& & x_1& \mathrm{}& x_6& u_1& y_1& \mathrm{}& y_4& \\ x_1& (\mathrm{}& N_{(1,4)}^{\frac{1}{3}}& 0& 0& 0& 0& 0& 0& )\mathrm{}\\ \mathrm{}& 0& \mathrm{}& 0& 0& 0& 0& 0\\ x_6& 0& 0& N_{(1,4)}^{\frac{1}{3}}& 0& 0& 0& 0\\ u_1& 0& 0& 0& N_{(1,4)}^{\frac{2}{3}}& 0& 0& 0\\ y_1& 0& 0& 0& 0& N_{(1,4)}^{\frac{2}{3}}& 0& 0\\ \mathrm{}& 0& 0& 0& 0& 0& \mathrm{}& 0\\ y_4& 0& 0& 0& 0& 0& 0& N_{(1,4)}^{\frac{2}{3}}\end{array}$$
So $`\sqrt{g}=N_{(1,4)}^{\frac{2}{3}}`$ and we see that the gauge equation is satisfied by $`F_{x_1\mathrm{}x_6\widehat{r}_a}=_{\widehat{r}_a}N_{(1,4)}^1`$ in the same manner as the field strength associated to the $`S2`$-brane in equations (40) and (41). The Einstein equations are also satisfied in the same way as the $`S2`$-brane solution. We first note that
$$F^{x_1\mathrm{}x_6\widehat{r}}F_{x_1\mathrm{}x_6\widehat{r}}=N_{(1,4)}^{\frac{2}{3}}N_{(1,4)}^2(_{u_1}N_{(1,4)})^2+4N_{(1,4)}^{\frac{2}{3}}N_{(1,4)}^2(_{y_a}N_{(1,4)})^2$$
(47)
Let us now confirm that the Einstein equations in $`(1,10,)`$ are satisfied for our ansatz (4).
$`({}_{x_i}{}^{x_i}){\displaystyle \frac{1}{2.7!}}(7`$ $`F^{x_i\mu _1\mathrm{}\mu _6}F_{x_i\mu _1\mathrm{}\mu _6}{\displaystyle \frac{2}{3}}F^{\mu _1\mathrm{}\mu _7}F_{\mu _1\mathrm{}\mu _7}\widehat{\delta }_{}^{x_i}{}_{x_i}{}^{})`$
$`={\displaystyle \frac{1}{2.7!}}(7.6!{\displaystyle \frac{2.7!}{3}})F^{x_1\mathrm{}x_6\widehat{r}}F_{x_1\mathrm{}x_6\widehat{r}}\widehat{\delta }_{}^{x_i}{}_{x_i}{}^{}`$ (48)
$`=N_{(1,4)}^{\frac{2}{3}}\{_{u_1}_{u_1}\mathrm{ln}N_{(1,4)}^{\frac{1}{6}}_{y_a}_{y_a}\mathrm{ln}N_{(1,4)}^{\frac{1}{6}}\widehat{\delta }_{}^{y_a}{}_{y_a}{}^{}\}\widehat{\delta }_{}^{x_i}{}_{x_i}{}^{}`$
$`R_{}^{x_i}{}_{x_i}{}^{}`$
$`({}_{u_a}{}^{u_a}){\displaystyle \frac{1}{2.7!}}(7`$ $`F^{u_1\mu _1\mathrm{}\mu _6}F_{u_1\mu _1\mathrm{}\mu _6}{\displaystyle \frac{2}{3}}F^{\mu _1\mathrm{}\mu _7}F_{\mu _1\mathrm{}\mu _7})`$
$`={\displaystyle \frac{1}{2.7!}}(7.6!F^{x_1\mathrm{}x_6u_1}F_{x_1\mathrm{}x_6u_1}{\displaystyle \frac{2.7!}{3}}F^{x_1\mathrm{}x_6\widehat{r}}F_{x_1\mathrm{}x_6\widehat{r}})`$ (49)
$`=N_{(1,4)}^{\frac{2}{3}}N_{(1,4)}^2\{{\displaystyle \frac{1}{6}}(_{u_1}N_{(1,4)})^2+{\displaystyle \frac{4}{3}}(_{y_a}N_{(1,4)})^2\}`$
$`R_{}^{u_a}{}_{u_a}{}^{}`$
$`({}_{y_a}{}^{y_a}){\displaystyle \frac{1}{2.7!}}(7`$ $`F^{y_a\mu _1\mathrm{}\mu _6}F_{y_a\mu _1\mathrm{}\mu _6}{\displaystyle \frac{2}{3}}F^{\mu _1\mathrm{}\mu _7}F_{\mu _1\mathrm{}\mu _7})`$
$`={\displaystyle \frac{1}{2.7!}}(7.6!F^{x_1\mathrm{}x_6y_a}F_{x_1\mathrm{}x_6y_a}{\displaystyle \frac{2.7!}{3}}\widehat{\delta }_{}^{y_a}{}_{y_a}{}^{}F^{x_1\mathrm{}x_6\widehat{r}}F_{x_1\mathrm{}x_6\widehat{r}})`$ (50)
$`=N_{(1,4)}^{\frac{2}{3}}N_{(1,4)}^2\{{\displaystyle \frac{1}{3}}(_{u_1}N_{(1,4)})^2+{\displaystyle \frac{5}{6}}(_{y_a}N_{(1,4)})^2\}\widehat{\delta }_{}^{y_a}{}_{y_a}{}^{}`$
$`R_{}^{y_a}{}_{y_a}{}^{}`$
### B.3 The $`Spp`$-Wave
The $`Spp`$-wave solution arises from considering the lowest weight in the weight chain whose highest weight has root $`\beta =\alpha _1+\alpha _2+\mathrm{}\alpha _{10}`$, with a choice of local sub-algebra such that the temporal coordinate of $`M`$-theory is not one of the two distinguished coordinates as it is in the $`pp`$-wave solution. The line element derived from the group element (1) using $`E_\beta =iK_{}^{1}{}_{2}{}^{}`$ is a solution of the vacuum Einstein equations and is,
$$ds^2=(1K)dx_1^2+(1+K)dx_2^22iKdx_1dx_2+d\mathrm{\Omega }_{(1,8)}^2$$
(51)
Where $`d\mathrm{\Omega }_{(1,8)}^2=du_1^2+dy_1^2+\mathrm{}dy_8^2`$, and $`K=K(u_1,y_1,\mathrm{}y_8)`$. This $`Spp`$-wave metric is the line element expected from a double Wick rotation of the $`pp`$-wave solution. A further Wick rotation would give a $`pp`$-wave solution of the $`M`$-theory in $`(2,9)`$. It is non-static and has wavefronts that progress in the spacelike directions transverse to $`\{x_1,x_2\}`$.
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# 1 Introduction
## 1 Introduction
Relativistic approaches of the nuclear physics becomes important for high momentum phenomena, in particular for those associated with spin observables . Furthermore, the phenomena where pions appear as in many nuclear physics phenomena, the relativistic treatment is essential, since the fundamental chiral symmetry is related to the relativistic nature of particles .
The starting point of the relativistically covariant theory is the Bethe-Salpeter equation for the two nucleon system . In order to overcome a difficulty of solving the integral equation, a separable interaction is often employed, primarily as a mathematical manipulation . Ignoring a possible dependence on the total momentum $`P=p_1+p_2=p_1^{}+p_2^{}`$, where $`p_1`$ and $`p_2`$ are the momenta for the initial two nucleons, while $`p_1^{}`$ and $`p_2^{}`$ for the final state ones, the interaction is given as a function of the relative momenta $`p=(p_1p_2)/2`$, $`p^{}=(p_1^{}p_2^{})/2`$,
$`V_{\mathrm{sep}}(p^{},p)=\lambda g(p^{})g(p).`$ (1)
This rank I separable potential is non-local and is very much different from the widely used one-boson-exchange potential (OBEP) which is, to the leading order, given as a function of the momentum transfer $`q=p^{}p`$ and is local:
$`V_{\mathrm{OBEP}}(p^{},p)=V(q).`$ (2)
In their form of Eq. (1) and (2), they can not be equivalent, but, instead, higher rank separable interactions may be used to generate the OBEP when infinitely many terms are introduced. Practically, a finite rank (usually up to rank three) potential is used with a finite number of parameters determined in the phase shift analysis .
In this paper, we investigate whether the parameters of the separable potential may be related to those of the OBEP, since the latter is considered to be physically more fundamental, at least for longer range part of the NN interaction. By doing this, we expect that the separable potential can be understood with physics ground, not just a mathematically convenient tool.
In section 2, we briefly formulate the Bethe-Salpeter equation and provide an analytic solution when a rank I separable potential is used for a single channel problem of $`l=0`$. The rank I separable potential and OBEP are then related in the long wave length approximation. In section 3, we compare the phase shifts calculated from the two potentials. Final section is devoted to conclusions of the present work.
## 2 The BS equation with a separable interaction
Let us consider a single channel equation for $`{}_{}{}^{1}S_{0}^{}`$. This is sufficient for our present qualitative discussions. After angular integration, the BS equation is given by
$`T(p^{},p;s)=V(p^{},p)+{\displaystyle \frac{i}{4\pi ^3}}{\displaystyle 𝑑k_0k^2𝑑k\frac{V(p^{},k)T(k,p;s)}{\left(\frac{\sqrt{s}}{2}e_k+iϵ\right)^2k_0^2}},`$ (3)
where $`T(p^{},p;s)`$ is the $`T`$-matrix, $`V(p^{},p)`$ the interaction kernel, $`s=(p_1+p_2)^2`$ and $`e_k=\sqrt{\stackrel{}{k}^2+M_N^2}`$ with $`M_N`$ being the mass of the nucleon . The momentum variables are for four momentum, e.g., $`p=(p_0,\stackrel{}{p})`$, etc. In Eq. (3), we have considered the equation in the center of mass system. The rank I separable ansatz assumes to write the interaction
$`V_{\mathrm{sep}}(p^{},p)=\lambda g(p^{})g(p)`$ (4)
with a coupling constant $`\lambda `$ and a function $`g(p)`$ as a scalar function of $`p`$ and $`p^{}`$. The $`T`$-matrix is then obtained in the separable form as
$`T(p^{},p;s)`$ $`=`$ $`\tau (s)g(p^{})g(p),`$
$`\tau (s)`$ $`=`$ $`{\displaystyle \frac{1}{\frac{1}{\lambda }+h(s)}},`$ (5)
where
$`h(s)={\displaystyle \frac{i}{4\pi ^3}}{\displaystyle 𝑑k_0k^2𝑑k\frac{g(k)^2}{\left(\frac{\sqrt{s}}{2}e_k+iϵ\right)^2k_0^2}}.`$ (6)
Phase shifts are then given by the relation
$`T(p^{},p;s)={\displaystyle \frac{16\pi }{\sqrt{s}\sqrt{s4M_{N}^{}{}_{}{}^{2}}}}e^{i\delta }\mathrm{sin}\delta ,`$ (7)
where the relative momenta are $`p=(0,\stackrel{}{p})`$, $`p^{}=(0,\stackrel{}{p}^{})`$ for on-shell nucleons. Finally the phase shift $`\delta (s)`$ can be presented by
$`\mathrm{cot}\delta (s)={\displaystyle \frac{\lambda ^1+\mathrm{Re}(h(s))}{\mathrm{Im}(h(s))}}.`$ (8)
Now important terms of the OBEP can be written as
$`V_\sigma (q)=g_\sigma ^2{\displaystyle \frac{1}{q^2+m_\sigma ^2}}\left({\displaystyle \frac{\mathrm{\Lambda }_\sigma ^2m_\sigma ^2}{\mathrm{\Lambda }_\sigma ^2q^2}}\right)^2,V_\omega (q)=g_\omega ^2{\displaystyle \frac{1}{q^2+m_\omega ^2}}\left({\displaystyle \frac{\mathrm{\Lambda }_\omega ^2m_\omega ^2}{\mathrm{\Lambda }_\omega ^2q^2}}\right)^2,`$ (9)
$`V_\pi (q)={\displaystyle \frac{g_\pi ^2}{4M_N^2}}{\displaystyle \frac{q^2}{q^2+m_\pi ^2}}\left({\displaystyle \frac{\mathrm{\Lambda }_\pi ^2m_\pi ^2}{\mathrm{\Lambda }_\pi ^2q^2}}\right)^2,V_\rho (q)={\displaystyle \frac{g_\rho ^2}{2M_N^2}}{\displaystyle \frac{q^2}{q^2+m_\rho ^2}}\left({\displaystyle \frac{\mathrm{\Lambda }_\rho ^2m_\rho ^2}{\mathrm{\Lambda }_\rho ^2q^2}}\right)^2,`$ (10)
where $`q=pp^{}=(0,\stackrel{}{p}\stackrel{}{p}^{})`$. Here the masses $`m_\alpha `$, the coupling constants $`g_\alpha `$ and the cutoff parameters $`\mathrm{\Lambda }_\alpha (\alpha =\sigma ,\omega ,\pi ,\rho )`$ are given in Ref. , and are summarized in Table 1. In Eqs. (9) and (10) we picked up the dominant piece of the one boson exchange potential. The higher order terms are proportional to the initial and final relative momenta $`p`$ and $`p^{}`$. For the $`\rho `$-exchange potential, we use only the tensor coupling term, where the correction from the vector term is about 5 $`\%`$. The coupling strength given in Table 1(in the second row) produces only the $`f`$-coupling in Ref. .
For the parameterization of the separable potential, we assume the Yukawa function for $`g(p)`$ with the same mass parameter $`m`$ as in the OBEP, $`g(p)=1/(p^2m_b^2)`$. Then we try to impose that $`V_{\mathrm{sep}}`$ equals $`V_b`$ in the long wave length limit <sup>1</sup><sup>1</sup>1The relations shown here differ from those of Ref , where extra factor of $`4\pi ^2`$ was erroneously included. :
$`V_{\mathrm{sep}}(0,0)=V_b(0,0),`$ (11)
which determines the strength $`\lambda `$. In this way, we have a separable potential approximately related to the OBEP
$`V_{\mathrm{sep}}(p^{},p)=\lambda _b{\displaystyle \frac{1}{p^2m_b^2}}{\displaystyle \frac{1}{p^2m_b^2}},`$ (12)
where $`\lambda _b`$ are given by
$`\lambda _\sigma =g_\sigma ^2m_\sigma ^2\left(1{\displaystyle \frac{m_\sigma ^2}{\mathrm{\Lambda }_\sigma ^2}}\right)^2,\lambda _\omega =g_\omega ^2m_\omega ^2\left(1{\displaystyle \frac{m_\omega ^2}{\mathrm{\Lambda }_\omega ^2}}\right)^2,`$ (13)
$`\lambda _\pi ={\displaystyle \frac{g_\pi ^2m_\pi ^4}{4M_N^2}}\left(1{\displaystyle \frac{m_\pi ^2}{\mathrm{\Lambda }_\pi ^2}}\right)^2,\lambda _\rho ={\displaystyle \frac{g_\rho ^2m_\rho ^4}{2M_N^2}}\left(1{\displaystyle \frac{m_\rho ^2}{\mathrm{\Lambda }_\rho ^2}}\right)^2.`$ (14)
In the case of $`\pi `$ and $`\rho `$, we excluded $`q^2`$ dependence in the numerator of the Eqs. (15), otherwise we can not determine the $`\lambda `$ parameter. It corresponds to excluding the $`\delta `$-function term in the $`r`$ space, namely for the Eq. (11) we have used
$`V_\pi (q)={\displaystyle \frac{g_\pi ^2}{4M_N^2}}{\displaystyle \frac{m_\pi ^2}{q^2+m_\pi ^2}}\left({\displaystyle \frac{\mathrm{\Lambda }_\pi ^2m_\pi ^2}{\mathrm{\Lambda }_\pi ^2q^2}}\right)^2,V_\rho (q)={\displaystyle \frac{g_\rho ^2}{2M_N^2}}{\displaystyle \frac{m_\rho ^2}{q^2+m_\rho ^2}}\left({\displaystyle \frac{\mathrm{\Lambda }_\rho ^2m_\rho ^2}{\mathrm{\Lambda }_\rho ^2q^2}}\right)^2.`$ (15)
The numerical values of the $`\lambda ^{}`$s are also given in the Table 1 (last column).
## 3 Comparison of phase shifts
We have calculated phase shifts using the BS Eqs. (3) – (6), which are compared with those obtained from the OBEP. Since we make the comparison at relatively low energy region, it is sufficient to solve the Schrödinger equation. Here we compare various phase shifts calculated by using a single term corresponding to $`\sigma `$-, $`\omega `$-, $`\pi `$\- or $`\rho `$-exchange potentials. The resulting phase shifts are shown in Figs. 1. For later use, we show the one boson exchange potential of the $`{}_{}{}^{1}S_{0}^{}`$ channel as a function of $`r`$ in Figs. 2, where various terms of the OBEP are shown separately. The thick solid line in Fig. 2-(a) is the total potential including the $`\sigma `$, $`\omega `$, $`\pi `$ and $`\rho `$ exchange potentials
Now we discuss the phase shifts calculated from each meson exchange potential.
* Fig. 1-(a) shows the phase shifts calculated from the potentials of the $`\sigma `$ channel as functions of $`T_{lab}`$, the kinetic energy in the laboratory frame. Here $`T_{lab}`$ is related to $`s`$ by
$`T_{lab}={\displaystyle \frac{s4M_N^2}{2M_N}}.`$ (16)
The thick solid line represents the phase shift for the separable potential, and the thin solid line for the OBEP. As shown in Fig. 1-(a), both phase shifts start from 180 degrees, indicating a strong attraction as accommodating one bound state. Indeed, as shown in Fig. 2-(a), the depth of the $`\sigma `$-exchange potential of the OBEP reaches about 200 MeV at 0.75 fm. The strong attraction of the OBEP causes the raising behavior at $`T_{lab=0}`$, which turns to decreasing at $`T_{lab}=20`$ MeV. On the other hand, the separable potential can not be that strongly attractive. This can be checked by analyzing the scattering matrix of Eq. (5), which will be discussed later. In fact, the phase shift approaches the upper limit as indicated by the dashed line in Fig. 1-(a) in the limit $`\lambda \mathrm{}`$. Furthermore we have checked that it is not possible to reproduce the strong attraction of the $`\sigma `$-exchange potential, whatever $`m_b`$ value of the separable potential we choose. The fact that the separable potential can not be too strong has been discussed previously .
* Fig. 1-(b) shows the phase shifts calculated from the potentials of the $`\omega `$ channel as functions of $`T_{lab}`$. The phase shifts calculated from the two interactions (Separable and OBEP) show repulsive nature as it starts from 0 degree and decreases as the energy increases. As shown in Fig. 2-(a), the repulsive force of the $`\omega `$-exchange potential is very strong. As shown in Fig. 1-(b), the result of the separable potential resembles that of OBEP. However, that strong repulsion of the OBEP can not be reproduced by the separable potential, which is similar to the case of the $`\sigma `$-exchange potential. Once again, as shown by the dashed line in Fig. 1-(b) there is a lower bound of the phase shift of the separable potential in the limit $`\lambda \mathrm{}`$. However if we allow $`m_b`$ to change, it is possible to reproduce the phase shift of the $`\omega `$-exchange potential by using a parameter set of, for instance, $`m_b=630`$ MeV and $`\lambda =81.5\times 10^6`$ MeV<sup>2</sup>. To make $`m_b`$ small corresponds to the increase of repulsion. At this point, we recognize that the physical meaning of the mass parameter $`m_b`$ in the separable potential is different from that in the OBEP.
* Fig. 1-(c) shows the phase shifts calculated from the potentials of the $`\pi `$ channel as functions of $`T_{lab}`$. The phase shifts calculated from the two interactions look very different. Fig. 2-(b) shows that $`\pi `$-exchange potential is attractive at long distances $`r\underset{}{>}`$ 1 fm and repulsive at middle and short distances $`r\underset{}{<}`$ 1 fm. Therefore the phase shift calculated from OBEP starts from 0 degree, raising at first, then turns to decrease at $`T_{lab}5`$ MeV and becomes repulsive at $`T_{lab}20`$ MeV. Due to the form factor, the one-pion-term of the OBEP here is written as a sum of the long range attraction and the short range repulsion. On the other hand, the phase shift of the separable interaction is weakly attractive. Because there is only one term in the rank I separable potential, the contributions of attraction and repulsion can not be reproduced simultaneously. In particular, the coupling strength $`\lambda _\pi `$ determined from the relation (14) at low momentum region ($`p=p^{}=0`$) is too attractive, which therefore can not reproduce the repulsive behavior at higher $`T_{lab}`$. However, one can fit the repulsive behavior at higher energies by changing the range parameter $`m_\pi `$ and the coupling constant $`\lambda `$. For instance if we choose $`m_\pi =\mathrm{\Lambda }_\pi `$$`=`$1300 MeV and $`\lambda =135\times 10^6`$ MeV<sup>2</sup>, we can reproduce the phase shift at around $`T_{lab}`$ 200 MeV as indicated by the dashed line of Fig. 1-(c).
* Fig. 1-(d) shows the phase shifts calculated from the potentials of the $`\rho `$ channel as functions of $`T_{lab}`$. The phase shifts calculated from the two interactions look very different. With $`\lambda =37.2\times 10^6`$ MeV<sup>2</sup> which is determined by Eq. (14), the result of the separable potential is too attractive, such that it generates one bound state and the phase shift starts from 180 degrees. In contrast, as shown in Fig. 2-(b) the $`\rho `$-exchange piece of the OBEP is attractive at long distances $`r\underset{}{>}`$ 0.6 fm and repulsive at middle and short distances $`r\underset{}{<}`$ 0.6 fm. The attractive interaction here, however, is not very large due to the cancellation by the repulsive component. Therefore the phase shift calculated from OBEP starts raising from 0 degree, turns to decrease at $`T_{lab}40`$ MeV, and change into repulsion at $`T_{lab}160`$MeV. Just as in the $`\pi `$ case, we can re-fit the strength of the separable potential $`\lambda _\rho `$. By reducing the strength by about factor 8, $`\lambda =4.52\times 10^6`$ MeV<sup>2</sup>, we obtain the phase shift of the separable potential as shown by the dashed line of Fig. 1-(d), which looks rather close to the result of OBEP.
These results show that it is difficult to reproduce the phase shifts of each terms of the one-boson-exchange potential separately by the rank I separable potential when we use the parameters determined in the long wave length limit. As explained above in detail, the separable potential can not be stronger than a certain strength both for attractive and repulsive cases if we do not change the $`m_b`$ parameter.
The fact that the separable potential can not be stronger than a certain strength may be understood from Eq. (5) where the factor $`1/\lambda `$ vanishes in the limit $`|\lambda |\mathrm{}`$. Interestingly, in this limit there is no distinction between attractive ($`\lambda \mathrm{}`$) and repulsive ($`\lambda +\mathrm{}`$) interactions. In order to find the maximum strength of the separable potential, we plot in Fig. 4 the real and imaginary parts of $`h(s)`$ from which we can calculate the phase shift by using Eq. (8). The result is shown in Fig. 4. The real part monotonically decreases from 0.276 GeV<sup>-2</sup>, while the imaginary part starts from 0, reaches the maximum value at some $`s`$ and turns to decrease monotonically. This behavior resembles what is familiar in the non-relativistic scattering theory where the phase shift varies from 0 to 180 degrees when there is one bound state. In a relativistic theory, however, a naive argument in the non-relativistic theory can not be applied, since in the large $`s`$ region particle production may occur and the discussion within a fixed particle number can not be applied. In the present separable potential model, the treatment will break down at and beyond $`s=4(M_N+m_b)^2`$ where an unphysical pole of mass $`m_b`$ appears. In our calculation of the phase shift, in order to determine the initial value $`\delta (T_{lab}=0)`$, for the attractive interaction we increased $`\lambda `$ gradually from a small value and verified that there is a jump from $`\delta (T_{lab}=0)=0`$ to $`\delta (T_{lab}=0)=180`$ degrees at certain strength of $`\lambda `$ only once. Therefore, we conclude that the maximum strength of the separable potential in our method is what allows one bound state for an attractive interaction. Similarly for the repulsive case, it is also possible to show that there is the maximum strength of the interaction if $`m_b`$ is fixed.
Now turning to the full result of the $`{}_{}{}^{1}S_{0}^{}`$ channel, as indicated in Fig. 3, the separable potential can reproduce rather well the result of the total nuclear force of OBEP when we take $`\lambda =0.294\times 10^6`$ MeV<sup>2</sup> and $`m_b=224`$ MeV . The very strong attractive and repulsive forces of the $`\sigma `$\- and $`\omega `$-exchange potentials are largely canceled, yielding a rather mild nuclear force. This is the reason that the separable potential for nuclear reaction have been successful.
## 4 Summary
We have studied the relation between the rank I separable potential for the covariant Bethe-Salpeter equation and the one-boson-exchange potential (OBEP). Individual channels of $`\sigma `$-, $`\omega `$-, $`\pi `$\- and $`\rho `$-exchanges were investigated separately. As a result, it turned out that the rank I separable potential could not reproduce the phase shift calculated from each component of the OBEP when we use the parameters determined in the long wave length limit. As for the $`\sigma `$ channel, where the potential is strongly attractive, we could not reproduce the phase shift of OBEP even if we take the limit $`\lambda \mathrm{}`$ and we change $`m_b`$ parameter. Similarly as for the $`\omega `$ channel with strong attraction, the separable potential could not reproduce again the phase shift of OBEP even in the limit $`\lambda \mathrm{}`$. However we could reproduce the strong repulsion, if we change the $`m_b`$ parameter. These observations imply that the physical meaning of the mass parameters in the separable potential and OBEP are different. The mass parameter of the OBEP represents the interaction range of a local potential, while that of the separable potential could mimic, for instance, the range of a non-local interaction. The non-locality of the nuclear force is related to the structure of the nucleon at short ranges $`r\underset{}{<}`$ 0.5 fm . Concerning the $`\pi `$ and $`\rho `$ channels, where the potential consists of attraction at long distances and repulsion at short distances, the rank I separable potential could not reproduce the mixed nature of the interaction, although the interaction strengths are not as strong as the $`\sigma `$ and $`\omega `$ channels. Despite the above fact, the rank I separable potential can reproduce the experimental data of the $`{}_{}{}^{1}S_{0}^{}`$ phase shift up to the energy $`T_{lab}\underset{}{>}200`$ MeV where a mild attractive interaction dominates.
These results show that the rank I separable potential is not suited to the description of very strong attraction. For instance, phase shifts calculated from the separable potential can not become larger than 180 degrees, no matter how large the attraction coupling constant takes. Rather, the separable potential can describe relatively mild attraction and all repulsion. In the realistic nuclear force, such a mild strength is obtained by the sum of the strongly attractive $`\sigma `$-exchange and the strongly repulsive $`\omega `$-exchange potentials.
In this work, we have shown that the separable potential works well for the two nucleon system if parameters are chosen suitably, although the decomposition into components of physical OBEP does not make sense. In a sense, the different nature of the two potential should have been expected. The main purpose of the present paper was to see whether it is possible to make physical meaning of the separable potential in comparison with the OBEP by using a simple parameterization of one term of rank I. In order to perform a good description of phenomena in the covariant Bethe-Salpeter formalism, we can introduce a higher rank form. Such a work is now in progress, where the use of the improved rank one ansatz and of higher rank interactions are tested .
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# Chan-Paton factors and Higgsing from Vacuum String Field Theory
## 1 Introduction
Open String Field Theory (OSFT), , is a candidate for a non–perturbative definition of string theory. Central to it there is the concept of background independence: physics should be independent on the vacuum we decide to expand the theory on. Of course there can be vacua on which the theory might look simpler than in any other vacuum therefore, once such a vacuum is identified, the theory should take its simplest form around it. As the dynamical degrees of freedom of OSFT are open strings, a class of its vacua should contain all possible configurations of D–branes. Bosonic D–branes are unstable due to the omnipresence of the open string tachyon, hence they decay classically to a vacuum with no D–branes at all. The nature of this vacuum (tachyon vacuum) is universal, in the sense that it is independent of the details of the various BCFT’s that describe the original configuration of branes. For this reason there is strong expectation that OSFT on the tachyon vacuum should be able to describe all open string physics and, possibly, the related closed string physics that arises non–perturbatively through the classical decay of the branes, see for review and references therein. To overcome the difficulties in finding exact analytic solutions of OSFT, a model (Vacuum String Field Theory) has been conjectured by Rastelli Sen and Zwiebach which is a proposal for OSFT at the tachyon vacuum, . This model highly simplifies the form of the OSFT action by replacing the usual BRST operator with a $`c`$–midpoint insertion. VSFT has passed many tests. Classical solutions relative to different branes have been obtained, and the right ratios of tensions have been reproduced . The small on shell perturbations of the solution relative to the D25–brane have been shown to cover the open string spectrum on a D25–brane . Time dependent solutions interpolating from the D25–brane vacuum to the tachyon vacuum have been proposed, . However all these successes have to be contrasted with the singular nature of VSFT coming from the factorization of matter and ghost degrees of freedom. Hints on how such a singular nature arises have been given in : matter ghost factorization arises as a singular limit of a reparametrization of the worldsheet which shrinks the whole string to its midpoint. For this reason all of the classical solutions of VSFT exhibit a singular behavior at the midpoint and, as a consequence, observables are obtained by regularization of these singularities. Although a completely consistent regularization scheme has not been given yet (but see for interesting developments) the leading matter–ghost factorized form of VSFT is very powerful and, in all the above–mentioned cases, it gives rise unambiguously to the right physics, once midpoint singularities are correctly regularized. In this note we give a simple description of open strings states living on a set of $`N`$ D–branes. When the branes are coincident we encounter in the spectrum $`N^2`$ massless vectors, giving rise to a $`U(N)`$ gauge symmetry. This symmetry is part of the huge gauge symmetry of VSFT when one considers matter–ghost factorized gauge transformations. The Chan Paton factors arises from particular combinations of left/right excitations on the sliver, that takes the form the generalized Laguerre polynomials discovered in , see also . This $`U(N)`$ structure is dynamically generated (it is an intrinsic part of a classical solution) and there is no need to add it by hand as in first quantized string theory or even in usual OSFT. In this sense background independence is manifest.
Using the translation operator $`\text{e}^{ix\widehat{p}}`$ we construct an array of D24–branes and analyze its small on shell fluctuations. We show that open strings stretched between parallel branes at different positions are obtained by translating differently the left and right part of the classical solution. This is possible because the lump projector is left/right factorized. Of course this operation is ambiguous for what concerns the midpoint, since it does not have a left/right decomposition. Indeed we show that a naive use of left/right orthogonality cannot give rise to the correct shift in the mass formula, proportional to the distance<sup>2</sup> between two D–branes. By using wedge–state regularization we show that in the sliver limit there is a non vanishing contribution which is completely localized at the midpoint and gives rise to the correct shift in the mass formula. The mechanism is that of a twist anomaly, , which has proven to be crucial for obtaining the spectrum of strings around a single D25–brane and to give the correct ratio of D–branes.
The note is organized as follows. After a brief review of the simplest classical solutions, the sliver and the lump, , in section 3 we review the construction of orthogonal Neumann (i.e. zero momentum) projectors using half string vectors and Laguerre polynomials, given in . In the following section we show how open string states with the correct Chan–Paton factors arise on a classical solution given by the sum of $`N`$ such orthogonal projectors. In section 5 we consider a superposition of D24–branes and show how the transverse fluctuations are obtained, by simply exciting the transverse part of the classical solution with generic oscillators. Most of these are pure gauge but the excitation given by the midpoint oscillator cannot be gauged away. In section 6 we higgs the above system by translating the D24–branes at a certain distance from one another, thus creating an array in the transverse direction. We analyze the fluctuations of such a classical solution and show that the mass shift for strings stretched between different branes is generated by a twist anomaly. We further show that the wave–functional for such states is not continuous at the midpoint but gives rise to the expected change in the boundary conditions, living the midpoint position in the target space undetermined.
Some computations involving well known formulas for computing $``$–products on Neumann–type solutions are summarized in the appendix.
## 2 D–branes as projectors
The leading order VSFT action is
$$𝒮(\mathrm{\Psi })=K\left(\frac{1}{2}\mathrm{\Psi }|𝒬|\mathrm{\Psi }+\frac{1}{3}\mathrm{\Psi }|\mathrm{\Psi }\mathrm{\Psi }\right)$$
(1)
where
$$𝒬=\frac{1}{2i}\left(c(i)c(i)\right)$$
(2)
given the purely ghost form of the kinetic operator, the classical solutions can be matter/ghost factorized
$$\mathrm{\Psi }=\mathrm{\Psi }_m\mathrm{\Psi }_g$$
(3)
so the equation of motions factorizes too
$`𝒬\mathrm{\Psi }_g`$ $`=`$ $`\mathrm{\Psi }_g_g\mathrm{\Psi }_g`$ (4)
$`\mathrm{\Psi }_m`$ $`=`$ $`\mathrm{\Psi }_m_m\mathrm{\Psi }_m`$ (5)
One can thus consider a universal ghost solution and concentrate in finding projectors of the matter star algebra. The operator definition of the $`_m`$ product is
$${}_{123}{}^{}V_3|\mathrm{\Psi }_1_{1}^{}|\mathrm{\Psi }_2_2=_3\mathrm{\Psi }_1_m\mathrm{\Psi }_2|,$$
(6)
where $`V_3|`$ is the three string vertex, see . In the following we need both translationally invariant (D25–branes) and non-translationally invariant (Dk–branes) solutions. For simplicity we will consider the sliver and the lump, . The former is translationally invariant and is defined by
$$|\mathrm{\Xi }=𝒩e^{\frac{1}{2}a^{}Sa^{}}|0,a^{}Sa^{}=\underset{n,m=1}{\overset{\mathrm{}}{}}a_n^\mu S_{nm}a_m^\nu \eta _{\mu \nu }$$
(7)
where $`S=CT`$ and
$$T=\frac{1}{2X}(1+X\sqrt{(1+3X)(1X)})$$
(8)
where $`C`$ is the twist matrix and $`X=CV^{11}`$ is the diagonal Neumann coefficient, see for details.
To represent lower dimensional brane we need to break translation invariance in the transverse direction, this is usually done by passing to the oscillator basis
$`a_0^{(r)\alpha }={\displaystyle \frac{1}{2}}\sqrt{b}\widehat{p}^{(r)\alpha }i{\displaystyle \frac{1}{\sqrt{b}}}\widehat{x}^{(r)\alpha },a_0^{(r)\alpha }={\displaystyle \frac{1}{2}}\sqrt{b}\widehat{p}^{(r)\alpha }+i{\displaystyle \frac{1}{\sqrt{b}}}\widehat{x}^{(r)\alpha },`$ (9)
where $`b`$ is a free parameter and $`\alpha `$, $`\beta `$ denote transverse indices.
The transverse part of the vertex in this new basis becomes
$`|V_{3,}^{}=Ke^E^{}|\mathrm{\Omega }_b`$ (10)
with
$`K=\left({\displaystyle \frac{\sqrt{2\pi b^3}}{3(V_{00}+b/2)^2}}\right)^{\frac{k}{2}},E^{}={\displaystyle \frac{1}{2}}{\displaystyle \underset{r,s=1}{\overset{3}{}}}{\displaystyle \underset{M,N0}{}}a_M^{(r)\alpha }V_{MN}^{{}_{}{}^{}rs}a_N^{(s)\beta }\eta _{\alpha \beta }`$ (11)
where $`M,N`$ denote the couple of indices $`\{0,m\}`$ and $`\{0,n\}`$, respectively. The coefficients $`V_{MN}^{{}_{}{}^{}rs}`$ are given in Appendix B of .
The lump solution $`|\mathrm{\Xi }_k^{}`$ has the form (7) with $`S`$ along the parallel directions and $`S`$ replaced by $`S^{}`$ along the perpendicular ones. In turn $`S^{}=CT^{}`$ and $`T^{}`$ has the same form as $`T`$ eq.(8) with $`X`$ replaced by $`X^{}`$.
## 3 $`N`$ coincident D25-branes
There are several ways to construct coincident branes solutions in VSFT, the one we are going to use is in terms of Laguerre polynomials, explicitly given in .
Consider a left string vector $`\xi _n^\mu `$, such that
$`\rho _R\xi ^\mu `$ $`=`$ $`0`$ (12)
$`\rho _L\xi ^\mu `$ $`=`$ $`\xi ^\mu `$ (13)
The $`\rho _{R,L}`$ operators project into the the right/left Hilbert space at zero momentum, see .
With this half string vector it is possible to excite left-right symmetrically a string configuration, using the operator
$$𝐱=(a_\mu ^{},\xi ^\mu )(a_\nu ^{},C\xi ^\nu )=y\stackrel{~}{y}$$
(14)
where $`(,)`$ means inner product in level space and the operators $`\stackrel{~}{y}`$ $`y`$ are identified with right/left excitations. The half string vector $`\xi `$ is normalized by the following condition and definition
$`(\xi _\mu ,{\displaystyle \frac{1}{1T^2}}\xi ^\mu )`$ $`=`$ $`1`$ (15)
$`(\xi _\mu ,{\displaystyle \frac{T}{1T^2}}\xi ^\mu )`$ $`=`$ $`\kappa `$ (16)
where $`T=CS`$ is the Sliver Neumann coefficient, (8).
For every choice of $`\xi `$ satisfying 15, we can construct an infinite family of orthogonal projectors (D–branes) given by
$$|\mathrm{\Lambda }_n=(\kappa )^nL_n\left(\frac{𝐱}{\kappa }\right)|\mathrm{\Xi }$$
(17)
where $`L_n(x)`$ is the n-th Laguerre polynomial. These states obey the remarkable properties
$`|\mathrm{\Lambda }_n|\mathrm{\Lambda }_m`$ $`=`$ $`\delta _{nm}|\mathrm{\Lambda }_m`$ (18)
$`\mathrm{\Lambda }_n|\mathrm{\Lambda }_m`$ $`=`$ $`\delta _{nm}\mathrm{\Xi }|\mathrm{\Xi }`$ (19)
Due to these properties, once the sliver is identified with a single D–brane, a stack of $`N`$ D-branes can be given by
$$|N=\underset{n=0}{\overset{N1}{}}|\mathrm{\Lambda }_n$$
(20)
From (19) we further get that the $`bpz`$ norm of such a solution is $`N`$–times the one of the sliver.
So far we have considered left-right symmetric projectors which are in one to one correspondence with type 0 Laguerre polynomial, there are however non left-right symmetric states corresponding to generalized Laguerre polynomials, they are given by,
$`|\mathrm{\Lambda }_{nm}=\sqrt{{\displaystyle \frac{n!}{m!}}}\kappa ^m(iy)^{nm}L_m^{nm}\left({\displaystyle \frac{x}{\kappa }}\right)|\mathrm{\Xi }nm`$ (21)
$`|\mathrm{\Lambda }_{nm}=\sqrt{{\displaystyle \frac{m!}{n!}}}\kappa ^n(i\stackrel{~}{y})^{nm}L_n^{mn}\left({\displaystyle \frac{x}{\kappa }}\right)|\mathrm{\Xi }mn`$ (22)
and obey the properties <sup>1</sup><sup>1</sup>1Another realization of this algebra is given in
$`|\mathrm{\Lambda }_{nm}|\mathrm{\Lambda }_{pq}=\delta _{mp}|\mathrm{\Lambda }_{nq}`$ (23)
$`\mathrm{\Lambda }_{nm}|\mathrm{\Lambda }_{pq}=\delta _{mp}\delta _{nq}\mathrm{\Xi }|\mathrm{\Xi }`$ (24)
note in particular that $`|\mathrm{\Lambda }_n=|\mathrm{\Lambda }_{nn}`$. With these states we can implement partial–isometry–like operations, see also . Consider indeed
$`|+`$ $`=`$ $`{\displaystyle \underset{n=0}{\overset{N1}{}}}|\mathrm{\Lambda }_{n+1,n}`$ (25)
$`|+`$ $`=`$ $`{\displaystyle \underset{n=0}{\overset{N1}{}}}|\mathrm{\Lambda }_{n,n+1}`$ (26)
It’s trivial to see that
$`|+|+`$ $`=`$ $`|N`$ (27)
$`|+|+`$ $`=`$ $`|N|\mathrm{\Xi }`$ (28)
Note that any of the previous states can be obtained starting from the sliver by star products
$`|\mathrm{\Lambda }_{nm}=(|+)_{}^n|\mathrm{\Xi }(|+)_{}^m`$ (29)
We have in particular
$`|+|\mathrm{\Xi }`$ $`=`$ $`0`$ (30)
$`|\mathrm{\Xi }|+`$ $`=`$ $`0`$
As a final remark it is worth noting that the partial isometry that relates projectors to projectors is actually a $``$–rotation and hence a (matter ghost factorized) gauge transformation. We have indeed
$`\mathrm{\Lambda }_{nn}=\text{e}^{\frac{\pi }{2}(\mathrm{\Lambda }_{nm}\mathrm{\Lambda }_{mn})}\mathrm{\Lambda }_{mm}\text{e}^{\frac{\pi }{2}(\mathrm{\Lambda }_{nm}\mathrm{\Lambda }_{mn})}`$ (31)
as can be easily checked from (23)
## 4 $`U(N)`$ open strings
Let’s recall that the (matter–ghost factorized) open string cohomology around a (matter–ghost factorized) classical solution $`|\mathrm{\Psi }`$ is given by the following conditions
$`|\varphi `$ $`=`$ $`|\varphi |\mathrm{\Psi }+|\mathrm{\Psi }|\varphi `$ (32)
$`|\varphi `$ $``$ $`|\mathrm{\Lambda }|\mathrm{\Psi }|\mathrm{\Psi }|\mathrm{\Lambda }`$ (33)
The first representing $`𝒬_\mathrm{\Psi }`$–closed states while the second gauges away $`𝒬_\mathrm{\Psi }`$ exact–states.<sup>2</sup><sup>2</sup>2These conditions actually cover only the ghost–matter factorized cohomology
In the case of $`N`$–coincident D25–branes the classical solution is given by (20).
As multiple D–branes are obtained starting from the sliver by multiple $``$–products via (29), a generic open string state on the sliver can acquire a Chan–Paton factor $`(i,j)Adj[U(N)]`$ in the same way.
Let $`|\{g\},p`$ be an on–shell open string state on the sliver, identified by the collection of polarization tensors $`\{g\}`$ and momentum $`p`$. The Chan–Paton structure is simply given by
$`|(i,j);\{g\},p=(|+)_{}^i|\{g\},p(|+)_{}^j`$ (34)
There is a subtlety here, related to twist anomaly, , and the consequent breakdown of $``$–associativity. Indeed the expression (34) is ambiguous in the overall normalization in front: it depends on how the various star products involved are nested. This is so because all the states we are considering are constructed on the sliver, which fails to satisfy unambiguously its equation of motion when states at non zero momentum enter the game. Consider for simplicity the Hata–Kawano tachyon state,
$`|p=𝒩\text{e}^{(ta^{}+ix)p}|\mathrm{\Xi }=𝒩^{}\text{e}^{ip\widehat{X}\left(\frac{\pi }{2}\right)}|\mathrm{\Xi }`$ (35)
$`t=3{\displaystyle \frac{T^2T+1}{1+T}}v_0`$ (36)
this state satisfies (weakly) the linearized equation of motion (LEOM) with the sliver state
$`|p|\mathrm{\Xi }=|\mathrm{\Xi }|p=\text{e}^{Gp^2}|p`$ (37)
$`G=\mathrm{log2}p^2=1`$ (38)
The quantity $`G`$ gets a non vanishing value from the region very near $`k=0`$ in the continuous basis, where some of the remarkable properties, encoding associativity, between Neumann coefficients breaks down due to singularities that are regulated in a non associative way (like level truncation). Indeed (37) violates associativity if, as is the case, $`G0`$
$`(|p|\mathrm{\Xi })|\mathrm{\Xi }|p(|\mathrm{\Xi }|\mathrm{\Xi })`$ (39)
Just to fix a convention (and stressing once more that the only ambiguity is in the overall normalization) we decide to do first all the star products at zero momentum (that do not develop twist anomaly) and, as the last operation, multiply the result with the state at definite momentum $`|\{g\},p`$
Now we show that (34) satisfies the LEOM
$`|(i,j);\{g\},p=|(i,j);\{g\},p|N+|N|(i,j);\{g\},p`$ (40)
using (30) we get the relations
$`|N(|+)_{}^i=(|+)_{}^i|Ni`$ (41)
$`(|+)_{}^j|N=|Nj(|+)_{}^j`$ (42)
which allow to write the LEOM as
$`|(i,j);\{g\},p=(|+)_{}^i(|\{g\},p|Nj+|Ni|\{g\},p)(|+)_{}^j`$
It is proven in the appendix that
$`|\{g\},p|\mathrm{\Lambda }_{n1}=|\mathrm{\Lambda }_{n1}|\{g\},p=0`$ (43)
once the following conditions on the half string vector $`\xi ^\mu `$ are satisfied
$`(t,{\displaystyle \frac{1}{1\pm T}}\xi ^\mu )`$ $`=`$ $`0`$ (44)
$`(g_{\mu \nu _1\mathrm{}\nu _n},{\displaystyle \frac{1}{1\pm T}}\xi ^\mu )`$ $`=`$ $`0`$ (45)
The first condition states that the half string vector $`\xi ^\mu `$ should be “orthogonal” to the on–shell tachyon vector $`t=3\frac{T^2T+1}{1+T}v_0`$, this just constrains $`2D`$ components of $`\xi ^\mu `$ out of $`D\mathrm{}1`$, and as such is easy to implement. <sup>3</sup><sup>3</sup>3These conditions are actually not needed if we represent the tachyon state as $`\text{e}^{ip\widehat{X}\left(\frac{\pi }{2}\right)}|\mathrm{\Xi }`$ since, up to overall normalizations, midpoint insertions commutes with the star product; the role of such conditions is to avoid extra terms when we use the CBH formula to pass to the oscillator expression $`\text{e}^{(ta^{}+ix)p}|\mathrm{\Xi }`$. The second condition implies the “orthogonality” of the half string vector $`\xi ^\mu `$ with any of the on–shell polarization tensors $`\{g\}`$. We know from previous works, , that the level components of the polarization vectors should be identified with the $`k=0`$ eigenvector(s), of the continuous basis at zero momentum, . Again, in order for (45) to be satisfied, it is sufficient to ask that the components $`\xi ^\mu (k)`$ vanishes fast enough at $`k=0`$, see for explicit realizations of this condition.
Given (43) it follows directly that
$`|\{g\},p|Nj+|Ni|\{g\},p=|\{g\},p|\mathrm{\Xi }+|\mathrm{\Xi }|\{g\},p`$ (46)
hence the LEOM simplifies to
$`|(i,j);\{g\},p=(|+)_{}^i(|\{g\},p|\mathrm{\Xi }+|\mathrm{\Xi }|\{g\},p)(|+)_{}^j`$ (47)
This ensures the on–shellness of the state $`|(i,j);\{g\},p`$ once this is true for $`|\{g\},p`$. We thus recover $`N^2`$ kinematical copies of the spectrum on a single D–brane. Note that the left/right structure of these states is the same as a $`U(N)`$ double line notation, as the relations (23) certify. It should be noted that this Chan–Paton structure does not sit at the endpoints of the string, , but is rather “diluted” on the string halves. This can be traced back to the singular field redefinition that should relate OSFT with VSFT, see the conclusions.
## 5 $`N`$ coincident D24-branes
A system of $`N`$ coincident D24–branes can be represented by<sup>4</sup><sup>4</sup>4Other possibilities, for example putting the Laguerre polynomials on the codimension, are related to this by partial isometry and hence, due to (31), should be gauge equivalent
$`|N=\left({\displaystyle \underset{n=0}{\overset{N}{}}}|\mathrm{\Lambda }_n\right)|\mathrm{\Xi }^{}`$ (48)
where the state $`|\mathrm{\Xi }^{}`$ is the lump solution given in . The open string string sector with Lorentz indices coming from the 25 dimensional world volume ($`N^2`$ tachyons, $`U(N)`$–gluons, etc…) is exactly as in the previous section. In addition there are the physical states coming from transverse excitations. These states are given by exciting the transverse part of the classical solution $`|\mathrm{\Xi }^{}`$ with oscillators. The Chan–Paton degrees of freedom are encoded in the worldvolume part of the state as in the previous section. For example the transverse scalars are given by<sup>5</sup><sup>5</sup>5Note that once the relations (44) are implemented one can recast the Chan Paton indices directly on the classical solution and then act with oscillators to build onshell fluctuation
$`|(ij);g^{25},p_{}=𝒩\text{e}^{(ta^{}+ix)p_{}}|\mathrm{\Lambda }_{ij}g^{25}a_{25}^{}|\mathrm{\Xi }^{}`$ (49)
It’s easy to verify that these states satisfy the LEOM iff $`p_{}^2=0`$, we have indeed
$`|(ij);g^{25},p_{}`$ $`=`$ $`|N|(ij);g^{25},p_{}+|(ij);g^{25},p_{}|N`$ (50)
$`=`$ $`2^{p_{}^2}𝒩\text{e}^{(ta^{}+ix)p_{}}|\mathrm{\Lambda }_{ij}(\rho _L^{}+\rho _R^{})g^{25}a_{25}^{}|\mathrm{\Xi }^{}`$ (51)
The $`\rho _{L,R}^{}`$ are the left/right projectors with zero modes, see .
The level vector $`g^{25}`$ is completely arbitrary, but only its midpoint part is not pure gauge. Indeed, exactly as in , we can try to gauge away any of the states (49)
$`\delta |(ij);g^{25},p_{}=|Q_{ij}|N|N|Q_{ij}`$ (52)
where
$`|Q_{ij}=\text{e}^{(ta^{}+ix)p_{}}|\mathrm{\Lambda }_{ij}u^{25}a_{25}^{}|\mathrm{\Xi }^{}`$ (53)
We have
$`\delta |(ij);g^{25},p_{}=|(ij);(\rho _L^{}\rho _R^{})u^{25},p_{}`$ (54)
Thus the state is pure gauge if
$`g^{25}=(\rho _L^{}\rho _R^{})u^{25}`$ (55)
It is well known that the operator $`(\rho _L^{}\rho _R^{})`$ just change the twist parity of a given level vector. In the diagonal basis all vectors are paired except the one corresponding to $`k=0`$ that is only twist even, (at least if we restrict ourselves to vectors that have a non vanishing overlap with Fock–space vectors, see ). Thus the gauge transformation (54) gauges away all the components of $`g^{25}`$ except the $`k=0`$ one, which is the midpoint. One can construct higher transverse excitations by applying more and more transverse oscillators as in . Again only the $`k=0`$ oscillator(s) are not gauge trivial.
## 6 Higgsing
Now we want to “higgs” the previous system of $`N`$ coincident D24–branes to an array of $`N`$ D24–branes, displaced of a distance $`\mathrm{}`$ from one another in the transverse dimension $`y`$. This system is obtained by multiple translation of the previous classical solution (48).
$`|N^{(\mathrm{})}={\displaystyle \underset{n=0}{\overset{N1}{}}}\left(|\mathrm{\Lambda }_n\text{e}^{in\mathrm{}\widehat{p}}|\mathrm{\Xi }^{}\right)`$ (56)
As in it is very convenient to pass to the oscillator basis by
$`\widehat{p}={\displaystyle \frac{1}{\sqrt{b}}}\left(a_0+a_0^{}\right)`$ (57)
and to define the level vector
$`\beta _N={\displaystyle \frac{i\mathrm{}}{\sqrt{b}}}\left(1T^{}\right)_{0N}`$ (58)
The transverse part of the n-th D24–brane in (56) can thus be written as
$`|\mathrm{\Xi }_n^{}=\text{e}^{in\mathrm{}\widehat{p}}|\mathrm{\Xi }^{}=\text{e}^{\frac{n^2}{2}(\beta ,\frac{1}{1T^{}}\beta )+n\beta a^{}}|\mathrm{\Xi }^{}`$ (59)
As proven in we have
$`|\mathrm{\Xi }_n^{}|\mathrm{\Xi }_m^{}`$ $`=`$ $`\delta _{nm}|\mathrm{\Xi }_n^{}`$ (60)
$`\mathrm{\Xi }_n^{}|\mathrm{\Xi }_m^{}`$ $`=`$ $`\delta _{nm}\mathrm{\Xi }^{}|\mathrm{\Xi }^{}`$ (61)
We recall here that the orthogonality condition comes from a divergence at $`k=0`$ of the continuous basis of the primed Neumann matrices. Indeed, up to unimportant contributions, we have used the identification, see
$`\delta _{nm}=\mathrm{exp}\left[(nm)^2(\beta ,{\displaystyle \frac{1}{1+T^{}}}\beta )\right]=\mathrm{exp}\left[{\displaystyle \frac{\mathrm{}^2}{b}}(nm)^2\left({\displaystyle \frac{(1T^{})^2}{1+T^{}}}\right)_{00}\right]`$ (62)
Note that we don’t really need to use different projectors on the worldvolume as the degeneracy is lifted by the different space–translations of the various projectors, however one can still use the $`|\mathrm{\Lambda }_n`$’s in order to maintain the orthogonality as $`\mathrm{}0`$.
Now we come to the spectrum.
Type $`(n,n)`$ strings (the ones stretched between the same D–brane) are obtained by translation of strings on a single D24–brane
$`|(n,n);\{g\},p_{}=\text{e}^{in\mathrm{}p_{}}|\{g\},p_{}^{(n)}`$ (63)
where $`|\{g\},p_{}^{(n)}`$ is an on–shell state of the previous section constructed on $`|\mathrm{\Lambda }_{nn}|\mathrm{\Xi }^{}`$. Thus we get $`N`$ copies of the spectrum of a single D24–branes: $`N`$ tachyons, $`N`$ massless vectors etc… This gives a $`U(1)^N`$ gauge symmetry.
The situation changes when we want to consider strings stretched between two different D–branes. In this case we expect that a shift in the mass formula is generated, proportional to the square of the distance between the two branes. In order to construct $`(i,j)`$ states we have to translate the state $`|\mathrm{\Xi }^{}`$ differently with respect its left/right degrees of freedom. We use the following identification for the left/right momentum
$`\widehat{p}`$ $`=`$ $`\widehat{p}_L+\widehat{p}_R`$ (64)
$`\widehat{p}_{L,R}`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{b}}}\left(\rho _{L,R}^{}a+\rho _{L,R}^{}a^{}\right)_0`$ (65)
We then consider the state
$`\text{e}^{in\mathrm{}\widehat{p}_Lim\mathrm{}\widehat{p}_R}|\mathrm{\Xi }^{}\text{e}^{(n\beta _L+m\beta _R)a^{}}|\mathrm{\Xi }^{}=|\mathrm{\Xi }_{nm}^{}`$ (66)
where we have defined
$`\beta _{L,R}=\rho _{L,R}^{}\beta `$ (67)
The $`\rho ^{}`$ projectors obey the following properties up to midpoint subtleties, see later
$`\rho _L^{}+\rho _R^{}=1`$ (68)
$`(\rho _{L,R}^{})^2=\rho _{L,R}^{}`$ (69)
$`\rho _{L,R}^{}\rho _{R,L}^{}=0`$ (70)
If we naively use these properties, using the formulas of , it is easy to prove that
$`|\mathrm{\Xi }_{nm}^{}|\mathrm{\Xi }_{pq}^{}=\text{e}^{\frac{1}{2}\left((nm+pqnq)\beta \frac{1}{1T^{}}\beta +(mpnmpq+nq)\beta \frac{1}{1T^2}\beta \right)}|\mathrm{\Xi }_{nq}^{}`$ (71)
We can then normalize the above states in order to have
$`|\widehat{\mathrm{\Xi }}_{nm}^{}|\widehat{\mathrm{\Xi }}_{pq}^{}=\delta _{mp}|\widehat{\mathrm{\Xi }}_{nq}^{}`$ (72)
where
$`|\widehat{\mathrm{\Xi }}_{nm}^{}=\text{e}^{\frac{1}{4}\beta \frac{n^2+m^2+2nmT^{}}{1T^2}\beta }\text{e}^{(n\beta _L+m\beta _R)a^{}}|\mathrm{\Xi }^{}`$ (73)
Note that this normalization is quite formal as the quantity $`\beta \frac{1}{1+T}\beta `$ is actually divergent, this is not however a real problem as open string states are not normalized by the LEOM’s, moreover it should be noted that even if these left/right non–symmetric states have a vanishing normalization, they give rise to non vanishing objects (the projectors) by $``$ product.
Consider now, for simplicity, the “tachyon” state stretched from the i-th brane to the j-th. The corresponding state is given by
$`|(ij);p_{}=𝒩\text{e}^{(ta^{}+ix)p_{}}|\mathrm{\Xi }|\widehat{\mathrm{\Xi }}_{ij}^{}`$ (74)
Using (72) it is easy to see that the above state satisfies the LEOM
$`|(ij);p_{}=|(ij);p_{}|N^{(\mathrm{})}+|N^{(\mathrm{})}|(ij);p_{}=2^{p_{}^2+1}|(nm);p_{}`$ (75)
with $`p_{}^2=1`$, that is we don’t get the usual mass shift proportional to the distance<sup>2</sup> between the two D–branes.
However the algebra (72) is not quite correct. To elucidate this point it is worth considering the components of the level vector $`\beta `$ in the continuous part of the diagonal basis of the primed Neumann matrices, see for details. <sup>6</sup><sup>6</sup>6There are, of course, also the contributions from the discrete spectrum, but they are not singular for $`0<b<\mathrm{}`$ We have
$`\beta (k)={\displaystyle \frac{i\mathrm{}}{\sqrt{b}}}(1+\text{e}^{\frac{\pi |k|}{2}})V_0(k)`$ (76)
where $`V_0(k)`$ is the zero component of the normalized eigenvector of the continuous basis,
$`V_0(k)=\sqrt{{\displaystyle \frac{bk}{4\mathrm{sinh}\frac{\pi k}{2}}}}\left[4+k^2\left(\mathrm{}F_c(k){\displaystyle \frac{b}{4}}\right)^2\right]^{\frac{1}{2}}`$ (77)
The $`\beta `$ vector is finite at $`k=0`$,
$`\beta (0)={\displaystyle \frac{i\mathrm{}}{\sqrt{2\pi }}},`$ (78)
hence its left/right decomposition is not well defined. This implies that it is not correct to consider the quantity
$`\gamma =(\beta _L,{\displaystyle \frac{1}{1+T^{}}}\beta _R)={\displaystyle \frac{\mathrm{}^2}{b}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}\theta (k)\theta (k){\displaystyle \frac{\left(1+\text{e}^{\frac{\pi |k|}{2}}\right)^2}{1\text{e}^{\frac{\pi |k|}{2}}}}V_0(k)^2`$ (79)
as vanishing since it is formally indeterminate (it is “$`0\mathrm{}`$”). Assuming that $`\gamma `$ is non vanishing one easily obtains that the algebra (72) gets modified to
$`|\widehat{\mathrm{\Xi }}_{nm}^{}|\widehat{\mathrm{\Xi }}_{pq}^{}=\delta _{mp}\text{e}^{\frac{1}{4}\left[\left(np\right)^2+\left(mq\right)^2\right]\gamma }|\widehat{\mathrm{\Xi }}_{nq}^{}`$ (80)
Taking this modification into account we obtain
$`|(nm);p_{}|N+|N|(nm);p_{}=2^{p_{}^2+1+\frac{1}{4}(nm)^2\frac{\gamma }{\mathrm{log2}}}|(nm);p_{}`$ (81)
that gives the mass formula
$`p_{}^2=1+{\displaystyle \frac{1}{4}}(nm)^2{\displaystyle \frac{\gamma }{\mathrm{log2}}}`$ (82)
We recall that the mass for such a state should be given by ($`\alpha ^{}=1`$)
$`p_{}^2=1\left({\displaystyle \frac{\mathrm{\Delta }y_{nm}}{2\pi }}\right)^2=1\left({\displaystyle \frac{(nm)\mathrm{}}{2\pi }}\right)^2`$ (83)
The two formulas agrees iff
$`\gamma =(\beta _L,{\displaystyle \frac{1}{1+T^{}}}\beta _R)={\displaystyle \frac{\mathrm{}^2}{\pi ^2}}\mathrm{log2}`$ (84)
To verify this identity we need to regularize the ambiguous expression (79). We do it by substituting the lump Neumann coefficient $`T^{}`$ with the wedge states one $`T_N^{}`$. We remind that, see
$`T_N^{}`$ $`=`$ $`{\displaystyle \frac{T^{}+(T^{})^{N1}}{1(T^{})^N}}`$ (85)
$`T_N^{}T_N^{}`$ $`=`$ $`X^{}+(X_+^{},X_{}^{})(1𝒯_{}^{}{}_{N}{}^{}^{})^1𝒯_{}^{}{}_{N}{}^{}\left(\begin{array}{c}X_{}^{}\\ X_+^{}\end{array}\right)=T_{2N1}^{}`$ (86)
We have<sup>7</sup><sup>7</sup>7Note that the $`T^{}`$ in the denominator of (79) is actually obtained by the projector equation $`T^{}T^{}=T^{}`$ that is violated in wedge–state regularization
$`\gamma =(\beta ,\rho _L(T^{}){\displaystyle \frac{1}{1+T^{}T^{}}}\rho _R(T^{})\beta )=\underset{N\mathrm{}}{lim}\left(\beta \rho _L(T_N^{}){\displaystyle \frac{1}{1+T_{2N}^{}}}\rho _R(T_N^{})\beta \right)`$ (87)
where we have used the $``$–multiplication between wedge states
$`T_N^{}T_N^{}=T_{2N1}^{}T_{2N}^{},N1`$ (88)
The matrices $`T_N^{}`$ gets contributions from the continuous and the discrete spectrum but only the continuous spectrum is relevant in the large $`N`$ limit, moreover it is only the region infinitesimally near the point $`k=0`$ that really contributes. We have
$`\gamma =\left({\displaystyle \frac{i\mathrm{}}{\sqrt{2\pi }}}\right)^2\underset{N\mathrm{}}{lim}{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑k\rho _L^{(N)}(k){\displaystyle \frac{1}{1+t_{2N}(k)}}\rho _R^{(N)}(k)`$ (89)
with
$`t_N(k)={\displaystyle \frac{\text{e}^{\frac{\pi k}{2}}+\left(\text{e}^{\frac{\pi k}{2}}\right)^{N1}}{1\left(\text{e}^{\frac{\pi k}{2}}\right)^N}}`$ (90)
$`\rho _L^{(N)}(k)=1{\displaystyle \frac{1}{1+\left(\text{e}^{\frac{\pi k}{2}}\right)^{N1}}}`$ (91)
$`\rho _R^{(N)}(k)={\displaystyle \frac{1}{1+\left(\text{e}^{\frac{\pi k}{2}}\right)^{N1}}}`$ (92)
where we have used the expression of $`\rho _{L,R}^{}`$ in terms of the sliver matrix and the $``$ Neumann coefficients, , and their (continuous) eigenvalues, . Let’s evaluate the integral in the large $`N`$ limit ($`x=\frac{\pi k}{2},y=Nx`$)
$`{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑k\rho _L^{(N)}(k){\displaystyle \frac{1}{1+t_{2N}(k)}}\rho _R^{(N)}(k)`$ (93)
$`=`$ $`{\displaystyle \frac{2}{\pi }}{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑x{\displaystyle \frac{\text{e}^{Nx}\left(\text{e}^{Nx}1\right)}{\left(1\text{e}^x\right)\left(1+\text{e}^{Nx}\right)\left(1+\text{e}^{2Nx}\right)}}+O\left({\displaystyle \frac{1}{N}}\right)`$
$`=`$ $`{\displaystyle \frac{2}{N\pi }}{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑y{\displaystyle \frac{\text{e}^y\left(\text{e}^y1\right)}{\left(1\text{e}^{\frac{y}{N}}\right)\left(1+\text{e}^y\right)\left(1+\text{e}^{2y}\right)}}+O\left({\displaystyle \frac{1}{N}}\right)`$
$`=`$ $`{\displaystyle \frac{2}{\pi }}{\displaystyle _{\mathrm{}}^{\mathrm{}}}{\displaystyle \frac{dy}{y}}{\displaystyle \frac{\text{e}^y\left(\text{e}^y1\right)}{\left(1+\text{e}^y\right)\left(1+\text{e}^{2y}\right)}}+O\left({\displaystyle \frac{1}{N}}\right)`$
$`=`$ $`{\displaystyle \frac{2}{\pi }}\mathrm{log2}+O\left({\displaystyle \frac{1}{N}}\right)`$
so we get the right value of $`\gamma `$. Note that, when $`N\mathrm{}`$, the integrand completely localizes at $`k=0`$, as claimed before.
As a last remark we would like to discuss about the position of the midpoint in the transverse direction for the states $`|\mathrm{\Xi }_{nm}^{}`$. We know that for the usual lump solution we have,
$`\widehat{X}\left({\displaystyle \frac{\pi }{2}}\right)|\mathrm{\Xi }^{}=0`$ (94)
That is the lump functional has support on string states in which the midpoint is constrained to live on the worldvolume. This is interpreted as a Dirichlet condition, see also . Moreover, since we have
$`[\widehat{X}\left({\displaystyle \frac{\pi }{2}}\right),p]=[x_0,p]=i,`$ (95)
it is immediate to see that
$`\widehat{X}\left({\displaystyle \frac{\pi }{2}}\right)\text{e}^{in\widehat{p}\mathrm{}}|\mathrm{\Xi }^{}=n\mathrm{}\text{e}^{in\widehat{p}\mathrm{}}|\mathrm{\Xi }^{}`$ (96)
So shifted branes undergoes a consistent change in boundary conditions.
The operator $`\widehat{X}\left(\frac{\pi }{2}\right)`$ is proportional to the $`k=0`$ position operator,
$`\widehat{X}\left({\displaystyle \frac{\pi }{2}}\right)=2\sqrt{\pi }\widehat{x}_{k=0}`$ (97)
It’s easy to check that we have the following commutation relations
$`[2\sqrt{\pi }\widehat{x}_k,p]`$ $`=`$ $`2i\sqrt{{\displaystyle \frac{2\pi }{b}}}V_0(k)`$ (98)
$`[2\sqrt{\pi }\widehat{x}_k,p_L]`$ $`=`$ $`2i\sqrt{{\displaystyle \frac{2\pi }{b}}}V_0(k)\theta (k)`$ (99)
$`[2\sqrt{\pi }\widehat{x}_k,p_R]`$ $`=`$ $`2i\sqrt{{\displaystyle \frac{2\pi }{b}}}V_0(k)\theta (k)`$ (100)
That allows to write
$`\underset{k0^{}}{lim}2\sqrt{\pi }\widehat{x}_k\text{e}^{i(n\widehat{p}_L+m\widehat{p}_R)\mathrm{}}|\mathrm{\Xi }^{}`$ $`=`$ $`n\mathrm{}\text{e}^{i(n\widehat{p}_L+m\widehat{p}_R)\mathrm{}}|\mathrm{\Xi }^{}`$ (101)
$`\underset{k0^+}{lim}2\sqrt{\pi }\widehat{x}_k\text{e}^{i(n\widehat{p}_L+m\widehat{p}_R)\mathrm{}}|\mathrm{\Xi }^{}`$ $`=`$ $`m\mathrm{}\text{e}^{i(n\widehat{p}_L+m\widehat{p}_R)\mathrm{}}|\mathrm{\Xi }^{}`$ (102)
The string functional relative to this state is not continuous at the midpoint, this is the reason why the correct mass shell condition comes out from a twist anomaly. In the singular representation of VSFT in which the whole interior of a string is contracted to the midpoint, , these properties reproduce the expected change in the left/right boundary conditions, and show that the point $`k=0`$ naturally accounts for D–branes moduli.
## 7 Conclusions
In this paper we have addressed the problem of how Chan–Paton degrees of freedom arise on the physical excitations of multiple D–branes. For coincident D–branes we have proved that these degrees of freedom can be encoded in the algebra of Laguerre polynomials, discovered in , we expect however that other (gauge equivalent) descriptions can be given. We have further shown that if we consider an array of parallel displaced D24–branes, the expected higgsing $`U(N)U(1)^N`$ takes place. The shift in the mass formula is due to a twist anomaly, a contribution which is completely localized at the point $`k=0`$ of the continuous spectrum of the Neumann matrices with zero modes. This is the same phenomenon, discovered in , that gives rise to the correct mass formula for open string states on a single D25–brane. This result confirms once more that observables in VSFT are associated to midpoint subtleties. This fact is hardly a surprise given that the field redefinition relating OSFT on the tachyon vacuum to VSFT completely shrinks the body of a string to its midpoint and “resolves” the endpoints into the left/right halves, . This, we believe, is the reason why Chan–Paton factors are to be found in left/right excitations of VSFT classical solutions, and not in the endpoints, see .
It would be interesting to give a BCFT description of our construction. In particular the properties (101, 102) suggest that stretched states should be given by the insertion on the boundary of the lump surface state of a boundary changing vertex operator, in a way similar to , but with one insertion more . Another interesting development would be to switch on interactions in order to explicitly realize the $`U(N)`$ dynamical symmetry.
We hope we have given a clear picture on how some topics related to background independence naturally emerge from the string field algebra. We believe this to be relevant in understanding the elusive nature of closed string states around the tachyon vacuum.
###### Acknowledgments.
I thank Loriano Bonora for useful discussions and encouragement during the various stages of this work, I thank Frank Ferrari for a stimulating discussion and Ghasem Exirifard for help in integrations. This research was supported by the Italian MIUR under the program “Teoria dei Campi, Superstringhe e Gravità”.
## Appendix
## Appendix A Proof of (4.12)
A general open string state on the sliver $`|\mathrm{\Xi }`$ can be obtained differentiating the generating state, see
$`|\varphi _\beta =\text{e}^{(tp+\beta )a^{}}|\mathrm{\Xi }\text{e}^{ipx}`$ (103)
where
$`t=3{\displaystyle \frac{T^2T+1}{1+T}}v_0`$ (104)
is the on–shell tachyon vector, and $`\beta _\mu `$ is a level–Lorentz vector.
The $`|\mathrm{\Lambda }_n`$’s can be generated by the state ,
$`|\mathrm{\Xi }_\lambda =\text{e}^{\lambda a^{}}|\mathrm{\Xi }`$ (105)
General formulas of and allows to compute
$`|\varphi _\beta |\mathrm{\Xi }_\lambda `$ $`=`$ $`\text{e}^{Gp^2+A_{LR}(\beta ,\lambda )}\text{e}^{(tp+\rho _L\beta \rho _R\lambda )a^{}}|\mathrm{\Xi }\text{e}^{ipx}`$ (106)
$`|\mathrm{\Xi }_\lambda |\varphi _\beta `$ $`=`$ $`\text{e}^{Gp^2+A_{RL}(\beta ,\lambda )}\text{e}^{(tp+\rho _R\beta \rho _L\lambda )a^{}}|\mathrm{\Xi }\text{e}^{ipx}`$ (107)
where
$`A_{LR}(\beta ,\lambda )`$ $`=`$ $`{\displaystyle \frac{1}{2}}(\beta ,{\displaystyle \frac{T}{1T^2}}\beta )+(\beta ,{\displaystyle \frac{\rho _RT\rho _L}{1T^2}}C\lambda ){\displaystyle \frac{1}{2}}(\lambda ,{\displaystyle \frac{T}{1T^2}}\lambda )`$
$`+p\left((t,{\displaystyle \frac{T}{1T^2}}\beta )(t,{\displaystyle \frac{\rho _LT\rho _R}{1T^2}}\lambda )+(t,{\displaystyle \frac{\rho _R+T\rho _L}{1T^2}}\beta )(t,{\displaystyle \frac{\rho _L\rho _R}{1T^2}}\lambda )\right)`$
and
$`A_{RL}(\beta ,\lambda )=A_{LR}(\beta ,\lambda )|_{\rho _L\rho _R,\rho _R\rho _L}`$ (109)
We can restrict the polarization $`\beta _n^\mu `$ to the $`k=0`$ component, indeed every physical excitation of the tachyon wave function $`\text{e}^{tpa^{}+ipx}|\mathrm{\Xi }`$ should be localized there, see .
Therefore is not restrictive to ask
$`(\beta ,f(T)\xi )=0`$ (110)
once the half string vector $`\xi ^\mu `$ vanishes rapidly enough at $`k=0`$.
Moreover asking for (44) to be satisfied, using (23), it is easy to show that
$`|\{g\},p|N+|N|\{g\},p=|\{g\},p|\mathrm{\Xi }+|\mathrm{\Xi }|\{g\},p,`$ (111)
as claimed
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# Hydrodynamic scaling limit of continuum solid-on-solid model
## 1. Introduction
A process of great technological importance, molecular beam epitaxy (MBE), is used to manufacture computer chips and semiconductor devices, see Barabasi and Stanley . In general the chip is constructed by spraying a beam of atoms on a flat surface. There are three phenomena that take place in the construction process namely deposition, diffusion and desorption. The atoms arrive or deposit on the surface and do not stick on the first contact point, but diffuse or walk on the surface. When an atom reaches the edge of another wandering atom, the two atoms meet or glue together, forming islands. Sometimes an atom may jump out of the surface. Smaller islands may develop into larger islands affecting the roughness of the surface on the macroscopic scale. A rough surface does not have very good contact properties and engineers would like to understand the basics mechanisms affecting the morphology in general.
How deposition and desorption affect the morphology of the surface, it is quite well understood. It is of great importance to know how the profile of the surface evolves on the macroscopic scale if the atoms that make up the surface diffuse. We assume that no atoms arrive or leave the surface.
The present paper discusses a model for surface diffusion the continuum solid-on-solid model, known also as the fourth-order Ginzburg-Landau model. The system has a complex interaction that is not a nearest-neighbor interaction and has two conserved quantities in a nonperiodic box, namely the total slope of the surface and the linear mean of the surface slope. Under the assumption that the molecules of the surface follow the dynamics of the continuum solid-on-solid model, we will prove that the dynamics of the surface slope profile on the macroscale is a fourth-order nonlinear equation. Since the model is nongradient, the derivation of the limit is not trivial and we shall use the method developed in Varadhan and further extended in Quastel , Varadhan and Yau . Even though we have an interaction that is not of nearest-neighbor type, the microscopic current still decomposes as the Laplacian of the slope field and the fluctuations. Usually in nearest-neighbor models after subtracting the fluctuations from the current we end up with the gradient of some local function. We use the result that for continuum solid-on-solid model the space of exact functions has codimension one inside the space of closed functions, see Savu .
We also include a discussion of a perturbation of continuum solid-on-solid model, where the evolution of the surface is driven by both diffusion and electric field. The electric field will add one extra second-order term to the final nonlinear equation.
The continuum solid-on-solid model is an approximation of the discrete solid-on-solid model and hence is not considered a truly microscopic model. Unfortunately at the time of writing of this paper, we don’t have the required techniques to solve the discrete solid-on-solid model. However, partial rigorous results and numerical analysis are available for the discrete case, see Krug, Doobs and Majaniemi .
Finally the paper is organized as follows: section 2 contains the description of continuum solid-on-solid model and of the model for surface electromigration, the statement of the main results and a note on similar models considered in the literature so far; section 3 outlines the proof of the main result, the computation of the scaling limit of continuum solid-on-solid model; section 4 shows how the final nonlinear equation is identified; in section 5 we calculate the asymptotics of central limit variance, the main ingredient used in section 4 to prove that the microscopic current can be replaced by a multiple of the field Laplacian; in section 6 we calculate the scaling limit of the modified model to incorporate surface electromigration.
## 2. The models
The continuum solid-on-solid model, is a 1-dimensional lattice system with continuous order parameter, used to study the evolution of an interface. The model describes the movement of the interface at the mesoscopic level and hence is not considered a truly microscopic model that captures all the aspects of the particle motion, but it has the advantage of being more suitable for computations.
The height model. Let $`𝕋`$ be the torus represented as the interval $`[0,1]`$ with $`0`$ and $`1`$ identified. For each positive integer $`N`$, there are $`N`$ scaled periodic lattice points located at sites $`i/N`$ in $`𝕋`$, $`i=1,\mathrm{},N`$. We shall denote by $`h_i(t)`$ the height of the surface at the site $`i/N`$, at time $`t`$. Also because of the periodicity of the lattice points, $`h_{N+1}(t)=h_1(t)`$. The energy function $`H_N(h)`$ of a height configuration $`h`$ is chosen to be invariant under the global translation $`h_i(t)h_i(t)+c`$ and has the form
$$H_N(h)=\underset{i=1}{\overset{N}{}}V(h_{i+1}h_i).$$
In this paper we shall assume that the potential is quadratic $`V(x)=x^2`$. We shall require the evolution in time of the surface to preserve the sum of the heights, to have as invariant measure the infinite mass measure $`e^{H_N(h)}dh`$, and to be reversible. All these properties are satisfied by the solution of the stochastic differential system:
$$dh_i(t)=\frac{N^4}{2}(w_iw_{i1})dt+N^2(\sqrt{a_i}dB_i\sqrt{a_{i1}}dB_{i1}),1iN.$$
Above, $`a_i(h)=a(h_ih_{i1},h_{i+1}h_i,h_{i+2}h_{i+1})`$ where $`a`$ is a function with bounded continuous first derivatives that satisfies $`0<1/a^{}a(x_1,x_0,x_1)a^{}<\mathrm{}`$. The $`n`$ copies of the Brownian motion $`B_i,i=1,\mathrm{},n`$ are independent. We define the instantaneous current $`w_i(h)`$ of particles over the bond $`\{i,i+1\}`$
(2.1) $`w_i(h)`$ $`=`$ $`(_1a2_0a+_1a)(h_ih_{i1},h_{i+1}h_i,h_{i+2}h_{i+1})+`$
$`+a_i(h)(V^{}(h_ih_{i1})2V^{}(h_{i+1}h_i)+V^{}(h_{i+2}h_{i+1})).`$
Because the height process preserves the sum of the heights of the surface, this dynamics models surface diffusion.
The slope process. A crucial property of the dynamics of the heights is the gauge property, namely the dynamics invariance under the action of the group $`G`$ of translation in the $`(1,\mathrm{},1)`$ direction,
$$G=\{T:^N^N|T(x_1,\mathrm{},x_N)=(x_1+c,\mathrm{},x_N+c),c\}.$$
Hence there exists an induced dynamics on the quotient space $`^N/G`$ of equivalence classes. A representative of an equivalence class is the slope configuration $`x_i(t)=h_{i+1}(t)h_i(t)`$, $`1iN1`$ and $`x_N(t)=h_1(t)h_N(t)`$. Note that $`_{i=1}^Nx_i(t)=0`$.
As a function of the slope configuration of the surface, the energy becomes $`H_N(x)=_{i=1}^NV(x_i)`$.
In the sequel, we shall study the slope process rather than the height process. The slope process is reversible and has as equilibrium distribution, the product probability measure $`d\nu _N^{\mathrm{eq}}=e^{H_N(x)}/Z^Ndx`$. Below we write down the stochastic differential system, the generator, and the Dirichlet form of the slope process:
(2.2)
$$dx_i(t)=\frac{N^4}{2}(w_{i+1}2w_i+w_{i1})dt+N^2(\sqrt{a_{i1}}dB_{i1}2\sqrt{a_i}dB_i+\sqrt{a_{i+1}}dB_{i+1}),$$
(2.3)
$$N^4L_N(f)=\frac{N^4}{2}\underset{i=1}{\overset{N}{}}a_i(x)(_{i+1}2_i+_{i1})^2f+w_i(x)(_{i+1}2_i+_{i1})f,$$
(2.4)
$$N^4D_N(f)=\frac{N^4}{2}_^N\underset{i=1}{\overset{N}{}}a_i(x)(_{i+1}f2_if_{i1}f)^2(x)\frac{e^{H_N(x)}}{Z^N}dx.$$
The factor $`N^4`$ in the generator with the lattice spacing of $`1/N`$ represents the scaling of space and time. This scaling is needed to observe a nontrivial motion in the limit.
The diffusion (2.2) is driven in the direction of the linear vector fields $`X_i=_{i+1}2_i+_{i1}`$, $`1iN`$, therefore is not ergodic in the whole space $`^N`$. It becomes ergodic when restricted to the hyperplane $`x_1+\mathrm{}+x_N=N\overline{x}`$ of average slope $`\overline{x}`$. The unique equilibrium probability measure of the dynamics restricted to the hyperplane is the conditional probability $`e^{H_N(x)}/Z^Ndx`$ given $`x_1+\mathrm{}+x_N=N\overline{x}`$.
Instantaneously, the slope profile of the surface decreases at some site $`i`$ twice as much as increases at the adjacent sites $`i1`$, $`i+1`$. We see that any update of the slope configuration affects the slopes at three sites. This type of interaction, known as the three site interaction, is quite complex and has been very rarely studied, so far.
We should note an integration by parts property of the current $`w_i`$, for all bounded and smooth functions $`f`$:
(2.5)
$$E^{\mathrm{eq}}[w_if]=E^{\mathrm{eq}}[a_i(_{i+1}2_i+_{i1})f].$$
The expected value above is with respect to the equilibrium measure $`d\nu _N^{\mathrm{eq}}`$.
The dynamics in a nonperiodic box. Although we are concerned with the study of the slope process defined on a periodic space, we make use of a similar slope process evolving in a box with nonperiodic boundary. Below we describe this new dynamics and its main properties. Suppose $`\mathrm{\Lambda }`$ is a box with finitely many sites of the lattice $``$. The infinitesimal generators,
(2.6)
$$L_\mathrm{\Lambda }(f)=\frac{1}{2}\underset{i\mathrm{\Lambda },i+1,i1\mathrm{\Lambda }}{}a_i(x)(_{i+1}2_i+_{i1})^2f+w_i(x)(_{i+1}2_i+_{i1})f,$$
respectively,
(2.7)
$$L_{\mathrm{}}(f)=\frac{1}{2}\underset{i}{}a_i(x)(_{i+1}2_i+_{i1})^2f+w_i(x)(_{i+1}2_i+_{i1})f,$$
produce two diffusion processes.
In the case of the first dynamics, generated by $`L_\mathrm{\Lambda }`$, there is no transport of particles over the boundary of the box, so the box does not have periodic boundary. If the box $`\mathrm{\Lambda }=[l,l]`$ or $`\mathrm{\Lambda }=[il,i+l]`$ we use the shorter notation $`L_l`$, respectively, $`L_{i,l}`$ for the generator $`L_\mathrm{\Lambda }`$. The second dynamics generated by $`L_{\mathrm{}}`$ is a dynamics on the infinite lattice $``$. The dynamics $`L_{i,l}`$ preserves the average slope, $`y_{i,l}^1=\frac{x_{il}+\mathrm{}+x_{i+l}}{2l+1}`$, and the linear mean of the slopes, $`y_{i,l}^2=\frac{(l)x_{il}+\mathrm{}+lx_{i+l}}{l(l+1)}`$, inside the box $`[il,i+l]`$. The second conserved quantity, $`y_{i,l}^2`$, should be understood as a boundary condition that is preserved in time. The linear mean slope, $`y_{i,l}^2`$, is conserved in time because we have a model for surface diffusion and the total height of the surface does not changes as time passes by. We will use two different notation $`\overline{x}_{i,l}`$ and $`y_{i,l}^1`$ for the mean slope of the field in the box centered at $`i`$, of size $`l`$. Also $`y_{i,l}`$ stands for the vector of the conserved quantities, $`(y_{i,l}^1,y_{i,l}^2)`$. As a convention we will drop the subscript $`i`$, meaning the center of the box, if the box is centered at the origin.
Equilibrium measures. We proceed to describe next the equilibrium measures of the dynamics that we have introduced. For the dynamics $`L_{i,l}`$, the grand canonical measure is the product probability measure $`\nu _{\alpha ,i,l}^{gc}=_{j=il}^{i+l}\frac{e^{\alpha x_jV(x_j)}}{Z(\alpha )}dx_j,`$ whereas the canonical measure is the conditional probability measure $`\nu _{y,i,l}^c=\nu _{\alpha ,i,l}^{gc}(|y_{i,l})`$ given the level set
$$\{x^{2l+1}|\frac{x_{il}+\mathrm{}+x_{i+l}}{2l+1}=y_{i,l}^1,\frac{(l)x_{il}+\mathrm{}+lx_{i+l}}{l(l+1)}=y_{i,l}^2\}.$$
The canonical measure is the unique stationary probability measure for the restricted dynamics on this set. The Dirichlet forms of the operators $`L_l`$, $`L_{i,l}`$, with respect to the grand canonical measure with $`\alpha =0`$, are denoted by $`D_l(f)`$, respectively, $`D_{i,l}(f)`$. The Dirichlet form of the operator $`L_l`$ with respect to the canonical measure is denoted by $`D_{\nu _{y,l}^c}(f)`$. For $`L_{\mathrm{}}`$, the product measures $`\nu _\alpha ^{gc}=_i\frac{e^{\alpha x_iV(x_i)}}{Z(\alpha )}dx_i`$ are equilibrium measures. We shall use the notation $`\nu _N^{\mathrm{eq}}`$ for the product probability measure $`_{i=1}^N\frac{e^{V(x_i)}}{Z}dx_i`$.
A model for surface electromigration. We consider also a perturbation of the continuum solid-on-solid model. The new system describes the evolution of a one-dimensional surface driven by both surface diffusion and surface electromigration. The surface electromigration refers to the motion of atoms on a solid surface that is caused by an electric current in the material. The electric field interacts with the atoms of the surface as the wind blows the sand particles, and a ripple pattern is observed in the long run. Electromigration along interfaces is believed to play a crucial role in the failure of metallic circuits, see Schimschak and Krug , for further details.
We assume the electric field is a continuous function $`E(t,\theta )`$ defined on $`[0,T]\times 𝕋`$. As before $`x_i`$ represents the slope of the surface at the site $`i/N`$. The generator of the system on a periodic lattice, that incorporates the action of the electric field, is
(2.8)
$$N^4L_{N,E}(f)=N^4L_N(f)+\frac{N^2}{2}\underset{i=1}{\overset{N}{}}E(t,\frac{i}{N})a(x_{i1},x_i,x_{i+1})(_{i1}2_i+_{i+1})f.$$
Hydrodynamic scaling limit of the models. We shall call $`P_{N,T}^{\mathrm{neq}}`$ and $`P_{N,T}^{\mathrm{eq}}`$ the law up to time $`T`$ of the slope process (2.2) started in some nonequilibrium distribution $`\nu _N^{\mathrm{neq}}`$, equilibrium distribution $`\nu _N^{\mathrm{eq}}`$, respectively. The law up to time $`T`$ of the slope process (2.8), driven by the electric field, started in the nonequilibrium measure $`\nu _N^{\mathrm{neq}}`$ shall be called $`P_{N,E,T}^{\mathrm{neq}}`$.
Under $`P_{N,T}^{\mathrm{neq}}`$ and $`P_{N,T}^{\mathrm{eq}}`$ the random variable
(2.9)
$$\pi _N(t)=\frac{1}{N}\left(x_1(t)\delta _{\frac{1}{N}}+\mathrm{}+x_N(t)\delta _{\frac{N}{N}}\right)$$
has distributions $`Q_{N,T}^{\mathrm{neq}}`$ and $`Q_{N,T}^{\mathrm{eq}}`$, respectively. We also refer to the random variable (2.9) as the empirical distribution. Every realization of this random variable is a measure-valued continuous path and $`Q_{N,T}^{\mathrm{neq}}`$ is a distribution on the space $`𝒳=_{(l0)}C([0,T],_l)`$. The space $`𝒳`$ is endowed with the inductive limit topology, the strongest topology that makes all the inclusions of $`C([0,T],_l)`$ continuous. The space of signed measures $`_l`$, with total variation not exceeding $`l`$, is a metrizable space with the weak topology.
We say that a model has hydrodynamic scaling limit if under certain assumptions the sequence of laws of empirical distributions has a limit that is supported on the solution of an initial value problem.
###### Definition 2.1.
A sequence of initial distributions $`\nu _N^{\mathrm{neq}}`$ on $`^N`$ is said to correspond to the macroscopic slope profile $`m_0L^1(𝕋)`$ if the random variable $`\pi _N`$ converges weakly in probability to $`\delta _{m_0(\theta )d\theta }`$, i.e., for any continuous function $`\varphi C(𝕋)`$ and any $`ϵ>0`$,
$$\underset{N\mathrm{}}{lim\; sup}\nu _N^{\mathrm{neq}}\left\{\left|\frac{1}{N}\underset{i=1}{\overset{N}{}}\varphi \left(\frac{i}{N}\right)x_i_𝕋\varphi (\theta )m_0(\theta )𝑑\theta \right|>ϵ\right\}=0.$$
Our first result of the paper says that the continuum solid-on solid model has a scaling limit and provide the form of the limiting evolution equation. More precisely,
###### Theorem 2.1.
Assume that the potential $`V(x)`$ is equal to $`x^2/2`$. Let $`m_0L^1(𝕋)`$ be a macroscopic slope profile such that $`_𝕋m_0(\theta )𝑑\theta =0`$. Assume that the sequence of initial distributions $`\{\nu _N^{\mathrm{neq}}\}_N`$ corresponds to the profile $`m_0`$ and the initial relative entropy $`H(\nu _N^{\mathrm{neq}}|\nu _N^{\mathrm{eq}})`$ is of order $`𝒪(N)`$. Then the sequence of probability measures $`\{Q_{N,T}^{\text{neq}}\}_{N0}`$ is tight in $`𝒳`$.
Any possible limit, $`Q_T`$, of a convergent subsequence of $`\{Q_{N,T}^{\text{neq}}\}_{N0}`$ is concentrated on the weak solutions $`(m(t,\theta )d\theta )_{t[0,T]}`$ of the Cauchy problem with periodic boundary conditions
(2.10)
$$_tm=\frac{1}{2}_\theta ^2(\widehat{a}(m)_\theta ^2m),m(0,\theta )=m_0(\theta ),\theta 𝕋.$$
The transport coefficient $`\widehat{a}`$ is a nonrandom continuous function on $``$, given by the following variational formula,
(2.11)
$$\widehat{a}(\alpha )=\underset{g}{inf}E_{\nu _\alpha ^{gc}}\left[a(x_1,x_0,x_1)\left(1+(_{i+1}2_i+_{i1})\left(\underset{j}{}\tau ^jg\right)\right)^2\right].$$
The infimum on the line above is taken over all local functions $`g(x_s,\mathrm{},x_s,\overline{y}_{0,l}^1)`$ of the slope configuration. The shift $`\tau ^j`$ acts on configurations $`(\tau ^jx)_k=x_{k+j}`$ and on local functions, $`(\tau ^jg)(x)=g(\tau ^jx)`$. The expectation $`E_{\nu _\alpha ^{gc}}`$ is with respect to the grand canonical measure $`\nu _\alpha ^{gc}`$.
Note. It can be shown that the transport coefficient $`\widehat{a}`$ is a continuous, bounded above and below function (see Kipnis and Landim ). Also we expect that the methods of Landim, Olla and Varadhan can show that $`\widehat{a}`$ is smooth but we will not pursue it here.
Note. By a weak solution of the Cauchy problem (2.10) we mean a path $`m𝒳`$ such that for each time $`T0`$, the value of the path $`m`$ at the time $`T`$ is a Lebesgue absolutely continuous measures on $`𝕋`$, satisfying the energy estimate
(2.12)
$$_0^T_𝕋(_\theta ^2m(t,\theta ))^2𝑑\theta 𝑑t<\mathrm{}.$$
Moreover for each $`0T<\mathrm{}`$ and for each test function $`\varphi C^{1,2}([0,T]\times 𝕋)`$,
$$_𝕋m(T,\theta )\varphi (T,\theta )𝑑\theta _𝕋m(0,\theta )\varphi (0,\theta )𝑑\theta _0^T_𝕋m(s,\theta )_s\varphi (s,\theta )d\theta ds+$$
$$+\frac{1}{2}_0^T_𝕋\widehat{a}(m)_\theta ^2m(s,\theta )_\theta ^2\varphi (s,\theta )d\theta ds=0.$$
If we assume that the initial condition $`m_0`$ has the property that $`_𝕋m_0(\theta )𝑑\theta =0`$ then for each time $`t0`$ the solution satisfies $`_𝕋m(t,\theta )𝑑\theta =0`$.
Uniqueness of weak solutions of the Cauchy problem (2.10), that satisfy the energy estimate (2.12) has not been proved yet. If the transport coefficient $`\widehat{a}`$ does not depend on the field $`m`$, the uniqueness of the Cauchy problem is known and can be found in Eidelman book .
As will be explained later the fluctuation dissipation equation for the continuum solid-on-solid model follows from the direct sum decomposition of a Hilbert space to be defined next.
We define the Hilbert space of closed functions, $`𝒞_X`$ to be the space of those $`\xi L^2(d\nu _\alpha ^{gc})`$ that satisfy in the weak sense the equations $`X_i(\tau ^j\xi )=X_j(\tau ^i\xi )`$ for all integers $`i`$ and $`j`$. It is not hard to see that a subspace of $`𝒞_X`$ is the closed linear span in $`L^2(d\nu _\alpha ^{gc})`$ of functions $`\xi _g=X_0(_j\tau ^jg)`$, where $`g`$ is a bounded local function with bounded first derivatives. Even though the infinite sum $`_j\tau ^jg`$ does not make sense, the function $`\xi _g`$ is well defined because the vector field kills all but finitely many terms of the infinite sum. We shall call this space the space of exact functions and we shall use the notation $`_X`$. The space of exact functions has codimension one inside the space of closed functions. In this paper we do not include the proof of this result, since it is very technical and is not of probabilistic nature, however we include the statement, see Lemma 2.1. Lemma 2.1 is discussed in Savu .
###### Lemma 2.1.
Let $`\mathrm{𝟏}`$ denote the constant function $`1`$. The direct sum decomposition holds:
(2.13)
$$𝒞_X=\mathrm{𝟏}_X.$$
The second result of the paper proves that the model for surface electromigration (2.8) has a scaling limit as well, and calculates the limiting evolution equation.
###### Theorem 2.2.
Suppose the hypothesis of Theorem 2.1 are satisfied. Then the sequence of probability measures $`\{Q_{N,E,T}^{\mathrm{neq}}\}_{N0}`$ is tight in $`𝒳`$, and any possible limit $`Q_{E,T}`$ is supported on the weak solutions of the Cauchy problem
(2.14)
$$_tm=\frac{1}{2}_\theta ^2(\widehat{a}(m)(_\theta ^2m+E)),m(0,\theta )=m_0(\theta ),\theta 𝕋.$$
The transport coefficient $`\widehat{a}`$ is given by the same variational formula as in the statement of Theorem 2.1.
Similar models. The continuum solid-on-solid model belongs to a large class of Ginzburg-Landau models. The slope model is of nongradient type and has an unusual dynamics because two neighboring exchanges occur always simultaneously. Nishikawa , and Bertini, Olla and Landim have investigated the hydrodynamic scaling limit in higher dimension of the gradient version (i.e., $`a`$ is a constant function) of the slope model. They have found that on the macroscale the interface follows a fourth-order nonlinear evolution equation.
Another similar model, the second-order Ginzburg-Landau model, where the sum of the heights is not conserved, was the subject of extensive discussions in the literature: the hydrodynamic scaling limit was derived by Fritz and Guo, Papanicolau and Varadhan for the gradient version, and by Varadhan for the nongradient case, whereas the nonequilibrium fluctuations have been proved by Chang and Yau . The second-order Ginzburg-Landau model and the continuum solid-on-solid model correspond to Glauber, respectively, Kawasaki dynamics in the context of interacting particle systems. As expected, a different dynamics at the mesoscopic level causes different dynamics at the macroscopic level, a second-order parabolic differential equation in the case of second-order Ginzburg-Landau model versus a fourth-order parabolic differential equation for the continuum solid-on-solid model.
## 3. Hydrodynamic scaling limit of continuum solid-on-solid model
In this section we give a sketch of the main result, Theorem 2.1. We follow a standard scheme to derive the hydrodynamic scaling limit of the continuum solid-on-solid model. The existence of the limit follows from the tightness of the sequence of probability measures $`\{Q_{N,T}^{\mathrm{neq}}\}_N`$. Let $`Q_T`$ be the limit of some weakly convergent subsequence of $`\{Q_{N,T}^{\mathrm{neq}}\}_N`$. We proceed to characterize the limit $`Q_T`$, showing that it is supported on continuous paths with certain regularity property, known as the energy estimate (2.12). The most involved part of the argument is the identification of the possible weak limit $`Q_T`$ as some probability measure supported on the weak solution of the Cauchy problem (2.10). We shall make the assumption that the sequence of initial distributions $`\{\nu _N^{\mathrm{neq}}\}_N`$ corresponds to some macroscopic profile $`m_0L^1(𝕋)`$.
Note on notations. Throughout the paper we shall make use of the shorter notation $`lim\; sup_{z_1i_1,\mathrm{},z_ni_n}f(z_1,\mathrm{},z_n)`$ for the sequence of limits
$`lim\; sup_{z_1i_1}\mathrm{}lim\; sup_{z_ni_n}f(z_1,dots,z_n)`$.
Tightness. As a consequence of Prohorov theorem and Arzela-Ascoli theorem, the tightness of the sequence $`Q_{N,T}^{\mathrm{neq}}`$ follows from the next two Lemmas.
###### Lemma 3.1.
For any test function $`\varphi C^2(𝕋)`$, any finite time $`T`$, and any $`ϵ>0`$,
(3.1)
$$\underset{\delta 0,N\mathrm{}}{lim\; sup}P_{N,T}^{\mathrm{neq}}\left\{\underset{|st|\delta ,\mathrm{\hspace{0.33em}0}s,tT}{sup}\left|\frac{1}{N}\underset{i=1}{\overset{N}{}}\varphi \left(\frac{i}{N}\right)x_i(t)\frac{1}{N}\underset{i=1}{\overset{N}{}}\varphi \left(\frac{i}{N}\right)x_i(s)\right|>ϵ\right\}=0.$$
(3.2)
$$\underset{l\mathrm{},N\mathrm{}}{lim\; sup}P_{N,T}^{\mathrm{neq}}\left\{\underset{0tT}{sup}\frac{1}{N}\underset{i=1}{\overset{N}{}}|x_i(t)|>l\right\}=0.$$
Proof. The first convergence (3.1) follows from Garsia-Rodemich-Rumsey inequality. To prove the second convergence (3.2) we can use estimates on the moment generating function of hitting time of the diffusion process (2.2). The reader may consult Kipnis and Landim or Guo, Papanicolau and Varadhan or Savu for a complete proof.
It is interesting to note that the stronger superexponential estimates can be established for the process in equilibrium,
(3.3)
$$\underset{\delta 0,N\mathrm{}}{lim\; sup}\frac{1}{N}\mathrm{log}P_{N,T}^{\mathrm{eq}}\left\{\underset{0t,sT,ts\delta }{sup}\left|_s^t\frac{1}{N}\underset{i=1}{\overset{N}{}}N^2w_i(u)\varphi \left(\frac{i}{N}\right)du\right|ϵ\right\}=\mathrm{}.$$
(3.4)
$$\underset{l\mathrm{},N\mathrm{}}{lim\; sup}\frac{1}{N}\mathrm{log}P_{N,T}^{\mathrm{eq}}\left\{\underset{0tT}{sup}\frac{1}{N}\underset{i=1}{\overset{N}{}}|x_i(t)|l\right\}=\mathrm{}.$$
Energy estimate. Any limiting point $`Q_T`$ of the measure-valued sequence $`\{Q_{N,T}^{\mathrm{neq}}\}_N`$ is supported on paths $`\mu 𝒳`$ such that at each time $`t`$, $`\mu (t)`$ is a Lebesgue absolutely continuous measure on the torus $`𝕋`$ with density $`m(t,\theta )`$. Moreover for each finite time $`T`$, the density $`m(t,\theta )`$ satisfies the energy estimate
(3.5)
$$_0^T_𝕋(_\theta ^2m(t,\theta ))^2𝑑\theta 𝑑t<\mathrm{}.$$
We note that the energy estimate (3.5) is equivalent to the inequality
(3.6)
$$\underset{\varphi C^{1,2}([0,T]\times 𝕋)}{sup}_0^T_𝕋(2m(t,\theta )_\theta ^2\varphi C\varphi ^2)𝑑\theta 𝑑t<\mathrm{},$$
where $`C`$ is a constant not depending on $`\varphi `$.
The entropy inequality (4.10), Feynman-Kac formula (4.12), and Lemma 4.2 can be used to derive the estimate
$$\underset{l,N\mathrm{}}{lim\; sup}\underset{\varphi C^{1,2}([0,T]\times 𝕋)}{sup}E^{\mathrm{neq}}[\frac{1}{N}\underset{i=1}{\overset{N}{}}_0^T2\varphi (t,\frac{i}{N})N^2\mathrm{Av}_{j=il+1}^{i+l1}(\mathrm{\Delta }x)_j$$
(3.7)
$$𝕍_{i,l}(\mathrm{\Delta }x,y)\varphi ^2(t,\frac{i}{N})dt]<\mathrm{}.$$
Here, the cylinder function $`\mathrm{\Delta }x`$ is the discrete Laplacian of the slope field, $`x_12x_0+x_1`$, and the variance $`𝕍_{i,l}(\mathrm{\Delta }x,y)`$ is defined later in section 4, see definition LABEL:. As will be proved in Lemma 5.2, the variance, $`𝕍_{i,l}(\mathrm{\Delta }x,y)`$, has a uniform-in-$`y`$ lower bound, therefore the estimate (3.6) follows from (3.7), after integrating by parts. Moreover, we can conclude that at any time $`t`$ the weak second derivative of the measure $`\mu (t)`$ is in $`L^2(𝕋,d\theta )`$ and hence $`\mu (t)`$ is absolutely continuous with respect to Lebesgue measure on the torus $`𝕋`$.
Identification of the equation. That any limiting point of the measure-valued sequence $`\{Q_{N,T}^{\mathrm{neq}}\}_{N>0}`$ is supported on the weak solutions of (2.10) follows if the event corresponding to the violation of the limiting equation has probability zero in the limit.
For each test function $`\varphi C^{1,2}([0,T]\times 𝕥)`$ that are differentiable in time and twice differentiable in space, and each finite time $`T`$ we define the function
(3.8) $`V(t)={\displaystyle \frac{1}{N}}{\displaystyle \underset{i=1}{\overset{N}{}}}\varphi (t,{\displaystyle \frac{i}{N}})x_i(t){\displaystyle \frac{1}{N}}{\displaystyle \underset{i=1}{\overset{N}{}}}\varphi (0,{\displaystyle \frac{i}{N}})x_i(0)`$
$`{\displaystyle _0^t}{\displaystyle \frac{1}{N}}{\displaystyle \underset{i=1}{\overset{N}{}}}_s\varphi (s,{\displaystyle \frac{i}{N}})\overline{x}_{i,aN}(s)ds+{\displaystyle \frac{1}{2}}{\displaystyle _0^t}{\displaystyle \frac{1}{N}}{\displaystyle \underset{i=1}{\overset{N}{}}}\varphi ^{\prime \prime }(s,{\displaystyle \frac{i}{N}})\widehat{a}(\overline{x}_{i,aN}(s))\times `$
$`\times b^2(\overline{x}_{ibN,cN}(s)2\overline{x}_{i,cN}(s)+\overline{x}_{i+bN,cN}(s))ds,`$
and the event
(3.9)
$$O_{a,b,c,ϵ}^\varphi =\left\{\underset{0tT}{sup}|V(t)|>ϵ\right\}.$$
We will prove in the next sections that for each $`ϵ>0`$ we have
(3.10)
$$\underset{a,b,c0,N\mathrm{}}{lim\; sup}P_{N,T}^{\mathrm{neq}}(O_{a,b,c,ϵ}^\varphi )=0.$$
The proof of the result (3.10) is complicated and is divided into several steps. The function that defines this event can be written as a sum of functions, see the beginning of section 4. We will deal separately with each function in the sum and show that it converges to zero in probability.
## 4. Identification of the limiting equation
In this section we establish that the event (3.9) is negligible in the limit, and hence any weak limit $`Q_T`$ is supported on the solutions of the Cauchy problem (2.10). To save space, we are suppressing the time dependence of the test function $`\varphi `$ that defines the event (3.9).
Note on notations. Assume $`f`$ is some local function. We denote by $`\mathrm{Av}_{j=il}^{i+l}\tau ^jf`$ the average of shifts of $`f`$, namely
$$\frac{\tau ^{il}f+\mathrm{}+\tau ^{i+l}f}{2l+1}$$
.
We write $`V(t)=`$
(4.1)
$`={\displaystyle \frac{1}{N}}{\displaystyle \underset{i=1}{\overset{N}{}}}\varphi \left({\displaystyle \frac{i}{N}}\right)x_i(t){\displaystyle \frac{1}{N}}{\displaystyle \underset{i=1}{\overset{N}{}}}\varphi \left({\displaystyle \frac{i}{N}}\right)x_i(0){\displaystyle _0^t}{\displaystyle \frac{1}{N}}{\displaystyle \underset{i=1}{\overset{N}{}}}\varphi \left({\displaystyle \frac{i}{N}}\right)N^4L_N(x_i)ds+`$
(4.2)
$`+{\displaystyle \frac{1}{2}}{\displaystyle _0^t}{\displaystyle \frac{1}{N}}{\displaystyle \underset{i=1}{\overset{N}{}}}\left\{N^2\left[\varphi \left({\displaystyle \frac{i1}{N}}\right)2\varphi \left({\displaystyle \frac{i}{N}}\right)+\varphi \left({\displaystyle \frac{i+1}{N}}\right)\right]\varphi ^{\prime \prime }\left({\displaystyle \frac{i}{N}}\right)\right\}w_ids+`$
(4.3) $`+{\displaystyle \frac{1}{2}}{\displaystyle _0^t}{\displaystyle \frac{1}{N}}{\displaystyle \underset{i=1}{\overset{N}{}}}\varphi ^{\prime \prime }\left({\displaystyle \frac{i}{N}}\right)N^2(w_i\mathrm{Av}_{j=il_1}^{i+l_1}w_j)ds+`$
(4.4) $`+{\displaystyle _0^t}{\displaystyle \frac{1}{N}}{\displaystyle \underset{i=1}{\overset{N}{}}}\varphi ^{\prime \prime }\left({\displaystyle \frac{i}{N}}\right)N^2[\mathrm{Av}_{j=il_1}^{i+l_1}w_j\widehat{a}(\overline{x}_{i,l})\mathrm{Av}_{j=il_1}^{i+l_1}(\mathrm{\Delta }x)_j`$
$`\mathrm{Av}_{j=il_1}^{i+l_1}\tau ^jL_{\mathrm{}}f_r]ds+`$
(4.5) $`+{\displaystyle \frac{1}{2}}{\displaystyle _0^t}{\displaystyle \frac{1}{N}}{\displaystyle \underset{i=1}{\overset{N}{}}}\varphi ^{\prime \prime }\left({\displaystyle \frac{i}{N}}\right)N^2L_{\mathrm{}}(\mathrm{Av}_{j=il_1}^{i+l_1}\tau ^jf_r)ds`$
(4.6) $`+{\displaystyle \frac{1}{2}}{\displaystyle _0^t}{\displaystyle \frac{1}{N}}{\displaystyle \underset{i=1}{\overset{N}{}}}\varphi ^{\prime \prime }({\displaystyle \frac{i}{N}})N^2[\widehat{a}(\overline{x}_{i,l})\mathrm{Av}_{j=il_1}^{i+l_1}(\mathrm{\Delta }x)_j`$
$`\widehat{a}(\overline{x}_{i,aN})b^2(\overline{x}_{ibN,cN}2\overline{x}_{i,cN}+\overline{x}_{i+bN,cN})]ds.`$
We proceed to prove that each term in the sum above converges to $`0`$ in probability.
Martingale estimate (the term (4.1)). We call $`M_N(t)`$ the term (4.1). From Ito formula we know that the process $`\{M_N(t)\}_{t0}`$ is a martingale and
$$E^{\mathrm{neq}}[M_N^2(T)]=E^{\mathrm{neq}}\left[_0^T\frac{1}{N^2}\underset{i=1}{\overset{N}{}}N^4\left[\varphi \left(\frac{i1}{N}\right)2\varphi \left(\frac{i}{N}\right)+\varphi \left(\frac{i+1}{N}\right)\right]^2a_ids\right].$$
The test function $`\varphi `$ is chosen to have continuous second derivatives, therefore $`E^{\mathrm{neq}}[M_N^2(T)]`$ is of order $`𝒪(\frac{1}{N})`$. We can use the Doob’s inequality
$$P_{N,T}^{\mathrm{neq}}\left\{\underset{0tT}{sup}|M_N(t)|ϵ\right\}\frac{1}{ϵ^2}E[M_N^2(T)],ϵ>0,$$
to conclude that the martingale is negligible in the limit, i.e.,
$$\underset{N\mathrm{}}{lim}P_{N,T}^{\mathrm{neq}}\left\{\underset{0tT}{sup}|M_N(t)|ϵ\right\}=0,ϵ>0.$$
The term (4.2). A straightforward computation involving the Chebyshev inequality and the entropy inequality (4.10) proves that the term (4.2) converges in probability to $`0`$. The test function needs to have continuous fourth derivative.
A technical Lemma. We shall prove a Lemma that reduces the problem of establishing the negligibility of an event to finding that the largest eigenvalue of a Schrödinger operator is negative. For an operator $`A:`$ acting on a Hilbert space $``$ we denote by $`\mathrm{supspec}_{}A`$ the largest value in the spectrum of $`A`$.
###### Lemma 4.1.
Let $`\{x(t)\}_{t0}`$ be the slope process with generator (2.3). Under the assumption
(4.7)
$$\underset{N\mathrm{}}{lim\; sup}\mathrm{supspec}_{L^2(\nu _N^{\mathrm{eq}})}\left(\alpha g+\frac{N^4}{N}L_N\right)=$$
$$=\underset{N\mathrm{}}{lim\; sup}\underset{\rho ,E^{\mathrm{eq}}[\rho ^2]=1}{sup}\left[\alpha E^{\mathrm{eq}}[g\rho ^2]\frac{N^4}{N}D_N(\rho )\right]0,\alpha 0$$
it follows that the event $`\left\{\left|_0^Tg(x(s))𝑑s\right|ϵ\right\}`$ has negligible probability or:
(4.8)
$$\underset{N\mathrm{}}{lim}P_{N,T}^{\mathrm{neq}}\left\{\left|_0^Tg(x(t))𝑑t\right|ϵ\right\}=0,ϵ>0.$$
Proof. We can use Chebyshev inequality to reduce the proof of (4.8) to
(4.9)
$$\underset{N\mathrm{}}{lim}E^{\mathrm{neq}}\left[\left|_0^Tg(x(t))𝑑t\right|\right]=0.$$
Since we do not have much information about the initial nonequilibrium distribution we use the entropy inequality to replace the nonequilibrium distribution in (4.9) by the equilibrium distribution.
Before we continue, we remind that given two probability measures $`\nu `$ and $`\mu `$ on the same probability space such that $`\nu `$ is absolutely continuous with respect to $`\mu `$, we define the relative entropy of $`\nu `$ with respect to $`\mu `$ by $`H(\nu |\mu )=E_\mu \left[\frac{d\nu }{d\mu }\mathrm{log}\frac{d\nu }{d\mu }\right]`$ where $`\frac{d\nu }{d\mu }`$ is the Radon-Nikodym derivative of $`\nu `$ relative to $`\mu `$. The entropy $`H(\nu |\mu )`$, always a positive quantity, is the optimal constant that makes the entropy inequality
(4.10)
$$E_\nu [f]\frac{1}{\alpha }\left\{H(\nu |\mu )+\mathrm{log}E_\mu [e^{\alpha f}]\right\}.$$
true for any bounded, measurable function $`f`$ and $`\alpha >0`$. A trivial consequence of the entropy inequality (4.10) helps us to estimate $`\nu (A)`$, where $`A`$ is some event,
(4.11)
$$\nu (A)\frac{\mathrm{log}(2)+H(\nu |\mu )}{\mathrm{log}(1+\frac{1}{\mu (A)})}$$
In our context we use the entropy inequality for the distribution of the process started in nonequilibrium and the distribution of the process started in the equilibrium. We have,
$$E^{\mathrm{neq}}\left[\left|_0^Tg(x(t))𝑑t\right|\right]=E^{\mathrm{neq}}\left[\frac{1}{N\alpha }\left|_0^TN\alpha g(x(t))𝑑t\right|\right]$$
$$\frac{1}{N\alpha }H(\nu _N^{\mathrm{neq}}|\nu _N^{\mathrm{eq}})+\frac{1}{N\alpha }\mathrm{log}E^{\mathrm{eq}}\left[\mathrm{exp}\left(\alpha \left|_0^TNg(x(t))𝑑t\right|\right)\right]$$
$$\frac{C}{\alpha }+\frac{1}{N\alpha }\mathrm{log}E^{\mathrm{eq}}\left[\mathrm{exp}\left(N\alpha _0^Tg(x(t))𝑑t\right)+\mathrm{exp}\left(N\alpha _0^Tg(x(t))𝑑t\right)\right].$$
As a consequence of the inequality
$$\mathrm{log}(a+b)\mathrm{max}(\mathrm{log}(2a),\mathrm{log}(2b))\mathrm{log}(2)+\mathrm{max}(\mathrm{log}(a),\mathrm{log}(b))$$
for two positive numbers $`a`$ and $`b`$, (4.9) follows as soon as we have
$$\underset{N\mathrm{}}{lim\; sup}\frac{1}{N}\mathrm{log}E^{\mathrm{eq}}\left[\mathrm{exp}\left(N\alpha _0^Tg(x(t))𝑑t\right)\right]0,\alpha 0.$$
A trivial consequence of Feynman-Kac proves our Lemma,
(4.12)
$$\frac{1}{N}\mathrm{log}E^{\mathrm{eq}}\left[\mathrm{exp}\left(N\alpha _0^Tg(x(t))𝑑t\right)\right]T\mathrm{supspec}_{L^2(\nu _N^{\mathrm{eq}})}(\alpha g+\frac{N^4}{N}L_N).$$
The microscopic current $`w`$ can be replaced by local average of currents (the term (4.3)). We shall show that the term (4.3) converges to zero in probability. As a consequence the current $`w`$ is replaced by a local average of $`w`$, that is closer to a deterministic value. $`w`$ by itself is a single fluctuating random variable.
We check that the hypothesis of Lemma 4.1 is valid for the function $`g_{l,N}=\frac{1}{N}_{i=1}^N\varphi (\frac{i}{N})N^2(w_i\mathrm{Av}_{j=il}^{i+l}w_j)`$. Recall that if the test function $`\varphi `$ has continuous second-order derivative the quantity
$$\mathrm{Av}_{j=il}^{i+l}\varphi (\frac{j}{N})\varphi (\frac{i}{N})=\frac{\varphi (\frac{il}{N})\varphi (\frac{i}{N})+\mathrm{}+\varphi (\frac{i+l}{N})\varphi (\frac{i}{N})}{2l+1}$$
is of order $`𝒪(\frac{l^2}{N^2})`$. Now, on integrating by parts twice, it follows for a fixed function $`\rho `$, with $`E^{\mathrm{eq}}[\rho ^2]=1`$, that,
$$\left|E^{\mathrm{eq}}\left[\rho ^2\frac{1}{N}\underset{i=1}{\overset{N}{}}\varphi (\frac{i}{N})N^2(w_i\mathrm{Av}_{j=il}^{i+l}w_j)\right]\right|^2=$$
$$=4\left|E^{\mathrm{eq}}\left[\frac{N^2}{N}\underset{i=1}{\overset{N}{}}\sqrt{a_i}X_i(\rho )\sqrt{a_i}\left(\mathrm{Av}_{j=il}^{i+l}\varphi (\frac{j}{N})\varphi (\frac{i}{N})\right)\rho \right]\right|^2$$
$$=4E^{\mathrm{eq}}\left[\underset{i=1}{\overset{N}{}}\frac{N^4}{N^2}a_i\left(\mathrm{Av}_{j=il}^{i+l}\varphi (\frac{j}{N})\varphi (\frac{i}{N})\right)^2\rho ^2\right]D_N(\rho )C\frac{l^4}{N^4}\frac{N^4}{N}D_N(\rho ).$$
Therefore,
$$\underset{l,N\mathrm{}}{lim\; sup}\underset{\rho ,E^{\mathrm{eq}}[\rho ^2]=1}{sup}\frac{\left(E^{\mathrm{eq}}\left[\rho ^2\frac{N^2}{N}_{i=1}^N\left(\mathrm{Av}_{j=il}^{i+l}\varphi (\frac{j}{N})\varphi (\frac{i}{N})\right)w_i\right]\right)^2}{D_N(\rho )}\frac{N}{N^4}=0,$$
and hence (4.7) is satisfied. Moreover,
$$\underset{l,N\mathrm{}}{lim}P_{N,T}^{\mathrm{neq}}\left\{\left|_0^T\frac{1}{N}\underset{i=1}{\overset{N}{}}\varphi (\frac{i}{N})N^2(w_i\mathrm{Av}_{j=il}^{i+l}w_j)dt\right|ϵ\right\}=0,ϵ>0.$$
Inserting the fluctuations (the term (4.5)). Let $`f(x_s,\mathrm{},x_s,\overline{x}_{0,l})C^2(^{2s+1})`$ be a local function that depends on the slope configuration in a box of size $`s`$ and on the mean slope in a box of a large size $`l`$. We use the notation $`l_1=l\sqrt{l}`$ for a slightly smaller $`l`$. We want to show that the fluctuations approach zero, in the limit or that
(4.13)
$$\underset{l,N\mathrm{}}{lim}P_{N,T}^{\mathrm{neq}}\left\{\left|_0^T\frac{N^2}{N}\underset{i=1}{\overset{N}{}}\varphi \left(\frac{i}{N}\right)L_N(\mathrm{Av}_{j=il_1}^{i+l_1}\tau ^jf)dt\right|ϵ\right\}=0,ϵ>0.$$
We apply Ito formula,
$$_0^T\frac{N^2}{N}\underset{i=1}{\overset{N}{}}\varphi \left(\frac{i}{N}\right)L_N(\mathrm{Av}_{j=il_1}^{i+l_1}\tau ^jf)dt=$$
(4.14)
$$=\frac{1}{NN^2}\underset{i=1}{\overset{N}{}}\varphi \left(\frac{i}{N}\right)\left[(\mathrm{Av}_{j=il_1}^{i+l_1}\tau ^jf)(x(T))(\mathrm{Av}_{j=il_1}^{i+l_1}\tau ^jf)(x(0))\right]+\frac{1}{N^2}M_N(t).$$
The summand in (4.14) converges to zero, as the function $`f`$ is bounded. The second part of (4.14) approaches zero because the $`L^2`$ norm of $`\frac{M_N(T)}{N^2}`$ is of order $`𝒪(\frac{l^2}{N})`$, as we can see below. Let $`g=\frac{1}{N}_{i=1}^N\varphi (\frac{i}{N})(\mathrm{Av}_{j=il_1}^{i+l_1}\tau ^jf)(x)`$, then
$$E^{\mathrm{neq}}\left[\frac{M_N(T)^2}{N^4}\right]=E^{\mathrm{neq}}\left[_0^TL_Ng2gL_Ngdt\right]=E^{\mathrm{neq}}\left[_0^T\underset{i=1}{\overset{N}{}}a_i(X_ig)^2dt\right]=$$
$$=\frac{1}{N^2}E^{\mathrm{neq}}\left[_0^T\underset{i=1}{\overset{N}{}}a_i\left(\underset{k=1}{\overset{N}{}}\varphi \left(\frac{k}{N}\right)X_i(\mathrm{Av}_{j=kl_1}^{k+l_1}\tau ^jf)\right)^2dt\right].$$
The function $`\mathrm{Av}_{j=kl_1}^{k+l_1}\tau ^jf`$ is a cylinder function that depends just on the sites $`n`$ such that $`klnk+l`$. Therefore the vector field $`X_i`$ is zero when acting on most of the summands inside $`\mathrm{Av}_{j=kl_1}^{k+l_1}\tau ^jf`$. There are no more than $`2l`$ sites $`k`$ such that $`X_i(\mathrm{Av}_{j=kl_1}^{k+l_1}\tau ^jf)0`$. We put all these arguments together to conclude that $`E^{\mathrm{neq}}[M^2(t)/N^4]Cl^2/N.`$
Replacing the current by the Laplacian of the slope field (the term (4.4)). In our model the instantaneous current, $`w`$, can not be written as the discrete Laplacian $`\tau h2h\tau ^1h`$ of some local function $`h`$, thus we use the method of Varadhan for computing the hydrodynamic scaling limit of our model. The main idea is that the current decomposes as
(4.15)
$$w=\widehat{a}(\overline{x}_l)\mathrm{\Delta }x+L_{\mathrm{}}f$$
for a suitable coefficient $`\widehat{a}(\overline{x}_l)`$, where $`\overline{x}_l`$ is the average slope in a cube centered at the origin of microscopic side $`l`$. The equation (4.15) is known in the literature as the fluctuation-dissipation equations. We have explained before that terms of the form $`L_{\mathrm{}}f`$ have no effect on the macroscopic scale. A new feature is characteristic to our model due to the complex interaction of the system: after filtering off the fluctuations from the current we are left with the Laplacian of some function and not a gradient, as happened for models previously considered in the literature. The precise meaning of (4.15) is given below.
Our aim is to prove the existence of a sequence $`\{f_r\}_{r0}`$ of local functions and of the transport coefficient $`\widehat{a}`$ such that
$`\underset{r,l,N\mathrm{}}{lim\; sup}P_{N,T}^{\mathrm{neq}}\left\{\right|{\displaystyle _0^T}{\displaystyle \frac{N^2}{N}}{\displaystyle \underset{i=1}{\overset{N}{}}}\varphi \left({\displaystyle \frac{i}{N}}\right)[\mathrm{Av}_{j=il_1}^{i+l_1}w_j\widehat{a}(\overline{x}_{i,l})\mathrm{Av}_{j=il_1}^{i+l_1}(\mathrm{\Delta }x)_j`$
(4.16) $`\mathrm{Av}_{j=il_1}^{i+l_1}\tau ^jL_{\mathrm{}}f_r]dt|ϵ\}=0.`$
The local function $`f_r(x_s,\mathrm{},x_s,\overline{x}_{0,l})`$ depends on two arguments, the mean slope $`\overline{x}_{0,l}`$ in a box of size $`l`$ and the slope field inside a box of size $`s`$, the size $`s`$ being much smaller than $`l`$. We can assume that the operator $`L_{\mathrm{}}`$ does not act on the first argument $`\overline{x}_{0,l}`$, since we can show that the action of the operator $`L_{\mathrm{}}`$ at the boundary sites is negligible. To be more precise, $`\mathrm{Av}_{j=il_1}^{i+l_1}\tau ^jL_{\mathrm{}}f_r`$ stands for
$$\frac{L_{\mathrm{}}f_r(\overline{x}_{i,l},\tau _{il_1}x)+\mathrm{}+L_{\mathrm{}}f_r(\overline{x}_{i,l},\tau _{i+l_1}x)}{2l_1+1}.$$
Note that the function $`\mathrm{Av}_{j=il_1}^{i+l_1}\tau ^jL_{\mathrm{}}f_r`$ depends just on the value of the field inside the box centered at $`i`$ and of size $`l`$.
As before we use Lemma 4.1 to conclude that the event (4) has negligible probability in the limit. An additional difficulty shows up. If $`\lambda _ϵ`$ is the largest eigenvalue of the perturbation, $`L+ϵW`$, of a negative operator $`L`$ with principal eigenvalue $`0`$, the eigenvalue $`\lambda _ϵ`$ has the formal series expansion,
$$\lambda _ϵ=0+ϵE_\nu [W]+ϵ^2<W,(L)^1W>_\nu +𝒪(ϵ^3),$$
Hence if the potential $`W`$ has mean zero, one expects $`lim_{ϵ0}\lambda _ϵϵ^2=<W,(L)^1W>_\nu `$. Fortunately for suitable potential $`W`$ the central limit variance $`<W,(L)^1W>_\nu `$ converges to zero.
We shall need in our context, a particular result about the largest eigenvalue of a perturbation operator, whose proof is found in Quastel \[Qua2\].
###### Lemma 4.2.
Let $`W`$ be a real potential that satisfies
$$<u,Wu>_\nu l^{1/2}D_l(u)^{1/2}||u||_2,l<Cϵ^{2/5}$$
for some $`C`$ small enough, or
$$W_{\mathrm{}}C,K(Cϵ)^{1/5}.$$
Provided that the generator $`l^4L_l`$ has spectral gap of order one the following estimate holds
(4.17)
$$ϵ^2l^5\mathrm{supspec}_{L^2(\nu )}(l^4L_lϵl^5W)l<W,(L_l)^1W>_\nu +𝒪(1).$$
Spectral gap. Indeed the generator $`l^4L_l`$ of our model, defined by (2.6), has a spectral gap of order $`1`$. The proof is standard by the method of Bakry-Emery (see Chang and Yau or Deuchel and Stroock ).
The operator $`L_l`$ is an unbounded operator defined on the subspace $`C_0^{\mathrm{}}(^{2l+1})`$ of the Hilbert space $`L^2(\nu _{\alpha ,l}^{gc})`$ and is negative definite, with spectrum included in the negative semiaxis of the real line. $`0`$ is an eigenvalue of the operator $`L_l`$ but the eigenspace corresponding to this eigenvalue is quite large, being infinite dimensional.
It is not hard to see that we can write the Hilbert space $`L^2(\nu _{\alpha ,l}^{gc})`$ as the direct sum $`_{y^2}L^2(\nu _{y,l}^c)`$. Moreover because of the ergodicity of the dynamics (2.6) on the level sets of the function $`y_{0,l}=(y_{0,l}^1,y_{0,l}^2)`$, we know that the eigenspace corresponding to zero of the restriction of the operator $`L_l`$ onto each Hilbert subspace $`L^2(\nu _{y,l}^c)`$, is one dimensional. The next eigenvalue of the restriction $`L_l|_{L^2(\nu _{y,l}^c)}`$ is a negative number. The distance between the largest eigenvalue and the next largest eigenvalue of the operator $`L_l|_{L^2(\nu _{y,l}^c)}`$ is called the spectral gap of the operator, because it is the gap in the spectrum of the operator.
###### Lemma 4.3.
(Spectral gap) There is a universal constant $`C`$ that does not depend on the conserved quantities $`y_{0,l}=(y_{0,l}^1,y_{0,l}^2)^2`$ such that
(4.18)
$$\frac{C}{l^4}E_{\nu _{y,l}^c}[\rho ^2]<(L_l)\rho ,\rho >_{\nu _{y,N}^c}$$
for any mean-zero function $`\rho `$, $`E_{\nu _{y,l}^c}[\rho ]=0`$.
Since the generator $`L_l`$ of our model has a spectral gap, see Lemma 4.3, we can introduce the central limit theorem variance in our context.
###### Definition 4.1.
Suppose we have a cylinder function $`f(x_s,\mathrm{}x_s)`$ such that
$`E_{\nu _{y,s}^c}[f]=0`$ for all possible values of $`y^2`$. Recall that $`\nu _{y,s}^c`$ is the canonical measure in a box centered at $`0`$ and size $`s`$ defined in section 2. The central limit theorem variance of $`f`$ on the box $`\mathrm{\Lambda }_{i,l}`$ is defined to be:
(4.19)
$$𝕍_{i,l}(f,y)=2(2l)<\mathrm{Av}_{j=il_1}^{i+l_1}\tau ^jf,(L_{i,l})^1(\mathrm{Av}_{j=il_1}^{i+l_1}\tau ^jf)>_{\nu _{y,i,l}^c}.$$
If the box $`\mathrm{\Lambda }_{i,l}`$ is centered at $`0`$ then we use the shorter notation $`𝕍_l(f,y)`$ for the CLT-variance. At this point we stress that $`𝕍_{i,l}(f,y)`$ is a local function depending on the field inside of $`\mathrm{\Lambda }_{i,l}`$, more precisely depending on the conserved quantities $`y_{i,l}=(y_{i,l}^1,y_{i,l}^2)`$, the mean slope and the linear mean of the slope field.
At this point we stress that $`𝕍_{i,l}(f,y)`$ is a local function depending on the field inside of $`\mathrm{\Lambda }_{i,l}`$, more precisely depending on the conserved quantities $`y_{i,l}=(y_{i,l}^1,y_{i,l}^2)`$, the mean slope and the linear mean of the slope field.
The strategy is to give a bound for the largest eigenvalue of a perturbation operator in terms of the CLT-variance $`𝕍_{i,l}(f,y)`$. Extra care must be taken because the CLT-variance $`𝕍_{i,l}(f,y)`$ is uniformly-in-y small on bounded sets and not on unbounded sets.
The canonical measure $`\nu _{i,l,y}^c`$ for our model has been obtained by conditioning the grand canonical measure on the configurations with fixed mean slope and fixed linear mean slope in a box centered at $`i`$ and of size $`l`$. The second conditioning makes the canonical measure not to have identical marginals. Actually the expected values of the marginals depends linearly on the site. However the finite-dimensional marginals of the canonical measure converges toward the finite-dimensional marginals of the grand canonical distribution as the size of the box approaches infinity, see Lemma 5.1. To benefit of this fact we will replace the CLT-variance $`𝕍_l(f,y)`$ in a box of size $`l`$ with its expectation $`E^{\mathrm{eq}}[𝕍_l(f,y)|y_k]`$ with respect to the canonical measure $`\nu _{y,k}^c`$ in a box of larger size $`k`$. We let $`k`$ go first to infinity. We formalize below.
Let us define the function g as
(4.20)
$$g=\frac{N^2}{N}\underset{i=1}{\overset{N}{}}\varphi \left(\frac{i}{N}\right)\left[\mathrm{Av}_{j=il_1}^{i+l_1}w_j\widehat{a}(\overline{x}_{i,l})\mathrm{Av}_{j=il_1}^{i+l_1}(\mathrm{\Delta }x)_j\mathrm{Av}_{j=il_1}^{i+l_1}\tau ^jL_{\mathrm{}}f_r\right]$$
Thanks to Lemma 4.1 the event (4) has negligible probability if
$$\underset{r,l,N\mathrm{}}{lim\; sup}\mathrm{supspec}_{L^2(\nu _N^{\mathrm{eq}})}(g+2\beta \frac{N^4}{N}L_N)0,\beta >0.$$
We write the operator $`g+2\beta \frac{N^4}{N}L_N`$ as a sum of operators and we estimate the size of the principal eigenvalue of each operator in the sum.
(4.21)
$$g+2\beta \frac{N^4}{N}L_N=\mathrm{\Omega }_1+\mathrm{\Omega }_2+\mathrm{\Omega }_3+\mathrm{\Omega }_4$$
where,
$`\mathrm{\Omega }_1`$ $`=`$ $`g{\displaystyle \frac{1}{\beta N}}{\displaystyle \underset{i=1}{\overset{N}{}}}\varphi \left({\displaystyle \frac{i}{N}}\right)^2𝕍_{i,l}(w\widehat{a}(\overline{x}_{i,l})\mathrm{\Delta }xL_{\mathrm{}}f_r,y)+\beta {\displaystyle \frac{N^4}{N}}L_N`$
$`\mathrm{\Omega }_2`$ $`=`$ $`{\displaystyle \frac{1}{\beta N}}{\displaystyle \underset{i=1}{\overset{N}{}}}\varphi \left({\displaystyle \frac{i}{N}}\right)^2[𝕍_{i,l}(w\widehat{a}(\overline{x}_{i,l})\mathrm{\Delta }xL_{\mathrm{}}f_r,y)`$
$`E^{\mathrm{eq}}[𝕍_{i,l}(w\widehat{a}(\overline{x}_{i,l})\mathrm{\Delta }xL_{\mathrm{}}f_r,y)|y_{i,k}]]+\beta {\displaystyle \frac{N^4}{N}}L_N`$
$`\mathrm{\Omega }_3`$ $`=`$ $`{\displaystyle \frac{1}{\beta N}}{\displaystyle \underset{i=1}{\overset{N}{}}}\varphi \left({\displaystyle \frac{i}{N}}\right)^2E^{\mathrm{eq}}[𝕍_{i,l}(w\widehat{a}(\overline{x}_{i,l})\mathrm{\Delta }xL_{\mathrm{}}f_r,y)|y_{i,k}]\mathrm{𝟏}_{|y_{i,k}|\delta }`$
$`\mathrm{\Omega }_4`$ $`=`$ $`{\displaystyle \frac{1}{\beta N}}{\displaystyle \underset{i=1}{\overset{N}{}}}\varphi \left({\displaystyle \frac{i}{N}}\right)^2E^{\mathrm{eq}}[𝕍_{i,l}(w\widehat{a}(\overline{x}_{i,l})\mathrm{\Delta }xL_{\mathrm{}}f_r,y)|y_{i,k}]\mathrm{𝟏}_{|y_{i,k}|\delta }.`$
The operators $`\mathrm{\Omega }_2`$, $`\mathrm{\Omega }_3`$ and $`\mathrm{\Omega }_4`$ are understood as multiplication operators.
The operator $`\mathrm{\Omega }_1`$. Assume that $`M`$ is an upper bound for the test function $`|\varphi |`$. Define
$$g_{i,l}=\mathrm{Av}_{j=il_1}^{i+l_1}w_j\widehat{a}(\overline{x}_{i,l})\mathrm{Av}_{j=il_1}^{i+l_1}(\mathrm{\Delta }x)_j\mathrm{Av}_{j=il_1}^{i+l_1}\tau ^jL_{\mathrm{}}f_r.$$
We have
$`\mathrm{supspec}_{L^2(\nu _N^{\mathrm{eq}})}(\mathrm{\Omega }_1)`$
$`\underset{|\lambda |M}{sup}\underset{\rho ,E^{\mathrm{eq}}[\rho ^2]=1}{sup}[E^{\mathrm{eq}}[\lambda N^2g_{0,l}\rho ^2{\displaystyle \frac{\lambda ^2}{\beta }}𝕍_l(w\widehat{a}(\overline{x}_l)\mathrm{\Delta }xL_{\mathrm{}}f_r,y)\rho ^2]`$
$`{\displaystyle \frac{\beta N^4}{2l+1}}D_l(\rho )]`$
$`\underset{|\lambda |M}{sup}\underset{y_l^2}{sup}[\underset{\rho ,E^{\mathrm{eq}}[\rho ^2|y_l]=1}{sup}[E^{\mathrm{eq}}[\lambda N^2g_{0,l}\rho ^2]{\displaystyle \frac{\beta N^4}{2l+1}}D_{\nu _{y,l}^c}(\rho )]`$
$`{\displaystyle \frac{\lambda ^2}{\beta }}𝕍_l(w\widehat{a}(\overline{x}_l)\mathrm{\Delta }xL_{\mathrm{}}f_r,y)].`$
Integrating by parts we can show that there is a constant $`C`$, not depending on $`l`$ and the values of the conserved quantities, such that for each density $`\rho `$:
$$E^{\mathrm{eq}}[g_{0,l}\rho ^2|y_l]C\frac{\sqrt{D_{\nu _{y,l}^c}(\rho )}}{\sqrt{2l+1}}.$$
Hence the hypothesis of Lemma 4.2 is satisfied and
$$\underset{N\mathrm{}}{lim\; sup}\underset{\rho ,E^{\mathrm{eq}}[\rho ^2|y_l]=1}{sup}\left[E^{\mathrm{eq}}[\lambda N^2g_{0,l}\rho ^2]\frac{\beta N^4}{2l+1}D_{\nu _{y,l}^c}(\rho ^2)\right]=$$
$$=\frac{1}{(2l+1)l^4}\underset{N\mathrm{}}{lim\; sup}N^4\underset{\rho ,E^{\mathrm{eq}}[\rho ^2|y_l]=1}{sup}\left[E^{\mathrm{eq}}[\frac{\lambda (2l+1)l^4}{N^2}g_{0,l}\rho ^2]\beta l^4D_{\nu _{y,l}^c}(\rho ^2)\right]$$
$$\frac{\lambda ^2}{\beta }𝕍_l(w\widehat{a}(\overline{x}_l)\mathrm{\Delta }xL_{\mathrm{}}f_r,y).$$
The convergence on the line above is uniform over the set of all possible values of the conserved quantity $`y_l=(y_l^1,y_l^2)^2`$, therefore the principal eigenvalue of the operator $`\mathrm{\Omega }_1`$ becomes negative as $`N\mathrm{}`$.
The operator $`\mathrm{\Omega }_2`$. Let us call
$$v_{i,k}=V_{i,l}(w\widehat{a}(\overline{x}_{i,l})\mathrm{\Delta }xL_{\mathrm{}}f_r,y)E^{eq}[V_{i,l}(w\widehat{a}(\overline{x}_{i,l})\mathrm{\Delta }xL_{\mathrm{}}f_r,y)|y_{i,k}].$$
We observe that $`\mathrm{\Omega }_2`$ contains $`v_{i,k}`$ without being multiplied with a factor $`N^2`$. Then,
$$\mathrm{supspec}_{L^2(\nu _N^{\mathrm{eq}})}(\mathrm{\Omega }_2)=\mathrm{supspec}_{L^2(\nu _N^{\mathrm{eq}})}\left[\frac{1}{\beta N}\underset{i}{}\varphi \left(\frac{i}{N}\right)^2v_{i,k}+\beta \frac{N^4}{N}L_N\right]$$
$$\underset{|\lambda |M}{sup}\underset{y_k^2}{sup}[\underset{\rho ,E^{\mathrm{eq}}[\rho ^2|y_k]=1}{sup}[E^{\mathrm{eq}}\left[\frac{\lambda ^2}{\beta }v_{0,k}\rho ^2\right]\frac{\beta N^4}{2k+1}D_{\nu _{y,k}^c}(\rho )]=$$
$$=\frac{1}{(2k+1)k^4}\underset{|\lambda |M}{sup}\underset{y_k^2}{sup}N^4[\underset{\rho ,E^{\mathrm{eq}}[\rho ^2|y_k]=1}{sup}[E^{\mathrm{eq}}\left[\frac{\lambda ^2(2k+1)k^4}{N^4\beta }v_{0,k}\rho ^2\right]\beta k^4D_{\nu _{y,k}^c}(\rho )].$$
Moreover, $`v_{i,k}`$ is a bounded function (see Lemma 5.2) and $`E^{\mathrm{eq}}[v_{0,k}|y_k]=0`$. We can apply Lemma 4.2.
$$\underset{N\mathrm{}}{lim\; sup}N^4[\underset{\rho ,E^{\mathrm{eq}}[\rho ^2|y_k]=1}{sup}[E^{\mathrm{eq}}\left[\frac{\lambda ^2(2k+1)k^4}{N^4\beta }v_{0,k}\rho ^2\right]\beta k^4D_{\nu _{y,k}^c}(\rho )]0.$$
The convergence on the line above is uniform over all possible values of the conserved quantities $`y_l=(y_l^1,y_l^2)`$; therefore $`lim\; sup_N\mathrm{}\mathrm{supspec}_{L^2(\nu _N^{\mathrm{eq}})}(\mathrm{\Omega }_2)0.`$
The operator $`\mathrm{\Omega }_3`$. We refer to Lemma 4.1. Rather than proving that
$`lim\; sup_{\delta ,r,l,k,N\mathrm{}}\mathrm{supspec}_{L^2(\nu _N^{\mathrm{eq}})}(\mathrm{\Omega }_3)=0`$ we will show
(4.22)
$$\underset{\delta ,r,l,k,N\mathrm{}}{lim\; sup}\frac{1}{N}\mathrm{log}E^{\mathrm{eq}}\left[\mathrm{exp}^{_0^TN\mathrm{\Omega }_3}\right]ds=0,$$
where
$$\mathrm{\Omega }_3=\frac{1}{\beta N}\underset{i=1}{\overset{N}{}}\varphi \left(\frac{i}{N}\right)^2E^{\mathrm{eq}}[V_l(w\widehat{a}(\overline{x}_l)\delta xL_{\mathrm{}}f_r,y)|y_k]\mathrm{𝟏}_{|y_k|\delta }.$$
The equation (4.22) follows from the estimations
$$\frac{1}{N}\mathrm{log}E^{\mathrm{eq}}\left[\mathrm{exp}\left(N_0^T\mathrm{\Omega }_3𝑑s\right)\right]\frac{1}{N}\mathrm{log}E^{\mathrm{eq}}\left[\frac{1}{T}_0^T\mathrm{exp}(NT\mathrm{\Omega }_3)𝑑s\right]=$$
$$=\frac{1}{N}\mathrm{log}E^{\mathrm{eq}}\left[\mathrm{exp}\left(\underset{i=1}{\overset{N}{}}\varphi \left(\frac{i}{N}\right)^2E^{\mathrm{eq}}[𝕍_l(w\widehat{a}(\overline{x}_l)\mathrm{\Delta }xL_{\mathrm{}}f_r,y)|y_k]\mathrm{𝟏}_{|y_k|\delta }\right)\right]$$
$$\mathrm{log}E_{\nu _0^{gc}}[\mathrm{exp}(M^2E^{\mathrm{eq}}[𝕍_l(w\widehat{a}(\overline{x}_l)\mathrm{\Delta }xL_{\mathrm{}}f_r,y)|y_k]\mathrm{𝟏}_{|y_k|\delta }])].$$
The dominated convergence theorem can be applied (see the note at the end of section 5) and gives us,
$$\underset{\delta 0}{lim}\mathrm{log}E_{\nu _0^{gc}}[\mathrm{exp}(M^2E^{\mathrm{eq}}[𝕍_l(w\widehat{a}(\overline{x}_l)\mathrm{\Delta }xL_{\mathrm{}}f_r,y)|y_k]\mathrm{𝟏}_{|y_k|\delta }])].$$
The operator $`\mathrm{\Omega }_4`$. The main purpose of sections 5 is to show that there exists a sequence of functions $`\{f_r\}_{r0}`$ such that for any $`\delta >0`$.
$$\underset{r,l,k\mathrm{}}{lim\; sup}\underset{|y_k|\delta }{sup}E^{\mathrm{eq}}[𝕍_l(w\widehat{a}(\overline{x}_l)\mathrm{\Delta }xL_{\mathrm{}}f_r,y)|y_k]=0,\mathrm{and}$$
$$\underset{r}{sup}\underset{y_k^2}{sup}E^{\mathrm{eq}}[𝕍_l(w\widehat{a}(\overline{x}_l)\mathrm{\Delta }xL_{\mathrm{}}f_r,y)|y_k]\mathrm{}.$$
The term (4.6) It follows that the term (4.6) converges to zero in probability from the following two lemmas. We refer the reader to Bertini, Olla and Landim or Guo, Papanicolau and Varadhan for the proof. The proof uses mainly the entropy inequality (4.10) and consequence of Feynman-Kac formula (4.12).
###### Lemma 4.4.
(Local ergodicity) Let $`f`$ be a cylinder function. Define $`\stackrel{~}{f}:`$ to be the function $`\stackrel{~}{f}(\alpha )=E_{\nu _\alpha ^{gc}}[f]`$. Then for any $`\delta >0`$, $`\varphi :𝕋`$ a smooth function,
$$\underset{k,N\mathrm{}}{lim\; sup}P_{N,T}^{\mathrm{neq}}\left\{_0^t\right|\frac{1}{N}\underset{i=1}{\overset{N}{}}\varphi \left(\frac{i}{N}\right)(\tau ^if(x(s))\stackrel{~}{f}(\mathrm{Av}_{j=ik}^{i+k}x_j(s)))|ds\delta \}=0.$$
###### Lemma 4.5.
(Two-Block estimate) For any continuous function $`g:`$, let
$$F_{k,a,N}=\left\{_0^t\mathrm{Av}_{i=1}^N\mathrm{Av}_{j=iaN}^{i+aN}(g(\mathrm{Av}_{l=ik}^{i+k}x_l(s))g(\mathrm{Av}_{l=jk}^{j+k}x_l(s)))^2𝑑s\delta \right\}.$$
Then for any $`\delta >0`$,
$$\underset{k\mathrm{},a0,N\mathrm{}}{lim\; sup}P_{N,T}^{\mathrm{neq}}\{F_{k,a,N}\}=0.$$
## 5. Computation of the central limit theorem variances
In this section we compute the value of the limit:
$$\underset{l\mathrm{}}{lim\; sup}\underset{|\alpha |\delta }{sup}E_{\nu _\alpha ^{gc}}[𝕍_l(f,y)]$$
for a particular class of cylinder functions $`f`$ to be described later.
We shall need a result that relates canonical and grand canonical measures, known as the equivalence of ensemble. The equivalence of ensemble says that asymptotically the marginal in a fixed box of the canonical measure is the marginal of the grand canonical measure.
###### Lemma 5.1.
(Equivalence of ensemble) Let $`f(x_s,\mathrm{},x_s)`$ be a bounded, local function. Then, for any $`ϵ>0`$, there exist $`N_1`$ and $`\delta _1>0`$ such that
$$|E_{\nu _{y,k}^c}[f]E_{\nu _\alpha ^{gc}}[f]|ϵ$$
as long as $`k>N_1`$, $`l>N_1`$ and $`|y_k^1\alpha |<\delta _1`$, $`y_k^2`$.
We shall need a new notation, namely $`X_0^{}`$ for the adjoint of the vector field $`X_0=_12_0+_1`$ with respect to the inner-product of the Hilbert space $`L^2(ad\nu _N^{\mathrm{eq}})`$. The adjoint is given by the formula,
$$X_0^{}(h)=X_0(ah)(V^{}(x_1)2V^{}(x_0)+V^{}(x_1))ah.$$
Note that the current $`w`$ is equal to $`X_0^{}(2)`$, and the slope Laplacian $`\mathrm{\Delta }x`$ is equal to $`X_0^{}(\frac{1}{a})`$.
###### Lemma 5.2.
Let $`h(x_s,\mathrm{},x_s)`$ be a bounded, cylinder function such that $`f=X_0^{}(h)L^2(\nu _{\alpha ,s}^{gc})`$ for all $`\alpha `$. Then there exists $`C(h)<\mathrm{}`$ that depends just on the function $`h`$ such that
(5.1)
$$\underset{l,y_{0,l}^2}{sup}\frac{1}{l}<(L_l)^1(\underset{|j|l\sqrt{l}}{}\tau ^jf),\underset{|i|l\sqrt{l}}{}\tau ^if>_{\nu _{y,l}^c}C(h).$$
In addition if the function $`h`$ is bounded away from zero, $`hC>0`$, we have the lower bound, $`C_1(h)>0`$ that depends just on $`h`$,
(5.2)
$$C_1(h)\underset{l,y_{0,l}^2}{sup}\frac{1}{l}<(L_l)^1(\underset{|j|l\sqrt{l}}{}\tau ^jf),\underset{|i|l\sqrt{l}}{}\tau ^if>_{\nu _{y,l}^c}.$$
Proof. Call $`𝕍_l(f,y)=\frac{1}{l}<(L_l)^1(_{|j|l\sqrt{l}}\tau ^jf),_{|i|l\sqrt{l}}\tau ^if>_{\nu _{y,l}^c}`$. We use the variational formula
(5.3)
$$𝕍_l(f,y)=\underset{u}{sup}\frac{<u,_{|j|l\sqrt{l}}\tau ^jf>_{\nu _{y,l}^c}^2}{lD_{\nu _{y,l}^c}(u)}.$$
We have
$$<u,\underset{|j|l\sqrt{l}}{}\tau ^jf>_{\nu _{y,l}^c}^2=<u,\underset{|j|l\sqrt{l}}{}X_j^{}(\tau ^jh)>_{\nu _{y,l}^c}^2=E_{\nu _{y,l}^c}\left[\underset{|j|l\sqrt{l}}{}a_jX_j(u)\tau ^jh\right]^2$$
$$D_{\nu _{y,l}^c}(u)E_{\nu _{y,l}^c}\left[\underset{|j|l\sqrt{l}}{}a_j(\tau ^jh)^2\right]2lC(h)D_{\nu _{y,l}^c}(u),$$
therefore the upper bound for the variance is established.
For the lower bound, we may chose a particular function $`u`$ in (5.3), namely $`u=_{i=l}^l\frac{i^2}{2}x_i`$. Note that $`X_i(u)=1`$ for any $`l+1il1`$. It follows,
$$\frac{<u,_{|j|l\sqrt{l}}\tau ^jf>_{\nu _{y,l}^c}^2}{lD_{\nu _{y,l}^c}(u)}=\frac{E_{\nu _{y,l}^c}\left[_{|j|l\sqrt{l}}a_j\tau ^jh\right]^2}{lE_{\nu _{y,l}^c}[_{|j|l1}a_j]}C_1(h).$$
We have just used that $`a`$, and $`h`$ are functions bounded away from zero, and $`a`$ is a bounded function.
Note. Lemma 5.2 is true for any local function $`f=_{i=s}^sX_j^{}(h_j)`$, where the functions $`h_j`$ are bounded, local functions.
###### Definition 5.1.
For a bounded, local functions $`f`$ such that $`E_{\nu _{y,l}^c}[f]=0`$, for all possible values of $`y^2`$, we define the semi-norm:
$$f_\alpha ^2=\underset{l\mathrm{},y_k^1\alpha }{lim\; sup}\frac{1}{l}E_{\nu _{y,k}^c}[<(L_l)^1(\underset{|i|l\sqrt{l}}{}\tau ^if),\underset{|j|l\sqrt{l}}{}\tau ^jf>_{\nu _{y,l}^c}].$$
We saw in Lemma 5.2 that $`f_\alpha `$ is a finite number as long as $`f`$ is equal to $`X_0^{}(h)`$, where $`h`$ is a bounded, local function. By polarization we can extend the semi-norm $`_\alpha `$ to a semi-inner product $`,_\alpha `$. For the remaining part of this section we compute the value of the semi-norm $`f_\alpha `$ for certain functions $`f`$.
###### Lemma 5.3.
Assume that $`g`$ is a bounded cylinder function with bounded first derivatives, $`f=L_{\mathrm{}}g`$, $`w=X_0(a)(V^{}(x_1)2V^{}(x_0)+V^{}(x_1))a`$ and $`\mathrm{\Delta }x=x_12x_0+x_1`$, then the following identities hold,
* $`L_{\mathrm{}}g_\alpha ^2=E_{\nu _\alpha ^{gc}}\left[a(x_1,x_0,x_1)\left(X_0\left(_j\tau ^jg\right)\right)^2\right],`$
* $`w_\alpha ^2=4E_{\nu _\alpha ^{gc}}[a(x_1,x_0,x_1)],`$
* $`L_{\mathrm{}}g,w_\alpha =2E_{\nu _\alpha ^{gc}}[a(x_1,x_0,x_1)X_0\left(_j\tau ^jg\right)],`$
* $`L_{\mathrm{}}g,\mathrm{\Delta }x_\alpha =0,`$
* $`w,\mathrm{\Delta }x_\alpha =4.`$
Proof. One checks directly using equivalence of ensemble lemma 5.1 and the asymptotic shift invariance of $`\nu _{y,k}^c`$ that the relations a)-e) hold. In particular, for d) and e), it is important to notice that $`\mathrm{\Delta }x=x_12x_0+x_1=X_0^{}(\frac{1}{a})`$.
###### Definition 5.2.
Define the Hilbert space $`_\alpha `$ to be the closed linear span in $`L^2(ad\nu _\alpha ^{gc})`$ of the function $`1`$ and functions $`\xi _g=X_0(_j\tau ^jg)`$, where $`g`$ is a bounded local function with bounded first derivatives.
It is not hard to see that if $`f`$ is equal to $`X_0^{}(h)`$ then,
$$f,L_{\mathrm{}}g_\alpha =E_{\nu _\alpha ^{gc}}\left[aProj__\alpha (2h)X_0\left(\underset{j}{}\tau ^jg\right)\right],f,w_\alpha =E_{\nu _\alpha ^{gc}}[aProj__\alpha (2h)2].$$
On both lines above $`Proj__\alpha `$ stands for the projection opertor in the subspace $`_\alpha `$. We are left to calculate $`f_\alpha `$ for a function $`f=X_0^{}(h)`$. As we will show in the following lemma, $`f_\alpha =E_{\nu _\alpha ^gc}[a(Proj__\alpha (2h))^2]`$.
###### Lemma 5.4.
Suppose $`f`$ is equal to $`X_0^{}(h)`$, where $`h`$ is a bounded cylinder function. The value of the semi-norm of $`f`$,
$$f_\alpha =E_{\nu _\alpha ^{gc}}[a(Proj__\alpha (2h))^2].$$
Proof. Using Cauchy-Schwartz inequality it follows that,
$$f_\alpha =\underset{l\mathrm{},y_k^1\alpha }{lim\; inf}\frac{1}{l}E_{\nu _{y,k}^c}[<(L_l)^1(\underset{|i|l\sqrt{l}}{}\tau ^if),\underset{|j|l\sqrt{l}}{}\tau ^jf>_{\nu _{y,l}^c}]$$
$$E_{\nu _\alpha ^{gc}}[a(Proj__\alpha (2h))^2].$$
Consider $`g=_{|j|l\sqrt{l}}\tau ^jf`$, where $`f=X_0^{}(h)`$ then
$$\frac{1}{l}<(L_l)^1(\underset{|i|l\sqrt{l}}{}\tau ^if),\underset{|j|l\sqrt{l}}{}\tau ^jf>_{\nu _{y,l}^c}=\underset{\rho ,D_{\nu _{y,l}^c}(\rho )}{sup}\frac{<_{|j|l\sqrt{l}}\tau ^jf>_{\nu _{y,l}^c}^2}{l}=$$
$$=\frac{<\rho _l,_{|j|l\sqrt{l}}\tau ^jf>_{\nu _{y,l}^c}^2}{2l^2}=\frac{<\rho _l,_{|j|l\sqrt{l}}X_j^{}(\tau ^jh)>_{\nu _{y,l}^c}^2}{2l^2}=$$
$$=\frac{1}{2l^2}E_{\nu _{y,l}^c}\left[\underset{|j|l\sqrt{l}}{}a(x_j,x_j,x_j)X_j(\rho _l)\tau ^jh\right]^2.$$
Above $`\rho _l`$ is some maximizer function with $`D_{\nu _{y,l}^c}(\rho _l)=2l`$.
Call $`u_l=\frac{1}{2l}_{|j|l\sqrt{l}}\tau ^j(X_j\rho _l)`$.
$$E_{\nu _\alpha ^{gc}}[a(u_l)^2]E_{\nu _\alpha ^{gc}}\left[\frac{a}{2l}\underset{|j|l\sqrt{l}}{}(\tau ^j(X_j\rho _l))^2\right]=\frac{1}{l}D_{\nu _\alpha ^{gc}}(\rho _l)2.$$
The above inequality shows that $`\{u_l\}_l`$ is a bounded sequence in $`L^2(ad\nu _\alpha ^{gc})`$ and hence has a weakly convergent subsequence in $`L^2(ad\nu _\alpha ^{gc})`$. Let $`u`$ be the weak limit of a convergent subsequence of $`\{u_l\}_l`$. It follows that
$$\underset{l\mathrm{},y_k^1\alpha }{lim\; sup}\frac{1}{l}E_{\nu _{y,k}^c}[<(L_l)^1(\underset{|i|l\sqrt{l}}{}\tau ^if),\underset{|j|l\sqrt{l}}{}\tau ^jf>_{\nu _{y,l}^c}]$$
$$2E_{\nu _\alpha ^{gc}}[auh]^2=2E_{\nu _\alpha ^{gc}}[auProj__\alpha h]^2$$
$$2E_{\nu _\alpha ^{gc}}[au^2]E_{\nu _\alpha ^{gc}}[a(Proj__\alpha h)^2]E_{\nu _\alpha ^{gc}}[a(Proj__\alpha 2h)^2].$$
The key point that has allowed us to write the above inequalities is that the function $`u`$ has the property $`X_a(\tau ^bu)=X_b(\tau ^au)`$ for all integers $`a`$ and $`b`$, and a function with this property belongs to $`_\alpha `$, (see Lemma 2.1 or Savu , ). The next estimate shows that $`X_a(\tau ^bu)=X_b(\tau ^au)`$ is valid in a weak sense. Consider a smooth test function $`\varphi `$, with bounded first derivatives. Assume $`a>b`$. We have
$$<X_a(\tau ^bu_l),\varphi >_{\nu _\alpha ^{gc}}=<\frac{1}{2l}\underset{|j|l\sqrt{l}}{}\tau ^{bj}(\rho _l),X_b^{}X_a^{}\varphi >_{\nu _\alpha ^{gc}},$$
and
$`<X_a(\tau ^bu_l)X_b(\tau ^au_l),\varphi >_{\nu _\alpha ^{gc}}=`$
$`=<{\displaystyle \frac{1}{2l}}{\displaystyle \underset{j=b+1+l\sqrt{l}}{\overset{a+l\sqrt{l}}{}}}\tau ^j(\rho _l){\displaystyle \frac{1}{2l}}{\displaystyle \underset{j=bl+\sqrt{l}}{\overset{a1l+\sqrt{l}}{}}}\tau ^j(\rho _l),X_b^{}X_a^{}\varphi >_{\nu _\alpha ^{gc}}=`$
$`=<{\displaystyle \frac{1}{2l}}{\displaystyle \underset{j=b+1+l\sqrt{l}}{\overset{a+l\sqrt{l}}{}}}\tau ^jX_{bj}(\rho _l)+{\displaystyle \frac{1}{2l}}{\displaystyle \underset{j=bl+\sqrt{l}}{\overset{a1l+\sqrt{l}}{}}}\tau ^jX_{bj}(\rho _l),X_a^{}\varphi >_{\nu _\alpha ^{gc}}.`$
By Cauchy-Schwartz inequality, we obtain
$`<X_a(\tau ^bu_l)X_b(\tau ^au_l),\varphi >_{\nu _\alpha ^{gc}}^2`$
$`E_{\nu _\alpha ^{gc}}\left[\left({\displaystyle \frac{1}{2l}}{\displaystyle \underset{j=b+1+l\sqrt{l}}{\overset{a+l\sqrt{l}}{}}}\tau ^jX_{bj}(\rho _l)+{\displaystyle \frac{1}{2l}}{\displaystyle \underset{j=bl+\sqrt{l}}{\overset{a1l+\sqrt{l}}{}}}\tau ^jX_{bj}(\rho _l)\right)^2\right]E_{\nu _\alpha ^{gc}}[(X_a^{}\varphi )^2]`$
$`{\displaystyle \frac{C(ab)}{(2l)^2}}\left({\displaystyle \underset{j=bal+\sqrt{l}}{\overset{1l+\sqrt{l}}{}}}E_{\nu _\alpha ^{gc}}[(X_j\rho _l)^2]+{\displaystyle \underset{j=ba+1+l\sqrt{l}}{\overset{l\sqrt{l}}{}}}E_{\nu _\alpha ^{gc}}[(X_j\rho _l)^2]\right)`$
(5.4) $`{\displaystyle \frac{C(ab)}{(2l)^2}}D_{\nu _\alpha ^{gc}}(\rho _l){\displaystyle \frac{C(ab)}{2l}}.`$
As $`l`$ converges to infinity, the sequence $`\{u_l\}_l`$ approaches $`u`$ in the weak sense. Combining this fact with (5) we can establish $`X_a(\tau ^bu)=X_b(\tau ^au)`$.
Now we can conclude that if we have a local function $`f`$ such that $`f=X_0^{}(h)`$ then
$$f,f_\alpha =\underset{l\mathrm{},y_k^1\alpha }{lim}\frac{1}{l}E_{\nu _{y,k}^c}[<(L_l)^1(\underset{|i|l\sqrt{l}}{}\tau ^if),\underset{|j|l\sqrt{l}}{}\tau ^jf>_{\nu _{y,l}^c}]=$$
$$=E_{\nu _\alpha ^{gc}}[a(Proj__\alpha (2h))^2].$$
The results proved so far in this section allow us to conclude that if $`g`$ is a local function equal to $`X_0^{}(h)`$ and $`b(y_k^1)`$ is a coefficient that depends on the mean-slope in a box of size $`k`$ then
$$g+bw+L_{\mathrm{}}f,g+bw+L_{\mathrm{}}f_\alpha =\underset{l\mathrm{},y_k^1\alpha }{lim}E_{\nu _{y,k}^c}[𝕍_l(g+bw+L_{\mathrm{}}f,y)]=$$
$$=E_{a\nu _\alpha ^{gc}}[(Proj__\alpha (2h)+2b+X_0(\underset{i}{}\tau ^if))^2].$$
Since $`Proj__\alpha (2h)`$ is a function in the closed linear span of the functions $`1`$ and $`X_0(_i\tau ^if)`$ where $`f`$ is a local function it follows the existence of a sequence of functions $`\{f_r\}_{r>0}`$ and of a coefficient $`b=\frac{1}{\widehat{a}}`$ such that
$$\underset{r\mathrm{}}{lim}\mathrm{\Delta }x+bw+L_{\mathrm{}}f_r_\alpha =0.$$
###### Lemma 5.5.
Let $`a^{}>0`$ be an upper bound for the function $`a`$. For any $`ϵ>0`$ and $`\delta >0`$ there exists a smooth function $`f(x_l,\mathrm{},x_l,\alpha )`$ such that
$$\underset{|\alpha |\delta }{sup}wL_{\mathrm{}}f\widehat{a}(\alpha )\mathrm{\Delta }x_\alpha <ϵ,\mathrm{and}$$
$$\underset{\alpha }{sup}wL_{\mathrm{}}f\widehat{a}(\alpha )\mathrm{\Delta }x_\alpha <2a^{}.$$
Proof. We calculate
$$w+L_{\mathrm{}}f\widehat{a}(\alpha )\mathrm{\Delta }x_\alpha =4\left[E_{\nu _\alpha ^{\mathrm{gc}}}[a(1+X_0(\underset{i}{}\tau ^if))^2]\widehat{a}(\alpha )\right].$$
Let us introduce the notation $`_{l,B}`$ for the set of cylinder functions $`f(x_l,\mathrm{},x_l)`$ with $`f_{\mathrm{}}B`$ and $`_if_{\mathrm{}}B`$ for $`i=l,\mathrm{},l`$. Let $`A(f,\alpha )=E_{\nu _\alpha ^{\mathrm{gc}}}[a(1+X_0(_i\tau ^if))^2]\widehat{a}(\alpha )`$ and $`A_{l,B}(\alpha )=inf_{f_{l,B}}A(f,\alpha )`$. The function $`A_{l,B}(\alpha )`$ is upper semicontinuous and nonincreasing in $`l`$ and $`B`$, and for each $`\alpha `$, $`lim_{l,B\mathrm{}}A_{l,B}(\alpha )=0`$, therefore,
$$\underset{l,B\mathrm{}}{lim}\underset{|\alpha |\delta +1}{sup}A_{l,B}(\alpha )=0.$$
Therefore we can find for each $`\alpha `$ a cylinder function $`f(\alpha )`$ such that $`A(f(\alpha ),\alpha )ϵ`$, if $`|\alpha |\delta +1`$. We extend $`f`$ to be zero on $`|\alpha |>\delta +1`$. Then on $`|\alpha |>\delta +1`$ we have that $`A(f(\alpha ),\alpha )=E_{\nu _\alpha ^{\mathrm{gc}}}[a(x_1,x_0,x_1)]\widehat{a}(\alpha )`$ and $`A(f(\alpha ),\alpha )a^{}`$. The avoid the problem that $`f`$ is not smooth, we take the convolution of $`f`$ with a smoothing kernel $`\varphi `$. The required function is the convolution $`f\varphi `$. For a complete argument see Lemma 2.6 in Quastel .
Note. From the equivalence of ensemble lemma 5.1, we know that for any bounded cylinder function $`f`$ and any $`ϵ>0`$, there exit $`N_1`$ and $`\delta _1>0`$ such that
$$|E^{\mathrm{eq}}[𝕍_l(w\widehat{a}(\overline{x}_{i,l})\mathrm{\Delta }xL_{\mathrm{}}f,y)|y_k]w\widehat{a}(\alpha )\mathrm{\Delta }xL_{\mathrm{}}f_\alpha |ϵ,$$
as long as $`k>N_1`$, $`l>N_1`$, and $`|y_k^1\alpha |<\delta _1`$. Then Lemma 5.5 helps us to conclude that for any $`ϵ>0`$ there exists a bounded cylinder function $`f`$ such that
$$\underset{l,k\mathrm{}}{lim\; sup}\underset{|y_k|\delta }{sup}E^{\mathrm{eq}}[𝕍_l(w\widehat{a}(\overline{x}_l)\mathrm{\Delta }xL_{\mathrm{}}f,y)|y_k]<2ϵ,\mathrm{and}$$
$$\underset{y_k^2}{sup}E^{\mathrm{eq}}[𝕍_l(w\widehat{a}(\overline{x}_l)\mathrm{\Delta }xL_{\mathrm{}}f_r,y)|y_k]3a^{}.$$
$`a^{}`$ on the line above is the upper bound for the function $`a`$ in the hypothesis of Lemma 5.5.
Properties of the transport coefficient. The transport coefficient $`\widehat{a}`$ has been defined as the unique real number such that
(5.5)
$$\underset{f}{inf}w\widehat{a}(\alpha )\mathrm{\Delta }xL_{\mathrm{}}f_\alpha ^2=0.$$
It is important to notice that the transport coefficient is also given by the formula:
$$\widehat{a}(\alpha )=\frac{w,\mathrm{\Delta }x_\alpha }{\mathrm{\Delta }x,\mathrm{\Delta }x_\alpha }=\frac{E_{\nu _\alpha ^{gc}}[2\frac{2}{a}a]}{E_{\nu _\alpha ^{gc}}[(Proj__\alpha (\frac{2}{a}))^2a]}=\frac{1}{E_{\nu _\alpha ^{gc}}[(Proj__\alpha (\frac{1}{a}))^2a]}=$$
$$=\frac{1}{sup_{\gamma ,g}E_{\nu _\alpha ^{gc}}[(Proj_{\gamma +X_0(_j\tau ^jg)}\frac{1}{a})^2a]}=\underset{\gamma ,g}{inf}\frac{E_{\nu _\alpha ^{gc}}[(\gamma +X_0(_j\tau ^jg))^2a]}{E_{\nu _\alpha ^{gc}}[\frac{1}{a}(\gamma +X_0(_j\tau ^jg))a]^2}=$$
(5.6)
$$=\underset{g}{inf}E_{\nu _\alpha ^{gc}}\left[a(x_1,x_0,x_1)(1+X_0(\underset{j}{}\tau ^jg))^2\right].$$
## 6. Hydrodynamic limit of the model for surface electromigration
In this section we will show that the model for surface electromigration 2.8 has a hydrodynamic scaling limit as well. We will prove Theorem 2.2. The proof is inspired by Quastel . From Cameron-Martin-Girsanov formula we can find the Radon-Nikodym derivative
$`{\displaystyle \frac{dP_{N,E,T}^{\mathrm{neq}}}{dP_{N,T}^{\mathrm{eq}}}}`$ $`=`$ $`{\displaystyle \frac{d\nu _N^{\mathrm{neq}}}{d\nu _N^{\mathrm{eq}}}}\mathrm{exp}({\displaystyle \frac{1}{2}}{\displaystyle \underset{i=1}{\overset{N}{}}}{\displaystyle _0^T}E(t,{\displaystyle \frac{i}{N}})\sqrt{a(x_{i1},x_i,x_{i+1})}dB_i`$
$`{\displaystyle \frac{1}{8}}{\displaystyle \underset{i=1}{\overset{N}{}}}{\displaystyle _0^T}E^2(t,{\displaystyle \frac{i}{N}})a(x_{i1},x_i,x_{i+1})dt).`$
and the relative entropy
$$H(P_{N,E,T}^{\mathrm{neq}}|P_{N,T}^{\mathrm{eq}})=H(\nu _N^{\mathrm{neq}}|\nu _N^{\mathrm{eq}})+E_{P_{N,E,T}^{\mathrm{neq}}}\left[\frac{1}{8}\underset{i=1}{\overset{N}{}}_0^TE^2(t,\frac{i}{N})a(x_{i1},x_i,x_{i+1})𝑑t\right].$$
Hence there exists a constant $`C`$ such that $`H(P_{N,E,T}^{\mathrm{neq}}|P_{N,T}^{\mathrm{eq}})CN`$. We can use the inequality (4.11) and the superexponential estimates (3.3), (3.4) to conclude that the sequence of probability measures $`\{Q_{N,E,T}^{\mathrm{neq}}\}_{N0}`$ is tight in the topology of the space $`𝒳`$. Denote by $`Q_{E,T}`$ a weak limit of a subsequence of $`\{Q_{N,E,T}^{\mathrm{neq}}\}_{N0}`$.
We proceed further to identify the limit $`Q_{E,T}`$. To save space we ignore the time dependence of the test function $`\varphi `$.
$`\frac{1}{N}_{i=1}^N\varphi (\frac{i}{N})x_i(T)\frac{1}{N}_{i=1}^N\varphi (\frac{i}{N})x_i(0)=`$
$`={\displaystyle _0^T}{\displaystyle \frac{1}{2N}}{\displaystyle \underset{i=1}{\overset{N}{}}}\left(N^2w_i+E(t,{\displaystyle \frac{i}{N}})a_i\right)\varphi ^{\prime \prime }({\displaystyle \frac{i}{N}})dt+M_N(T)=`$
$`={\displaystyle _0^T}{\displaystyle \frac{N^2}{2N}}{\displaystyle \underset{i=1}{\overset{N}{}}}(w_i\widehat{a}(\overline{x}_{i,k})(\mathrm{\Delta }x)_iL_{\mathrm{}}\tau ^if_r)\varphi ^{\prime \prime }({\displaystyle \frac{i}{N}})dt+M_N(T)`$
(6.1) $`+{\displaystyle _0^T}{\displaystyle \frac{N^2}{2N}}{\displaystyle \underset{i=1}{\overset{N}{}}}L_{N,E}\tau ^if_r\varphi ^{\prime \prime }({\displaystyle \frac{i}{N}})dt+`$
(6.2) $`+{\displaystyle _0^T}{\displaystyle \frac{1}{2N}}{\displaystyle \underset{i=1}{\overset{N}{}}}\left[E(t,{\displaystyle \frac{i}{N}})(a_i\widehat{a}(\overline{x}_{i,k}))+N^2(L_{\mathrm{}}\tau ^if_rL_{N,E}\tau ^if_r)\right]\varphi ^{\prime \prime }({\displaystyle \frac{i}{N}})dt`$
$`+{\displaystyle _0^T}{\displaystyle \frac{1}{2N}}{\displaystyle \underset{i=1}{\overset{N}{}}}\widehat{a}(\overline{x}_{i,k})\left[N^2(\mathrm{\Delta }x)_i+E(t,{\displaystyle \frac{i}{N}})\right]\varphi ^{\prime \prime }({\displaystyle \frac{i}{N}})dt`$
Above $`\{M_N(t)\}_{t0}`$ is a martingale and as in section 4 we can prove that is negligible in the limit. Because of the entropy inequality (4.10) and the superexponential estimate
$`\underset{r,l,N\mathrm{}}{lim\; sup}{\displaystyle \frac{1}{N}}\mathrm{log}E^{\mathrm{eq}}[\mathrm{exp}\left(\right|{\displaystyle _0^T}N^2{\displaystyle \underset{i=1}{\overset{N}{}}}[\mathrm{Av}_{j=il_1}^{i+l_1}w_j\widehat{a}(\overline{x}_{i,l})\mathrm{Av}_{j=il_1}^{i+l_1}(\mathrm{\Delta }x)_j`$
$`\mathrm{Av}_{j=il_1}^{i+l_1}\tau ^jL_{\mathrm{}}f_r]ds\left|\right)]0,`$
the event on line (4) is negligible under $`P_{N,E,T}^{\mathrm{neq}}`$. For the new model that we have the negligible fluctuations are $`L_{N,E}f`$ and not $`L_{\mathrm{}}f`$ (i.e., the term (6.1) has no contribution towards the limit of the model). The contribution coming from nontrivial fluctuations $`L_{N,E}fL_{\mathrm{}}f`$ are gathered in the coefficient in front of the vector field $`E`$ in the nonlinear equation (2.14). Because of this reason, the coefficient turns out to be the transport coefficient $`\widehat{a}`$, defined by the variational formula (2.11) and not $`\stackrel{~}{a}(\alpha )=E_{\nu _\alpha ^{gc}}[a]`$
We shall show that the sequence of smooth local functions $`\{f_r\}_{r0}`$ introduced in the Note right after the Lemma 5.5 can be chosen to have the additional property that the term (6.2) is negligible in the limit. Note that the term (6.2) is negligible if $`lim\; sup_{r,k,N\mathrm{}}P_{N,E,T}^{\mathrm{neq}}\{|F(T)|>ϵ\}=0`$ where,
(6.3)
$$F(T)=_0^T\frac{1}{2N}\underset{i=1}{\overset{N}{}}E(t,\frac{i}{N})\left[a_ig_i\widehat{a}(\overline{x}_{i,k})\right]dt,g_i=1+X_i\left(\underset{j}{}\tau ^jf_r\right)$$
We can write
(6.6) $`F(T)`$ $`=`$ $`{\displaystyle _0^T}{\displaystyle \frac{1}{2N}}{\displaystyle \underset{i=1}{\overset{N}{}}}E(t,{\displaystyle \frac{i}{N}})(a_ig_iE_{\nu _{\overline{y}_{i,k}^1}^{gc}}[a_ig_i])+`$
$`+{\displaystyle _0^T}{\displaystyle \frac{1}{2N}}{\displaystyle \underset{i=1}{\overset{N}{}}}E(t,{\displaystyle \frac{i}{N}})(E_{\nu _{\overline{y}_{i,k}^1}^{gc}}[a_ig_i]\widehat{a}(\overline{x}_{i,k}))\mathrm{𝟏}_{|\overline{y}_{i,k}^1|\delta }+`$
$`+{\displaystyle _0^T}{\displaystyle \frac{1}{2N}}{\displaystyle \underset{i=1}{\overset{N}{}}}E(t,{\displaystyle \frac{i}{N}})(E_{\nu _{\overline{y}_{i,k}^1}^{gc}}[a_ig_i]\widehat{a}(\overline{x}_{i,k}))\mathrm{𝟏}_{|\overline{y}_{i,k}^1|\delta }`$
Recall that in the proof of Lemma 5.5 we introduced $`A(f,\alpha )`$ to be the difference $`E_{\nu _\alpha ^{\mathrm{gc}}}[a(1+X_0(_i\tau ^if))^2]\widehat{a}(\alpha )`$. We will show that we can modify the sequence $`\{f_r\}_{r0}`$ to have the properties
(6.7)
$$\underset{\alpha ,r}{sup}A(f_r(\alpha ),\alpha )a^{},\underset{\alpha ,r}{sup}|E_{\nu _\alpha ^{gc}}[ag_0]\widehat{a}(\alpha )|5a^{},$$
(6.8)
$$\underset{r\mathrm{}}{lim}\underset{|\alpha |\delta }{sup}A(f_r(\alpha ),\alpha )=0,\mathrm{and}$$
(6.9)
$$\underset{r\mathrm{}}{lim}\underset{|\alpha |\delta }{sup}|E_{\nu _\alpha ^{gc}}[ag_0]\widehat{a}(\alpha )|=0.$$
We write,
$$E_{\nu _\alpha ^{gc}}[ag_0]\widehat{a}(\alpha )=(E_{\nu _\alpha ^{gc}}[ag_0]E_{\nu _\alpha ^{gc}}[ag_0^2])+(E_{\nu _\alpha ^{gc}}[ag_0^2]\widehat{a}(\alpha ))=$$
$$=E_{\nu _\alpha ^{gc}}[ag_0(g_01)]+A(f_r(\alpha ),\alpha ).$$
Since $`|E_{\nu _\alpha ^{gc}}[ag_0(g_01)]|4a^{}`$ and $`A(f_r(\alpha ),\alpha )a^{}`$ uniform in $`\alpha `$ and $`r`$, (6.7) follows. Consider $`f^{}(\alpha )`$ to be the minimizer of $`A(f,\alpha )`$ among local functions of $`x_l`$ through $`x_l`$. $`f^{}`$ has the property that for any $`f_r`$,
$$E_{\nu _\alpha ^{gc}}\left[a\left(1+X_0(\underset{j}{}\tau ^jf^{})\right)X_0(\underset{j}{}\tau ^jf_r)\right]=0.$$
Now, thanks to Cauchy-Schwartz inequality,
$$|E_{\nu _\alpha ^{gc}}[ag_0(g_01)]|=\left|E_{\nu _\alpha ^{gc}}\left[aX_0(\underset{j}{}\tau ^jf_r)X_0(\underset{j}{}\tau ^j(f^{}f_r))\right]\right|$$
$$E_{\nu _\alpha ^{gc}}[a(X_0(\underset{j}{}\tau ^jf_r))^2]^{1/2}E_{\nu _\alpha ^{gc}}[a(X_0(\underset{j}{}\tau ^j(f^{}f_r)))^2]^{1/2}.$$
We modify $`f_r`$ such that
$$\underset{|\alpha |\delta }{sup}E_{\nu _\alpha ^{gc}}\left[a\left(X_0(\underset{j}{}\tau ^j(f^{}f_r))\right)^2\right]\frac{ϵ^2}{6a^{}}.$$
To conclude (6.7) is negligible because of local ergodicity lemma 4.4, (6.8) is negligible because for any fixed $`\delta `$ we have (6.9) and (6.9) shrinks to zero once $`\delta `$ becomes very large. Also it is important to have property (6.7). Our result follows.
## Acknowledgement
I would like to thank Professor Jeremy Quastel for proposing me this problem and for introducing me to interacting particle systems.
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# Six-Dimensional Gauge Theory on the Chiral Square
## 1 Introduction
Six-dimensional (6D) gauge theories have intriguing properties that make them potentially relevant for extensions of the standard model of particle physics . For example, the global $`SU(2)_W`$ gauge anomaly cancels only in the case where the number of quark and lepton generations is a multiple of three . Furthermore, simple compactifications of two dimensions preserve a discrete symmetry which is a subgroup of the 6D Lorentz group, such that the neutrino masses are forced to be of the Dirac type and proton decay is adequately suppressed even when baryon number is maximally violated at the TeV scale .
As with any theory that has fermions in more than four dimensions, a major constraint imposed by the observed properties of quarks and leptons is that the compactification allows the existence of chiral fermions in the effective 4D theory obtained by integration over the extra-dimensional coordinates. The simplest compactification of two extra dimensions, namely on a torus, does not satisfy this constraint. A toroidal compactification, which can be viewed as a parallelogram with the opposite sides being identified, leads to 4D fermions that are vector-like with respect to any gauge symmetry. The next simplest compactification, a parallelogram with the adjacent sides being identified, automatically leaves at most a single 4D fermion of definite chirality as the zero mode of any chiral 6D fermion . Identifying adjacent sides requires these to have the same length $`L`$, so that the parallelogram has to be a rhombus in this case. For simplicity, we consider the most symmetric compactification of this type: a square. As pointed out in , this configuration is naturally preferred by radion and moduli stabilization mechanisms, a simple example of which is the stabilization by quantum, Casimir-like effects. We refer to this compactification as the “chiral square”.
In this paper we study 6D gauge theories compactified on the chiral square. There are various questions that we address: what are the boundary conditions for gauge fields on the chiral square? what kind of gauge fixing conditions may be imposed in the 6D theory? what is the spectrum of Kaluza-Klein (KK) modes for gauge fields? how does the Higgs mechanism work if the heavy KK gauge fields acquire masses both from compactification and from spontaneous breaking via a vacuum expectation value (VEV)? what are the interactions of the KK modes? Related studies of gauge theories in six dimensions have appeared in .
Before tackling gauge fields, let us recapitulate some properties of the chiral square derived in Ref. . Figure 1 shows a chiral square, with adjacent sides identified. This space has the topology of a sphere, but has a flat metric except for conical singularities at $`(0,0)`$ and $`(L,L)`$ (deficit angle of $`3\pi /2`$ each) and a third one at the identified points $`(0,L)(L,0)`$ (deficit angle of $`\pi `$).
From this point of view there is nothing special going on at the sides that are being glued together. It is nevertheless useful to formulate this compactification by considering the above square region on the $`x^4x^5`$ plane, together with boundary conditions on the edges of the square that encode the identification of adjacent sides. In particular, field values should be equal at identified points, modulo possible symmetries of the theory. The physics at identified points is identical if the Lagrangians take the same values for any field configuration:
$`(x^\mu ,y,0)=(x^\mu ,0,y),`$
$`(x^\mu ,y,L)=(x^\mu ,L,y),`$ (1)
for any $`y[0,L]`$. For a free field $`\mathrm{\Phi }`$, this requirement is consistent with
$$\mathrm{\Phi }(x^\mu ,y,0)=e^{i\theta }\mathrm{\Phi }(x^\mu ,0,y),$$
(2)
for an arbitrary phase $`\theta `$, provided one requires a “smoothness” condition on the derivatives normal to the “edges” of the square
$`_5\mathrm{\Phi }|_{(x^4,x^5)=(y,0)}=e^{i\theta }_4\mathrm{\Phi }|_{(x^4,x^5)=(0,y)}.`$ (3)
Similar relations should be imposed at $`(L,y)(y,L)`$. However, it was found in that this system admits nontrivial solutions only when $`\theta `$ takes one of the four discrete values $`n\pi /2`$, for $`n=0,1,2,3`$. Only those fields that satisfy boundary conditions corresponding to $`n=0`$ admit a zero-mode. Furthermore, when considering 6D Weyl fermions, one finds that their 4D left- and right-handed chiralities obey boundary conditions corresponding to integers that differ by one: $`n_Ln_R=\pm 1`$, where the sign depends on the 6D chirality. Hence, fermions propagating on this space lead necessarily to a chiral low-energy theory: at most one of the left- or right-handed chiralities has a zero mode. This compactification is equivalent to the $`T^2/Z_4`$ orbifold .
The chiral square possesses a discrete $`Z_8`$ symmetry that acts on the right- and left-handed components of 6D Weyl spinors as
$`\mathrm{\Psi }_{\pm R}(x^\mu ,x^4,x^5)`$ $``$ $`e^{i(n_R^\pm \pm 1/2)\pi /2}\mathrm{\Psi }_{\pm R}(x^\mu ,x^4,x^5),`$
$`\mathrm{\Psi }_{\pm L}(x^\mu ,x^4,x^5)`$ $``$ $`e^{i(n_L^\pm 1/2)\pi /2}\mathrm{\Psi }_{\pm L}(x^\mu ,x^4,x^5),`$ (4)
where $`+`$ or $``$ label the 6D chirality, and $`n_L^\pm `$, $`n_R^\pm `$ label the boundary conditions satisfied by $`\mathrm{\Psi }_{\pm L}`$ and $`\mathrm{\Psi }_{\pm R}`$, respectively. Note that each KK tower transforms as a single entity under the $`Z_8`$, i.e. the symmetry commutes with KK number. This $`Z_8`$ symmetry is at the heart of the Dirac nature of neutrinos and the suppression of baryon number violation.
A KK-parity can also be naturally imposed on these scenarios and ensures that the lightest KK particle (LKP) is a viable dark matter candidate . This $`Z_2`$ symmetry distinguishes among KK modes and acts as
$$\mathrm{\Phi }^{(j,k)}(x^\mu )(1)^{j+k}\mathrm{\Phi }^{(j,k)}(x^\mu ),$$
(5)
where $`\mathrm{\Phi }`$ stands for a field of any spin, and $`j,k`$ are integers labeling the KK level. The KK-parity has a geometrical interpretation as a rotation by $`\pi `$ about the center of the chiral square. In particular, KK-parity requires that localized operators at $`(0,0)`$ and $`(L,L)`$ be identical. Localized operators at $`(0,L)(L,0)`$ have coefficients that are, in general, unrelated to those on the previous two conical singularities.
In section 2 we give the appropriate boundary conditions for gauge fields propagating on the chiral square. We concentrate on those boundary conditions that preserve a zero mode, i.e. we do not study the breaking of gauge symmetries by boundary conditions. Next we turn to deriving the self-interactions of the KK modes in the mass eigenstate basis for non-Abelian gauge fields (section 3), gauge interactions of fermions (section 4) and scalars (section 5). In section 6 we analyze spontaneously broken 6D gauge symmetries. We also address in detail the gauge fixing procedure, with and without breaking by nonzero VEV’s, discuss the associated ghost action and corresponding KK decomposition, isolate the linear combinations of scalars that provide the longitudinal polarizations for massive gauge fields, and identify the additional scalars coming from the extra dimensional components of the gauge fields. We summarize and conclude in section 7.
## 2 Abelian gauge fields
Let us first study a 6D Abelian gauge field, $`A^\alpha (x^\beta )`$ with $`\alpha ,\beta =0,1,\mathrm{}5`$, whose propagation in the $`x^4,x^5`$ plane is restricted to a square of size $`L`$ ($`x^\nu `$, $`\nu =0,1,2,3`$ are the Minkowski spacetime coordinates). The action has the usual quadratic form in the gauge field strength, $`F^{\alpha \beta }`$,
$$S=d^4x_0^L𝑑x^4_0^L𝑑x^5\left(\frac{1}{4}F_{\alpha \beta }F^{\alpha \beta }+_{GF}\right),$$
(1)
and includes a gauge fixing term, $`_{GF}`$, which we choose such that the mixings of $`A_\mu `$, $`\mu =0,1,2,3`$, with $`A_4`$ and $`A_5`$ vanish:
$$_{GF}=\frac{1}{2\xi }\left[\frac{}{}_\mu A^\mu \xi \left(_4A_4+_5A_5\right)\right]^2,$$
(2)
where $`\xi `$ is the gauge fixing parameter.
The action may also include localized kinetic terms at the fixed points $`(0,0)`$, $`(L,L)`$ and $`(0,L)(L,0)`$. Such terms are generated radiatively, as was explicitly shown in the cases of 5D theories and of 6D theories compactified on the $`T^2/Z_2`$ orbifold , and are phenomenologically important . In this work we assume that they are sufficiently small that they can be taken into account perturbatively. Therefore, Eq. (1) is our starting point for the KK decomposition. We defer a detailed study of localized terms for future work .
### 2.1 Boundary conditions
Integrating by parts, the action (1) can be written as
$`S`$ $`=`$ $`{\displaystyle }d^4x{\displaystyle _0^L}dx^4{\displaystyle _0^L}dx^5\{{\displaystyle \frac{1}{4}}F_{\mu \nu }F^{\mu \nu }{\displaystyle \frac{1}{2\xi }}(_\mu A^\mu )^2+{\displaystyle \frac{1}{2}}[(_4A_\mu )^2+(_5A_\mu )^2]`$ (3)
$`+{\displaystyle \frac{1}{2}}\left[(_\mu A_4)^2+(_\mu A_5)^2\xi (_4A_4+_5A_5)^2(_4A_5_5A_4)^2\right]`$
$`+_4\left[A_4_\mu A^\mu \right]+_5\left[A_5_\mu A^\mu \right]\text{}\}.`$
The last two terms are surface terms that are important in determining the possible boundary conditions.
The variation of $`S`$ with respect to $`A^\alpha `$ must vanish everywhere in the bulk, leading to the following field equations:
$`^\mu F_{\mu \nu }+{\displaystyle \frac{1}{\xi }}_\mu _\nu A^\mu `$ $`=`$ $`\left(_4^2+_5^2\right)A_\nu ,`$
$`\left(_\mu ^\mu \xi _4^2_5^2\right)A_4`$ $`=`$ $`(\xi 1)_4_5A_5,`$ (4)
$`\left(_\mu ^\mu _4^2\xi _5^2\right)A_5`$ $`=`$ $`(\xi 1)_4_5A_4.`$
Furthermore, we require the surface terms in $`\delta S`$ to vanish everywhere on the boundary:
$`{\displaystyle }d^4x\{{\displaystyle _0^L}dx^4[{\displaystyle \frac{}{}}F_{5\mu }\delta A^\mu +F_{45}\delta A_4+(_\mu A^\mu \xi _4A_4\xi _5A_5)\delta A_5]|_{x^5=0}^{x^5=L}`$
$`+{\displaystyle _0^L}dx^5[{\displaystyle \frac{}{}}F_{4\mu }\delta A^\mu F_{45}\delta A_5+(_\mu A^\mu \xi _4A_4\xi _5A_5)\delta A_4]|_{x^4=0}^{x^4=L}\}=0.`$ (5)
This leads to boundary conditions that guarantee a well-defined, self-adjoint problem, which in turn ensures that the differential operators in Eqs. (4) possess a complete set of orthogonal eigenfunctions. Possible localized terms in the original action would give additional contributions to Eq. (5), and thus would correspond to a modification of the boundary conditions. Requiring that the boundary contributions to $`\delta S`$ vanish also guarantees that there is no flow of charges, such as energy or momentum, across the boundary, or that any flow is explicitly associated with localized terms that act as sources for the corresponding charges. As already mentioned, we do not consider localized terms in what follows.
As discussed in Section 1, we consider the case where the points $`(y,0)`$ and $`(0,y)`$ are identified in the sense that the Lagrangians at these points are equal, and likewise $`(y,L)`$ and $`(L,y)`$ are identified, for any $`0yL`$. Given that any matter field $`\mathrm{\Phi }`$ satisfies the boundary conditions (2) and (3), and analogous relations at the boundaries $`x^4=L`$ and $`x^5=L`$, the requirement that the boundary conditions are compatible with the gauge symmetry implies
$`D_\mu \mathrm{\Phi }|_{(x^4,x^5)=(y,0)}=e^{i\theta }D_\mu \mathrm{\Phi }|_{(x^4,x^5)=(0,y)},`$
$`D_4\mathrm{\Phi }|_{(x^4,x^5)=(y,0)}=e^{i\theta }D_5\mathrm{\Phi }|_{(x^4,x^5)=(0,y)},`$
$`D_5\mathrm{\Phi }|_{(x^4,x^5)=(y,0)}=e^{i\theta }D_4\mathrm{\Phi }|_{(x^4,x^5)=(0,y)}.`$ (6)
The first and second equations are derived by differentiating Eq. (2) with respect to $`x^\mu `$ and $`y`$, respectively, and then replacing partial derivatives by covariant ones,
$$D_\alpha =_\alpha ig_6A_\alpha ,$$
(7)
where the 6D gauge coupling, $`g_6`$, has mass dimension $`1`$. The last equation in (6) is obtained directly from Eq. (3).
Eq. (6) implies that the boundary conditions are invariant under 6D gauge transformations only if $`A_\mu `$, $`A_4`$ and $`A_5`$ satisfy the “folding” identifications
$`A_\mu (y,0)`$ $`=`$ $`A_\mu (0,y),`$
$`A_4(y,0)`$ $`=`$ $`A_5(0,y),`$ (8)
$`A_5(y,0)`$ $`=`$ $`A_4(0,y),`$
and the same relations between fields at $`(y,L)`$ and $`(L,y)`$. These boundary conditions have also been derived in .
Alternatively, one can understand the sign in the last equation of (8) by recalling that the folding boundary condition (2) is closely related to rotations by $`\pi /2`$ about the origin of the larger square $`L<x^4,x^5<L`$. Under $`(x^4,x^5)(x^5,x^4)`$, the gauge field satisfies the covariant transformation law $`(A_4,A_5)(A_5,A_4)`$, and identifying boundary points that differ by such a rotation leads to Eq. (8).
In the presence of the folding identifications (8), Eq. (5) implies that either the gauge field values are fixed on the boundary ($`\delta A_\alpha =0`$), or else
$`F_{\mu 5}(y,0)`$ $`=`$ $`F_{\mu 4}(0,y),`$
$`F_{45}(y,0)`$ $`=`$ $`F_{45}(0,y),`$ (9)
$`\left[{\displaystyle \frac{}{}}_\mu A^\mu \xi \left(_4A_4+_5A_5\right)\right]_{(x^4,x^5)=(y,0)}`$ $`=`$ $`\left[{\displaystyle \frac{}{}}_\mu A^\mu \xi \left(_4A_4+_5A_5\right)\right]_{(x^4,x^5)=(0,y)}.`$
In the latter case, differentiating Eqs. (8) with respect to $`y`$ and combining with Eq. (9) we find some constraints on the $`x^4`$ and $`x^5`$ derivatives of $`A_\alpha `$ on adjacent sides of the square. For $`A_\mu `$, the full set of boundary conditions reads
$`A_\mu (y,0)`$ $`=`$ $`A_\mu (0,y),`$
$`_4A_\mu |_{(x^4,x^5)=(y,0)}`$ $`=`$ $`_5A_\mu |_{(x^4,x^5)=(0,y)},`$ (10)
$`_5A_\mu |_{(x^4,x^5)=(y,0)}`$ $`=`$ $`_4A_\mu |_{(x^4,x^5)=(0,y)},`$
and analogous relations at $`(y,L)`$ and $`(L,y)`$. The conditions on the $`A_4`$ and $`A_5`$ components of the gauge field are more conveniently expressed in terms of the fields $`A_\pm =A_4\pm iA_5`$:
$`A_\pm (y,0)`$ $`=`$ $`iA_\pm (0,y),`$
$`_4A_\pm |_{(x^4,x^5)=(y,0)}`$ $`=`$ $`i_5A_\pm |_{(x^4,x^5)=(0,y)},`$ (11)
$`_5A_\pm |_{(x^4,x^5)=(y,0)}`$ $`=`$ $`\pm i_4A_\pm |_{(x^4,x^5)=(0,y)}.`$
### 2.2 Kaluza-Klein decomposition
Without loss of generality, we may expand the fields $`A_\mu `$, $`A_\pm `$ in terms of complete sets of functions satisfying the boundary conditions (10) and (11). Using the complete sets of functions satisfying these boundary conditions found in , we may write
$`A_\mu (x^\nu ,x^4,x^5)`$ $`=`$ $`{\displaystyle \frac{1}{L}}\left[A_\mu ^{(0,0)}(x^\nu )+{\displaystyle \underset{j1}{}}{\displaystyle \underset{k0}{}}f_0^{(j,k)}(x^4,x^5)A_\mu ^{(j,k)}(x^\nu )\right],`$
$`A_+(x^\nu ,x^4,x^5)`$ $`=`$ $`{\displaystyle \frac{1}{L}}{\displaystyle \underset{j1}{}}{\displaystyle \underset{k0}{}}f_3^{(j,k)}(x^4,x^5)A_+^{(j,k)}(x^\nu ),`$
$`A_{}(x^\nu ,x^4,x^5)`$ $`=`$ $`{\displaystyle \frac{1}{L}}{\displaystyle \underset{j1}{}}{\displaystyle \underset{k0}{}}f_1^{(j,k)}(x^4,x^5)A_{}^{(j,k)}(x^\nu ),`$ (12)
where $`j`$ and $`k`$ are integers, $`A_\mu ^{(j,k)}`$ and $`A_\pm ^{(j,k)}`$ are canonically normalized KK modes, and the KK functions $`f_n`$ with $`n=0,1,2,3`$ are:
$`f_{0,2}^{(j,k)}(x^4,x^5)={\displaystyle \frac{1}{1+\delta _{j,0}}}\left[\mathrm{cos}\left({\displaystyle \frac{jx^4+kx^5}{R}}\right)\pm \mathrm{cos}\left({\displaystyle \frac{kx^4jx^5}{R}}\right)\right],`$
$`f_{1,3}^{(j,k)}(x^4,x^5)=i\mathrm{sin}\left({\displaystyle \frac{jx^4+kx^5}{R}}\right)\mathrm{sin}\left({\displaystyle \frac{kx^4jx^5}{R}}\right).`$ (13)
These satisfy the two-dimensional Klein-Gordon equation,
$`\left(_4^2+_5^2+M_{j,k}^2\right)f_n^{(j,k)}(x^4,x^5)=0,`$ (14)
where
$$M_{j,k}^2\frac{j^2+k^2}{R^2},$$
(15)
with $`R=L/\pi `$, and are normalized so that
$$\frac{1}{L^2}_0^L𝑑x^4_0^L𝑑x^5\left[f_n^{(j,k)}(x^4,x^5)\right]^{}f_n^{(j^{},k^{})}(x^4,x^5)=\delta _{j,j^{}}\delta _{k,k^{}}.$$
(16)
Note that $`f_1^{(j,k)}=f_3^{(j,k)}`$, so that the explicit minus sign in the expansion of $`A_+`$ shown in Eq. (12) leads to $`A_{}^{(j,k)}=A_+^{(j,k)}`$. Derivatives along $`x^4`$ or $`x^5`$ acting on the KK functions satisfy
$$_\pm f_n^{(j,k)}(x^4,x^5)=ir_{j,\pm k}M_{j,k}f_{n1}^{(j,k)}(x^4,x^5),$$
(17)
where
$$_\pm =_4\pm i_5$$
(18)
and $`r_{j,k}`$ are complex phases that depend only on the KK-numbers:
$$r_{j,k}\frac{j+ik}{\sqrt{j^2+k^2}}.$$
(19)
Before inserting the KK expansions (12) into the action (3), note that
$`_4A_5_5A_4`$ $`=`$ $`{\displaystyle \frac{1}{L}}{\displaystyle \underset{j1}{}}{\displaystyle \underset{k0}{}}M_{j,k}A_H^{(j,k)}(x^\nu )f_0^{(j,k)}(x^4,x^5),`$
$`_4A_4+_5A_5`$ $`=`$ $`{\displaystyle \frac{1}{L}}{\displaystyle \underset{j1}{}}{\displaystyle \underset{k0}{}}M_{j,k}A_G^{(j,k)}(x^\nu )f_0^{(j,k)}(x^4,x^5),`$ (20)
where we defined at each KK level two real scalar fields, $`A_H^{(j,k)}`$ and $`A_G^{(j,k)}`$, by
$$A_\pm ^{(j,k)}=r_{j,\pm k}\left(A_H^{(j,k)}iA_G^{(j,k)}\right).$$
(21)
The explicit factors of $`M_{j,k}`$ in Eqs. (20) ensure that the scalars $`A_H^{(j,k)}`$ and $`A_G^{(j,k)}`$ are canonically normalized. It is clear that $`A_H^{(j,k)}`$ correspond to excitations which are invariant under 6D gauge transformations: $`A_\alpha A_\alpha +_\alpha \chi /g_6`$, where $`\chi `$ is a gauge parameter that, like $`A_\mu `$ in Eq. (12), has an expansion in terms of $`f_0^{(j,k)}`$. The orthogonal excitations, $`A_G^{(j,k)}`$, shift under such a gauge transformation and correspond to the Nambu-Goldstone modes eaten by the massive 4D fields, $`A_\mu ^{(j,k)}`$.
After integrating over $`x^4`$ and $`x^5`$, we find the 4D Lagrangian for the KK modes (one can check that the last two terms in Eq. (3) give no contribution). The gauge boson of level $`(0,0)`$ remains massless, while the gauge bosons of level $`(j,k)`$ with $`j1`$ appear as massive vector particles in an $`R_\xi `$ gauge, with a Lagrangian
$$\frac{1}{4}F_{\mu \nu }^{(j,k)}F^{(j,k)\mu \nu }+\frac{1}{2}M_{j,k}^2\left(A_\mu ^{(j,k)}\right)^2\frac{1}{2\xi }\left(^\mu A_\mu ^{(j,k)}\right)^2,$$
(22)
where $`F_{\mu \nu }^{(j,k)}=_\mu A_\nu ^{(j,k)}_\nu A_\mu ^{(j,k)}`$ is the 4D field strength of level $`(j,k)`$, and for the zero-mode one just sets $`M_{0,0}=0`$. At each $`(j,k)`$ level with $`j1`$ one finds that $`A_H^{(j,k)}`$ and $`A_G^{(j,k)}`$, as defined in Eqs. (20) and (21), are mass eigenstates in the gauge defined by Eq. (2), and are described by the following terms in the 4D Lagrangian:
$$\frac{1}{2}\left[\left(_\mu A_H^{(j,k)}\right)^2M_{j,k}^2\left(A_H^{(j,k)}\right)^2+\left(_\mu A_G^{(j,k)}\right)^2\xi M_{j,k}^2\left(A_G^{(j,k)}\right)^2\right].$$
(23)
One can explicitly check that the field equations (4), when expressed in terms of $`A_H^{(j,k)}`$ and $`A_G^{(j,k)}`$, are satisfied. In the unitary gauge, $`\xi \mathrm{}`$, only the $`A_H^{(j,k)}(x)`$ scalars propagate with masses $`M_{j,k}`$, whereas the fields $`A_G^{(j,k)}(x)`$ are the Nambu-Goldstone bosons eaten by the $`A_\mu ^{(j,k)}(x)`$ KK gauge bosons.
## 3 Non-Abelian gauge fields
The boundary conditions, KK decomposition, and identification of the physical states in the case of non-Abelian gauge fields are analogous to the ones discussed in Section 2 for Abelian fields. In this section we present the self-interactions of the KK modes associated with non-Abelian gauge fields, and then we study the ghost fields required by gauge fixing.
### 3.1 Self-Interactions
The kinetic term of a 6D non-Abelian gauge field, $`A_\alpha ^a`$, where $`a`$ labels the generators of the adjoint representation, is given by
$$\frac{1}{4}F_{\alpha \beta }^aF^{a\alpha \beta }=\frac{1}{4}\left(F_{\mu \nu }^aF^{a\mu \nu }2F_{+\mu }^aF_{}^{a\mu }\right)+\frac{1}{8}\left(F_+^a\right)^2.$$
(1)
The gauge field strengths introduced here are defined by
$`F_{\mu \nu }^a=_\mu A_\nu ^a_\nu A_\mu ^a+g_6f^{abc}A_\mu ^bA_\nu ^c,`$
$`F_{\pm \mu }^a=_\pm A_\mu ^a_\mu A_\pm ^a+g_6f^{abc}A_\pm ^bA_\mu ^c=F_{4\mu }^a\pm iF_{5\mu }^a,`$
$`F_+^a=_+A_{}^a_{}A_+^a+g_6f^{abc}A_+^bA_{}^c=2iF_{45}^a,`$ (2)
where $`f^{abc}`$ are the group structure constants.
The trilinear terms included in Eq. (1) are given by
$$\frac{g_6}{4}f^{abc}\{4A^{a\mu }A_\nu ^b_\mu A^{c\nu }+[A_{}^aA_+^b_+A_{}^c+2A_{}^aA^{b\mu }(_\mu A_+^c_+A_\mu ^c)+\mathrm{H}.\mathrm{c}.]\}.$$
(3)
Upon KK decomposition and integration over the $`x^4`$ and $`x^5`$ coordinates, these give rise in the 4D Lagrangian to trilinear interactions (proportional to the 4D gauge coupling, $`g_4=g_6/L`$) among spin-1 modes,
$$g_4f^{abc}\delta _{0,0,0}^{(j_1k_1)(j_2k_2)(j_3k_3)}A_\mu ^{(j_1,k_1)a}A_\nu ^{(j_2,k_2)b}^\mu A^{(j_3,k_3)c\nu },$$
(4)
as well as trilinear interactions involving one, two or three scalar gauge fields:
$`{\displaystyle \frac{g_4}{2}}f^{abc}\delta _{1,3,0}^{(j_1k_1)(j_2k_2)(j_3k_3)}A_{}^{(j_1,k_1)a}`$ $`[ir_{j_2,k_2}M_{j_2,k_2}A_\mu ^{(j_2,k_2)b}A^{(j_3,k_3)c\mu }\left(^\mu A_+^{(j_2,k_2)b}\right)A_\mu ^{(j_3,k_3)c}`$ (5)
$`+{\displaystyle \frac{i}{2}}r_{j_3,k_3}M_{j_3,k_3}A_+^{(j_2,k_2)b}A_{}^{(j_3,k_3)c}]+\mathrm{H}.\mathrm{c}.`$
Here we used Eq. (17) to express the $`_\pm `$ derivatives in terms of the KK masses.
We have introduced the following notation:
$$\delta _{n_1,\mathrm{},n_r}^{(j_1,k_1)\mathrm{}(j_r,k_r)}\frac{1}{L^2}_0^L𝑑x^4_0^L𝑑x^5f_{n_1}^{(j_1,k_1)}\mathrm{}f_{n_r}^{(j_r,k_r)}.$$
(6)
This integral was computed for $`r=3`$ in Ref. , and the result in the particular cases relevant for trilinear interactions read
$`\delta _{n,\overline{n},0}^{(j_1,k_1)(j_2,k_2)(j_3,k_3)}={\displaystyle \frac{1}{2\left(1+\delta _{j_1,0}\right)\left(1+\delta _{j_2,0}\right)\left(1+\delta _{j_3,0}\right)}}[\text{}7\delta _{j_1,0}\delta _{j_2,0}\delta _{j_3,0}+\delta _{j_1,j_3j_2}\delta _{k_1,k_3k_2}`$
$`+i^n\left(\delta _{j_1,j_3+k_2}\delta _{k_1,k_3j_2}+\delta _{j_1,k_2k_3}\delta _{k_1,j_3j_2}\right)+(i)^n\left(\delta _{j_1,j_3k_2}\delta _{k_1,j_2+k_3}+\delta _{j_1,k_3k_2}\delta _{k_1,j_2j_3}\right)`$
$`+(1)^n(\delta _{j_1,j_2+j_3}\delta _{k_1,k_2+k_3}+\delta _{j_1,j_2j_3}\delta _{k_1,k_2k_3}+\delta _{j_1,j_2k_3}\delta _{k_1,j_3+k_2}+\delta _{j_1,j_2+k_3}\delta _{k_1,k_2j_3})\text{}],`$
(7)
where $`\overline{n}n\mathrm{mod}4`$, and we have used the fact that $`j=0`$ implies $`k=0`$.
The quartic terms included in Eq. (1),
$$\frac{g_6^2}{4}f^{abc}f^{ade}\left(A_\mu ^bA_\nu ^cA^{d\mu }A^{e\nu }2A_\mu ^bA_+^cA^{d\mu }A_{}^e\frac{1}{2}A_+^bA_{}^cA_+^dA_{}^e\right),$$
(8)
lead to the following quartic interactions of the KK modes:
$`{\displaystyle \frac{g_4^2}{4}}f^{abc}f^{ade}`$ $`[\text{}\delta _{0,0,0,0}^{(j_1,k_1)\mathrm{}(j_4,k_4)}A_\mu ^{(j_1,k_1)b}A_\nu ^{(j_2,k_2)c}A^{(j_3,k_3)d\mu }A^{(j_4,k_4)e\nu }`$ (9)
$`+\mathrm{\hspace{0.17em}2}\delta _{3,1,0,0}^{(j_1,k_1)\mathrm{}(j_4,k_4)}A_+^{(j_1,k_1)c}A_{}^{(j_2,k_2)e}A_\mu ^{(j_3,k_3)b}A^{(j_4,k_4)d\mu }`$
$`{\displaystyle \frac{1}{2}}\delta _{3,1,3,1}^{(j_1,k_1)\mathrm{}(j_4,k_4)}A_+^{(j_1,k_1)b}A_{}^{(j_2,k_2)c}A_+^{(j_3,k_3)d}A_{}^{(j_4,k_4)e}].`$
Of particular interest for phenomenolgy are vertices involving at least a zero mode, and for those one can use Eq. (7) because
$$\delta _{n,\overline{n},0,0}^{(j_1,k_1)(j_2,k_2)(j_3,k_3)(0,0)}=\delta _{n,\overline{n},0}^{(j_1,k_1)(j_2,k_2)(j_3,k_3)}.$$
(10)
In order to extract the interactions of the physical states, the $`A_\pm ^{(j,k)a}`$ complex scalars have to be replaced in Eqs. (5) and (9) by the two real scalar fields, $`A_H^{(j,k)a}`$ and $`A_G^{(j,k)a}`$, as prescribed by Eq. (21). These heavy 4D scalars are in the adjoint representation of the non-Abelian gauge group, so we will refer to them as “spinless adjoints”.
### 3.2 Ghost fields
For completeness we now turn to determining the ghost action associated with the gauge fixing term given by the obvious adaptation of Eq. (2) to non-Abelian fields. This arises from inserting in the path integral a factor
$$𝒟\chi \delta [G(A)]\text{det}\left(\frac{\delta G(A^\chi )}{\delta \chi }\right)=1,$$
(11)
where $`\chi `$ is the gauge transformation parameter, $`A^\chi `$ is the transformed gauge field, and the gauge fixing condition is
$$G(A^a)=^\mu A_\mu ^a\xi \left(_4A_4^a+_5A_5^a\right)\omega ^a,$$
(12)
for arbitrary functions $`\omega ^a`$. After integrating with a Gaussian weight over $`\omega ^a`$ one recovers Eq. (2). Since
$$(A_\alpha ^a)^\chi =A_\alpha ^a+\frac{1}{g_6}D_\alpha \chi ^a,$$
(13)
with $`D_\alpha `$ the covariant derivative in the adjoint representation, we find
$$\frac{\delta G(A^\chi )}{\delta \chi }=\frac{1}{g_6}\left[_\mu D^\mu \xi \left(_4D_4+_5D_5\right)\right].$$
(14)
These terms in the Lagrangian may be taken into account by including a ghost term in the 6D action, given by
$$\overline{c}^a\left[_\mu D^\mu \xi \left(_4D_4+_5D_5\right)\right]c^a,$$
(15)
where $`c^a`$ is an anti-commuting 6D scalar in the adjoint representation of the gauge group.
The above procedure did not take into account the compactification of the two extra dimensions and the associated boundary conditions. The free part of the 6D ghost Lagrangian (15) is
$$\overline{c}^a\left[_\mu ^\mu \xi \left(_4^2+_5^2\right)\right]c^a,$$
(16)
which up to the factor $`\xi `$ is the same as for a scalar. It follows that the KK expansion for the ghost fields can be written as
$$c^a(x^\mu ,x^4,x^5)=\frac{1}{L}\underset{(j,k)}{}c^{(j,k)a}(x^\mu )f_{n_c}^{(j,k)}(x^4,x^5),$$
(17)
where the $`f_{n_c}^{(j,k)}`$ belong to one of the sets of KK functions defined in Eqs. (2.2). The KK modes $`c^{(j,k)a}(x^\mu )`$ have mass $`\sqrt{\xi }M_{j,k}`$. It would appear that $`n_c`$ can take any of the integral values $`0,1,2,3`$. However, for $`n_c=1,2,3`$, the ghost fields lack the zero modes necessary to ensure the gauge invariance of amplitudes involving the zero-mode gauge fields. More generally, since $`A_\mu (x^\mu ,x^4,x^5)`$ satisfies boundary conditions with $`n=0`$, the necessary relations between the couplings involving only gauge fields and those involving ghosts and gauge fields can be satisfied only if the ghosts satisfy the same boundary conditions as $`A_\mu `$. Therefore, only the value $`n_c=0`$ is allowed. The zero mode $`c^{(0,0)}`$ is the ghost required by 4D gauge invariance.
After integrating by parts Eq. (15), the interactions between gauge bosons and ghost fields are given by
$$g_6f^{abc}(^\mu \overline{c}^a)A_\mu ^bc^c\frac{1}{2}\xi g_6f^{abc}\left[(_+\overline{c}^a)A_{}^b+(_{}\overline{c}^a)A_+^b\right]c^c,$$
(18)
where $`A_\pm `$ were defined before Eq. (11) and $`_\pm `$ were defined in (18). It is worth pointing out that the boundary conditions for $`c^a`$ given in Eq. (17), together with the boundary conditions for $`A_4`$, $`A_5`$ discussed in section 2, ensure that one can freely integrate by parts Eq. (18) without generating surface terms.
Using the KK expansions (17), as well as Eq. (20), Eq. (18) leads to KK interactions between ghosts and gauge fields:
$$g_4f^{abc}\delta _{0,0,0}^{(j_1,k_1)(j_2,k_2)(j_3,k_3)}\left(^\mu \overline{c}^{(j_1,k_1)a}\right)A_\mu ^{(j_2,k_2)b}c^{(j_3,k_3)c},$$
(19)
and between ghosts and the “spinless adjoints”:
$$\frac{i}{2}g_4f^{abc}\xi r_{j_1,k_1}M_{j_1,k_1}\delta _{3,1,0}^{(j_1,k_1)(j_2,k_2)(j_3,k_3)}\overline{c}^{(j_1,k_1)a}A_{}^{(j_2,k_2)b}c^{(j_3,k_3)c}+\mathrm{H}.\mathrm{c}.,$$
(20)
where we used Eq. (17) to simplify the derivative, and $`A_\pm ^{(j,k)}`$ are given in terms of the mass eigenstates $`A_H^{(j,k)}`$ and $`A_G^{(j,k)}`$ by Eq. (21).
## 4 Gauge interactions of fermions
The 6D chiral fermions have four degrees of freedom, and the Dirac equation can be written using a set of six anti-commuting matrices, $`\mathrm{\Gamma }^\alpha `$. We use here the following $`8\times 8`$ representation:
$$\mathrm{\Gamma }^\mu =\gamma ^\mu \sigma ^0,\mathrm{\Gamma }^{4,5}=i\gamma _5\sigma ^{1,2}$$
(1)
where $`\gamma ^\mu `$, are the 4D $`\gamma `$-matrices, $`\gamma _5=i\gamma ^0\gamma ^1\gamma ^2\gamma ^3`$, $`\sigma ^0`$ is the $`2\times 2`$ unit matrix and $`\sigma ^i`$ are the Pauli matrices. The 6D fermions have $`+`$ or $``$ chirality, defined by the eigenvalue of the $`\overline{\mathrm{\Gamma }}=\gamma _5\sigma ^3`$ chirality operator: $`\overline{\mathrm{\Gamma }}\mathrm{\Psi }_\pm =\pm \mathrm{\Psi }_\pm `$. A 6D chiral fermion includes both 4D chiralities: $`\gamma _5\sigma ^0\mathrm{\Psi }_{\pm _L}=\mathrm{\Psi }_{\pm _L}`$ and $`\gamma _5\sigma ^0\mathrm{\Psi }_{\pm _R}=\mathrm{\Psi }_{\pm _R}`$.
It is useful to define
$$\mathrm{\Gamma }^\pm =\frac{1}{2}\left(\mathrm{\Gamma }^4\pm i\mathrm{\Gamma }^5\right),$$
(2)
and to observe that
$`\mathrm{\Gamma }^4P_LP_\pm =i\mathrm{\Gamma }^5P_LP_\pm ,`$
$`\mathrm{\Gamma }^4P_RP_\pm =\pm i\mathrm{\Gamma }^5P_RP_\pm ,`$ (3)
where $`P_{L,R}`$, $`P_\pm `$ project on the respective 4D and 6D chiralities, from which it follows that
$$\mathrm{\Gamma }^+\mathrm{\Psi }_{+_L}=\mathrm{\Gamma }^{}\mathrm{\Psi }_{+_R}=\mathrm{\Gamma }^{}\mathrm{\Psi }__L=\mathrm{\Gamma }^+\mathrm{\Psi }__R=0.$$
(4)
The fermion kinetic term can then be written in terms of 4D chiral components as follows:
$$i\overline{\mathrm{\Psi }}_\pm \mathrm{\Gamma }^\alpha D_\alpha \mathrm{\Psi }_\pm =i\left(\overline{\mathrm{\Psi }}_{\pm _L}\mathrm{\Gamma }^\mu D_\mu \mathrm{\Psi }_{\pm _L}+\overline{\mathrm{\Psi }}_{\pm _R}\mathrm{\Gamma }^\mu D_\mu \mathrm{\Psi }_{\pm _R}+\overline{\mathrm{\Psi }}_{\pm _L}\mathrm{\Gamma }^\pm D_{}\mathrm{\Psi }_{\pm _R}+\overline{\mathrm{\Psi }}_{\pm _R}\mathrm{\Gamma }^{}D_\pm \mathrm{\Psi }_{\pm _L}\right).$$
(5)
Here $`D_\pm =D_4\pm iD_5`$ and $`D_\alpha `$ is the covariant derivative defined in Eq. (7) with either $`A_\alpha `$ being an Abelian gauge field, or $`A_\alpha =T^aA_\alpha ^a`$ where $`A_\alpha ^a`$ is a non-Abelian gauge field and $`T^a`$ are the gauge group generators corresponding to the representation of $`\mathrm{\Psi }_\pm `$.
The KK expansions of a 6D fermion of + chirality with a left-handed zero mode component are ,
$`\mathrm{\Psi }_{+_L}`$ $`=`$ $`{\displaystyle \frac{1}{L}}\left[\mathrm{\Psi }_{+_L}^{(0,0)}(x^\nu )+{\displaystyle \underset{j1}{}}{\displaystyle \underset{k0}{}}f_0^{(j,k)}(x^4,x^5)\mathrm{\Psi }_{+_L}^{(j,k)}(x^\nu )\right]\left(\begin{array}{c}1\\ 0\end{array}\right),`$ (8)
$`\mathrm{\Psi }_{+_R}`$ $`=`$ $`{\displaystyle \frac{i}{L}}{\displaystyle \underset{j1}{}}{\displaystyle \underset{k0}{}}r_{j,k}f_3^{(j,k)}(x^4,x^5)\mathrm{\Psi }_{+_R}^{(j,k)}(x^\nu )\left(\begin{array}{c}0\\ 1\end{array}\right),`$ (11)
while those containing a right-handed zero mode are,
$`\mathrm{\Psi }_{+_L}`$ $`=`$ $`{\displaystyle \frac{i}{L}}{\displaystyle \underset{j1}{}}{\displaystyle \underset{k0}{}}r_{j,k}^{}f_1^{(j,k)}(x^4,x^5)\mathrm{\Psi }_{+_L}^{(j,k)}(x^\nu )\left(\begin{array}{c}1\\ 0\end{array}\right),`$ (14)
$`\mathrm{\Psi }_{+_R}`$ $`=`$ $`{\displaystyle \frac{1}{L}}\left[\mathrm{\Psi }_{+_R}^{(0,0)}(x^\nu )+{\displaystyle \underset{j1}{}}{\displaystyle \underset{k0}{}}f_0^{(j,k)}(x^4,x^5)\mathrm{\Psi }_{+_R}^{(j,k)}(x^\nu )\right]\left(\begin{array}{c}0\\ 1\end{array}\right).`$ (17)
The phases of the KK functions for left- and right-handed fermions are correlated, but one can choose an overall phase such that the functions $`f_n(x^4,x^5)`$ with $`n=0,1,3`$ are given in Eq. (2.2). For a fermion of $``$ chirality, the same equations are valid except for an interchange of the $`L`$ and $`R`$ indices.
After integrating the 6D fermion kinetic term shown in Eq. (5) over $`x^4`$ and $`x^5`$, the 4D Lagrangian includes the usual kinetic terms for all KK modes, mass terms of the type $`M_{j,k}\overline{\mathrm{\Psi }}_{\pm _L}^{(j,k)}\mathrm{\Psi }_{\pm _R}^{(j,k)}+\mathrm{H}.\mathrm{c}.`$, and interactions among KK modes. The latter ones, in the case of a fermion with $`+`$ chirality having a left-handed zero mode, include interactions of the $`\mathrm{\Psi }_{+L}^{(j,k)}`$ fermions with a spin-1 KK mode,
$$g_4\delta _{0,0,0}^{(j_1,k_1)(j_2,k_2)(j_3,k_3)}\overline{\mathrm{\Psi }}_{+L}^{(j_1,k_1)}A_\mu ^{(j_2,k_2)}\gamma ^\mu \mathrm{\Psi }_{+L}^{(j_3,k_3)},$$
(18)
interactions of the $`\mathrm{\Psi }_{+R}^{(j,k)}`$ fermions with a spin-1 KK mode,
$$g_4\delta _{1,0,3}^{(j_1,k_1)(j_2,k_2)(j_3,k_3)}r_{j_1,k_1}^{}r_{j_3,k_3}\overline{\mathrm{\Psi }}_{+R}^{(j_1,k_1)}A_\mu ^{(j_2,k_2)}\gamma ^\mu \mathrm{\Psi }_{+R}^{(j_3,k_3)},$$
(19)
and gauge interactions of the fermions with the spinless adjoints,
$$ig_4\delta _{0,1,3}^{(j_1,k_1)(j_2,k_2)(j_3,k_3)}r_{j_2,k_2}^{}r_{j_3,k_3}\overline{\mathrm{\Psi }}_{+L}^{(j_1,k_1)}\left(A_H^{(j_2,k_2)}+iA_G^{(j_2,k_2)}\right)\mathrm{\Psi }_{+R}^{(j_3,k_3)}+\mathrm{H}.\mathrm{c}.$$
(20)
The $`\delta _{n,\overline{n},0}^{(j_1k_1)(j_2,k_2)(j_3k_3)}`$ coefficients with $`n+\overline{n}=0`$ are given in Eq. (7), the complex numbers $`r_{j,k}`$ are given in Eq. (19), and $`g_4=g_6/L`$ is the 4D gauge coupling.
In the case of a fermion with $`+`$ chirality having a right-handed zero mode, the spin-1 KK modes have interactions with $`\mathrm{\Psi }_{+R}^{(j,k)}`$ given by
$$g_4\delta _{0,0,0}^{(j_1,k_1)(j_2,k_2)(j_3,k_3)}\overline{\mathrm{\Psi }}_{+R}^{(j_1,k_1)}A_\mu ^{(j_2,k_2)}\gamma ^\mu \mathrm{\Psi }_{+R}^{(j_3,k_3)},$$
(21)
and interactions with $`\mathrm{\Psi }_{+L}^{(j,k)}`$ given by
$$g_4\delta _{3,0,1}^{(j_1,k_1)(j_2,k_2)(j_3,k_3)}r_{j_1,k_1}r_{j_3,k_3}^{}\overline{\mathrm{\Psi }}_{+L}^{(j_1,k_1)}A_\mu ^{(j_2,k_2)}\gamma ^\mu \mathrm{\Psi }_{+L}^{(j_3,k_3)},$$
(22)
while the gauge interactions of the spinless adjoints with the fermions are given by
$$ig_4\delta _{0,3,1}^{(j_1,k_1)(j_2,k_2)(j_3,k_3)}r_{j_2,k_2}r_{j_3,k_3}^{}\overline{\mathrm{\Psi }}_{+R}^{(j_1,k_1)}\left(A_H^{(j_2,k_2)}iA_G^{(j_2,k_2)}\right)\mathrm{\Psi }_{+L}^{(j_3,k_3)}+\mathrm{H}.\mathrm{c}.$$
(23)
A 6D fermion of $``$ chirality has the same gauge interactions as $`\mathrm{\Psi }_+`$ except for an interchange of the 4D chiralities. More explicitly, if $`\mathrm{\Psi }_R`$ has a zero mode, then the gauge interactions of $`\mathrm{\Psi }_R`$ and $`\mathrm{\Psi }_L`$ are as in Eqs. (18)-(20) with an interchange of the $`L`$ and $`R`$ indices. If $`\mathrm{\Psi }_L`$ has a zero mode, then $`\mathrm{\Psi }_R`$ and $`\mathrm{\Psi }_L`$ have gauge interactions given by Eqs. (21)-(23) with $`L`$ and $`R`$ interchanged.
## 5 Gauge interactions of scalars
Consider now a 6D scalar field $`\mathrm{\Phi }`$ transforming under a certain nontrivial representation of a gauge symmetry, with an action given by
$$S_\mathrm{\Phi }=d^4x_0^L𝑑x^4_0^L𝑑x^5\left[\left(D_\alpha \mathrm{\Phi }\right)^{}D^\alpha \mathrm{\Phi }M_\mathrm{\Phi }^2\mathrm{\Phi }^{}\mathrm{\Phi }\frac{\lambda _6}{2}\left(\mathrm{\Phi }^{}\mathrm{\Phi }\right)^2\right],$$
(1)
where $`\lambda _6`$ is a parameter of mass dimension $`2`$, and $`D_\alpha `$ is the covariant derivative associated with a gauge field $`A_\alpha ^a`$, as in Eq. (7). As in the previous section, we use the short-hand notation $`A_\alpha =T^aA_\alpha ^a`$ where $`T^a`$ are the gauge group generators corresponding to the representation of $`\mathrm{\Phi }`$.
The KK decomposition of the scalar has been derived in :
$$\mathrm{\Phi }(x^\mu ,x^4,x^5)=\frac{1}{L}\underset{(j,k)}{}\mathrm{\Phi }^{(j,k)}(x^\mu )f_n^{(j,k)}(x^4,x^5),$$
(2)
with $`n=0,1,2`$, or $`3`$. The scalar KK modes $`\mathrm{\Phi }^{(j,k)}`$ have masses
$$M_\mathrm{\Phi }^{(j,k)}=\sqrt{M_\mathrm{\Phi }^2+M_{j,k}^2},$$
(3)
where $`M_{j,k}`$ are the KK masses given in Eq. (15).
Using the KK decomposition of gauge fields given in Eq. (12), and integrating over the extra dimensional coordinates, we find that the 4D Lagrangian includes interactions of two KK scalars with one spin-1 KK field,
$$ig_4\delta _{\overline{n},0,n}^{(j_1,k_1)(j_2,k_2)(j_3,k_3)}\mathrm{\Phi }^{(j_1,k_1)}A_\mu ^{(j_2,k_2)}^\mu \mathrm{\Phi }^{(j_3,k_3)}+\mathrm{H}.\mathrm{c}.,$$
(4)
as well as interactions of two KK scalars with two spin-1 KK fields,
$$g_4^2\delta _{\overline{n},0,0,n}^{(j_1,k_1)\mathrm{}(j_4,k_4)}\mathrm{\Phi }^{(j_1,k_1)}A_\mu ^{(j_2,k_2)}A^{(j_3,k_3)\mu }\mathrm{\Phi }^{(j_4,k_4)}.$$
(5)
In particular, the interactions of the $`A_\mu ^{(0,0)}`$ fields are dictated by 4D gauge invariance.
The 4D Lagrangian also includes interactions of two $`\mathrm{\Phi }^{(j,k)}`$ scalars with one of the $`A_H^{(j,k)}`$ and $`A_G^{(j,k)}`$ spinless adjoints,
$`{\displaystyle \frac{g_4}{2}}M_{j_3,k_3}\mathrm{\Phi }^{(j_1,k_1)}[\text{}(\delta _{\overline{n},1,n1}^{(j_1,k_1)(j_2,k_2)(j_3,k_3)}r_{j_2,k_2}^{}r_{j_3,k_3}\delta _{\overline{n},3,n+1}^{(j_1,k_1)(j_2,k_2)(j_3,k_3)}r_{j_2,k_2}r_{j_3,k_3}^{})A_H^{(j_2,k_2)}`$ (6)
$`+i(\delta _{\overline{n},3,n+1}^{(j_1,k_1)(j_2,k_2)(j_3,k_3)}r_{j_2,k_2}r_{j_3,k_3}^{}+\delta _{\overline{n},1,n1}^{(j_1,k_1)(j_2,k_2)(j_3,k_3)}r_{j_2,k_2}^{}r_{j_3,k_3})A_G^{(j_2,k_2)}]\mathrm{\Phi }^{(j_3,k_3)}+\mathrm{H}.\mathrm{c}.,`$
\[here we replaced the derivatives using Eq. (2.17)\], interactions of two $`\mathrm{\Phi }^{(j,k)}`$ scalars with two of the spinless adjoints,
$$g_4^2r_{j_2,k_2}r_{j_3,k_3}^{}\delta _{\overline{n},3,1,n}^{(j_1,k_1)\mathrm{}(j_4,k_4)}\mathrm{\Phi }^{(j_1,k_1)}\left(A_H^{(j_2,k_2)}iA_G^{(j_2,k_2)}\right)\left(A_H^{(j_3,k_3)}+iA_G^{(j_3,k_3)}\right)\mathrm{\Phi }^{(j_4,k_4)},$$
(7)
and, finally, self-interactions of the $`\mathrm{\Phi }^{(j,k)}`$ scalars
$$\frac{\lambda _4}{2}\delta _{\overline{n},n,\overline{n},n}^{(j_1,k_1)\mathrm{}(j_4,k_4)}\mathrm{\Phi }^{(j_1,k_1)}\mathrm{\Phi }^{(j_2,k_2)}\mathrm{\Phi }^{(j_3,k_3)}\mathrm{\Phi }^{(j_4,k_4)},$$
(8)
where $`\lambda _4=\lambda _6/L^2`$ is the 4D quartic coupling.
## 6 Spontaneous symmetry breaking
We now discuss the case where the gauge symmetry is broken by the VEV of a 6D scalar $`\mathrm{\Phi }`$. The action is given in Eq. (1), with $`M_\mathrm{\Phi }^2<0`$. By adding an irrelevant constant, the terms in the 6D Lagrangian that include $`\mathrm{\Phi }`$ can be rewritten as
$$_\mathrm{\Phi }=\left|D_\alpha \mathrm{\Phi }\right|^2\frac{\lambda _6}{2}\left(\mathrm{\Phi }^{}\mathrm{\Phi }\frac{1}{2}v_6^2\right)^2,$$
(1)
where $`v_6>0`$ is the 6D VEV, which has mass dimension two and is defined by the relation
$$M_\mathrm{\Phi }^2=\frac{1}{2}\lambda _6v_6^2.$$
(2)
Let us focus on the case where the gauge group is $`SU(N)`$ and the scalar $`\mathrm{\Phi }`$ is in the fundamental representation of the gauge group. The formulae we derive below also apply (with trivial modifications) to $`SO(N)`$ or $`U(1)`$ gauge groups. In order to analyze the spectrum and interactions in the presence of the 6D VEV, we parameterize $`\mathrm{\Phi }`$ as
$$\mathrm{\Phi }=\frac{1}{\sqrt{2}}(v_6+h)U_\eta \widehat{\mathrm{\Phi }}_0,$$
(3)
where $`\widehat{\mathrm{\Phi }}_0=(0,\mathrm{},0,1)`$ defines the direction of the VEV, $`h`$ is the single 6D real scalar that is orthogonal to the Nambu-Goldstone modes, and $`U_\eta `$ is an unitary matrix that depends on the 6D Nambu-Goldstone bosons, $`\eta ^a`$:
$$U_\eta =e^{i\eta /v_6},$$
(4)
where $`\eta =\eta ^aX^a`$, with $`X^a`$ the broken generators. The sum over $`a`$ includes the $`2N1`$ generators of the coset $`SU(N)/SU(N1)`$. In terms of $`h`$ and $`\eta `$, the Lagrangian (1) is given by
$$_\mathrm{\Phi }=\frac{1}{2}(_\alpha h)^2\frac{\lambda _6}{8}h^2\left(2v_6+h\right)^2+\frac{1}{2}(v_6+h)^2\widehat{\mathrm{\Phi }}_0^{}\left|U_\eta ^{}_\alpha U_\eta ig_6U_\eta ^{}A_\alpha U_\eta \right|^2\widehat{\mathrm{\Phi }}_0.$$
(5)
where $`A_\alpha =A_\alpha ^aT^a`$, with $`a`$ running over all group generators $`T^a`$.
### 6.1 Physical states
Expanding $`U_\eta `$ in powers of $`\eta `$, we find the terms in $`_\mathrm{\Phi }`$ that are quadratic in $`\eta `$ or $`A_\alpha `$:
$$\frac{1}{2}v_6^2\widehat{\mathrm{\Phi }}_0^{}\left(\frac{1}{v_6}_\alpha \eta g_6A_\alpha \right)^2\widehat{\mathrm{\Phi }}_0.$$
(6)
After integration by parts, these become
$$\frac{1}{2}\widehat{\mathrm{\Phi }}_0^{}\left[(_\alpha \eta )^2+g_6v_6\{\eta ,_\alpha A^\alpha \}+g_6^2v_6^2A_\alpha A^\alpha \right]\widehat{\mathrm{\Phi }}_0,$$
(7)
where $`\{\mathrm{},\mathrm{}\}`$ is the anticommutator. We see that $`\eta ^a`$ mixes with the broken components of both $`A_\mu ^a`$ and $`_4A_4^a+_5A_5^a`$, where the latter includes the KK modes $`A_G^{(j,k)a}`$ that are eaten by the spin-1 modes in the limit of zero VEV \[see Eq. (20)\]. In order to make the physics more transparent, we modify the gauge fixing term (2) to
$$_{GF}^v=\frac{1}{2\xi }\left[\frac{}{}_\mu A^{\mu a}\xi \left(_4A_4^a+_5A_5^a\right)+\xi g_6v_6\eta ^a\right]^2,$$
(8)
where it is understood that $`\eta ^a`$ is non-vanishing only along the direction of the broken generators. The terms involving the unbroken components of $`A_\mu ^a`$ remain as in Eq. (2). Since $`\widehat{\mathrm{\Phi }}_0^{}\{\eta ,_\alpha A^\alpha \}\widehat{\mathrm{\Phi }}_0=2\eta ^a_\alpha A^{\alpha a}`$, where $`a`$ runs over the broken generators, the crossed terms in Eq. (8) involving $`\eta ^a`$ and $`A_\mu ^a`$ cancel the corresponding terms in (7). The remaining terms in Eqs. (7) and (8), involving $`\eta `$, $`A_4`$ and $`A_5`$, are then
$$\frac{1}{2}\widehat{\mathrm{\Phi }}_0^{}\left[(_\alpha \eta )^2g_6^2v_6^2A_+A_{}\xi g_6^2v_6^2\eta ^2+g_6v_6(\xi 1)\{\eta ,_4A_4+_5A_5\}\right]\widehat{\mathrm{\Phi }}_0,$$
(9)
together with the second line in Eq. (3). Here we used $`A_\alpha ^aA^{\alpha a}=A_\mu ^aA^{\mu a}A_+^aA_{}^a`$, where $`A_\pm ^a=A_4\pm iA_5`$.
We turn next to the KK decomposition. This is achieved by using the KK expansions for $`A_\mu ^a`$ and $`A_\pm ^a`$, Eqs. (12), supplemented by
$`h(x^\mu ,x^4,x^5)`$ $`=`$ $`{\displaystyle \frac{1}{L}}{\displaystyle \underset{(j,k)}{}}h^{(j,k)}(x^\mu )f_0^{(j,k)}(x^4,x^5),`$
$`\eta ^a(x^\mu ,x^4,x^5)`$ $`=`$ $`{\displaystyle \frac{1}{L}}{\displaystyle \underset{(j,k)}{}}\eta ^{(j,k)a}(x^\mu )f_0^{(j,k)}(x^4,x^5).`$ (10)
In assuming a constant VEV for $`\mathrm{\Phi }`$ we have implicitly imposed that it satisfy boundary conditions with $`n=0`$, a fact we used in Eqs. (10).
Using Eqs. (20) as well as the orthogonality of the KK functions $`f_n^{(j,k)}`$, it is now straightforward to obtain the new terms in the KK action. One finds new mass contributions proportional to the VEV for the broken components of $`A_\mu ^{(j,k)a}`$ and $`A_H^{(j,k)a}`$ modes, where the latter ones are defined in Eq. (21):
$$\frac{1}{2}g_4^2v_4^2\widehat{\mathrm{\Phi }}_0^{}\left(A_\mu ^{(j,k)}A^{(j,k)\mu }+A_H^{(j,k)}A_H^{(j,k)}\right)\widehat{\mathrm{\Phi }}_0.$$
(11)
These, together with the mass contributions due to momentum in extra dimensions, shown in Eqs. (22) and (23), lead to masses for these modes given by
$$M_A^{(j,k)}=\sqrt{M_{j,k}^2+g_4^2v_4^2},$$
(12)
with $`M_{j,k}`$ defined by Eq. (15). We chose to express these results in terms of the 4D gauge coupling, $`g_4=g_6/L`$, and the 4D VEV, $`v_4=v_6L`$, with mass dimension one. Also, from Eq. (5), the masses of the $`h^{(j,k)}`$ real scalars are given by
$$M_h^{(j,k)}=\sqrt{M_{j,k}^2+\lambda _4v_4^2},$$
(13)
where $`\lambda _4=\lambda _6/L^2`$ is the 4D quartic coupling.
Finally, the mass terms involving $`\eta ^{(j,k)}`$ and $`A_G^{(j,k)}`$ \[see Eqs. (9) and (23)\] mix these modes, such that, for the components along the broken generators, the physical states are given by the linear combinations
$`\stackrel{~}{A}_G^{(j,k)a}`$ $`=`$ $`{\displaystyle \frac{1}{M_A^{(j,k)}}}\left(M_{j,k}A_G^{(j,k)a}+g_4v_4\eta ^{(j,k)a}\right),`$
$`\stackrel{~}{\eta }^{(j,k)a}`$ $`=`$ $`{\displaystyle \frac{1}{M_A^{(j,k)}}}\left(M_{j,k}\eta ^{(j,k)a}g_4v_4A_G^{(j,k)a}\right),`$ (14)
for $`j>0`$ and $`k0`$, while for $`j=k=0`$ we define
$$\stackrel{~}{A}_G^{(0,0)a}=\eta ^{(0,0)a}.$$
(15)
This latter mode is eaten by the would-be zero-mode of the gauge KK tower. The free Lagrangian terms for these states are
$$\frac{1}{2}\widehat{\mathrm{\Phi }}_0^{}\left[\left(_\mu \stackrel{~}{\eta }^{(j,k)}\right)^2\left(M_A^{(j,k)}\stackrel{~}{\eta }^{(j,k)}\right)^2+\left(_\mu \stackrel{~}{A}_G^{(j,k)}\right)^2\xi M_{(j,k)}^2\left(\stackrel{~}{A}_G^{(j,k)}\right)^2\right]\widehat{\mathrm{\Phi }}_0,$$
(16)
while for the unbroken components it is given by
$$\frac{1}{2}\left[\left(_\mu A_G^{(j,k)a}\right)^2\xi M_{j,k}^2\left(A_G^{(j,k)a}\right)^2\right].$$
(17)
Note that under an infinitesimal 6D gauge transformation
$$A_\alpha ^aA_\alpha ^a+\frac{1}{g_6}_\alpha \chi ^a+\mathrm{},\eta ^a\eta ^a+v_6\chi ^a+\mathrm{},$$
(18)
one has $`_4A_4^a+_5A_5^a_4A_4^a+_5A_5^a+(_4^2+_5^2)\chi ^a/g_6+\mathrm{}`$, which after KK expansion, Eq. (20), translates at the linear level into
$$A_G^{(j,k)}A_G^{(j,k)}+\frac{1}{g_4}M_{j,k}\chi ^{(j,k)}+\mathrm{},\eta ^{(j,k)}\eta ^{(j,k)}+v_4\chi ^{(j,k)}+\mathrm{}.$$
(19)
In the above, the dimensionless gauge transformation parameter, $`\chi `$, has an expansion
$$\chi (x^\mu ,x^4,x^5)=\underset{(j,k)}{}\chi ^{(j,k)}(x^\mu )f_0^{(j,k)}(x^4,x^5),$$
(20)
and the $`\mathrm{}`$ stand for higher order terms that, in general, mix the KK levels. It is then clear that the inhomogeneus piece of the gauge transformation drops from $`\stackrel{~}{\eta }^{(j,k)a}`$, while the orthogonal combination, $`\stackrel{~}{A}_G^{(j,k)a}`$, shifts under the gauge transformation. $`\stackrel{~}{A}_G^{(j,k)}`$ is the generalization of the would-be Nambu-Goldstone boson of sections 2 and 3 to the case where the gauge symmetry is broken by a constant scalar VEV. Note that in the unitary gauge, $`\xi \mathrm{}`$, each KK level gauge boson eats a linear combination of $`A_4`$, $`A_5`$ and $`\eta `$.
As a result of the modification arising from $`v_6`$ in the gauge fixing term, Eq. (8), the ghost Lagrangian, Eq. (15), receives a new contribution
$$\xi g_6^2v_6^2\overline{c}^ac^a,$$
(21)
where $`a`$ runs over the broken components only. Hence, in the presence of the Higgs VEV, the KK masses associated with the broken components of the ghost fields are $`\sqrt{\xi }M_A^{(j,k)}`$, where $`M_A^{(j,k)}`$ was defined in Eq. (12), while those associated with the unbroken ones remain at $`\sqrt{\xi }M_{j,k}`$.
### 6.2 Trilinear and quartic couplings
The interactions among KK modes follow from the nonlinear terms in Eq. (5). These can be derived by expanding $`U_\eta `$ in a power series:
$`U_\eta ^{}_\alpha U_\eta `$ $`=`$ $`i_\alpha \widehat{\eta }{\displaystyle \frac{1}{2}}[_\alpha \widehat{\eta },\widehat{\eta }]{\displaystyle \frac{i}{3!}}[[_\alpha \widehat{\eta },\widehat{\eta }],\widehat{\eta }]+\mathrm{},`$
$`U_\eta ^{}A_\alpha U_\eta `$ $`=`$ $`A_\alpha +i[A_\alpha ,\widehat{\eta }]{\displaystyle \frac{1}{2}}[[A_\alpha ,\widehat{\eta }],\widehat{\eta }]{\displaystyle \frac{i}{3!}}[[[A_\alpha ,\widehat{\eta }],\widehat{\eta }],\widehat{\eta }]+\mathrm{},`$ (22)
where $`[\mathrm{},\mathrm{}]`$ is the commutator and $`\widehat{\eta }=\eta /v_6`$.
The interaction terms up to quadratic order in $`\eta `$ and $`A_\alpha `$ that appear in the Lagrangian $`_\mathrm{\Phi }`$ are
$$\frac{\lambda _6}{8}\left(4v_6h^3+h^4\right)+\frac{1}{2}\left(2v_6h+h^2\right)\widehat{\mathrm{\Phi }}_0^{}\left(_\alpha \widehat{\eta }g_6A_\alpha \right)^2\widehat{\mathrm{\Phi }}_0.$$
(23)
Higher order terms in $`\eta `$ and $`A_\alpha `$ include additional trilinear and quartic interactions involving the $`h`$ scalar,
$$\frac{i}{2}v_6\left(v_6+2h\right)\widehat{\mathrm{\Phi }}_0^{}\left\{\frac{1}{2}\{_\alpha \widehat{\eta },[^\alpha \widehat{\eta },\widehat{\eta }]\}+g_6\left[\widehat{\eta }(_\alpha \widehat{\eta })A^\alpha (_\alpha \widehat{\eta })A^\alpha \widehat{\eta }\right]\right\}\widehat{\mathrm{\Phi }}_0,$$
(24)
interactions of three $`\eta `$ fields with an $`A_\alpha `$,
$$\frac{g_6}{2}v_6^2\widehat{\mathrm{\Phi }}_0^{}\left[\text{}(_\alpha \widehat{\eta })[A^\alpha ,\widehat{\eta }]\widehat{\eta }\widehat{\eta }[A^\alpha ,\widehat{\eta }]_\alpha \widehat{\eta }+\widehat{\eta }(_\alpha \widehat{\eta })[\widehat{\eta },A^\alpha ][\widehat{\eta },A^\alpha ](_\alpha \widehat{\eta })\widehat{\eta }\right]\widehat{\mathrm{\Phi }}_0,$$
(25)
quartic interactions among the $`\eta `$ fields,
$$\frac{1}{2}v_6^2\widehat{\mathrm{\Phi }}_0^{}\{\text{}\frac{1}{3!}(\text{}(_\alpha \widehat{\eta })[[_\alpha \widehat{\eta },\widehat{\eta }],\widehat{\eta }]+\mathrm{H}.\mathrm{c}.)+\frac{1}{4}[_\alpha \widehat{\eta },\widehat{\eta }]^2\}\widehat{\mathrm{\Phi }}_0,$$
(26)
and interactions involving five or more fields.
Replacing the KK expansions, Eq. (23) leads to trilinear interactions involving the spin-1 modes
$$v_4\delta _{0,0,0}^{(j_1,k_1)(j_2,k_2)(j_3,k_3)}h^{(j_1,k_1)}\widehat{\mathrm{\Phi }}_0^{}\left(\frac{1}{v_4}_\mu \eta ^{(j_2,k_2)}g_4A_\mu ^{(j_2,k_2)}\right)\left(\frac{1}{v_4}^\mu \eta ^{(j_3,k_3)}g_4A^{(j_3,k_3)\mu }\right)\widehat{\mathrm{\Phi }}_0.$$
(27)
Additional trilinear couplings of a spin-1 mode, coming from Eq. (24), are given by
$$\frac{i}{2}g_4\delta _{0,0,0}^{(j_1,k_1)(j_2,k_2)(j_3,k_3)}\widehat{\mathrm{\Phi }}_0^{}\left[\eta ^{(j_1,k_1)}\left(^\mu \eta ^{(j_2,k_2)}\right)A_\mu ^{(j_3,k_3)}\left(^\mu \eta ^{(j_1,k_1)}\right)A_\mu ^{(j_2,k_2)}\eta ^{(j_3,k_3)}\right]\widehat{\mathrm{\Phi }}_0,$$
(28)
In the above two equations, one should express $`\eta ^{(j,k)}`$ in terms of the mass eigenstates $`\stackrel{~}{\eta }^{(j,k)}`$ and $`\stackrel{~}{A}_G^{(j,k)}`$ by inverting Eqs. (14):
$$\eta ^{(j,k)a}=\frac{1}{M_A^{(j,k)}}\left(M_{j,k}\stackrel{~}{\eta }^{(j,k)a}+g_4v_4\stackrel{~}{A}_G^{(j,k)a}\right),$$
(29)
except for the zero-mode, which is given by Eq. (15). There are also trilinear interactions among scalars which include at least one $`h^{(j,k)}`$ field and no derivatives:
$`v_4h^{(j_1,k_1)}`$ $`\{\text{}{\displaystyle \frac{1}{2}}\lambda _4\delta _{0,0,0}^{(j_1,k_1)(j_2,k_2)(j_3,k_3)}h^{(j_2,k_2)}h^{(j_3,k_3)}+r_{j_2,k_2}r_{j_3,k_3}^{}\delta _{0,3,1}^{(j_1,k_1)(j_2,k_2)(j_3,k_3)}`$ (30)
$`\times \widehat{\mathrm{\Phi }}_0^{}(g_4A_H^{(j_2,k_2)}+i\omega _{j_2,k_2}\stackrel{~}{\eta }^{(j_2,k_2)})(g_4A_H^{(j_3,k_3)}i\omega _{j_3,k_3}\stackrel{~}{\eta }^{(j_3,k_3)})\widehat{\mathrm{\Phi }}_0\text{}\},`$
where we defined the KK-number-dependent, dimensionless ratios
$$\omega _{j,k}=\frac{M_A^{(j,k)}}{v_4}.$$
(31)
The remaining trilinear couplings involve only $`\stackrel{~}{\eta }^{(j,k)}`$, $`A_H^{(j,k)}`$ and $`A_G^{(j,k)}`$ scalars, and can be derived from Eq. (24).
The quartic interactions include terms involving $`h^{(j,k)}`$ and the spin-1 fields,
$$\frac{1}{2}\delta _{0,0,0,0}^{(j_1,k_1)\mathrm{}(j_4,k_4)}h^{(j_1,k_1)}h^{(j_2,k_2)}\widehat{\mathrm{\Phi }}_0^{}\left(\frac{1}{v_4}_\mu \eta ^{(j_3,k_3)}g_4A_\mu ^{(j_3,k_3)}\right)\left(\frac{1}{v_4}^\mu \eta ^{(j_4,k_4)}g_4A^{(j_4,k_4)\mu }\right)\widehat{\mathrm{\Phi }}_0,$$
(32)
as well as the spinless fields $`h^{(j,k)}`$, $`A_H^{(j,k)}`$ and $`\stackrel{~}{\eta }^{(j,k)}`$:
$`h^{(j_1,k_1)}h^{(j_2,k_2)}`$ $`\{\text{}{\displaystyle \frac{1}{8}}\lambda _4\delta _{0,0,0,0}^{(j_1,k_1)\mathrm{}(j_4,k_4)}h^{(j_3,k_3)}h^{(j_4,k_4)}+r_{j_3,k_3}r_{j_4,k_4}^{}\delta _{0,0,3,1}^{(j_1,k_1)\mathrm{}(j_4,k_4)}`$ (33)
$`\times \widehat{\mathrm{\Phi }}_0^{}(g_4A_H^{(j_3,k_3)}+i\omega _{j_3,k_3}\stackrel{~}{\eta }^{(j_3,k_3)})(g_4A_H^{(j_4,k_4)}i\omega _{j_4,k_4}\stackrel{~}{\eta }^{(j_4,k_4)})\widehat{\mathrm{\Phi }}_0\text{}\}.`$
The higher order terms in Eqs. (24), (25) and (26) lead to additional quartic couplings among KK modes, which include at most a single 4D spin-1 field. These couplings are straightforward to derive but have long expressions so that we do not display all of them here.
We end this section by observing that the Abelian case can be recovered from the previous formulae by setting $`\widehat{\mathrm{\Phi }}_0=1`$, and considering a single gauge index $`a`$.
## 7 Conclusions
Gauge theories in six dimensions may be relevant for physics beyond the Standard Model provided the size of two dimensions is below $`10^{16}`$ cm. Theories of this type have been proposed in the past, with compactification scales ranging from the electroweak scale to the GUT scale. This paper presents the first in-depth study of the gauge interactions among KK modes. We have concentrated on the simplest compactification of two dimensions that leads to zero mode fermions of definite 4D chirality, namely the “chiral square” defined in Ref. .
After identifying a gauge fixing procedure appropriate for this compactification, we have determined a set of gauge invariant boundary conditions for gauge fields. The ensuing KK decomposition of a 6D gauge field $`A^\alpha `$, $`\alpha =0,1,\mathrm{},5`$, includes a tower of spin-1 modes, $`A_\mu ^{(j,k)}`$, that have a zero mode ($`j=k=0`$), and two towers of spin-0 modes that have no zero mode, $`A_4^{(j,k)}`$ and $`A_5^{(j,k)}`$, where the pair of KK numbers $`(j,k)`$ take the integer values $`j1`$, $`k0`$. The spin-1 zero mode is associated with the unbroken 4D gauge invariance, while the other spin-1 modes become heavy. Each nonzero spin-1 mode, $`A_\mu ^{(j,k)}`$, has a longitudinal polarization given by a linear combination of $`A_4^{(j,k)}`$ and $`A_5^{(j,k)}`$. The linear combination $`A_G^{(j,k)}`$ of spin-0 modes that play the role of Nambu-Goldstone boson eaten by the heavy spin-1 mode depends on the KK numbers, as shown in Eq. (21). The orthogonal combination, $`A_H^{(j,k)}`$, is gauge invariant and remains as an additional physical degree of freedom. Therefore, unlike the case of one extra dimension, where the extra components of the gauge field can be gauged away, any gauge field in two extra dimensions implies the existence of a tower of heavy spinless particles in the adjoint representation of the gauge group (“spinless adjoints”), whose interactions depend on the KK numbers.
The self-interactions of 6D non-Abelian gauge fields induce the following terms in the 4D Lagrangian involving KK modes: trilinear couplings of spin-1 modes, given in Eq. (4); couplings of two spin-1 modes and one spin-0 mode, and couplings of three spin-0 modes, given in Eq. (5); quartic couplings of spin-1 modes, couplings of two spin-1 modes and two spin-0 modes, and quartic couplings of spin-0 modes, given in Eq. (9); and finally, couplings of one spin-1 mode, or one spinless adjoints, and two modes of the ghost field, given in Eqs. (19) and (20), respectively.
The gauge interactions of a chiral 6D fermion induce couplings of a gauge field mode to fermion modes of both 4D chiralities. These depend on the 4D chirality of the zero-mode fermion. The spin-1 modes couple to fermion modes according to Eqs. (18), (19), (21), and (22), while the Yukawa couplings of the spinless adjoints to the fermion modes are given in Eqs. (20) and (23).
The gauge interactions of a 6D scalar field induce couplings of two scalar modes to one or two spin-1 modes, given in Eq. (4) and (5), and also to one or two spinless adjoints, given in Eqs. (6) and (7). A 6D quartic self-interaction of a scalar induces quartic couplings of the scalar modes, as in Eq. (8).
We have also studied the case where the gauge symmetry is broken by a the VEV of a 6D scalar that has a zero mode. In this case, all gauge KK modes receive a contribution to their mass from the spontaneous breaking. The longitudinal polarizations of the heavy KK modes are now given by a linear combination of $`A_4`$, $`A_5`$ and the scalar that acquires the VEV, as shown in Eq. (14). The longitudinal polarization of the would-be zero-mode gauge field is provided by the zero-mode of the additional spinless adjoint, as given in Eq. (15). We showed how the longitudinal modes and additional scalars can be identified by studying the transformation properties under 6D gauge transformations. We displayed the interactions of these scalars with the spin-1 modes and among themselves in Eqs. (27)-(33).
All the couplings among various KK modes mentioned above are induced at tree level by bulk interactions. Loop corrections generate operators localized at the three conical singularities of the chiral square. Even though these preserve KK parity, they lead to additional couplings among KK modes. These are perturbative corrections, but nevertheless may have important phenomenological consequences as they give rise to mixing among all KK modes belonging to the same tower. This effect is studied in a subsequent paper .
The detailed construction of 6D gauge theories with explicit boundary conditions presented here opens up various theoretical and model building avenues of research, including issues pertaining to symmetry breaking by boundary conditions , relation to other compactifications , the structure of gauge and gravitational anomalies on a compact space , and latticized or deconstructed versions of the chiral square compactification.
Given the constraints on the compactification scale of universal extra dimensions from electroweak measurements , flavor-changing processes or collider searches , are as low as a few hundred GeV, it would be particularly interesting to use the tools developed here for analyzing the phenomenological implications of the Standard Model in two universal extra dimensions compactified on the chiral square. That model is well motivated by the natural happenstances of proton stability, constraint on the number of fermion generations, and existence of a dark matter candidate. More generally, the spectra and interactions derived here are relevant for any extensions of the Standard Model in the context of 6D theories.
Acknowledgements: We have benefited from the hospitality of the Aspen Center for Physics at various stages of this project. G.B. acknowledges the support of the State of São Paulo’s Agency for the Promotion of Research (FAPESP), and the Brazilian National Council for Technological and Scientific Development (CNPq). The work of B.D. was supported by DOE under contract DE-FG02-92ER-40704.
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# Finding Exponential Product Formulas of Higher Orders
## 1 Introduction: Why do we need the exponential product formula?
First of all, we discuss as to why we have to treat the exponential operator and why we need an approximant in order to treat the exponential operator. The exponential operator appears in various fields of physics as a formal solution of the differential equation of the form
$$\frac{}{t}f(t)=f(t),$$
(5)
where $`f`$ is a function or a vector and $``$ is an operator or a matrix. Typical examples are the Schrödinger equation
$$\mathrm{i}\frac{}{t}\psi (x,t)=\psi (x,t)$$
(6)
(we put $`\mathrm{}=1`$ here and hereafter), the Hamilton equation
$$\frac{\mathrm{d}}{\mathrm{d}t}\left(\begin{array}{c}\stackrel{}{p}(t)\\ \stackrel{}{q}(t)\end{array}\right)=\left(\begin{array}{c}\stackrel{}{p}(t)\\ \stackrel{}{q}(t)\end{array}\right),$$
(7)
(see Eq. (14) below) and the diffusion equation with a potential
$$\frac{\mathrm{d}}{\mathrm{d}t}P(x,t)=P(x,t).$$
(8)
A solution of Eq. (5) is given in the form of the Green’s function as
$$f(t)=G(t;0)f(0)=\mathrm{e}^tf(0),$$
(9)
although it is only a formal solution; obtaining the Green’s function $`G(t;0)\mathrm{e}^t`$ is just as difficult as solving the equation (5) in any other way. Another important incident of the exponential operator is the partition function in equilibrium quantum statistical physics:
$$Z=Tr\mathrm{e}^\beta ,$$
(10)
where $``$ is a quantum Hamiltonian.
The exponential operator, however, is hard to compute in many interesting cases. The most straightforward way of computing the exponential operator $`\mathrm{e}^x`$ is to diagonalize the operator $``$. In quantum many-body problems, however, the basis of the diagonalized representation is often nontrivial, because we are typically interested in the Hamiltonian with two terms or more that are mutually non-commutative; for example, the Ising model in a transverse field,
$$=\underset{i,j}{}J_{ij}\sigma _i^z\sigma _j^z\mathrm{\Gamma }\underset{i}{}\sigma _i^x,$$
(11)
and the Hubbard model,
$$=t\underset{\sigma =,}{}\underset{i,j}{}\left(c_{i\sigma }^{}c_{j\sigma }+c_{j\sigma }^{}c_{i\sigma }\right)+U\underset{i}{}n_in_i.$$
(12)
In the first example (11), the quantization axis of the first term is the spin $`z`$ axis, while that of the second term is the spin $`x`$ axis. The two terms are therefore mutually non-commutative. In the second example (12), the first term is diagonalizable in the momentum space, whereas the second term is diagonalizable in the coordinate space. In both examples, each term is easily diagonalizable. Since one quantization axis is different from the other, the diagonalization of the sum of the terms becomes suddenly difficult.
The same situation arises in chaotic Hamilton dynamics. Consider a classical Hamiltonian
$$H(\stackrel{}{p},\stackrel{}{q})=K(\stackrel{}{p})+V(\stackrel{}{q}),$$
(13)
where $`K(\stackrel{}{p})`$ is the kinetic term and $`V(\stackrel{}{q})`$ is the potential term. The Hamilton equation is expressed in the form
$$\frac{\mathrm{d}}{\mathrm{d}t}\left(\begin{array}{c}\stackrel{}{p}(t)\\ \stackrel{}{q}(t)\end{array}\right)=\left(\begin{array}{c}\hfill \frac{\mathrm{d}}{\mathrm{d}\stackrel{}{q}}V(\stackrel{}{q})\\ \hfill \frac{\mathrm{d}}{\mathrm{d}\stackrel{}{p}}K(\stackrel{}{p})\end{array}\right)\left(\begin{array}{cc}& \widehat{V}\\ \widehat{K}& \end{array}\right)\left(\begin{array}{c}\stackrel{}{p}\\ \stackrel{}{q}\end{array}\right),$$
(14)
where the operators $`\widehat{K}`$ and $`\widehat{V}`$ are symbolic ones standing for the operations
$$\widehat{K}\stackrel{}{p}\frac{\mathrm{d}}{\mathrm{d}\stackrel{}{p}}K(\stackrel{}{p})\text{and}\widehat{V}\stackrel{}{q}\frac{\mathrm{d}}{\mathrm{d}\stackrel{}{q}}V(\stackrel{}{q}).$$
(15)
Although each operation of $`\widehat{K}`$ and $`\widehat{V}`$ is simple enough, the “Hamiltonian” operator
$$\left(\begin{array}{cc}& \widehat{V}\\ \widehat{K}& \end{array}\right)$$
(16)
is not easily tractable. This is because the kinetic part and the potential part,
$$𝒦\left(\begin{array}{cc}& \\ \widehat{K}& \end{array}\right)\text{and}𝒱\left(\begin{array}{cc}& \widehat{V}\end{array}\right),$$
(17)
do not commute with each other; see an example in Sect. 2.2 below.
To summarize this section, we frequently encounter the situation where the exponential operator of each term, $`\mathrm{e}^{xA}`$ and $`\mathrm{e}^{xB}`$, is easily obtained and yet the desired exponential operator $`\mathrm{e}^{x(A+B)}`$ is hard to come. This is the situation where the Trotter decomposition (1) becomes useful.
## 2 Why is the exponential product formula a good approximant?
We discussed in the previous section the importance of the exponential operator and the necessity of a way of treating it. We here discuss a remarkable advantage of the Trotter approximant to the exponential operator.
Let us first confirm that the Trotter approximant (1) is indeed a first-order approximant. By expanding the both sides of Eq. (1), we have
$`\mathrm{e}^{x(A+B)}`$ $`=`$ $`I+x(A+B)+{\displaystyle \frac{1}{2}}x^2(A+B)^2+O(x^3)`$ (18)
$`=`$ $`I+x(A+B)+{\displaystyle \frac{1}{2}}x^2\left(A^2+AB+BA+B^2\right)+O(x^3),`$
$`\mathrm{e}^{xA}\mathrm{e}^{xB}`$ $`=`$ $`\left(I+xA+{\displaystyle \frac{1}{2}}x^2A^2+O(x^3)\right)\left(I+xB+{\displaystyle \frac{1}{2}}x^2B^2+O(x^3)\right)`$ (19)
$`=`$ $`I+x(A+B)+{\displaystyle \frac{1}{2}}x^2\left(A^2+2AB+B^2\right)+O(x^3),`$
where $`I`$ is the identity operator. The difference between the two comes from the fact that in the approximant (19), the operator $`A`$ always comes on the left of the operator $`B`$. Hence we obtain
$$\mathrm{e}^{xA}\mathrm{e}^{xB}=\mathrm{e}^{x(A+B)+\frac{1}{2}x^2[A,B]+O(x^3)}.$$
(20)
In the actual application of the approximant, we divide the parameter $`x`$ into $`n`$ slices in the form
$$\left(\mathrm{e}^{\frac{x}{n}A}\mathrm{e}^{\frac{x}{n}B}\right)^n=\left[\mathrm{e}^{\frac{x}{n}(A+B)+\frac{1}{2}\left(\frac{x}{n}\right)^2[A,B]+O\left(\left(\frac{x}{n}\right)^3\right)}\right]^n=\mathrm{e}^{x(A+B)+\frac{1}{2}\frac{x^2}{n}[A,B]+O\left(\frac{x^3}{n^2}\right)}.$$
(21)
Thus the correction term vanishes in the limit $`n\mathrm{}`$. We refer to the integer $`n`$ as the Trotter number.
Now we discuss as to why we should be interested in generalizing the Trotter approximation. The Trotter approximant (1) and the generalized one (3), in fact, have a remarkable advantage over other approximants such as the frequently used one
$$\mathrm{e}^{x(A+B)}=I+x(A+B)+O(x^2).$$
(22)
The approximant of the form (3) conserves an important symmetry of the system in problems of quantum dynamics and Hamilton dynamics.
In problems of quantum dynamics, the exponential operator, or the Green’s function $`\mathrm{e}^{\mathrm{i}t}`$ is a unitary operator; hence the norm of the wave function does not change, which corresponds to the charge conservation. We here emphasize that the exponential product
$$\mathrm{e}^{\mathrm{i}tp_1A}\mathrm{e}^{\mathrm{i}tp_2B}\mathrm{e}^{\mathrm{i}tp_3A}\mathrm{}\mathrm{e}^{\mathrm{i}tp_MB}$$
(23)
is also a unitary operator. The perturbational approximant (22), on the other hand, does not conserve the norm of the wave function; in fact, the norm typically increases monotonically as the time passes as we demonstrate in Sect. 2.1 below.
In problems of Hamilton dynamics, the time evolution of the Hamilton system conserves the volume in the phase space $`\{\stackrel{}{p},\stackrel{}{q}\}`$, which is called the symplecticity in mathematics. The exponential product formula, in general, also has the symplecticity.
The time evolution of the Hamilton equation (14) is described by the exponential operator
$$\left(\begin{array}{c}\stackrel{}{p}(t)\\ \stackrel{}{q}(t)\end{array}\right)=\mathrm{e}^t\left(\begin{array}{c}\stackrel{}{p}(0)\\ \stackrel{}{q}(0)\end{array}\right),$$
(24)
where $``$ is the “Hamiltonian” operator (16). The Trotter decomposition approximates the time evolution with the operator
$$\mathrm{e}^t\left(\mathrm{e}^{\frac{t}{n}𝒦}\mathrm{e}^{\frac{t}{n}𝒱}\right)^n$$
(25)
with $`𝒦`$ and $`𝒱`$ given by Eq. (17). The operator $`\mathrm{e}^{\frac{t}{n}𝒦}`$ describes the time evolution over the time slice $`t/n`$ of a Hamilton system with only the kinetic energy $`K(p)`$. It thereby conserves the phase-space volume, so does the operator $`\mathrm{e}^{\frac{t}{n}𝒱}`$. The whole Trotter approximant therefore conserves the phase-space volume. This holds for any exponential product formula in the form (3) as well. Hence the exponential product formula, when used in the Hamilton dynamics, is sometimes called a symplectic integrator.
In equilibrium quantum statistical physics, the operator $`\mathrm{e}^\beta `$ does not have a particular symmetry except the symmetries of the Hamiltonian itself. The above advantage of the exponential product formula is hence lost when applied to numerical calculations of the partition function $`Z=Tr\mathrm{e}^\beta `$. In fact, in applying the higher-order decomposition (3) to the world-line quantum Monte Carlo simulation, some of the parameters $`\{p_1,p_2,\mathrm{},p_M\}`$ are negative, which causes the negative-sign problem in systems that usually do not have the negative-sign problem . The negative-sign problem is the problem that the Boltzmann weight of the system to be simulated becomes negative for some configurations.
Thanks to a recent development of the world-line quantum Monte Carlo simulation , the higher-order decomposition is not necessary anymore in some cases; the simulation is carried out in the limit $`n\mathrm{}`$ from the very beginning and hence the order of the correction term does not matter in such cases. See Appendix B for a brief review over the recent development.
### 2.1 Example: spin precession
The fact that the exponential product formula keeps the symmetry of the system is one of its remarkable advantages. In the present and next subsections, we demonstrate that this indeed affects numerical accuracy strongly. In the present subsection, we use a simple example of quantum dynamics, namely the spin precession.
Consider the simple Hamiltonian
$$=\sigma _z+\mathrm{\Gamma }\sigma _x=\left(\begin{array}{cc}1& \mathrm{\Gamma }\\ \mathrm{\Gamma }& 1\end{array}\right).$$
(26)
If we start the dynamics from the up-spin state
$$\psi (0)=\left(\begin{array}{c}1\\ 0\end{array}\right),$$
(27)
the spin precesses around the axis of the magnetic field $`\stackrel{}{H}=(\mathrm{\Gamma },0,1)`$ with the period
$$T=\frac{\pi }{\sqrt{1+\mathrm{\Gamma }^2}}.$$
(28)
Although it is easy to compute the dynamics exactly, we here use the Trotter approximant
$$G(t+\mathrm{\Delta }t;t)\mathrm{e}^{\mathrm{i}\mathrm{\Delta }t\sigma _z}\mathrm{e}^{\mathrm{i}\mathrm{\Delta }t\mathrm{\Gamma }\sigma _x}$$
(29)
and the perturbational approximant
$$G(t+\mathrm{\Delta }t;t)I\mathrm{i}\mathrm{\Delta }t=I\mathrm{i}\mathrm{\Delta }t(\sigma _z+\mathrm{\Gamma }\sigma _x).$$
(30)
The exact dynamics should conserve the energy expectation $``$. Figure. 1 shows the energy deviation due to the approximations.
The error in the energy of the Trotter approximation (29) oscillates periodically and never increases beyond the oscillation amplitude. The period of the oscillation in Fig. 1(a) is equal to that of the spin precession. We can understand this as follows: when the spin comes back to the original position after one cycle of the precession, it comes back accurately to the initial state (27) because of the unitarity of the Trotter approximation, and hence the oscillation.
In contrast, the error in the energy monotonically grows in the case of the perturbational approximant as is shown in Fig. 1(b). This is because the norm of the wave vector increases by the factor
$$1\mathrm{i}\mathrm{\Delta }t1+\mathrm{\Delta }t>1.$$
(31)
The remarkable difference between Fig. 1(a) and Fig. 1(b) thus comes from the fact that the Trotter approximant is unitary.
### 2.2 Example: symplectic integrator
We next demonstrate the Trotter decomposition (25) in an interesting example of chaotic dynamics. We again emphasize that keeping the symplecticity of the Hamilton dynamics has an important effect on numerical accuracy.
Let us first notice that the operators in Eq. (17) satisfy
$$𝒦^2=𝒱^2=0.$$
(32)
We therefore have
$`\mathrm{e}^{𝒦\mathrm{\Delta }t}\left(\begin{array}{c}\stackrel{}{p}\\ \stackrel{}{q}\end{array}\right)=\left(I+𝒦\mathrm{\Delta }t\right)\left(\begin{array}{c}\stackrel{}{p}\\ \stackrel{}{q}\end{array}\right)=\left(\begin{array}{c}\stackrel{}{p}\hfill \\ \stackrel{}{q}+\mathrm{\Delta }t\frac{\mathrm{d}}{\mathrm{d}\stackrel{}{p}}K(\stackrel{}{p})\hfill \end{array}\right),`$ (39)
$`\mathrm{e}^{𝒱\mathrm{\Delta }t}\left(\begin{array}{c}\stackrel{}{p}\\ \stackrel{}{q}\end{array}\right)=\left(I+𝒱\mathrm{\Delta }t\right)\left(\begin{array}{c}\stackrel{}{p}\\ \stackrel{}{q}\end{array}\right)=\left(\begin{array}{c}\stackrel{}{p}\mathrm{\Delta }t\frac{\mathrm{d}}{\mathrm{d}\stackrel{}{q}}V(\stackrel{}{q})\hfill \\ \stackrel{}{q}\hfill \end{array}\right).`$ (46)
Note that applying the two operators in the order $`\mathrm{e}^{𝒦\mathrm{\Delta }t}\mathrm{e}^{𝒱\mathrm{\Delta }t}`$ is different from applying them in the order $`\mathrm{e}^{𝒱\mathrm{\Delta }t}\mathrm{e}^{𝒦\mathrm{\Delta }t}`$; in the former, the update of $`\stackrel{}{q}`$ in the application of $`\mathrm{e}^{𝒦\mathrm{\Delta }t}`$ is done under the updated $`\stackrel{}{p}`$, whereas in the latter, it is done under $`\stackrel{}{p}`$ before the update.
Umeno and Suzuki demonstrated the use of symplectic integrators for chaotic dynamics of the system
$$K(\stackrel{}{p})=\frac{1}{2}(p_1{}_{}{}^{2}+p_2{}_{}{}^{2})\text{and}V(\stackrel{}{q})=\frac{1}{2}q_1{}_{}{}^{2}q_{2}^{}{}_{}{}^{2}.$$
(47)
The equipotential contour is given by $`|q_1q_2|=`$constant; hence the system is confined in the area surrounded by four hyperbolas as exemplified in Fig. 2(a).
The exact dynamics should conserve the energy. The Trotter approximation of the dynamics, Eq. (25), gives the energy fluctuation shown in Fig. 2(b). The energy, though deviates from the correct value sometimes, comes back after the deviation. In fact, the deviation occurs when the system goes into one of the four narrow valleys of the potential; it is suppressed again and again when the system comes back to the central area.
This is in striking contrast to the update due to the perturbational approximant
$$\left(\begin{array}{c}\stackrel{}{p}\\ \stackrel{}{q}\end{array}\right)\left(I+\mathrm{\Delta }t\right)\left(\begin{array}{c}\stackrel{}{p}\\ \stackrel{}{q}\end{array}\right)=\left(\begin{array}{c}\stackrel{}{p}\mathrm{\Delta }t\frac{\mathrm{d}}{\mathrm{d}\stackrel{}{q}}V(\stackrel{}{q})\\ \stackrel{}{q}+\mathrm{\Delta }t\frac{\mathrm{d}}{\mathrm{d}\stackrel{}{p}}K(\stackrel{}{p})\end{array}\right),$$
(48)
which yields the monotonic energy increase shown in Fig. 2(c). The reason of the difference between the approximants, though less apparent than in the case of the previous subsection, must be keeping the symplecticity, or the conservation of the phase-space volume.
## 3 Fractal decomposition
We emphasized in the previous section the importance of the exponential product formula. In the present section, we describe a way of constructing higher-order exponential product formulas recursively .
The easiest improvement of the Trotter formula (2) is the symmetrization:
$$S_2(x)\mathrm{e}^{\frac{x}{2}A}\mathrm{e}^{xB}\mathrm{e}^{\frac{x}{2}A}=\mathrm{e}^{x(A+B)+x^3R_3+x^5R_5+\mathrm{}}.$$
(49)
The symmetrized approximant has the property
$$S_2(x)S_2(x)=\mathrm{e}^{\frac{x}{2}A}\mathrm{e}^{xB}\mathrm{e}^{\frac{x}{2}A}\mathrm{e}^{\frac{x}{2}A}\mathrm{e}^{xB}\mathrm{e}^{\frac{x}{2}A}=I,$$
(50)
because of which the even-order terms vanish in the exponent of the right-hand side of Eq. (49). We can thereby promote the approximant (49) to a second-order approximant.
Now we introduce a way of constructing a symmetrized fourth-order approximant from the symmetrized second-order approximant (49). Consider a product
$`S(x)`$ $``$ $`S_2(sx)S_2((12s)x)S_2(sx)`$ (51)
$`=`$ $`\mathrm{e}^{\frac{s}{2}xA}\mathrm{e}^{sxB}\mathrm{e}^{\frac{1s}{2}xA}\mathrm{e}^{(12s)xB}\mathrm{e}^{\frac{1s}{2}xA}\mathrm{e}^{sxB}\mathrm{e}^{\frac{s}{2}xA},`$ (52)
where $`s`$ is an arbitrary real number for the moment. The expression (49) is followed by
$`S(x)`$ $`=`$ $`S_2(sx)S_2((12s)x)S_2(sx)`$ (53)
$`=`$ $`\mathrm{e}^{sx(A+B)+s^3x^3R_3+O(x^5)}\mathrm{e}^{(12s)x(A+B)+(12s)^3x^3R_3+O(x^5)}\mathrm{e}^{sx(A+B)+s^3x^3R_3+O(x^5)}`$
$`=`$ $`\mathrm{e}^{x(A+B)+\left[2s^3+(12s)^3\right]R_3+O(x^5)}.`$
(The readers should convince themselves by the Taylor expansion that the third-order correction in the exponent of the last line is just the sum of the third-order corrections in the exponents of the second line. This is not true for higher-order corrections.) Note that we arranged the parameters in the form $`\{s,12s,s\}`$ in Eq. (51) so that (i) the first-order term in the exponent of the last line of Eq. (53) should become $`x(A+B)`$ and (ii) the whole product $`S(x)`$ should be symmetrized, or should satisfy $`S(x)S(x)=I`$. Because of the second property, the even-order corrections vanish in the exponent of the last line of Eq. (53). Making the parameter $`s`$ a solution of the equation
$$2s^3+(12s)^3=0,\text{or}s=\frac{1}{2\sqrt[3]{2}}=1.351207191959657\mathrm{},$$
(54)
we promote the product (51) to a fourth-order approximant .
Following the same line of thought, we come up with another fourth-order approximant in the form
$`S_4(x)`$ $``$ $`S_2(s_2x)^2S_2((14s_2)x)S_2(s_2x)^2`$ (55)
$`=`$ $`\mathrm{e}^{\frac{s_2}{2}xA}\mathrm{e}^{s_2xB}\mathrm{e}^{s_2xA}\mathrm{e}^{s_2xB}\mathrm{e}^{\frac{13s_2}{2}xA}\mathrm{e}^{(14s_2)xB}\mathrm{e}^{\frac{13s_2}{2}xA}\mathrm{e}^{s_2xB}\mathrm{e}^{s_2xA}\mathrm{e}^{s_2xB}\mathrm{e}^{\frac{s_2}{2}xA},`$ (56)
where the parameter $`s_2`$ is a solution of the equation
$$4s_2{}_{}{}^{3}+(14s_2)^3=0,\text{or}s_2=\frac{1}{4\sqrt[3]{4}}=0.414490771794375\mathrm{}.$$
(57)
We can compare the fourth-order approximants (51) and (55) using the following diagram. Suppose that the exponential operator $`\mathrm{e}^{x(A+B)}`$ is a time-evolution operator from the time $`t=0`$ to the time $`t=x`$. In the product (51), the term $`S_2(sx)`$ on the right approximates the time evolution from $`t=0`$ to $`t=sx1.35x`$, the term $`S_2((12s)x)`$ in the middle approximates the time evolution from $`t=sx`$ to $`t=sx+(12s)x=(1s)x0.35x`$, and the term $`S_2(sx)`$ on the left approximates the time evolution from $`t=(1s)x`$ to $`t=(1s)x+sx=x`$. Let us express this time evolution as in Fig. 4(a).
The product (55) is similarly represented as in Fig. 4(b).
As is evident, the first product (51) has a part that goes into the “past,” or $`t<0`$. This can be problematic in some situations; in the diffusion from a delta-peak distribution, for example, there exists no “past” of the initial delta peak. The second product (55) does not have the problem and hence is recommended for general use.
Once we know how to construct the fourth-order approximant from the second-order approximant, the rest is quite straightforward . Following the construction (55), we construct the sixth-order approximant in the form
$`S_6(x)`$ $``$ $`S_4(s_4x)^2S_4((14s_4)x)S_4(s_4x)^2`$ (58)
$`=`$ $`\left(S_2(s_4s_2x)^2S_2(s_4(14s_2)x)S_2(s_4s_2x)^2\right)^2`$
$`\times S_2((14s_4)s_2x)^2S_2((14s_4)(14s_2)x)S_2((14s_4)s_2x)^2`$
$`\times \left(S_2(s_4s_2x)^2S_2(s_4(14s_2)x)S_2(s_4s_2x)^2\right)^2`$
with
$$4s_4^5+(14s_4)^5=0,\text{or}s_4=\frac{1}{4\sqrt[5]{4}}=0.373065827733272\mathrm{},$$
(59)
and further construct the eighth-order approximant in the form
$$S_8(x)S_6(s_6x)^2S_6((14s_6)x)S_6(s_6x)^2$$
(60)
with
$$4s_6^7+(14s_6)^7=0,\text{or}s_6=\frac{1}{4\sqrt[7]{4}}=0.359584649349992\mathrm{}.$$
(61)
These approximants are represented by the diagrams in Fig. 4. We can continue this recursive procedure, ending up with the exact time evolution, where the diagram ultimately becomes a fractal object. This is why the series of the approximants is called the fractal decomposition. It is an interesting thought that the back-and-forth time evolution in a fractal way reproduces the exact time evolution.
## 4 Time-ordered exponential
Before going into another way of constructing higher-order exponential product formulas, let us introduce, as an interlude, an important application of the exponential product formula. We show how to approximate the time-ordered exponential .
We have considered until now only the case where the operators $`A`$ and $`B`$ do not depend on $`x`$, or in other words, only the time evolution of a time-independent Hamiltonian. The fractal decomposition introduced in the previous section needs modification when applied to problems such as the quantum dynamics of a time-dependent Hamiltonian; in quantum annealing , for example, the transverse field $`\mathrm{\Gamma }`$ in the Hamiltonian (11) is changed in time.
The time-evolution operator of the quantum Hamiltonian
$$(t)=A(t)+B(t)$$
(62)
is not simply $`\mathrm{e}^{\mathrm{i}t}`$ but a time-ordered exponential in the form
$$G(t_2;t_1)=T\left[\mathrm{exp}\left(\mathrm{i}_{t_1}^{t_2}(s)𝑑s\right)\right].$$
(63)
It is quite well-known that
$$G_1(t+\mathrm{\Delta }t;t)\mathrm{e}^{\mathrm{i}\mathrm{\Delta }tA\left(t+\mathrm{\Delta }t\right)}\mathrm{e}^{\mathrm{i}\mathrm{\Delta }tB\left(t+\mathrm{\Delta }t\right)}$$
(64)
is an approximant of the first order of $`\mathrm{\Delta }t`$ and
$$G_2(t+\mathrm{\Delta }t;t)\mathrm{e}^{\frac{\mathrm{i}}{2}\mathrm{\Delta }tA\left(t+\frac{1}{2}\mathrm{\Delta }t\right)}\mathrm{e}^{\mathrm{i}\mathrm{\Delta }tB\left(t+\frac{1}{2}\mathrm{\Delta }t\right)}\mathrm{e}^{\frac{\mathrm{i}}{2}\mathrm{\Delta }tA\left(t+\frac{1}{2}\mathrm{\Delta }t\right)}$$
(65)
is an approximant of the second order. How do we construct higher-order approximants? We here show that a slight modification of the fractal decomposition gives the answer.
The key is to introduce a shift-time operator defined in
$$F(t)\mathrm{e}^{\mathrm{i}\mathrm{\Delta }t𝒯}G(t)=F(t+\mathrm{\Delta }t)G(t).$$
(66)
Note that the operator acts on the function on the left. The shift-time operator is expressed in the form
$$𝒯=\mathrm{i}\stackrel{}{\frac{}{t}}$$
(67)
in the case where $`F(t)`$ is an analytic function, but the definition (66) does not limit its use to the analytic case. If we have two shift-time operators, the result is
$`F(t)\mathrm{e}^{\mathrm{i}\mathrm{\Delta }t𝒯}G(t)\mathrm{e}^{\mathrm{i}\mathrm{\Delta }t𝒯}H(t)`$ $`=`$ $`F(t+\mathrm{\Delta }t)G(t)\mathrm{e}^{\mathrm{i}\mathrm{\Delta }t𝒯}H(t)`$ (68)
$`=`$ $`F(t+2\mathrm{\Delta }t)G(t+\mathrm{\Delta }t)H(t).`$
With the use of the shift-time operator, the time-ordered exponential (63) is transformed as
$$T\left[\mathrm{exp}\left(\mathrm{i}_t^{t+\mathrm{\Delta }t}(s)𝑑s\right)\right]=\mathrm{e}^{\mathrm{i}\mathrm{\Delta }t\left((t)+𝒯\right)}.$$
(69)
We can prove this by using the Trotter approximation as follows:
$`\mathrm{e}^{\mathrm{i}\mathrm{\Delta }t\left((t)+𝒯\right)}`$ $`=`$ $`\underset{n\mathrm{}}{lim}\left(\mathrm{e}^{\mathrm{i}\frac{\mathrm{\Delta }t}{n}(t)}\mathrm{e}^{\mathrm{i}\frac{\mathrm{\Delta }t}{n}𝒯}\right)^n`$ (70)
$`=`$ $`\underset{n\mathrm{}}{lim}\mathrm{e}^{\mathrm{i}\frac{\mathrm{\Delta }t}{n}(t)}\mathrm{e}^{\mathrm{i}\frac{\mathrm{\Delta }t}{n}𝒯}\mathrm{e}^{\mathrm{i}\frac{\mathrm{\Delta }t}{n}(t)}\mathrm{e}^{\mathrm{i}\frac{\mathrm{\Delta }t}{n}𝒯}\mathrm{}\mathrm{e}^{\mathrm{i}\frac{\mathrm{\Delta }t}{n}(t)}\mathrm{e}^{\mathrm{i}\frac{\mathrm{\Delta }t}{n}𝒯}`$
$`=`$ $`\underset{n\mathrm{}}{lim}\mathrm{e}^{\mathrm{i}\frac{\mathrm{\Delta }t}{n}\left(t+\mathrm{\Delta }t\right)}\mathrm{e}^{\mathrm{i}\frac{\mathrm{\Delta }t}{n}\left(t+\frac{n1}{n}\mathrm{\Delta }t\right)}\mathrm{}\mathrm{e}^{\mathrm{i}\frac{\mathrm{\Delta }t}{n}\left(t+\frac{1}{n}\mathrm{\Delta }t\right)}`$
$`=`$ $`T\left[\mathrm{exp}\left(\mathrm{i}{\displaystyle _t^{t+\mathrm{\Delta }t}}(s)𝑑s\right)\right].`$
Decomposing the Hamiltonian into two parts as in Eq. (62), we have now three parts in the exponent of the time-evolution operator as in
$$T\left[\mathrm{exp}\left(\mathrm{i}_t^{t+\mathrm{\Delta }t}(s)𝑑s\right)\right]=\mathrm{e}^{\mathrm{i}\mathrm{\Delta }t\left(A(t)+B(t)+𝒯\right)}.$$
(71)
We then approximate the exponential in the right-hand side of Eq. (71). The first-order approximant is given by
$`G_1(t+\mathrm{\Delta }t;t)`$ $`=`$ $`\mathrm{e}^{\mathrm{i}\mathrm{\Delta }tA(t)}\mathrm{e}^{\mathrm{i}\mathrm{\Delta }tB(t)}\mathrm{e}^{\mathrm{i}\mathrm{\Delta }t𝒯}`$ (72)
$`=`$ $`\mathrm{e}^{\mathrm{i}\mathrm{\Delta }tA(t+\mathrm{\Delta }t)}\mathrm{e}^{\mathrm{i}\mathrm{\Delta }tB(t+\mathrm{\Delta }t)}`$
and the second-order approximant is given by
$`G_2(t+\mathrm{\Delta }t;t)`$ $`=`$ $`\mathrm{e}^{\frac{\mathrm{i}}{2}\mathrm{\Delta }t𝒯}\mathrm{e}^{\frac{\mathrm{i}}{2}\mathrm{\Delta }tA(t)}\mathrm{e}^{\mathrm{i}\mathrm{\Delta }tB(t)}\mathrm{e}^{\frac{\mathrm{i}}{2}\mathrm{\Delta }tA(t)}\mathrm{e}^{\frac{\mathrm{i}}{2}\mathrm{\Delta }t𝒯}`$ (73)
$`=`$ $`\mathrm{e}^{\frac{\mathrm{i}}{2}\mathrm{\Delta }tA\left(t+\frac{1}{2}\mathrm{\Delta }t\right)}\mathrm{e}^{\mathrm{i}\mathrm{\Delta }tB\left(t+\frac{1}{2}\mathrm{\Delta }t\right)}\mathrm{e}^{\frac{\mathrm{i}}{2}\mathrm{\Delta }tA\left(t+\frac{1}{2}\mathrm{\Delta }t\right)}.`$
Higher-order approximants are given by the fractal decomposition of the three parts, $`A`$, $`B`$, and $`𝒯`$. The fractal decomposition of three parts is easily obtained by substituting
$$S_2(x)\mathrm{e}^{\frac{x}{2}A}\mathrm{e}^{\frac{x}{2}B}\mathrm{e}^{xC}\mathrm{e}^{\frac{x}{2}B}\mathrm{e}^{\frac{x}{2}A}=\mathrm{e}^{x(A+B+C)+O(x^3)}$$
(74)
for Eq. (49). The fourth-order approximant is thereby obtained as
$`G_4(t+\mathrm{\Delta }t;t)`$ $``$ $`\left(\mathrm{e}^{\frac{\mathrm{i}}{2}s_2\mathrm{\Delta }t𝒯}\mathrm{e}^{\frac{\mathrm{i}}{2}s_2\mathrm{\Delta }tA(t)}\mathrm{e}^{\mathrm{i}s_2\mathrm{\Delta }tB(t)}\mathrm{e}^{\frac{\mathrm{i}}{2}s_2\mathrm{\Delta }tA(t)}\mathrm{e}^{\frac{\mathrm{i}}{2}s_2\mathrm{\Delta }t𝒯}\right)^2`$ (75)
$`\times \mathrm{e}^{\frac{\mathrm{i}}{2}(14s_2)\mathrm{\Delta }t𝒯}\mathrm{e}^{\frac{\mathrm{i}}{2}(14s_2)\mathrm{\Delta }tA(t)}\mathrm{e}^{\mathrm{i}(14s_2)\mathrm{\Delta }tB(t)}\mathrm{e}^{\frac{\mathrm{i}}{2}(14s_2)\mathrm{\Delta }tA(t)}\mathrm{e}^{\frac{\mathrm{i}}{2}(14s_2)\mathrm{\Delta }t𝒯}`$
$`\times \left(\mathrm{e}^{\frac{\mathrm{i}}{2}s_2\mathrm{\Delta }t𝒯}\mathrm{e}^{\frac{\mathrm{i}}{2}s_2\mathrm{\Delta }tA(t)}\mathrm{e}^{\mathrm{i}s_2\mathrm{\Delta }tB(t)}\mathrm{e}^{\frac{\mathrm{i}}{2}s_2\mathrm{\Delta }tA(t)}\mathrm{e}^{\frac{\mathrm{i}}{2}s_2\mathrm{\Delta }t𝒯}\right)^2`$
$`=`$ $`\mathrm{e}^{\frac{\mathrm{i}}{2}s_2\mathrm{\Delta }tA\left(t+\frac{2s_2}{2}\mathrm{\Delta }t\right)}\mathrm{e}^{\mathrm{i}s_2\mathrm{\Delta }tB\left(t+\frac{2s_2}{2}\mathrm{\Delta }t\right)}\mathrm{e}^{\frac{\mathrm{i}}{2}s_2\mathrm{\Delta }tA\left(t+\frac{2s_2}{2}\mathrm{\Delta }t\right)}`$
$`\times \mathrm{e}^{\frac{\mathrm{i}}{2}s_2\mathrm{\Delta }tA\left(t+\frac{23s_2}{2}\mathrm{\Delta }t\right)}\mathrm{e}^{\mathrm{i}s_2\mathrm{\Delta }tB\left(t+\frac{23s_2}{2}\mathrm{\Delta }t\right)}\mathrm{e}^{\frac{\mathrm{i}}{2}s_2\mathrm{\Delta }tA\left(t+\frac{23s_2}{2}\mathrm{\Delta }t\right)}`$
$`\times \mathrm{e}^{\frac{\mathrm{i}}{2}s_2\mathrm{\Delta }tA\left(t+\frac{1}{2}\mathrm{\Delta }t\right)}\mathrm{e}^{\mathrm{i}s_2\mathrm{\Delta }tB\left(t+\frac{1}{2}\mathrm{\Delta }t\right)}\mathrm{e}^{\frac{\mathrm{i}}{2}s_2\mathrm{\Delta }tA\left(t+\frac{1}{2}\mathrm{\Delta }t\right)}`$
$`\times \mathrm{e}^{\frac{\mathrm{i}}{2}s_2\mathrm{\Delta }tA\left(t+\frac{3s_2}{2}\mathrm{\Delta }t\right)}\mathrm{e}^{\mathrm{i}s_2\mathrm{\Delta }tB\left(t+\frac{3s_2}{2}\mathrm{\Delta }t\right)}\mathrm{e}^{\frac{\mathrm{i}}{2}s_2\mathrm{\Delta }tA\left(t+\frac{3s_2}{2}\mathrm{\Delta }t\right)}`$
$`\times \mathrm{e}^{\frac{\mathrm{i}}{2}s_2\mathrm{\Delta }tA\left(t+\frac{s_2}{2}\mathrm{\Delta }t\right)}\mathrm{e}^{\mathrm{i}s_2\mathrm{\Delta }tB\left(t+\frac{s_2}{2}\mathrm{\Delta }t\right)}\mathrm{e}^{\frac{\mathrm{i}}{2}s_2\mathrm{\Delta }tA\left(t+\frac{s_2}{2}\mathrm{\Delta }t\right)}`$
with the coefficient $`s_2`$ given by Eq. (57).
## 5 Quantum analysis – Towards the construction of general decompositions –
In the last section before the summary, we discuss the calculus of the correction terms. In the fractal decomposition, we construct higher-order approximants recursively. Is it possible to construct higher-order approximants directly, not recursively? In fact, Ruth found (not systematically) a third-order formula
$$\mathrm{e}^{\frac{7}{24}xA}\mathrm{e}^{\frac{2}{3}xB}\mathrm{e}^{\frac{3}{4}xA}\mathrm{e}^{\frac{2}{3}xB}\mathrm{e}^{\frac{1}{24}xA}\mathrm{e}^{xB}=\mathrm{e}^{x(A+B)+O(x^4)},$$
(76)
which would not be found within the framework of the fractal decomposition.
For the purpose of finding higher-order formulas directly, we need to compute the correction terms in the exponent as
$$\mathrm{e}^{p_1xA}\mathrm{e}^{p_2xB}\mathrm{e}^{p_3xA}\mathrm{e}^{p_4xB}\mathrm{}\mathrm{e}^{p_MxB}=\mathrm{e}^{x(A+B)+x^2R_2+x^3R_3+\mathrm{}}.$$
(77)
This is one of the aims of the quantum analysis developed by one of the present authors (M.S.) . Then we can put the correction terms to zero up to a desired order and solve the set of non-linear simultaneous equations
$$R_2=0,R_3=0,\mathrm{},R_m=0,$$
(78)
thereby obtaining the parameters $`\{p_i\}`$.
### 5.1 Operator differential
The main feature of the quantum analysis is to introduce operator differential. In order to motivate the readers, suppose that we can write down an identity
$$\frac{\mathrm{d}}{\mathrm{d}x}f(A(x))=\frac{\mathrm{d}f(A)}{\mathrm{d}A}\frac{\mathrm{d}A(x)}{\mathrm{d}x},$$
(79)
where $`f(A)`$ is an operator functional. The derivative with respect to $`x`$ on the right-hand side is well-defined; for example, $`\mathrm{d}A(x)/\mathrm{d}x=B+2xC`$ for $`A(x)=xB+x^2C`$. Now, is it possible to define the differentiation $`\mathrm{d}f(A)/\mathrm{d}A`$?
Let us discuss as to what should be the definition of the operator differential in order for the identity (79) to hold. The definition of the $`x`$ derivative is expressed as
$$A(x+h)=A(x)+h\frac{\mathrm{d}A(x)}{\mathrm{d}x}+O(h^2).$$
(80)
The left-hand side of the identity (79) is given by the definition of the derivative as
$$\frac{\mathrm{d}}{\mathrm{d}x}f(A(x))=\underset{h0}{lim}\frac{f(A(x+h))f(A(x))}{h}=\underset{h0}{lim}\frac{f\left(A(x)+h\frac{\mathrm{d}A(x)}{\mathrm{d}x}\right)f(A(x))}{h}.$$
(81)
The identity (79) suggests that the operator differential $`\mathrm{d}f(A)/\mathrm{d}A`$ must be a hyperoperator that maps the operator $`\mathrm{d}A(x)/\mathrm{d}x`$ to the operator given by Eq. (81).
Thus we arrive at the definition of the operator differential within the framework of the quantum analysis : if we can express the operator given by
$$\mathrm{d}f(A)\underset{h0}{lim}\frac{f(A+h\mathrm{d}A)f(A)}{h}$$
(82)
in terms of a hyperoperator mapping from an arbitrary operator $`\mathrm{d}A`$ as in $`\mathrm{d}A\mathrm{d}f(A)`$, then we refer to the hyperoperator as an operator differential $`\mathrm{d}f(A)/\mathrm{d}A`$ and denote it in the form
$$\mathrm{d}f(A)=\frac{\mathrm{d}f(A)}{\mathrm{d}A}\mathrm{d}A.$$
(83)
We stress here that the operator differential $`\mathrm{d}f(A)/\mathrm{d}A`$ must be expressed in terms of $`A`$ and the commutation relation of $`A`$, or the “inner derivation”
$$\delta _A[A,],$$
(84)
but not in terms of the arbitrary operator $`\mathrm{d}A`$. The convergence of Eq. (82) is in the sense of the norm convergence which is uniform with respect to the arbitrary operator $`\mathrm{d}A`$.
Let us consider the application of the above in a simple example $`f(A)=A^2`$. The definition (82) is followed by
$`\mathrm{d}f(A)`$ $`=`$ $`\underset{h0}{lim}{\displaystyle \frac{(A+h\mathrm{d}A)^2A^2}{h}}=\underset{h0}{lim}{\displaystyle \frac{hA\mathrm{d}A+h\mathrm{d}AA+h^2(\mathrm{d}A)^2}{h}}`$ (85)
$`=`$ $`A\mathrm{d}A+\mathrm{d}AA=2A\mathrm{d}A(A\mathrm{d}A\mathrm{d}AA)`$
$`=`$ $`\left(2A\delta _A\right)\mathrm{d}A.`$
Thus we have
$$\frac{\mathrm{d}(A^2)}{\mathrm{d}A}=2A\delta _A.$$
(86)
If $`A=xB+x^2C`$, we use the result (86) for Eq. (79) and have
$`{\displaystyle \frac{\mathrm{d}}{\mathrm{d}x}}(xB+x^2C)^2`$ $`=`$ $`\left(2xB+2x^2C\delta _{xB+x^2C}\right)(B+2xC)`$ (87)
$`=`$ $`(2xB+2x^2C)(B+2xC)[xB+x^2C,B+2xC]`$
$`=`$ $`2xB^2+4x^2BC+2x^2CB+4x^3C^22x^2(BCCB)x^2(CBBC)`$
$`=`$ $`2xB^2+3x^2BC+3x^2CB+4x^3C^2,`$
which is indeed identical to the result of straightforward algebra.
We cannot see in this simple example any merit of the use of the quantum analysis. The readers should wait for more complicated examples given later in Sec. 5.3, where we show that the differential of exponential operators is given in terms of the inner derivation. The Lie algebra is defined by commutation relations, or the inner derivation; it is hence essential to obtain results in terms of the inner derivation, not in terms of naive expansions such as the right-hand side of Eq. (87).
### 5.2 Inner derivation
We here provide some of the important formulas of the inner derivation (84) as preparation for the next subsection, where we give the differential of exponential operators.
First, we have linearity: for any c-numbers $`a`$ and $`b`$, the inner derivation of the operators $`A`$ and $`B`$ satisfies
$$\delta _{aA+bB}=[aA+bB,]=a[A,]+b[B,]=a\delta _A+b\delta _B.$$
(88)
Any powers of the operator $`A`$ are commutable with the inner derivation of any powers of the same operator:
$$[A^m,\delta _{A^n}]=0,$$
(89)
because
$$A^m\delta _{A^n}B=A^m[A^n,B]=[A^n,A^mB]=\delta _{A^n}A^mB$$
(90)
for an arbitrary operator $`B`$ and any integers $`m`$ and $`n`$. We can generalize the identity (89) to the case of any analytic functions of the operator $`A`$:
$$[f(A),\delta _{g(A)}]=0,$$
(91)
where $`f(A)`$ and $`g(A)`$ are defined by the Taylor expansion as
$$f(A)=\underset{n=0}{\overset{\mathrm{}}{}}a_nA^n\text{and}g(A)=\underset{n=0}{\overset{\mathrm{}}{}}b_nA^n.$$
(92)
Next, we prove the identity
$$\delta _{f(A)g(A)}=f(A)\delta _{g(A)}+g(A)\delta _{f(A)}\delta _{g(A)}\delta _{f(A)}.$$
(93)
The proof is as follows: for an arbitrary operator $`B`$, we have
$`f(A)\delta _{g(A)}B+g(A)\delta _{f(A)}B\delta _{g(A)}\delta _{f(A)}B`$ $`=`$ $`f(A)[g(A),B]+g(A)[f(A),B][g(A),[f(A),B]]`$ (94)
$`=`$ $`f(A)[g(A),B]+[f(A),B]g(A)=[f(A)g(A),B]`$
$`=`$ $`\delta _{f(A)g(A)}B.`$
Note that we can rewrite the identity (93) as
$`\delta _{f(A)g(A)}`$ $`=`$ $`\delta _{g(A)}f(A)+g(A)\delta _{f(A)}\delta _{g(A)}\delta _{f(A)}`$ (95)
$`=`$ $`\delta _{g(A)}\left(f(A)\delta _{f(A)}\right)+g(A)\delta _{f(A)}`$
because of the identity (91). In the special case $`f(A)=A`$, we have
$$\delta _{Ag(A)}=\delta _{g(A)}\left(A\delta _A\right)+g(A)\delta _A.$$
(96)
With the repeated use of the identity (96), we then prove the identity
$$\delta _{A^n}=A^n\left(A\delta _A\right)^n$$
(97)
for any integer $`n`$. This is proved by means of mathematical induction. The identity (97) indeed holds for $`n=1`$. Assume now that
$$\delta _{A^{n1}}=A^{n1}\left(A\delta _A\right)^{n1}.$$
(98)
Then the identity (96) yields
$`\delta _{A^n}`$ $`=`$ $`\delta _{AA^{n1}}=\delta _{A^{n1}}\left(A\delta _A\right)+A^{n1}\delta _A=\left[A^{n1}(A\delta _A)^{n1}\right]\left(A\delta _A\right)+A^{n1}\delta _A`$ (99)
$`=`$ $`A^n\left(A\delta _A\right)^n.`$
An interesting and quite well-known identity is
$$\mathrm{e}^{xA}B\mathrm{e}^{xA}=\mathrm{e}^{x\delta _A}B.$$
(100)
We can prove this by differentiating the left-hand side by $`x`$. First, note that
$$\frac{\mathrm{d}}{\mathrm{d}x}\mathrm{e}^{xA}B\mathrm{e}^{xA}=\mathrm{e}^{xA}AB\mathrm{e}^{xA}\mathrm{e}^{xA}BA\mathrm{e}^{xA}=\mathrm{e}^{xA}[A,B]\mathrm{e}^{xA}.$$
(101)
We thereby have the following in each order of $`x`$:
$`{\displaystyle \frac{\mathrm{d}}{\mathrm{d}x}}\mathrm{e}^{xA}B\mathrm{e}^{xA}|_{x=0}`$ $`=`$ $`\mathrm{e}^{xA}[A,B]\mathrm{e}^{xA}|_{x=0}=\delta _AB,`$ (102)
$`{\displaystyle \frac{\mathrm{d}^2}{\mathrm{d}x^2}}\mathrm{e}^{xA}B\mathrm{e}^{xA}|_{x=0}`$ $`=`$ $`\mathrm{e}^{xA}[A,[A,B]]\mathrm{e}^{xA}|_{x=0}=\delta _{A}^{}{}_{}{}^{2}B,`$
$`\mathrm{},`$
which proves the identity (100). As a corollary, we obtain the following identity:
$$\mathrm{e}^{\delta _A}\mathrm{e}^{\delta _B}=\mathrm{e}^{\delta _\mathrm{\Phi }}\text{if}\mathrm{e}^A\mathrm{e}^B=\mathrm{e}^\mathrm{\Phi }.$$
(104)
The proof is straightforward; for an arbitrary operator $`C`$, we have
$$\mathrm{e}^{\delta _A}\mathrm{e}^{\delta _B}C=\mathrm{e}^A\mathrm{e}^BC\mathrm{e}^B\mathrm{e}^A=\mathrm{e}^\mathrm{\Phi }C\mathrm{e}^\mathrm{\Phi }=\mathrm{e}^{\delta _\mathrm{\Phi }}C.$$
(105)
### 5.3 Differential of exponential operators
We are now in a position to discuss the differential of exponential operators. We begin with the differential of the power of an operator, $`f(A)=A^n`$, a generalization of the identity (86). The result is
$$\frac{\mathrm{d}\left(A^n\right)}{\mathrm{d}A}=\frac{A^n\left(A\delta _A\right)^n}{\delta _A}=\frac{\delta _{A^n}}{\delta _A}.$$
(106)
An important comment is in order. The identity (106) does not claim that the inverse of $`\delta _A`$ is well-defined. In fact, the inner derivation $`\delta _A`$ in the denominator is canceled when we expand the numerator of the second expression. The denominator is well-defined only in such cases.
We use the identity (97) in the derivation of the identity (106). The definition (82) is followed by
$`\mathrm{d}f(A)`$ $`=`$ $`\underset{h0}{lim}{\displaystyle \frac{\left(A+h\mathrm{d}A\right)^nA^n}{h}}={\displaystyle \underset{j=1}{\overset{n}{}}}A^{j1}(\mathrm{d}A)A^{nj}`$ (107)
$`=`$ $`\left(nA^{n1}{\displaystyle \underset{j=1}{\overset{n}{}}}A^{j1}\delta _{A^{nj}}\right)\mathrm{d}A=\left\{nA^{n1}{\displaystyle \underset{j=1}{\overset{n}{}}}A^{j1}\left[A^{nj}\left(A\delta _A\right)^{nj}\right]\right\}\mathrm{d}A`$
$`=`$ $`{\displaystyle \underset{j=1}{\overset{n}{}}}A^{j1}\left(A\delta _A\right)^{nj}\mathrm{d}A={\displaystyle \frac{A^n\left(A\delta _A\right)^n}{A\left(A\delta _A\right)}}\mathrm{d}A`$
$`=`$ $`{\displaystyle \frac{A^n\left(A\delta _A\right)^n}{\delta _A}}\mathrm{d}A={\displaystyle \frac{\delta _{A^n}}{\delta _A}}\mathrm{d}A.`$
Note again that the transformation in the fourth line is well-defined only because the expansion of the numerator cancels the denominator.
We can generalize the identity (106) to any analytic functions defined by the Taylor expansion (92). The result is
$$\frac{\mathrm{d}f(A)}{\mathrm{d}A}=\frac{f(A)f\left(A\delta _A\right)}{\delta _A}=\frac{\delta _{f(A)}}{\delta _A}.$$
(108)
It is interesting to note that the operator differential or the quantum derivative is expressed by a difference form of hyperoperators. As a special case, we arrive at the identity
$$\frac{\mathrm{de}^A}{\mathrm{d}A}=\frac{\mathrm{e}^A\mathrm{e}^{A\delta _A}}{\delta _A}=\mathrm{e}^A\frac{1\mathrm{e}^{\delta _A}}{\delta _A}.$$
(109)
### 5.4 Example: Baker-Campbell-Hausdorff formula
We now use the formula (109) for the derivation of the Baker-Campbell-Hausdorff formula, or the derivation of higher-order terms of the exponent $`\mathrm{\Phi }(x)`$ given in
$$\mathrm{e}^{\mathrm{\Phi }(x)}=\mathrm{e}^{xA}\mathrm{e}^{xB}.$$
(110)
The differential of the left-hand side of Eq. (110) gives
$$\frac{\mathrm{d}}{\mathrm{d}x}\mathrm{e}^{\mathrm{\Phi }(x)}=\frac{\mathrm{de}^\mathrm{\Phi }}{\mathrm{d}\mathrm{\Phi }}\frac{\mathrm{d}\mathrm{\Phi }(x)}{\mathrm{d}x}=\mathrm{e}^{\mathrm{\Phi }(x)}\frac{1\mathrm{e}^{\delta _{\mathrm{\Phi }(x)}}}{\delta _{\mathrm{\Phi }(x)}}\frac{\mathrm{d}\mathrm{\Phi }(x)}{\mathrm{d}x}$$
(111)
owing to Eq. (109), while the differential of the right-hand side of Eq. (110) gives
$$\frac{\mathrm{d}}{\mathrm{d}x}\mathrm{e}^{xA}\mathrm{e}^{xB}=\mathrm{e}^{xA}A\mathrm{e}^{xB}+\mathrm{e}^{xA}\mathrm{e}^{xB}B=\mathrm{e}^{xA}\mathrm{e}^{xB}\left(\mathrm{e}^{xB}A\mathrm{e}^{xB}+B\right)=\mathrm{e}^{\mathrm{\Phi }(x)}\left(\mathrm{e}^{x\delta _B}A+B\right),$$
(112)
where we have used the identity (100). Equating the both sides, we have
$$\frac{\mathrm{d}\mathrm{\Phi }(x)}{\mathrm{d}x}=\frac{\delta _{\mathrm{\Phi }(x)}}{1\mathrm{e}^{\delta _{\mathrm{\Phi }(x)}}}\left(\mathrm{e}^{x\delta _B}A+B\right)=\frac{\delta _{\mathrm{\Phi }(x)}}{\mathrm{e}^{\delta _{\mathrm{\Phi }(x)}}1}\left(A+\mathrm{e}^{x\delta _A}B\right).$$
(113)
The second equality is due to the identity (104).
We can expand the right-hand side of Eq. (113) as follows. Note here that
$$\mathrm{e}^{\delta _{\mathrm{\Phi }(x)}}=\mathrm{e}^{x\delta _A}\mathrm{e}^{x\delta _B}\text{yields}\delta _{\mathrm{\Phi }(x)}=\mathrm{log}\left(\mathrm{e}^{x\delta _A}\mathrm{e}^{x\delta _B}\right).$$
(114)
Thus we transform Eq. (113) as
$$\frac{\mathrm{d}\mathrm{\Phi }(x)}{\mathrm{d}x}=\frac{\mathrm{log}\left(\mathrm{e}^{x\delta _A}\mathrm{e}^{x\delta _B}\right)}{\mathrm{e}^{x\delta _A}\mathrm{e}^{x\delta _B}1}\left(A+\mathrm{e}^{x\delta _A}B\right)=\underset{k=0}{\overset{\mathrm{}}{}}\frac{(1)^k}{k+1}\left(\mathrm{e}^{x\delta _A}\mathrm{e}^{x\delta _B}1\right)^k\left(A+\mathrm{e}^{x\delta _A}B\right).$$
(115)
We finally arrive at
$$\mathrm{\Phi }(x)=\underset{k=0}{\overset{\mathrm{}}{}}\frac{(1)^k}{k+1}_0^x\left(\mathrm{e}^{t\delta _A}\mathrm{e}^{t\delta _B}1\right)^k\left(A+\mathrm{e}^{t\delta _A}B\right)𝑑t.$$
(116)
It is very important to notice here that all the expansion terms are given by commutation relations. One of the merits of the quantum analysis is to be able to express the expansion in terms of commutation relations.
Let us derive, for example, the term of the third order of $`x`$, or the second order of $`t`$ of Eq. (116). Up to the second order, we have
$`\mathrm{e}^{t\delta _A}\mathrm{e}^{t\delta _B}1`$ $``$ $`t\left(\delta _A+\delta _B\right)+{\displaystyle \frac{t^2}{2}}\left(\delta _{A}^{}{}_{}{}^{2}+2\delta _A\delta _B+\delta _{B}^{}{}_{}{}^{2}\right)`$ (117)
$`=`$ $`t\delta _{A+B}+{\displaystyle \frac{t^2}{2}}\left(\delta _{A+B}^{}{}_{}{}^{2}+\delta _A\delta _B\delta _B\delta _A\right),`$
$`\left(\mathrm{e}^{t\delta _A}\mathrm{e}^{t\delta _B}1\right)^2`$ $``$ $`t^2\delta _{A+B}^{}{}_{}{}^{2},`$ (118)
and hence
$`\left(\mathrm{e}^{t\delta _A}\mathrm{e}^{t\delta _B}1\right)^0\left(A+\mathrm{e}^{t\delta _A}B\right)`$ $``$ $`\left(A+B\right)+t\delta _AB+{\displaystyle \frac{t^2}{2}}\delta _{A}^{}{}_{}{}^{2}B,`$ (119)
$`\left(\mathrm{e}^{t\delta _A}\mathrm{e}^{t\delta _B}1\right)^1\left(A+\mathrm{e}^{t\delta _A}B\right)`$ $``$ $`t\delta _{A+B}\left(A+B\right)`$ (120)
$`+{\displaystyle \frac{t^2}{2}}\left(\delta _{A+B}^{}{}_{}{}^{2}+\delta _A\delta _B\delta _B\delta _A\right)\left(A+B\right)+t^2\left(\delta _A+\delta _B\right)\delta _AB`$
$`=`$ $`{\displaystyle \frac{t^2}{2}}\left(\delta _A\delta _BA\delta _B\delta _AB\right)+t^2\left(\delta _A+\delta _B\right)\delta _AB,`$
$`\left(\mathrm{e}^{t\delta _A}\mathrm{e}^{t\delta _B}1\right)^2\left(A+\mathrm{e}^{t\delta _A}B\right)`$ $``$ $`t^2\delta _{A+B}^{}{}_{}{}^{2}\left(A+B\right)=0.`$ (121)
Summing up the second-order terms with the coefficient $`(1)^k/(k+1)`$, we have
$$\frac{t^2}{2}\delta _{A}^{}{}_{}{}^{2}B\frac{t^2}{4}\left(\delta _A\delta _BA\delta _B\delta _AB\right)\frac{t^2}{2}\left(\delta _A+\delta _B\right)\delta _AB=\frac{t^2}{4}\left(\delta _{A}^{}{}_{}{}^{2}B+\delta _{B}^{}{}_{}{}^{2}A\right),$$
(122)
which we integrate to obtain
$$\frac{x^3}{12}\left(\delta _{A}^{}{}_{}{}^{2}B+\delta _{B}^{}{}_{}{}^{2}A\right)=\frac{x^3}{12}\left([A,[A,B]]+[[A,B],B]\right).$$
(123)
### 5.5 Example: Ruth’s formula
We now extend the above computation to the exponential product
$$\mathrm{e}^{p_1xA}\mathrm{e}^{p_2xB}\mathrm{e}^{p_3xA}\mathrm{e}^{p_4xB}\mathrm{e}^{p_5xA}\mathrm{e}^{p_6xB}=\mathrm{e}^{\mathrm{\Phi }(x)}$$
(124)
and seek Ruth’s formula (76) as a specific solution of the general formula. We compute the second-order and third-order correction terms of $`\mathrm{\Phi }(x)`$, defined in
$$\mathrm{\Phi }(x)=x(A+B)+x^2R_2+x^3R_3+O(x^4),$$
(125)
and put $`R_2=R_3=0`$.
The same computation as from Eq. (111) through Eq. (116) produces
$`\mathrm{\Phi }(x)`$ $`=`$ $`{\displaystyle \underset{k=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(1)^k}{k+1}}{\displaystyle _0^x}\left(\mathrm{e}^{p_1t\delta _A}\mathrm{e}^{p_2t\delta _B}\mathrm{e}^{p_3t\delta _A}\mathrm{e}^{p_4t\delta _B}\mathrm{e}^{p_5t\delta _A}\mathrm{e}^{p_6t\delta _B}1\right)^k`$ (126)
$`\times \left(p_1A+\mathrm{e}^{p_1t\delta _A}p_2B+\mathrm{e}^{p_1t\delta _A}\mathrm{e}^{p_2t\delta _B}p_3A\mathrm{}\right)dt.`$
Note again that all the terms are given by commutation relations.
For the term $`k=0`$, we have up to the second order of $`x`$,
$`p_1A+\mathrm{e}^{p_1t\delta _A}p_2B+\mathrm{e}^{p_1t\delta _A}\mathrm{e}^{p_2t\delta _B}p_3A\mathrm{}`$ (127)
$``$ $`p_1A+\left(1+tp_1\delta _A+{\displaystyle \frac{t^2}{2}}p_{1}^{}{}_{}{}^{2}\delta _{A}^{}{}_{}{}^{2}\right)p_2B`$
$`+\left[1+t\left(p_1\delta _A+p_2\delta _B\right)+{\displaystyle \frac{t^2}{2}}\left(p_{1}^{}{}_{}{}^{2}\delta _{A}^{}{}_{}{}^{2}+2p_1p_2\delta _A\delta _B+p_{2}^{}{}_{}{}^{2}\delta _{B}^{}{}_{}{}^{2}\right)\right]p_3A+\mathrm{}`$
$`=`$ $`\left(p_1+p_3+p_5\right)A+\left(p_2+p_4+p_6\right)B`$
$`+t\left[p_1p_2\delta _AB+p_2p_3\delta _BA+\left(p_1+p_3\right)p_4\delta _AB+\left(p_2+p_4\right)p_5\delta _BA+\left(p_1+p_3+p_5\right)p_6\delta _AB\right]`$
$`+{\displaystyle \frac{t^2}{2}}[p_1^2p_2\delta _{A}^{}{}_{}{}^{2}B+p_2^2p_3\delta _{B}^{}{}_{}{}^{2}A+2p_1p_2p_3\delta _A\delta _BA+(p_1+p_3)^2p_4\delta _{A}^{}{}_{}{}^{2}B`$
$`+2p_2p_3p_4\delta _B\delta _AB+\left(p_2+p_4\right)^2p_5\delta _{B}^{}{}_{}{}^{2}A+2\left(p_1p_2+p_1p_4+p_3p_4\right)p_5\delta _A\delta _BA`$
$`+(p_1+p_3+p_5)^2p_6\delta _{A}^{}{}_{}{}^{2}B+2(p_2p_3+p_2p_5+p_4p_5)p_6\delta _B\delta _AB].`$
The zeroth-order term with respect to $`t`$ appears only here and hence we have the conditions
$$p_1+p_3+p_5=1\text{and}p_2+p_4+p_6=1.$$
(128)
Using Eq. (128) and the identity $`\delta _BA=\delta _AB`$, we can reduce the right-hand side of Eq. (127) as
$$A+B+t(12q)\delta _AB+\frac{t^2}{2}\left[(1q3r)\delta _{A}^{}{}_{}{}^{2}B+(q3s)\delta _{B}^{}{}_{}{}^{2}A\right]$$
(129)
with
$`q`$ $``$ $`p_2p_3+p_2p_5+p_4p_5,`$ (130)
$`r`$ $``$ $`p_1p_2p_3+p_1p_2p_5+p_1p_4p_5+p_3p_4p_5,`$ (131)
$`s`$ $``$ $`p_2p_3p_4+p_2p_3p_6+p_2p_5p_6+p_4p_5p_6.`$ (132)
For $`k=1`$, we first have
$$\mathrm{e}^{p_1t\delta _A}\mathrm{e}^{p_2t\delta _B}\mathrm{e}^{p_3t\delta _A}\mathrm{e}^{p_4t\delta _B}\mathrm{e}^{p_5t\delta _A}\mathrm{e}^{p_6t\delta _B}1t\delta _{A+B}+\frac{t^2}{2}\left[\delta _{A}^{}{}_{}{}^{2}+\delta _{B}^{}{}_{}{}^{2}+2(1q)\delta _A\delta _B+2q\delta _B\delta _A\right],$$
(133)
where we already used the conditions in Eq. (128). Applying Eq. (133) to Eq. (129) and dropping higher-order terms, we note that the first-order term vanishes and have
$`\left(\mathrm{e}^{p_1t\delta _A}\mathrm{e}^{p_2t\delta _B}\mathrm{e}^{p_3t\delta _A}\mathrm{e}^{p_4t\delta _B}\mathrm{e}^{p_5t\delta _A}\mathrm{e}^{p_6t\delta _B}1\right)\left(p_1A+\mathrm{e}^{p_1t\delta _A}p_2B+\mathrm{e}^{p_1t\delta _A}\mathrm{e}^{p_2t\delta _B}p_3A\mathrm{}\right)`$ (134)
$``$ $`t^2(12q)\delta _{A+B}\delta _AB+{\displaystyle \frac{t^2}{2}}\left[\delta _{A}^{}{}_{}{}^{2}+\delta _{B}^{}{}_{}{}^{2}+2(1q)\delta _A\delta _B+2q\delta _B\delta _A\right](A+B)`$
$`=`$ $`t^2(12q)\left(\delta _{A}^{}{}_{}{}^{2}B\delta _{B}^{}{}_{}{}^{2}A\right)+{\displaystyle \frac{t^2}{2}}\left[\delta _{A}^{}{}_{}{}^{2}B+\delta _{B}^{}{}_{}{}^{2}A2(1q)\delta _{A}^{}{}_{}{}^{2}B2q\delta _{B}^{}{}_{}{}^{2}A\right]`$
$`=`$ $`t^2\left({\displaystyle \frac{1}{2}}q\right)\left(\delta _{A}^{}{}_{}{}^{2}B\delta _{B}^{}{}_{}{}^{2}A\right).`$
The second-order term of $`t`$ in the term $`k=2`$ vanishes just as in Eq. (121). Thus we arrive at
$$\mathrm{\Phi }(x)=x(A+B)+\frac{x^2}{2}(12q)\delta _AB+\frac{x^3}{3!}\left[\left(\frac{1}{2}3r\right)\delta _{A}^{}{}_{}{}^{2}B+\left(\frac{1}{2}3s\right)\delta _{B}^{}{}_{}{}^{2}A\right]+O(x^4).$$
(135)
Putting the second-order and third-order terms to zero, we have a set of simultaneous equations of the parameters as
$`p_1+p_3+p_5`$ $`=`$ $`1,`$ (136)
$`p_2+p_4+p_6`$ $`=`$ $`1,`$ (137)
$`2q=2\left(p_2p_3+p_2p_5+p_4p_5\right)`$ $`=`$ $`1,`$ (138)
$`6r\stackrel{(\text{138})}{=}3\left(p_1+2p_3p_4p_5\right)`$ $`=`$ $`1,`$ (139)
$`6s\stackrel{(\text{138})}{=}3\left(2p_2p_3p_4+p_6\right)`$ $`=`$ $`1.`$ (140)
We can confirm that Ruth’s formula (76), or
$$p_1=\frac{7}{24},p_2=\frac{2}{3},p_3=\frac{3}{4},p_4=\frac{2}{3},p_5=\frac{1}{24},\text{and}p_6=1$$
(141)
is indeed a solution of the above set of simultaneous equations. With six variables for five equations, the solution is in fact a continuous line; Ruth’s solution (141) is just a point on the line. By adjusting the last variable $`p_6`$, we have the continuous solution shown in Fig. 6. (We can solve the set of equations with five parameters by putting $`p_6=0`$, but the solution is complex.)
### 5.6 Example: perturbational composition
We finally present an interesting exercise, motivated by the “perturbational composition” . Suppose that we apply a weak transverse field to an Ising spin. We ask what is the correction term in the exponent of the right-hand side of
$$\mathrm{e}^{\frac{x}{2}\gamma \sigma _x}\mathrm{e}^{x\sigma _z}\mathrm{e}^{\frac{x}{2}\gamma \sigma _x}=\mathrm{e}^{\mathrm{\Phi }(x,\gamma )}=\mathrm{e}^{x\left(\sigma _z+\gamma C_1(x)+O(\gamma ^2)\right)}.$$
(142)
Notice that we expand the exponent with respect to the perturbation parameter $`\gamma `$, not with respect to $`x`$ as in the preceding sections. The first-order perturbation term $`C_1(x)`$ in turn contains higher orders of $`x`$. We could explicitly compute the $`2\times 2`$ matrices on both sides of Eq. (142), expand them with respect to $`\gamma `$ and compare them term by term, but the quantum analysis provides a more elegant way of computation.
We differentiate the both sides of Eq. (142) with respect to $`\gamma `$:
$`{\displaystyle \frac{\mathrm{d}}{\mathrm{d}\gamma }}\mathrm{e}^{x\left(\sigma _z+\gamma C_1(x)+O(\gamma ^2)\right)}`$ $`=`$ $`{\displaystyle \frac{\mathrm{de}^\mathrm{\Phi }}{\mathrm{d}\mathrm{\Phi }}}{\displaystyle \frac{\mathrm{\Phi }(x,\gamma )}{\gamma }}=\mathrm{e}^\mathrm{\Phi }{\displaystyle \frac{1\mathrm{e}^{\delta _\mathrm{\Phi }}}{\delta _\mathrm{\Phi }}}{\displaystyle \frac{\mathrm{\Phi }}{\gamma }},`$ (143)
$`{\displaystyle \frac{\mathrm{d}}{\mathrm{d}\gamma }}\mathrm{e}^{\frac{x}{2}\gamma \sigma _x}\mathrm{e}^{x\sigma _z}\mathrm{e}^{\frac{x}{2}\gamma \sigma _x}`$ $`=`$ $`{\displaystyle \frac{x}{2}}\left(\sigma _x\mathrm{e}^{\frac{x}{2}\gamma \sigma _x}\mathrm{e}^{x\sigma _z}\mathrm{e}^{\frac{x}{2}\gamma \sigma _x}+\mathrm{e}^{\frac{x}{2}\gamma \sigma _x}\mathrm{e}^{x\sigma _z}\mathrm{e}^{\frac{x}{2}\gamma \sigma _x}\sigma _x\right)`$ (144)
$`=`$ $`{\displaystyle \frac{x}{2}}\mathrm{e}^\mathrm{\Phi }\left(\mathrm{e}^\mathrm{\Phi }\sigma _x\mathrm{e}^\mathrm{\Phi }+\sigma _x\right)={\displaystyle \frac{x}{2}}\mathrm{e}^\mathrm{\Phi }\left(\mathrm{e}^{\delta _\mathrm{\Phi }}+1\right)\sigma _x.`$
Equating the both sides, we have
$`{\displaystyle \frac{\mathrm{\Phi }}{\gamma }}=xC_1(x)+O(\gamma )`$ $`=`$ $`{\displaystyle \frac{x}{2}}{\displaystyle \frac{\delta _\mathrm{\Phi }}{1\mathrm{e}^{\delta _\mathrm{\Phi }}}}\left(1+\mathrm{e}^{\delta _\mathrm{\Phi }}\right)\sigma _x`$ (145)
$`=`$ $`{\displaystyle \frac{x}{2}}\left(x\delta _{\sigma _z}+O(\gamma )\right){\displaystyle \frac{1+\mathrm{e}^{\delta _\mathrm{\Phi }}}{1\mathrm{e}^{\delta _\mathrm{\Phi }}}}\sigma _x.`$
Putting $`\gamma =0`$, we have
$$C_1(x)=\frac{1}{2}x\delta _{\sigma _z}\frac{1+\mathrm{exp}\left(x\delta _{\sigma _z}\right)}{1\mathrm{exp}\left(x\delta _{\sigma _z}\right)}\sigma _x=\frac{1}{2}\underset{n=0}{\overset{\mathrm{}}{}}a_nx^n\delta _{\sigma _z}^{}{}_{}{}^{n}\sigma _x$$
(146)
owing to the fact $`\delta _\mathrm{\Phi }=x\delta _{\sigma _z}+O(\gamma )`$, where we have made the Taylor expansion
$$x\frac{1+\mathrm{e}^x}{1\mathrm{e}^x}=\underset{n=0}{\overset{\mathrm{}}{}}a_nx^n$$
(147)
with $`a_0=1`$. We also note that the function (147) is even with respect to $`x`$ and hence $`a_n=0`$ for odd integers $`n`$.
The right-hand side of Eq. (146) is explicitly calculated as follows. In each order, we have
$`\delta _{\sigma _z}\sigma _x`$ $`=`$ $`[\sigma _z,\sigma _x]=2\mathrm{i}\sigma _y,`$ (148)
$`\delta _{\sigma _z}^{}{}_{}{}^{2}\sigma _x`$ $`=`$ $`2\mathrm{i}[\sigma _z,\sigma _y]=4\sigma _x,`$ (149)
$`\delta _{\sigma _z}^{}{}_{}{}^{3}\sigma _x`$ $`=`$ $`4[\sigma _z,\sigma _x]=8\mathrm{i}\sigma _y,`$
$`\mathrm{},`$
or in general,
$$\delta _{\sigma _z}^{}{}_{}{}^{n}\sigma _x=\{\begin{array}{cc}2^n\mathrm{i}\sigma _y\hfill & \text{for odd }n,\hfill \\ 2^n\sigma _x\hfill & \text{for even }n.\hfill \end{array}$$
(151)
We substitute Eq. (151) for each even-order term of the right-hand side of Eq. (146) and arrive at
$$C_1(x)=\frac{1}{2}\underset{n=0}{\overset{\mathrm{}}{}}a_n(2x)^n\sigma _x=x\frac{1+\mathrm{e}^{2x}}{1\mathrm{e}^{2x}}\sigma _x=\left(x\mathrm{coth}x\right)\sigma _x.$$
(152)
(In the second equality, we used the Taylor expansion (147) in the reverse direction.)
In summary, we have
$$\mathrm{e}^{\frac{x}{2}\gamma \sigma _x}\mathrm{e}^{x\sigma _z}\mathrm{e}^{\frac{x}{2}\gamma \sigma _x}=\mathrm{e}^{x\left[\sigma _z+\gamma \left(x\mathrm{coth}x\right)\sigma _x+O(\gamma ^2)\right]}.$$
(153)
The coefficient $`x\mathrm{coth}x`$ behaves as shown in Fig. 6. We have $`x\mathrm{coth}x1`$ for small $`x`$ as is expected, but $`x\mathrm{coth}xx`$ for large $`x`$, and hence the first-order perturbation term grows as $`x^2`$.
## 6 Summary
In the present article, we have reviewed a continual effort on generalization of the Trotter formula to higher-order exponential product formulas. As was emphasized in Sect. 2, the exponential product formula is a good and useful approximant, particularly because it conserves important symmetries of the system dynamics.
We focused on two algorithms of constructing higher-order exponential product formulas. The first is the fractal decomposition, where we construct higher-order formulas recursively. The second is to make use of the quantum analysis, where we compute higher-order correction terms directly. As interludes, we also have described the decomposition of symplectic integrators, the approximation of time-ordered exponentials, and the perturbational composition. It is our hope that the readers find the present article a useful and tutorial “manual” when they numerically investigate dynamical systems. For more practical applications of the exponential product formulas, we refer the readers to the review articles found in Refs. .
## Appendix A Hybrid exponential product formula
We mention here another kind of the exponential product formula . Consider the Trotter approximant
$$\mathrm{e}^{xA}\mathrm{e}^{xB}=\mathrm{e}^{x(A+B)+\frac{1}{2}x^2[A,B]+O(x^3)}.$$
(154)
We can cancel out the second-order correction term in the form
$$\mathrm{e}^{xA}\mathrm{e}^{xB}\mathrm{e}^{\frac{1}{2}x^2[A,B]}=\mathrm{e}^{x(A+B)+O(x^3)}.$$
(155)
If, in some problems, the commutation relation $`[A,B]`$ is easily diagonalized, Eq. (155) may be a useful approximant.
A more complicated one is the fourth-order approximant
$$\mathrm{e}^{\frac{x^3}{432}[B,[A,B]]}S_a\left(\frac{x}{3}\right)S_b\left(\frac{x}{3}\right)S_a\left(\frac{x}{3}\right)\mathrm{e}^{\frac{x^3}{432}[B,[A,B]]}=\mathrm{e}^{x(A+B)+O(x^5)},$$
(156)
where
$$S_a(x)\mathrm{e}^{\frac{1}{2}xA}\mathrm{e}^{xB}\mathrm{e}^{\frac{1}{2}xA}\text{and}S_b(x)\mathrm{e}^{\frac{1}{2}xB}\mathrm{e}^{xA}\mathrm{e}^{\frac{1}{2}xB}.$$
(157)
In fact, the diffusion equation is described by
$$A=\frac{1}{2}\mathrm{\Delta }\text{and}B=V(\stackrel{}{q})$$
(158)
and we have
$$[B,[A,B]]=\left(V(\stackrel{}{q})\right)^2.$$
(159)
The above type of the exponential product formula was referred to as the hybrid exponential product formula. We do not give its details in this article, since commutation relations are not easily diagonalized except for a few specific problems.
## Appendix B World-line quantum Monte Carlo method
In the present appendix, we give a short review of the world-line quantum Monte Carlo method. The world-line quantum Monte Carlo method is to transform the partition function (10) of a quantum system $``$ into the partition function of a classical system by means of the path-integral representation and simulate the latter system. We explain the method using the transverse Ising model (11), or $`=A+B`$ with
$$A=\underset{i,j}{}J_{ij}\sigma _i^z\sigma _j^z\text{and}B=\mathrm{\Gamma }\underset{i}{}\sigma _i^x.$$
(160)
The starting point is the Trotter decomposition (21) of the partition function, namely the Suzuki-Trotter transformation , of the form:
$`Z`$ $`=`$ $`Tr\mathrm{e}^\beta =\underset{n\mathrm{}}{lim}Tr\left(\mathrm{e}^{\frac{\beta }{n}A}\mathrm{e}^{\frac{\beta }{n}B}\right)^n=\underset{n\mathrm{}}{lim}{\displaystyle \underset{\left\{\sigma _i\right\}}{}}\left\{\sigma _i^{(0)}\right\}\left|\left(\mathrm{e}^{\frac{\beta }{n}A}\mathrm{e}^{\frac{\beta }{n}B}\right)^n\right|\left\{\sigma _i^{(0)}\right\}`$ (161)
$`=`$ $`\underset{n\mathrm{}}{lim}{\displaystyle \underset{\left\{\sigma _i^{(0)}\right\}}{}}\left\{\sigma _i^{(0)}\right\}\left|\mathrm{e}^{\frac{\beta }{n}A}\mathrm{e}^{\frac{\beta }{n}B}\mathrm{e}^{\frac{\beta }{n}A}\mathrm{e}^{\frac{\beta }{n}B}\mathrm{}\mathrm{e}^{\frac{\beta }{n}B}\right|\left\{\sigma _i^{(0)}\right\}.`$
In the second line, we have taken the trace with respect to a complete basis set by using the spin $`z`$ axis as the quantization axis:
$$\sigma _k^z|\left\{\sigma _i^{(0)}\right\}=\sigma _k^{(0)}|\left\{\sigma _i^{(0)}\right\},$$
(162)
where the eigenvalue is $`\sigma _k^{(0)}=\pm 1`$. The meaning of the superscript $`(0)`$ becomes self-evident just below. We further insert the resolution of unity in between each pair of the exponential operators in the last line of Eq. (161), obtaining
$`Z=\underset{n\mathrm{}}{lim}{\displaystyle \underset{\left\{\sigma _i^{(0)}\right\}}{}}{\displaystyle \underset{\left\{\sigma _i^{(1)}\right\}}{}}{\displaystyle \underset{\left\{\sigma _i^{(2)}\right\}}{}}\mathrm{}{\displaystyle \underset{\left\{\sigma _i^{(n1)}\right\}}{}}\left\{\sigma _i^{(0)}\right\}\left|\mathrm{e}^{\frac{\beta }{n}A}\right|\left\{\sigma _i^{(0)}\right\}\left\{\sigma _i^{(0)}\right\}\left|\mathrm{e}^{\frac{\beta }{n}B}\right|\left\{\sigma _i^{(1)}\right\}`$ (163)
$`\times \left\{\sigma _i^{(1)}\right\}\left|\mathrm{e}^{\frac{\beta }{n}A}\right|\left\{\sigma _i^{(1)}\right\}\left\{\sigma _i^{(1)}\right\}\left|\mathrm{e}^{\frac{\beta }{n}B}\right|\left\{\sigma _i^{(2)}\right\}\mathrm{}\left\{\sigma _i^{(n1)}\right\}\left|\mathrm{e}^{\frac{\beta }{n}B}\right|\left\{\sigma _i^{(0)}\right\}.`$
In the above expression, we used the fact that the operator $`A`$ is diagonal in the representation of $`\{\sigma _i^{(m)}\}`$ and hence made the complete set on the both sides of each operator $`\mathrm{e}^{\frac{\beta }{n}A}`$ identical. In contrast, the operator $`\mathrm{e}^{\frac{\beta }{n}B}`$ has off-diagonal elements.
Let us calculate the matrix elements in Eq. (163). The matrix element of the operator $`\mathrm{e}^{\frac{\beta }{n}A}`$ is easy:
$$\left\{\sigma _i^{(m)}\right\}\left|\mathrm{e}^{\frac{\beta }{n}A}\right|\left\{\sigma _i^{(m)}\right\}=\mathrm{exp}\left(\frac{\beta }{n}\underset{i,j}{}J_{ij}\sigma _i^{(m)}\sigma _j^{(m)}\right).$$
(164)
This is because the operators $`\left\{\sigma _i^z\right\}`$ are all diagonal in the present representation as in Eq. (162). On the other hand, the operator $`\mathrm{e}^{\frac{\beta }{n}B}`$ has off-diagonal elements as well as diagonal ones in the following form:
$$\left\{\sigma _i^{(m)}\right\}\left|\mathrm{e}^{\frac{\beta }{n}B}\right|\left\{\sigma _i^{(m+1)}\right\}=\underset{i}{}\sigma _i^{(m)}\left|\mathrm{e}^{\frac{\beta \mathrm{\Gamma }}{n}\sigma _i^x}\right|\sigma _i^{(m+1)}$$
(165)
with each matrix element given by
$$\begin{array}{c}\\ \hfill \sigma _i^{(m)}\left|\mathrm{e}^{\frac{\beta \mathrm{\Gamma }}{n}\sigma _i^x}\right|\sigma _i^{(m+1)}=\end{array}\begin{array}{cc}& \begin{array}{cc}|\sigma _i^{(m+1)}=+1& |\sigma _i^{(m+1)}=1\end{array}\\ \begin{array}{c}\sigma _i^{(m)}=+1|\\ \\ \sigma _i^{(m)}=1|\end{array}& \left(\begin{array}{cc}\mathrm{cosh}\frac{\beta \mathrm{\Gamma }}{n}& \mathrm{sinh}\frac{\beta \mathrm{\Gamma }}{n}\\ & \\ \mathrm{sinh}\frac{\beta \mathrm{\Gamma }}{n}& \mathrm{cosh}\frac{\beta \mathrm{\Gamma }}{n}\end{array}\right)\end{array}.$$
(166)
These matrix elements are expressed in a single equation
$$\sigma _i^{(m)}\left|\mathrm{e}^{\frac{\beta \mathrm{\Gamma }}{n}\sigma _i^x}\right|\sigma _i^{(m+1)}=\mathrm{exp}\left(\gamma _n\sigma _i^{(m)}\sigma _i^{(m+1)}+\delta _n\right),$$
(167)
where the parameters $`\gamma _n`$ and $`\delta _n`$ are defined in
$$\mathrm{e}^{\gamma _n+\delta _n}=\mathrm{cosh}\frac{\beta \mathrm{\Gamma }}{n}\text{and}\mathrm{e}^{\gamma _n+\delta _n}=\mathrm{sinh}\frac{\beta \mathrm{\Gamma }}{n},$$
(168)
or more explicitly defined by
$$\gamma _n=\frac{1}{2}\mathrm{log}\mathrm{tanh}\frac{\beta \mathrm{\Gamma }}{n}\text{and}\delta _n=\frac{1}{2}\mathrm{log}\frac{1}{2}\mathrm{sinh}\frac{2\beta \mathrm{\Gamma }}{n}.$$
(169)
The expressions (164) and (167) give the partition function (163) in the form
$$Z=\underset{n\mathrm{}}{lim}\underset{\left\{\sigma _i^{(m)}\right\}}{}\mathrm{e}^{\beta _n}$$
(170)
with the resulting classical Hamiltonian
$$\beta _n\frac{\beta }{n}\underset{m=0}{\overset{n1}{}}\underset{i,j}{}J_{ij}\sigma _i^{(m)}\sigma _j^{(m)}+\gamma _n\underset{m=0}{\overset{n1}{}}\underset{i}{}\sigma _i^{(m)}\sigma _i^{(m+1)},$$
(171)
where we dropped a constant term due to $`\delta _n`$. Note that the periodic boundary conditions, $`\sigma _i^{(n)}\sigma _i^{(0)}`$, must be required in the second term of Eq. (171) because the trace operation in Eq. (161) demands it.
The classical Hamiltonian (171) is interpreted as follows (Fig. 7). Suppose that the original quantum system (160) is defined on a square lattice. The first term of Eq. (171) indicates that the two-dimensional system is replicated into $`n`$ layers with the intra-layer interaction reduced by $`n`$ times. The second term of Eq. (171) represents the inter-layer interactions. The coupling is $`\gamma _n/\beta `$ as defined in Eq. (171). Thus the quantum system on a square lattice is mapped to an Ising model on a cubic lattice. In general, a $`d`$-dimensional quantum system is mapped to a $`(d+1)`$-dimensional classical system. The additional axis is called the Trotter direction. The physical quantities of the quantum system can be estimated by Monte Carlo simulation of the mapped classical system. This is the basic idea of the world-line quantum Monte Carlo method .
We can use this mapping in order to study the quantum annealing . Suppose that we look for the ground state of the diagonal part $`A`$ of the system (160). Random exchange interactions $`\{J_{ij}\}`$ may produce many local minima that are only slightly above the ground state in energy but far apart from the ground state in the phase space. The simulated annealing, a well-known method of ground-state search, is often trapped in a local minima and does not reach the ground state. In quantum annealing, we use the transverse field $`\mathrm{\Gamma }`$ in order to induce tunneling from local minima to the ground state. We first apply the off-diagonal part $`B`$ of Eq. (160) strongly and turn it off gradually, hoping to end up with the ground state of the diagonal part $`A`$. This corresponds to a Monte Carlo simulation of the mapped classical system (171) with the intra-layer coupling $`\gamma _n`$ being infinitesimally weak at the beginning and infinitely strong at the end. Each layer of the system (171) is first independent of each other and is gradually frozen into an identical configuration, which we hope is the ground state.
An annoying problem inherent in the algorithm of the quantum Monte Carlo method is the systematic error due to the finite Trotter number $`n`$. It used to be that simulations were carried out for various finite values of $`n`$, quantities were estimated in each simulation, and then the limit $`n\mathrm{}`$ was taken in the process of the data analysis, which was called the Trotter extrapolation. A recent development of the quantum Monte Carlo method dramatically changed the situation. We here mention the development briefly; see Ref. for a tutorial and exhaustive review of the topic.
In the most recent quantum Monte Carlo algorithm, it is possible for some systems to take the Trotter limit before we set up the classical system for simulation. Taking the Trotter limit $`n\mathrm{}`$, we have a continuum Trotter axis (Fig. 9). (Note again that the boundary conditions are required in the Trotter direction.) The interaction is described as follows (Fig. 9). Instead of Ising spins on lattice points of a Trotter axis, we have up-spin domains and down-spin domains on the axis. Instead of intra-layer interactions between a pair of nearest-neighbor spins, we have parallel-spin areas and anti-parallel-spin areas. In Monte Carlo simulation, we update the up-spin domains and down-spin domains on the basis of the energy of the parallel-spin areas and anti-parallel-spin areas.
It is thus possible in such situations to carry out a simulation in the Trotter limit $`n\mathrm{}`$. Monte Carlo estimates of such a simulation are free of the systematic error of the order $`\beta ^2/n`$ in Eq. (21), and hence do not need the higher-order exponential product formula for such systems.
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# Gigantic transmission band edge resonance in periodic stacks of anisotropic layers
## 1 Introduction
The subject of this paper is the Fabry-Perot cavity resonance in periodic layered structures. This phenomenon, known for many decades, is also referred to as the *transmission band edge resonance*, because it occurs in the vicinity of photonic band edge frequencies in finite photonic crystals and is accompanied by sharp transmission peaks. This basic effect can occur in a finite periodic array of any two alternating materials with different refractive indices. It can also be realized in a waveguide environment, or in a finite array of coupled resonators. The only essential requirements are: (a) low absorption, (b) the appropriate number $`N`$ of unit cells in the periodic array, and (c) the presence of a frequency gap (a stop-band) in the frequency spectrum of the periodic array. The latter is a universal property of almost any lossless spatially periodic structure. This phenomenon has found numerous and diverse practical applications in optics. Below we briefly outline some basic features of the classical band edge resonance in periodic stacks of isotropic layers. We only highlight those points, which are necessary for the following comparative analysis of stacks involving anisotropic layers. More detailed information on the relevant aspects of electrodynamics of layered dielectric media can be found in an extensive literature on the subject (see, for example, , and references therein).
Consider a simplest periodic array of alternating dielectric layers $`A`$ and $`B`$, as shown in Fig. 1(a). The layers are made of transparent isotropic materials with different refractive indices $`n_A`$ and $`n_B`$. Such periodic stacks are also referred to as a $`1D`$ photonic crystal. The electromagnetic eigenmodes of the periodic structure in Fig. 1(a) are Bloch waves with a typical wave number/frequency diagram shown in Fig. 1(b).
Consider now a finite periodic stack composed of $`N`$ unit cells $`L`$ in Fig. 1(a). Such a stack is commonly referred to as a Fabry-Perot cavity. If the number $`N`$ of the double layers $`L`$ is significant, the stack periodicity causes coherent interference of light scattered by the layers interfaces. In practice, periodic stacks having as few as several periods $`L`$, can display almost total reflectivity at the band gap frequencies, provided that the refractive indices $`n_A`$ and $`n_B`$ of the adjacent layers differ significantly.
A typical frequency dependence of the finite stack transmittance is presented in Fig. 2. The frequency range shown includes the vicinity of the photonic band edge (BE) $`g`$ in Fig. 1(b). The sharp transmission peaks below the photonic band edge frequency $`\omega _g`$ correspond to transmission band edge resonances, also known as Fabry-Perot cavity resonances. At each resonance, the electromagnetic field inside the periodic stack is close to a standing wave composed of a forward and a backward Bloch eigenmodes with large and nearly equal amplitudes. The slab boundaries coincide with the standing wave nodes, where the forward and backward Bloch components interfere destructively, as illustrated in Figs. 3 and 4. The latter circumstance determines the wave numbers of the forward and backward Bloch components at the resonance frequencies
$$k_sk_g\pm \frac{\pi }{NL}s,s=1,2,\mathrm{},$$
(1)
where $`s`$ is the order number of the resonant peak in Fig. 2, and $`k_g`$ is the wave number corresponding to the photonic band edge. In our case, $`k_g=\pi /L`$. The resonance frequencies themselves can be expressed in terms of the dispersion relation $`\omega \left(k\right)`$ of the respective frequency band. Indeed, just below the photonic band edge $`g`$ in Fig. 1(b), the dispersion relation $`\omega \left(k\right)`$ can be approximated by the quadratic parabola
$$\omega \omega _g+\frac{\omega _g^{\prime \prime }}{2}\left(kk_g\right)^2,\text{ where }\omega _g^{\prime \prime }=\left(\frac{^2\omega }{k^2}\right)_{k=k_g}<0.$$
(2)
This relation together with (1) yield the frequencies of the resonant transmission peaks in Fig. 2
$$\omega _s\left(N\right)\omega _g+\frac{\omega _g^{\prime \prime }}{2}\left(\frac{\pi }{NL}s\right)^2,s=1,2,\mathrm{}$$
(3)
where $`\omega _g=\omega \left(k_g\right)`$ is the band edge frequency. For example, the transmission peak $`1`$ closest to the photonic band edge is located at
$$\omega _1\left(N\right)\omega _g+\frac{\omega _g^{\prime \prime }}{2}\left(\frac{\pi }{NL}\right)^2.$$
(4)
The dependence (3) is illustrated in Fig. 2.
Electromagnetic field distribution inside the stack at the frequency of band edge resonance is shown in Figs. 3 and 4 for the first two transmission resonances, respectively. For a given amplitude $`\mathrm{\Psi }_I`$ of the incident wave, the maximal field intensity $`\left|\mathrm{\Psi }\left(z\right)\right|^2`$ inside the slab depends on the number $`N`$ of the double layers in the stack and on the order number $`s`$ of the resonance peak in Fig. 2
$$\mathrm{max}\left|\mathrm{\Psi }\left(z\right)\right|^2\left|\mathrm{\Psi }_I\right|^2\left(\frac{N}{s}\right)^2.$$
(5)
The maximal field intensity (5) is proportional to squared thickness of the slab and, for a large $`N`$, is greatly enhanced compared to that of the incident light. By contrast, at the slab boundaries at $`z=0`$ and $`z=D=NL`$, the field amplitude $`\mathrm{\Psi }\left(z\right)`$ always remains comparable to $`\mathrm{\Psi }_I`$ to satisfy the electromagnetic boundary conditions (69). At frequencies outside the resonance transparency peaks, the field amplitude inside the stack drops sharply.
The exact definition of the physical values plotted in Figs. 3 and 4, as well as the numerical parameters of the periodic stacks used to generates these plots, are given in the Appendix.
The expressions (1) through ((5) are valid if $`N1`$ and only apply to the transmission resonances close enough to the photonic band edge. In further consideration, we will focus on the most powerful first resonance $`s=1`$, closest to the photonic band edge.
Above, we outlined some basic features of the transmission band edge resonance in finite stacks of isotropic layers. The question we would like to address in this paper is whether the presence of anisotropic layers in a finite periodic stack can qualitatively change the nature of the Fabry-Perot cavity resonance. We will show that, indeed, in periodic stacks involving anisotropic layers, the transmission resonance can be significantly stronger, compared to what is achievable with common periodic stacks of isotropic layers. For instance, in the periodic stack shown in Fig. 5, the field intensity associated with the transmission band edge resonances can be proportional to $`N^4`$, rather than $`N^2`$. The latter implies that a stack of $`N`$ anisotropic layers can perform as well as a common stack of $`N^2`$ isotropic layers. And this is a huge difference! The physical reason for this is that periodic stacks of anisotropic layers can support the kind of $`k\omega `$ diagrams that are impossible in stacks of isotropic layers. Specifically, a dispersion curve $`\omega \left(k\right)`$ of the stack in Fig. 5 can develop a degenerate band edge, as shown in Fig. 6(b). Just below the degenerate band edge $`d`$, the dispersion curve can be approximated as
$$\omega \omega _d+\frac{\omega _d^{\prime \prime \prime \prime }}{24}\left(kk_d\right)^4,\text{where }\omega _d^{\prime \prime \prime \prime }=\left(\frac{^4\omega }{k^4}\right)_{k=k_d}<0,$$
(6)
which implies a huge density of modes. And this is what makes all the difference compared to the case (2) of a regular band edge.
For convenience, the domain of definition of the Bloch wave number $`k`$ in Fig. 6 is chosen between $`0`$ and $`2\pi /L`$. Since the Bloch wave number is defined up to a multiple of $`2\pi /L`$, the representation in Fig. 6 is equivalent to that in Fig. 1(b). Note that the points $`a`$, $`g`$, and $`d`$ in Fig. 6(a) and 6(b) lie at the Brillouin zone boundary at $`k=\pi /L`$.
In Fig. 7 we present the transmission dispersion of the finite periodic stacks in Fig. 5 composed of 32 unit cells $`L`$ and having the $`k\omega `$ diagram in Fig. 6(b). The frequency range shown includes the degenerate band edge (DBE) at $`\omega =\omega _d`$. The field intensity distribution at the frequency of the first transmission resonance 1 is shown in Fig. 8(b). One can see that for a given $`N`$, the resonant field intensity in Fig. 8 is significantly larger compared to that of the vicinity of a regular band edge, shown in Fig. 3. Specifically, in the case of degenerate band edge
$$\mathrm{max}\left|\mathrm{\Psi }\left(z\right)\right|^2\left|\mathrm{\Psi }_I\right|^2\left(\frac{N}{s}\right)^4,$$
(7)
compared to the estimation (5) related to a regular band edge. The transmission bandwidth in the case of DBE appears to be much smaller.
In practice, the field amplitude associated with the transmission resonance is limited not only by the number of layers in the stack, but also by such factors as absorption, nonlinearity, imperfections of the periodic array, stack dimensions in the $`XY`$ plane, incident radiation bandwidth, etc. All else being equal, the stack with degenerate band edge can have much fewer layers and, therefore, can be much thinner compared to a regular stack of isotropic layers with similar performance. Much smaller dimensions can be very attractive for a variety of practical applications. Similar effect associated with degenerate band edge (6) of the $`k\omega `$ diagram can also be achieved in the waveguide environment, as well as in finite periodic arrays of coupled multiple-mode resonators. On the down side, the realization of $`k\omega `$ diagram having a dispersion curve with DBE, requires more sophisticated periodic arrays, such as the one shown in Fig. 5.
The rest of the paper is organized as follows.
In Section 2 we briefly outline the electrodynamics of periodic stratified media and examine the relation between anisotropy of the layers and the $`k\omega `$ diagram of the stack. We show that at the frequency of degenerate band edge $`d`$, the $`4\times 4`$ transfer matrix of a unit cell $`L`$ cannot be diagonalized or even reduced to the block-diagonal form. Based on this, we establish necessary symmetry conditions for a periodic stack to develop a degenerate band edge (6) and to display the peculiar resonance properties associated with it. We prove that such periodic stacks must have at least two misaligned anisotropic layers in a unit cell, as shown in the example in Fig. 5.
In Section 3, we consider the scattering problem for a finite periodic stack. We analyze the eigenmode composition of electromagnetic field at the frequency of transmission resonance near degenerate band edge. We show that in contrast to the case (5) of a regular band edge, in the degenerate band edge case (7) the resonance field inside the stack does not reduce to a superposition of forward and backward propagating eigenmodes. Instead, the contribution of evanescent eigenmodes becomes equally important and leads to much stronger dependence (7) of the resonance field intensity on the number of layers in the stack.
The physical and geometrical parameters of stacks used for numerical simulations are specified in the Appendix.
## 2 Electrodynamics of periodic stacks of anisotropic layers
This section starts with a brief description of some basic electrodynamic properties of periodic layered media composed of lossless anisotropic layers. Then we turn to the particular case of periodic stacks with degenerate band edge. The scattering problem for periodic finite stacks, including the Fabry-Perot cavity resonance in the vicinity of degenerate photonic band edge will be considered in the next section.
### 2.1 Transverse electromagnetic waves in stratified media
Our consideration is based on time-harmonic Maxwell equations in heterogeneous nonconducting media
$$\times \stackrel{}{𝐄}\left(\stackrel{}{r}\right)=i\frac{\omega }{c}𝐁\left(\stackrel{}{r}\right),\times \stackrel{}{𝐇}\left(\stackrel{}{r}\right)=i\frac{\omega }{c}𝐃\left(\stackrel{}{r}\right),$$
(8)
where electric and magnetic fields and inductions are related by linear constitutive equations
$$\stackrel{}{𝐃}\left(\stackrel{}{r}\right)=\widehat{\epsilon }\left(\stackrel{}{r}\right)\stackrel{}{𝐄}\left(\stackrel{}{r}\right),\stackrel{}{𝐁}\left(\stackrel{}{r}\right)=\widehat{\mu }\left(\stackrel{}{r}\right)\stackrel{}{𝐇}\left(\stackrel{}{r}\right).$$
(9)
In further consideration we assume that:
1. The direction of plane wave propagation coincide with the normal $`z`$ to the layers.
2. The second rank tensors $`\widehat{\epsilon }\left(\stackrel{}{r}\right)`$ and $`\widehat{\mu }\left(\stackrel{}{r}\right)`$ are dependent on a single Cartesian coordinate $`z`$, normal to the layers.
3. The $`z`$ direction is a two-fold symmetry axis of the stack, implying that the anisotropy axes of individual layers are either parallel, or perpendicular the $`z`$ direction. The latter defines the case of in-plane anisotropy.
Under the above restrictions, the normal field components $`E_z`$ and $`H_z`$ of electromagnetic wave are zeroes, and the system (8) of six time-harmonic Maxwell equations reduces to the following system of four ordinary linear differential equations for the transverse field components
$$\frac{}{z}\mathrm{\Psi }\left(z\right)=i\frac{\omega }{c}M\left(z\right)\mathrm{\Psi }\left(z\right),\text{ where }\mathrm{\Psi }\left(z\right)=\left[\begin{array}{c}E_x\left(z\right)\\ E_y\left(z\right)\\ H_x\left(z\right)\\ H_y\left(z\right)\end{array}\right].$$
(10)
The $`4\times 4`$ matrix $`M\left(z\right)`$ in (10) is referred to as the (reduced) Maxwell operator.
In a lossless nonmagnetic medium with in-plane anisotropy, the electric permittivity and magnetic permeability tensors have the form
$$\widehat{\epsilon }=\left[\begin{array}{ccc}\epsilon _{xx}& \epsilon _{xy}& 0\\ \epsilon _{xy}& \epsilon _{yy}& 0\\ 0& 0& \epsilon _{zz}\end{array}\right],\widehat{\mu }=\widehat{1}.$$
(11)
This yields the following explicit expression for the Maxwell operator $`M\left(z\right)`$ in (10)
$$M\left(z\right)=\left[\begin{array}{cccc}0& 0& 0& 1\\ 0& 0& 1& 0\\ \epsilon _{xy}& \epsilon _{yy}& 0& 0\\ \epsilon _{xx}& \epsilon _{xy}& 0& 0\end{array}\right],$$
(12)
where the components of the permittivity tensor may vary from layer to layer.
#### 2.1.1 The transfer matrix formalism
The Cauchy problem
$$\frac{}{z}\mathrm{\Psi }\left(z\right)=i\frac{\omega }{c}M\left(z\right)\mathrm{\Psi }\left(z\right),\mathrm{\Psi }\left(z_0\right)=\mathrm{\Psi }_0$$
(13)
for the reduced Maxwell equation (10) has a unique solution
$$\mathrm{\Psi }\left(z\right)=T(z,z_0)\mathrm{\Psi }\left(z_0\right),$$
(14)
where the $`4\times 4`$ matrix $`T(z,z_0)`$ is referred to as the *transfer matrix*. The transfer matrix (14) uniquely relates the values of time-harmonic electromagnetic field $`\mathrm{\Psi }`$ at any two points $`z`$ and $`z_0`$ of the stratified medium. From the definition (14), it follows that
$$T(z,z_0)=T(z,z^{})T(z^{},z_0),T(z,z_0)=T^1(z_0,z),T(z,z)=I.$$
(15)
The transfer matrix of a stack of layers is defined as
$$T_S=T(D,0),$$
where $`z=0`$ and $`z=D`$ are the stack boundaries. The greatest advantage of the transfer matrix formalism stems from the fact that the transfer matrix of an arbitrary stack is a sequential product of the transfer matrices $`T_m`$ of the constitutive layers
$$T_S=\underset{m}{}T_m.$$
(16)
If the individual layers $`m`$ are homogeneous, the corresponding single-layer transfer matrices $`T_m`$ can be explicitly expressed in terms of the respective Maxwell operators $`M_m`$
$$T_m=\mathrm{exp}\left(iD_mM_m\right),$$
(17)
where $`D_m`$ is the thickness of the $`m`$-th layer. The explicit expression for the Maxwell operator $`M_m`$ of a uniform dielectric layer with in-plane anisotropy is given by Eq. (12). Thus, Eq. (16) together with (17) and (12) provide an explicit analytical expression for the transfer matrix $`T_S`$ of a stack of dielectric layers with in-plane anisotropy.
The $`4\times 4`$ transfer matrix of an arbitrary lossless stratified medium displays the fundamental property of $`J`$-unitarity
$$T^{}=JT^1J,\text{ where }J=\left[\begin{array}{cccc}0& 0& 0& 1\\ 0& 0& 1& 0\\ 0& 1& 0& 0\\ 1& 0& 0& 0\end{array}\right],$$
(18)
which implies, in particular, that
$$\left|detT\right|=1.$$
(19)
Different versions of the $`4\times 4`$ transfer matrix formalism have been used in electrodynamics of stratified media composed of birefringent and/or gyrotropic layers for decades (see, for example, and references therein). In this paper we use exactly the same notations and terminology as in our previous publications on electrodynamics of stratified media.
### 2.2 Eigenmodes in periodic layered media
In a periodic layered medium, all material tensors are periodic functions of $`z`$, and so is the $`4\times 4`$ matrix $`M(z)`$ in (10). Usually, the solutions $`\mathrm{\Psi }_k\left(z\right)`$ of the reduced Maxwell equation (10) with the periodic $`M(z)`$ can be chosen in the Bloch form
$$\mathrm{\Psi }_k\left(z+L\right)=e^{ikL}\mathrm{\Psi }_k\left(z\right),$$
(20)
where the Bloch wave number $`k`$ is defined up to a multiple of $`2\pi /L`$. The definition (14) of the transfer matrix together with Eq. (20) yield
$$T(z+L,z)\mathrm{\Psi }_k\left(z\right)=e^{ikL}\mathrm{\Psi }_k\left(z\right).$$
(21)
Introducing the transfer matrix of a unit cell $`L`$
$$T_L=T(L,0),$$
(22)
we have from Eq. (21)
$$T_L\mathrm{\Phi }_k=e^{ikL}\mathrm{\Phi }_k,\text{ where }\mathrm{\Phi }_k=\mathrm{\Psi }_k\left(0\right).$$
(23)
Thus, the eigenvectors of the transfer matrix $`T_L`$ of the unit cell are uniquely related to the Bloch solutions $`\mathrm{\Psi }_k\left(z\right)`$ of the reduced Maxwell equation (10)
$$\mathrm{\Phi }_i=\mathrm{\Psi }_i\left(0\right),i=1,2,3,4.$$
(24)
The respective four eigenvalues
$$X_i=e^{ik_iL},i=1,2,3,4$$
(25)
of $`T_L`$ are the roots of the characteristic polynomial $`F_4\left(X\right)`$ of the forth degree
$$F_4\left(X\right)=det\left(T_LXI\right)=0.$$
(26)
Unit cell of a periodic stack can be chosen differently. For example, the choice $`A_1BA_2`$ specified in Fig. 5 is as good as $`A_2A_1B`$. Different choice of a unit cell corresponds to its shift in the $`z`$ direction and results in the following transformation of the respective transfer matrix $`T_L`$
$$T_L^{}=T(0,Z)T_LT(Z,0)=\left[T(Z,0)\right]^1T_LT(Z,0),$$
(27)
where $`Z`$ is the amount of the shift. The modified transfer matrix is similar to the original one and has the same set of eigenvalues.
For any given $`\omega `$, the characteristic equation (26) defines a set of four eigenvalues (25). Real $`k`$ (or, equivalently, $`\left|X\right|=1`$) correspond to propagating Bloch modes, while complex $`k`$ (or, equivalently, $`\left|X\right|1`$) correspond to evanescent modes.
The $`J`$-unitarity (18) of $`T_L`$ imposes the following restriction on its eigenvalues (25)
$$\left\{X_i^1\right\}\left\{X_i^{}\right\},i=1,2,3,4$$
(28)
or, equivalently
$$\{k_i\}\{k_i^{}\},i=1,2,3,4,$$
(29)
for any given $`\omega `$. In view of the relations (29) or (28), one can distinguish the following three different situations.
1. All four wave numbers are real
$$k_1k_1^{},k_2k_2^{},k_3k_3^{},k_4k_4^{}.$$
(30)
In the example in Fig. 6, this case relates to the frequency range
$$0<\omega <\omega _a.$$
(31)
In this case, all four Bloch eigenmodes are propagating.
2. Two wave numbers are real and the other two are complex
$$k_1=k_1^{},k_2=k_2^{},k_4=k_3^{},\text{where }k_3k_3^{},k_4k_4^{}.$$
(32)
This case relates to the frequency range
$$\omega _a<\omega <\omega _g,$$
(33)
in Fig. 6(a), or the frequency range
$$\omega _a<\omega <\omega _d,$$
(34)
in Fig. 6(b). In both cases (33) and (34), two of the four Bloch eigenmodes are propagating and the remaining two are evanescent with complex conjugated wave numbers.
3. All four wave numbers are complex
$$k_2=k_1^{},k_4=k_3^{},\text{where }k_1k_1^{},k_2k_2^{},k_3k_3^{},k_4k_4^{}.$$
(35)
This situation relates to a frequency gap, where all four Bloch eigenmodes are evanescent. In the example in Fig. 6, this case corresponds to
$$\omega _g<\omega ,\text{ or }\omega _d<\omega .$$
(36)
Notice that the case (B) of two propagating and two evanescent modes can only occur in periodic stacks of anisotropic layers. If all the layers in a unit cell are isotropic, the four Bloch eigenmodes are either all propagating or all evanescent, as is the case in Fig. 1. For a given frequency $`\omega `$, the four Bloch eigenmodes correspond to two different polarizations and two opposite directions of propagation.
#### 2.2.1 Non-Bloch solutions at stationary points of the $`k\omega `$ diagram
So far, we have considered only the cases where all four solutions for the reduced Maxwell equation (10) can be chosen in the Bloch form (20). All such cases fall into one of the following three categories: (A) all four eigenmodes are propagating, (B) two modes are propagating and the other two are evanescent, (C) all four eigenmodes are evanescent. This classification of the eigenmodes does not apply at the frequencies of stationary points on the $`k\omega `$ diagram, where the group velocity $`u`$ of some of the propagating modes vanishes
$$u=d\omega /dk=0.$$
(37)
At a stationary point of a dispersion curve, not all four solutions of the Maxwell equation (10) are Bloch waves, as defined in (20). Instead, some of the solutions can be algebraically diverging non-Bloch eigenmodes. Such eigenmodes can be essential for understanding the resonance effects in finite and semi-infinite periodic arrays.
Let us consider the situation at stationary points (37) on the $`k\omega `$ diagram in terms of the matrix $`T_L`$. Although at any given frequency $`\omega `$, the reduced Maxwell equation (10) has exactly four linearly independent solutions, it does not imply that the respective transfer matrix $`T_L`$ in (23) must have four linearly independent eigenvectors (23). Indeed, although the matrix $`T_L`$ is invertible, it is neither Hermitian, nor unitary and, therefore, may not be diagonalizable. Specifically, if the frequency approaches one of the stationary points (37), some of the eigenvectors $`\mathrm{\Phi }_k`$ in (23) become nearly parallel to each other. Eventually, as $`\omega `$ reaches the the stationary point value, the number of linearly independent eigenvectors $`\mathrm{\Phi }_k`$ becomes lesser than four, and the relation (24) does not apply at that point.
Examples of different stationary points (37) are shown in Fig. 6. Using these examples, let us take a closer look at these special frequencies.
At the frequencies $`\omega _g`$ of the photonic band edge $`g`$ in Fig. 6(a), the four eigenmodes include:
* one propagating mode with $`k=\pi /L`$ and zero group velocity,
* one non-Bloch eigenmode linearly diverging with $`z`$,
* a pair of evanescent modes with equal and opposite imaginary wave numbers.
At the frequencies $`\omega _a`$ corresponding to the point $`a`$ in Fig. 6, the four solutions of Eq. (10) include:
* the propagating mode with $`k=\pi /L`$ and zero group velocity,
* one non-Bloch eigenmode linearly diverging with $`z`$,
* a pair of propagating modes having equal and opposite group velocities and belonging to the dispersion curve other than the one containing the point $`a`$.
Of special interest here is the frequency $`\omega _d`$ of degenerate photonic band edge $`d`$ in Fig. 6(b). In this case, the four solutions of Eq. (10) include :
* the propagating mode with $`k=\pi /L`$ and zero group velocity,
* three non-Bloch eigenmodes diverging as $`z`$, $`z^2`$, and $`z^3`$, respectively.
At any particular frequency, the existence of non-Bloch eigenmodes can be directly linked to the canonical Jordan form of the respective transfer matrix $`T_L`$. Indeed, the $`4\times 4`$ matrix $`T_L`$, being invertible, can have one of the following five different canonical forms
$`\stackrel{~}{T}_1`$ $`=\left[\begin{array}{cccc}X_1& 0& 0& 0\\ 0& X_2& 0& 0\\ 0& 0& X_3& 0\\ 0& 0& 0& X_4\end{array}\right],`$ (42)
$`\stackrel{~}{T}_{21}`$ $`=\left[\begin{array}{cccc}X_1& 1& 0& 0\\ 0& X_1& 0& 0\\ 0& 0& X_3& 0\\ 0& 0& 0& X_4\end{array}\right],\stackrel{~}{T}_{22}=\left[\begin{array}{cccc}X_1& 1& 0& 0\\ 0& X_1& 0& 0\\ 0& 0& X_2& 1\\ 0& 0& 0& X_2\end{array}\right],`$ (51)
$`\stackrel{~}{T}_3`$ $`=\left[\begin{array}{cccc}X_1& 1& 0& 0\\ 0& X_1& 1& 0\\ 0& 0& X_1& 0\\ 0& 0& 0& X_2\end{array}\right],\stackrel{~}{T}_4=\left[\begin{array}{cccc}X& 1& 0& 0\\ 0& X& 1& 0\\ 0& 0& X& 1\\ 0& 0& 0& X\end{array}\right],`$ (60)
of which all but $`\stackrel{~}{T}_1`$ have non-trivial Jordan blocks and, therefore, are not diagonalizable. Each $`m\times m`$ Jordan block is associated with a single eigenvector of the respective $`T`$ \- matrix. Therefore, the total number of eigenvectors of different $`4\times 4`$ matrix in (51) is lesser than the number four of the solutions for the reduced Maxwell equation (10). The only exception is the diagonalizable case $`\stackrel{~}{T}_1`$, where the four $`T_L`$ eigenvectors correspond to four Bloch eigenmodes, as prescribed by (24). Generally, each nontrivial $`m\times m`$ Jordan block of the $`T`$ \- matrix is associated with the $`m`$ eigenmodes of Eq. (10), of which one is a propagating Bloch eigenmode with zero group velocity, and the other $`m1`$ are non-Bloch eigenmodes algebraically diverging with $`z`$. The details can be found in any course on linear algebra, for example, . Let us consider each case separately.
At general frequencies, different from those of stationary points (37), the $`4\times 4`$ matrix $`T_L`$ is always diagonalizable. Its canonical (diagonalized) form is trivial and coincides with $`\stackrel{~}{T}_1`$ in (51). The four eigenvectors are defined in (24). All the possibilities here reduce to one of the three cases (A), (B), or (C) described earlier in this Section. None of them involves non-Bloch solutions.
At the frequency $`\omega _g`$ of a regular photonic band edge in Fig. 6(a), the canonical Jordan form of the respective transfer matrix is $`\stackrel{~}{T}_{21}`$ in (51), where
$$X_1=1,X_3=X_3^{}=X_4^11.$$
The $`2\times 2`$ Jordan block relates to one propagating Bloch mode with $`k=\pi /L`$ and zero group velocity, and one non-Bloch linearly diverging eigenmode. The pair of real eigenvalues $`X_3`$ and $`X_4=X_3^1`$ relate to the pair of evanescent modes at $`\omega =\omega _g`$.
By contrast, at the frequency of photonic band edge $`g`$ in Fig. 1(b) related to the stack of isotropic layers, there is no evanescent modes. The canonical Jordan form of the respective transfer matrix coincides with $`\stackrel{~}{T}_{22}`$ in (51), where
$$X_1=X_2=1.$$
Each of the two $`2\times 2`$ identical Jordan blocks relates to one propagating Bloch mode with $`k=\pi /L`$ and zero group velocity, as well as one non-Bloch linearly diverging eigenmode. The two Jordan blocks of $`\stackrel{~}{T}_{22}`$ correspond to two different polarizations of light.
At the frequency $`\omega _a`$ in Fig. 6, the canonical Jordan form of the respective transfer matrix is $`\stackrel{~}{T}_{21}`$, where
$$X_1=1,X_3=X_4^{},\left|X_3\right|=\left|X_4\right|=1.$$
The double eigenvalue $`X_1=1`$ of the $`2\times 2`$ Jordan block relates to one propagating Bloch mode with $`k=\pi /L`$ and zero group velocity, and one non-Bloch linearly diverging eigenmode. The pair of complex eigenvalues $`X_3`$ and $`X_4=X_3^{}`$ relate to the pair of propagating modes having equal and opposite group velocities and belonging to the dispersion curve other than the one containing the point $`a`$.
The canonical form $`\stackrel{~}{T}_3`$ in (51) relates to a $`k\omega `$ diagram with stationary inflection point. Such a stationary point cannot be realized in a reciprocal periodic stack at normal light propagation .
Of particular interest here is the case (6) of degenerate band edge. At the degenerate band edge frequency $`\omega _d`$, the canonical Jordan form of the respective transfer matrix is $`\stackrel{~}{T}_4`$ in (51)
$$T_L\stackrel{~}{T}_4=\left[\begin{array}{cccc}X& 1& 0& 0\\ 0& X& 1& 0\\ 0& 0& X& 1\\ 0& 0& 0& X\end{array}\right],\text{ where }X=\pm 1$$
(61)
More specifically, in the case shown in Fig. 6(b)
$$X=X_d=e^{ik_dL}=1,\text{at }\omega =\omega _d.$$
The matrix (61) presents a single $`4\times 4`$ Jordan block and has a single eigenvector, corresponding to the propagating eigenmode with $`k=\pi /L`$ and zero group velocity. The other three solutions for the Maxwell equation (10) at $`\omega =\omega _d`$ are non-Bloch eigenmodes diverging as $`z`$, $`z^2`$, and $`z^3`$, respectively.
If the frequency $`\omega `$ deviates from the stationary point (37), the transfer matrix $`T_L`$ become diagonalizable with the canonical Jordan form $`\stackrel{~}{T}_1`$ in (51). The perturbation theory relating the non-Bloch eigenmodes at the frequency of degenerate band edge to the Bloch eigenmodes in the vicinity of this point is presented in .
### 2.3 Symmetry conditions for the existence of degenerate band edge
Not any periodic stack can develop a degenerate band edge, defined in (6). Some fundamental restrictions can be derived from symmetry considerations. These restrictions stem from the fact that at the frequency $`\omega _d`$ of degenerate band edge, the transfer matrix $`T_L`$ must have the Jordan canonical form (61). Such a matrix cannot be reduced to a block-diagonal form, let alone diagonalized. Therefore,
*\- a necessary condition for the existence of degenerate band edge is that the symmetry of the periodic array does not impose the reducibility of the transfer matrix* $`T_L`$ *to a block-diagonal form.*
The above condition does not imply that the transfer matrix $`T_L`$ must not be reducible to a block-diagonal form at *any* frequency $`\omega `$ on the $`k\omega `$ diagram. Indeed, at a general frequency $`\omega `$, the matrix $`T_L`$ is certainly reducible, and even diagonalizable. The strength of *the symmetry imposed* reducibility is that it leaves no room for exceptions, such as the frequency $`\omega _d`$ of degenerate band edge, where the transfer matrix $`T_L`$ must not be reducible to a block-diagonal form. Therefore, in the case of symmetry imposed reducibility, the very existence of degenerate band edge is ruled out. Observe that in most periodic layered structures, the symmetry of the periodic array *does require* the matrix $`T_L`$ to be similar to a block-diagonal matrix at all frequencies. In all these cases, the stack symmetry is incompatible with the existence of degenerate band edge on the $`k\omega `$ diagram.
Let us apply the above criterion to some specific cases.
In periodic stacks of isotropic layers, the Maxwell equations for the waves with the $`x`$\- and the $`y`$ \- polarizations are identical and decoupled, implying that the respective transfer matrix can be reduced to the block-diagonal form
$$\stackrel{~}{T}_L=\left[\begin{array}{cccc}T_{11}& T_{12}& 0& 0\\ T_{21}& T_{22}& 0& 0\\ 0& 0& T_{11}& T_{12}\\ 0& 0& T_{21}& T_{22}\end{array}\right].$$
(62)
The two identical blocks in (62) correspond to two different polarizations of light. The characteristic polynomial $`F_4(X)`$ of the block-diagonal matrix (62) factorizes into a product of two identical second degree polynomials related to electromagnetic waves with the $`x`$\- and the $`y`$ \- polarizations, respectively
$$F_4(X)=F_2(X)F_2(X).$$
(63)
The block-diagonal structure of the matrix (62) rules out the existence of degenerate band edge in periodic stacks of isotropic layers. In fact, the transfer matrix $`T_L`$ in this case can only have the following two canonical forms
$$\stackrel{~}{T}_1=\left[\begin{array}{cccc}X& 0& 0& 0\\ 0& X^1& 0& 0\\ 0& 0& X& 0\\ 0& 0& 0& X^1\end{array}\right],\stackrel{~}{T}_{22}=\left[\begin{array}{cccc}\pm 1& 1& 0& 0\\ 0& \pm 1& 0& 0\\ 0& 0& \pm 1& 1\\ 0& 0& 0& \pm 1\end{array}\right],$$
compatible with (62). The case $`\stackrel{~}{T}_1`$ relates to a general frequency, while the case $`\stackrel{~}{T}_{22}`$ relates to a photonic band edge, like the one shown in Fig. 1(b).
Let us now turn to the situation where all or some of the layers of the periodic stack are birefringent. The in-plane dielectric anisotropy (11) may allow for degenerate band edge on the $`k\omega `$ diagram, but not automatically.
Let us start with the simplest periodic array in which all anisotropic layers of the stack have aligned in-plane anisotropy. The term ”aligned” means that one can choose the directions of the in-plane Cartesian axes $`x`$ and $`y`$ so that the permittivity tensors in all layers are diagonalized simultaneously. In this setting, the Maxwell equations for the waves with the $`x`$\- and the $`y`$ \- polarizations are still separated, implying that the respective transfer matrix can be reduced to the block-diagonal form
$$\stackrel{~}{T}_L=\left[\begin{array}{cccc}T_{11}& T_{12}& 0& 0\\ T_{21}& T_{22}& 0& 0\\ 0& 0& T_{33}& T_{34}\\ 0& 0& T_{43}& T_{44}\end{array}\right].$$
(64)
The two blocks in (64) correspond to the $`x`$ and $`y`$ polarization of light. The forth degree characteristic polynomial of the block-diagonal matrix (64) factorizes into the product
$$F_4(X)=F_x(X)F_y(X),$$
(65)
where $`F_x(X)`$ and $`F_y(X)`$ are independent second degree polynomials related to electromagnetic waves with the $`x`$\- and the $`y`$ \- polarizations, respectively. The typical $`k\omega `$ diagram in this case will be similar to that shown in Fig. 6(a) with two separate curves related to two linear polarizations of light. Again, the block-diagonal structure of the matrix (64) rules out the existence of degenerate band edge in periodic stacks with aligned anisotropic layers. The transfer matrix $`T_L`$ in this case can only have the following two canonical forms
$$\stackrel{~}{T}_1=\left[\begin{array}{cccc}X_1& 0& 0& 0\\ 0& X_1^1& 0& 0\\ 0& 0& X_2& 0\\ 0& 0& 0& X_2^1\end{array}\right],\stackrel{~}{T}_{22}=\left[\begin{array}{cccc}\pm 1& 1& 0& 0\\ 0& \pm 1& 0& 0\\ 0& 0& X& 0\\ 0& 0& 0& X^1\end{array}\right],$$
compatible with (64). The case $`\stackrel{~}{T}_1`$ relates to a general frequency, while the case $`\stackrel{~}{T}_{22}`$ relates to a photonic band edge.
As we have seen, the presence of anisotropic layers may not necessarily lift the symmetry prohibition for the degenerate band edge, because the symmetry of the periodic array may still be incompatible with the canonical Jordan form (61). Generally if the space symmetry group $`G`$ of the layered structure includes a mirror plane $`m_{||}`$ parallel to the $`z`$ direction, this would guarantee the reducibility of the matrix $`T_L`$ to a block-diagonal form. Indeed, a standard line of reasoning gives that if $`m_{||}G`$, and the $`y`$ axis is chosen perpendicular to the mirror plane $`m_{||}`$, then the waves with the $`x`$\- and the $`y`$ \- polarizations have different parity with respect to the symmetry operation of reflection and, therefore, are decoupled. The latter leads to reducibility of the matrix $`T_L`$ to the block-diagonal form (64). Thus, a formal necessary condition for the existence of degenerate band edge on the $`k\omega `$ diagram can be written as follows
$$m_{||}G.$$
(66)
None of the common periodic layered structures satisfies this criterion and, therefore, none of them can develop the degenerate band edge. For example, even if anisotropic layers are present, but the anisotropy axes in all anisotropic layers are either aligned, or perpendicular to each other, the symmetry group $`G`$ of the stack still has the mirror plane $`m_{||}`$, which guarantees the separation of the $`x`$\- and the $`y`$ \- polarizations and the reducibility of the respective transfer matrix $`T_L`$ to the block-diagonal form (64). The only way to satisfy the condition (66) and, thereby, to allow for degenerate band edge on the $`k\omega `$ diagram, is to have at least two misaligned anisotropic layers in a unit cell with the misalignment angle being different from $`0`$ and $`\pi /2`$, as shown in the example in Fig. 5.
Observe that the presence of $`B`$ layers in the periodic array in Fig. 5 is also essential, unless the two anisotropic layers $`A_1`$ and $`A_2`$ have different thicknesses or are made of different anisotropic materials. In Fig. 5, the layers $`A_1`$ and $`A_2`$ differ only by their orientation in the $`xy`$ plane, but otherwise, they are identical. In such a case, if the $`B`$ layers are removed, the point symmetry group of the periodic stack in Fig. 5 rises from $`D_2`$ to $`D_{2h}`$ acquiring the glide mirror plane $`m_{||}`$. This, according to the criterion (66), imposes the reducibility of the matrix $`T_L`$ to the block-diagonal form (64), regardless of the misalignment angle between the adjacent $`A`$ layers. The symmetry imposed reducibility rules out the possibility of the degenerate band edge (6). The $`k\omega `$ diagram of the periodic stack in Fig. 5 with the $`B`$ layers removed is shown in Fig. 6(d).
In the numerical example considered in the next section, the $`B`$ layers are simply empty gaps of certain thickness $`D_B`$ between the adjacent double layers $`A_1A_2`$. The misalignment angle is chosen $`\pi /4`$. By changing the thickness $`D_B`$ of the gap $`B`$, one can change the $`k\omega `$ diagram of the periodic stack, as shown in Fig. 6. Similar effect can be achieved by changing the misalignment angle between the adjacent $`A`$ layers.
## 3 Transmission resonance in the vicinity of degenerate band edge
### 3.1 Scattering problem for periodic semi-infinite stack
To solve the scattering problem for a plane monochromatic wave incident on a finite stack of anisotropic layers we use the following standard approach based on the $`4\times 4`$ transfer matrix.
Let $`\mathrm{\Psi }_I\left(z\right)`$, $`\mathrm{\Psi }_R\left(z\right)`$, and $`\mathrm{\Psi }_P\left(z\right)`$, be the incident, reflected, and passed plane waves in vacuum. Allowing for general elliptic polarization of the waves, we have
$$\mathrm{\Psi }_I\left(z\right)=\left[\begin{array}{c}A_x\\ A_y\\ A_y\\ A_x\end{array}\right]e^{i\frac{\omega }{c}z},\mathrm{\Psi }_R\left(z\right)=\left[\begin{array}{c}R_x\\ R_y\\ R_y\\ R_x\end{array}\right]e^{i\frac{\omega }{c}z},\mathrm{\Psi }_P\left(z\right)=\left[\begin{array}{c}P_x\\ P_y\\ P_y\\ P_x\end{array}\right]e^{i\frac{\omega }{c}z}.$$
The respective domains of definition are
$$\begin{array}{c}z0\text{, for }\mathrm{\Psi }_I\left(z\right)\text{ and }\mathrm{\Psi }_R\left(z\right),\\ Dz\text{, for }\mathrm{\Psi }_P\left(z\right).\end{array}$$
(67)
where $`D`$ is the stack thickness. The field inside the stack is denoted by $`\mathrm{\Psi }_T\left(z\right)`$. Except for the stationary points (37) on the $`k\omega `$ diagram, $`\mathrm{\Psi }_T\left(z\right)`$ can be decomposed into a superposition of the four Bloch solutions (20) of the Maxwell equation (10)
$$\mathrm{\Psi }_T\left(z\right)=\underset{k}{}\mathrm{\Psi }_k\left(z\right),0zD.$$
(68)
This representation is meaningful if the periodic stack contains a significant number of unit cells $`L`$. Otherwise, if there are just a few layers in the stack, the representation (68) is formally valid, but not particularly useful.
The boundary conditions at the two slab/vacuum interfaces are
$$\mathrm{\Psi }\left(0\right)=\mathrm{\Psi }_I\left(0\right)+\mathrm{\Psi }_R\left(0\right),\mathrm{\Psi }\left(D\right)=\mathrm{\Psi }_P\left(D\right).$$
(69)
The transfer matrix $`T_S`$ of a periodic stack is
$$T_S=\left(T_L\right)^N.$$
(70)
The relation
$$\mathrm{\Psi }\left(D\right)=T_S\mathrm{\Psi }\left(0\right),$$
(71)
together with the pair of boundary conditions (69) allow to express both the reflected wave $`\mathrm{\Psi }_R`$ and the wave $`\mathrm{\Psi }_P`$ passed through the slab, in terms of a given incident wave $`\mathrm{\Psi }_I`$ and the elements of the transfer matrix $`T_S`$. This also gives the transmittance/reflectance coefficients of the slab defined as
$$\tau _D=\frac{S_P}{S_I}=\frac{\left|\mathrm{\Psi }_P\left(D\right)\right|^2}{\left|\mathrm{\Psi }_I\left(0\right)\right|^2},\rho _D=\frac{S_R}{S_I}=\frac{\left|\mathrm{\Psi }_R\left(0\right)\right|^2}{\left|\mathrm{\Psi }_I\left(0\right)\right|^2}.$$
(72)
where $`S=cW`$ is the Poynting vector of the respective wave. In the case of a lossless stack
$$\tau _D+\rho _D=1.$$
The field distribution $`\mathrm{\Psi }_T\left(z\right)`$ inside the slab is found using either of the following expressions
$$\mathrm{\Psi }_T\left(z\right)=T(z,0)\left[\mathrm{\Psi }_I\left(0\right)+\mathrm{\Psi }_R\left(0\right)\right]=T(z,D)\mathrm{\Psi }_P\left(D\right),0zD.$$
(73)
The above procedure is commonly used for the frequency-domain analysis of periodic and non-periodic layered structures involving anisotropic and/or gyrotropic layers (see, for example, ,,,, and references therein).
In addition to the field distribution inside the slab, we are also interested in its eigenmode composition. The latter is particularly important since it allows to explain the fundamental difference between Fabry-Perot resonance in the vicinity of a degenerate band edge and a similar resonance in the vicinity of a regular band edge. Throughout this section we consider only the first transmission resonance, closest to the respective band edge.
### 3.2 Field composition at the frequency of transmission resonance
In the case of transmission resonance in a periodic stack of isotropic layers, the resonance field inside the stack is a simple standing wave composed of two propagating Bloch modes with opposite group velocity (see Eqs. (1) through ((5) and the comments therein). The introduction of anisotropy in itself does not change qualitatively the resonance picture, as shown in Fig. 9. The only difference is that the electromagnetic field $`\mathrm{\Psi }_T\left(z\right)`$ inside the stack can now have both propagating and evanescent components. But at the frequency of a transmission resonance, the contribution of the evanescent components is negligible, as shown in Fig. 10. So, basically, one can still see the resonance field inside the stack as a simple standing wave composed of a pair of propagating modes with greatly enhanced amplitude, compared to that of the incident wave. The formulas (1) through ((5) still apply here, provided that the number $`N`$ of unit cells in the stack is not too small.
In Fig. 10(a) we show the squared amplitudes (intensity distributions)
$$\left|\mathrm{\Psi }_j\left(z\right)\right|^2,j=1,2,3,4,$$
(74)
of the individual Bloch components of the resulting field $`\mathrm{\Psi }_T\left(z\right)`$ at the frequency of the first transmission resonance near the regular band edge $`g`$ in Fig. 6(a). The numbers 1 and 2 designate the forward and backward propagating components of the standing wave. The evanescent contributions 3 and 4 are negligible. In Fig. 10(b) we show the squared amplitudes of the combined contribution of the pair of propagating waves ($`pr`$), and the combined contribution of the pair of evanescent waves ($`ev`$)
$$\left|\mathrm{\Psi }_{pr}\left(z\right)\right|^2=\left|\mathrm{\Psi }_1\left(z\right)+\mathrm{\Psi }_2\left(z\right)\right|^2,\left|\mathrm{\Psi }_{ev}\left(z\right)\right|^2=\left|\mathrm{\Psi }_3\left(z\right)+\mathrm{\Psi }_4\left(z\right)\right|^2.$$
(75)
Obviously,
$$\mathrm{\Psi }_T\left(z\right)\mathrm{\Psi }_{pr}\left(z\right).$$
(76)
Now let us turn to the case of transmission resonance in the vicinity of degenerate band edge $`d`$ in Fig. 6(b). There are several features that sharply distinguish this case from the similar transmission resonance near the regular band edge $`g`$ in Fig. 6(a).
First of all, for a given number $`N`$ of unit cells in the stack ($`N=32`$), the field intensity in Fig. 8(b) is by several orders of magnitude higher, compared to that in Fig. 9(b), in spite of the close similarity of all numerical parameters of the respective periodic structures (see the Appendix). In the case of transmission resonance in the vicinity of degenerate band edge (DBE), the resonance field intensity increases as $`N^4`$, while in the case of a regular band edge (BE), the field intensity is proportional to $`N^2`$.
The second distinction is that in the case of degenerate band edge (DBE), the field intensity near the slab boundaries at $`z=0`$ and $`z=D`$, increases as
$$\text{DBE case: }\left|\mathrm{\Psi }_T\left(z\right)\right|^2z^4,\left(Dz\right)^4.$$
(77)
By contrast, in the case of a regular band edge (BE), the field intensity near the stack boundaries rises at a much slower rate
$$\text{Regular BE case: }\left|\mathrm{\Psi }_T\left(z\right)\right|^2z^2,\left(Dz\right)^2,$$
(78)
which is characteristic of a regular standing wave composed of two propagating components.
Finally, in the case of a regular band edge, the transmission resonance field $`\mathrm{\Psi }_T\left(z\right)`$ is a standing wave composed of two propagating Bloch modes with opposite group velocities. The evanescent modes do not participate in the formation of the resonance field, as clearly seen in Fig. 10. By contrast, in the case of transmission resonance in the vicinity of DBE, the role of evanescent components in the formation of the resonance field is absolutely crucial. Indeed, as shown in Fig. 11, amplitude of the propagating and evanescent components are comparable in magnitude. More importantly, the combined contribution $`\mathrm{\Psi }_{pr}\left(z\right)`$ of the two propagating components does not even resemble a standing wave with the nodes at the slab boundary, as was the case in Fig. 10(b). Instead, comparing Fig. 8(b) and Fig. 11(b) we see that at the slab boundaries at $`z=0`$ and $`z=D`$. the propagating and evanescent components interfere destructively, almost canceling each other
$$\text{at }z=0\text{ and }z=D\text{}\mathrm{\Psi }_{pr}\left(z\right)\mathrm{\Psi }_{ev}\left(z\right),$$
while the individual Bloch components (74) remain huge. We remind that in all cases, the intensity of the incident wave $`\mathrm{\Psi }_I`$ is unity.
Qualitatively, such a bizarre resonance behavior in the vicinity of the degenerate band edge can be characterized as follows. As we already mentioned, at the frequency $`\omega _d`$ of degenerate band edge, three of the four solutions for the Maxwell equation (10) are non-Bloch eigenmodes diverging as $`z`$, $`z^2`$, $`z^3`$. Although the frequency $`\omega _1`$ of the transmission cavity resonance is slightly different from $`\omega _d`$ and, therefore, all four eigenmodes are Bloch waves, the close proximity of $`\omega _1`$ to $`\omega _d`$ causes all the abnormalities seen in Fig. 10. To describe these behavior in mathematically consistent way, one should start with the one solutions at $`\omega =\omega _d`$ as zero approximation. Then using the perturbation theory for the non-diagonalizable transfer matrix $`T_L\left(\omega _d\right)`$, one can derive an asymptotic theory for the case of large $`N`$. The perturbation theory for the degenerate band edge $`d`$ was developed in . It turns out that near the slab boundaries at $`zD`$ and $`\left(Dz\right)D`$, the field $`\mathrm{\Psi }_T\left(z\right)`$ at the transmission resonance frequency $`\omega _1`$ is well approximated by quadratically diverging non-Bloch eigenmode, corresponding to $`\omega =\omega _d`$. It qualitatively explains an extremely rapid growth (77) of the resonance field inside the slab as one moves away from either slab boundary. It also explains the sharp dependence (7) of the resonance field intensity on the number $`N`$ of unit cells in the stack.
Acknowledgment and Disclaimer: Effort of A. Figotin is sponsored by the Air Force Office of Scientific Research, Air Force Materials Command, USAF, under grant number FA9550-04-1-0359.
## 4 Appendix
### 4.1 Numerical parameters of layered arrays
The periodic array in Fig. 5. has three layers in a unit cell $`L`$, of which two ($`A_1`$ and $`A_2`$) are anisotropic and have the same thickness $`D_A`$. The third layer $`B`$ is isotropic with the thickness $`D_B`$. In our numerical simulations, the $`B`$ layers are empty gaps of variable thickness.
The anisotropic layers $`A_1`$ and $`A_2`$ are made of the same lossless dielectric material and have the same thickness. The respective dielectric permittivity tensor is
$$\widehat{\epsilon }=\left[\begin{array}{ccc}\epsilon _A+\delta \mathrm{cos}2\phi & \delta \mathrm{sin}2\phi & 0\\ \delta \mathrm{sin}2\phi & \epsilon _A\delta \mathrm{cos}2\phi & 0\\ 0& 0& \epsilon _{zz}\end{array}\right],$$
where $`\delta `$ describes the magnitude of in-plane anisotropy, while the angle $`\phi `$ defines the orientation of the anisotropy axes of the respective layer in the $`xy`$ plane. The most critical parameter of the periodic structure in Fig. 5 is the misalignment angle
$$\phi =\phi _1\phi _2$$
(79)
between the adjacent anisotropic layers $`A_1`$ and $`A_2`$. This angle determines the symmetry of the periodic array and, eventually, what kind of $`k\omega `$ diagram it can display. It is important for our purposes that the misalignment angle is different from $`0`$ and $`\pi /2`$. In our numerical simulations we set
$$\phi _1=0,\phi _2=\pi /4,$$
and use the following expressions for the dielectric permittivity of the $`A`$ layers
$$\widehat{\epsilon }_{A1}=\left[\begin{array}{ccc}\epsilon _A+\delta & 0& 0\\ 0& \epsilon _A\delta & 0\\ 0& 0& \epsilon _{zz}\end{array}\right],\widehat{\epsilon }_{A2}=\left[\begin{array}{ccc}\epsilon _A& \delta & 0\\ \delta & \epsilon _A& 0\\ 0& 0& \epsilon _{zz}\end{array}\right].$$
(80)
In our numerical simulations we set
$$\epsilon _A=\mathrm{13.\hspace{0.17em}61},\delta =\mathrm{12.\hspace{0.17em}4}.$$
The numerical value of $`\epsilon _{zz}`$ is irrelevant.
In the case of the periodic array of isotropic layers in Fig. 1(a), we set
$$\epsilon _A=3.78,\delta =0,D_A=D_B=0.5\times L.$$
### 4.2 Description of plots
In all plots of the field intensity distribution we, in fact, plotted the following physical quantity
$$\left|\mathrm{\Psi }\left(z\right)\right|^2=\stackrel{}{E}\left(z\right)\stackrel{}{E}^{}\left(z\right)+\stackrel{}{H}\left(z\right)\stackrel{}{H}^{}\left(z\right)_L,$$
(81)
which is the squared field amplitude averaged over local unit cell. The real electromagnetic energy density $`W\left(z\right)`$ is similar to $`\left|\mathrm{\Psi }\left(z\right)\right|^2`$. Both of them are strongly oscillating functions of the coordinate $`z`$ with the period of oscillations coinciding with the unit cell length $`L`$. Thus, the quantity (81) can be interpreted as the smoothed field intensity distribution, with the correction coefficient of the order of unity.
In all plots, the wave number $`k`$ and the frequency $`\omega `$ are expressed in units of $`L^1`$ and $`cL^1`$, respectively.
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# Security analysis of communication system based on the synchronization of different order chaotic systems
## 1 Introduction
In recent years, a considerable effort has been devoted to extend the chaotic communication applications to the field of secure communications. It has been noticed that there exists an interesting relationship between chaos and cryptography: many properties of chaotic systems have their corresponding counterparts in traditional cryptosystems, such as:
* Ergodicity and Confusion: The output has the same distribution for any input.
* Sensitivity to initial conditions/control parameter and Diffusion with a small change in the plaintext/secret key: A small deviation in the input can cause a large change at the output.
* Mixing property and Diffusion with a small change in one plain-block of the whole plaintext: A small deviation in the local area can cause a large change in the whole space.
* Deterministic dynamics and Deterministic pseudo-randomness: A deterministic process can cause a random-like (pseudo-random) behavior.
* Structure complexity and Algorithm (attack) complexity: A simple process has a very high complexity.
As a result of investigating the above relationships, a rich variety of chaos-based cryptosystems for end-to-end communications have been proposed , some of them fundamentally flawed by a lack of robustness and security .
Most analog chaos-based cryptosystems are secure communication schemes designed for noisy channels, based on the technique of chaos synchronization, first shown by Pecora and Carrol .
Reduced order synchronization is a new interesting topic which has recently drawn attention from several researchers . In it is shown that second order driven oscillators can be synchronized with canonical projection of a higher order chaotic system by means of non-linear feedback.
In a recent paper Bowong proposed a scheme based on reduced order synchronization and feedback with application to secure communications . The transmitter is a four-order chaotic oscillator, modulated by the plaintext and whose output is added to the plaintext as a masking signal. The receiver consists of a Duffing second-order system, that is enslaved to the transmitter by means of non-linear feedback of the error.
Bowong presented two examples based on the plaintext modulation of a chaotic oscillator and subsequent additive masking of the plaintext with the oscillator signal. The following equation system defines the transmitter operation:
$`\dot{x}_{1m}=`$ $`x_{2m},`$
$`\dot{x}_{2m}=`$ $`cx_{2m}d^2(x_{1m}x_{3m}x_{1m}x_{3m})+u(t),`$ (1)
$`\dot{x}_{3m}=`$ $`x_{4m},`$
$`\dot{x}_{4m}=`$ $`ax_{4m}x_{3m}b^2x_{3m}^3f(x_{1m}+x_{1m}^2/2)+e(t),`$
$`y_T=`$ $`x_{1m}+u(t),`$
with parameter values:
$$a=0.03,b=1,c=0.3,d=0.985,f=0.1,$$
(2)
where the function $`e(t)`$ was specified as $`e(t)=0.7\mathrm{cos}(t)`$ and being the term $`u(t)`$ the plaintext message. The signal $`y_T=x_{1m}+u(t)`$ constitutes the transmitted ciphertext to the receiver end. Actually, the above Eq. (1) is written as it should be, in spite of the erroneous formulation given in (14) of \[18, §4\].
The receiver was constructed as follows:
$`\dot{x}_{1s}=`$ $`x_{2s},`$
$`\dot{x}_{2s}=`$ $`\lambda x_{2s}\omega _0^2x_{1s}\gamma x_{1s}^3+K_1\mathrm{cos}(\omega _1t+\theta _1)+K_2\mathrm{cos}(\omega _2t+\theta _2)+\upsilon ,`$
$`y_s=`$ $`x_{1s},`$
with parameter values:
$$\lambda =1,\omega _0=10,\gamma =100,K_1=K_2=1,\omega _1=2,\omega _2=4,\theta _1=\theta _2=0,$$
(3)
and where $`\upsilon `$ is the feedback control law which forces the error $`e=x_{1s}x_{1m}`$ to converge exponentially to zero as $`t\mathrm{}`$.
The retrieved plaintext $`\widehat{u}(t)`$ is calculated as the difference between the ciphertext and the output of the reduced order system $`\widehat{u}(t)=y_Ty_s=u(t)e`$. It was shown that the transmitter-receiver system was capable of accurately retrieving the plaintext after an initial synchronization period of 10 seconds. Afterwards it was claimed that the system can be used for secure communications and some examples are provided.
In this letter it is shown that the proposed cryptosystem is insecure and two different procedures to break it are also presented: by high-pass filtering and by means of a simple reduced order intruder receiver.
## 2 Missing security analysis and system key specification
In , the author asserted that the scheme is applicable to secure communication. However, no analysis of security was included to support this claim. Furthermore, there is no mention of the secret key, when it is well known that a secure communication system cannot exist without a key. In it is not considered whether there should be a key in the proposed system, what it should consist of, what the available key space would be (how many different keys exist in the system), what precision to use, and how it would be managed. None of these elements should be neglected when describing a secure communication system .
Moreover, being the transmitter and the receiver implemented with different kind of systems, it is not explained how the encryption keys, if any, may be related to the corresponding decryption keys. Usually, in many chaotic cryptosystems the system parameters play the role of key, but it is not the case in , because the transmitter and the receiver do not make use of the same parameters.
## 3 Plaintext retrieval by a filtering attack
It was supposed for some time that chaotic masking was an adequate means for secure transmission, because chaotic systems present some properties as sensitive dependence on parameters and initial conditions, ergodicity, mixing, and dense periodic points. These properties make them similar to pseudorandom noise , which has been used traditionally as a masking signal for cryptographic purposes. The basic fundamental requirement of the pseudorandom noise used in cryptography is that its spectrum should be infinitely broad, flat and of much higher power density than the signal to be concealed. In other words, the plaintext power spectrum should be effectively buried into the pseudorandom noise power spectrum.
The secure application proposed in does not satisfy this condition. On the contrary, the spectrum of the signal generated by the chaotic oscillator is of narrow band, decaying very fast with increasing frequency, showing a power density much lower than the plaintext at the plaintext frequencies used. Hence it can not cope with a filtering attack intended to separate the masking signal and the plaintext.
To illustrate this fact we consider the the two examples in \[18, §4\] corresponding to the following plaintexts:
$`u_1(t)`$ $`=\mathrm{cos}(7t),`$
$`u_2(t)`$ $`=(1+\mathrm{sin}(0.2t))\mathrm{cos}(7t),`$
whose waveforms are illustrated in Fig. 1.
The transmitter proposed in was simulated with a four-order Runge-Kutta integration algorithm in MATLAB 6.5, with a step size of $`10^3`$. Fig. 2 illustrates the logarithmic power spectra, as a function of frequency, of the ciphertexts $`y_{T1}`$ and $`y_{T2}`$ when the plaintext signals $`u_1(t)`$ and $`u_2(t)`$ are encrypted, respectively, with the same parameter values previously described in (2). The power spectra were calculated using a 8192-point Discrete Fourier Transform with a sampling frequency of 32 Hz; previously, the analyzed signal segments were multiplied by a 4-term Blackman-Harris window , to avoid aliasing artifacts.
It can be seen in both examples that the plaintext signal components clearly emerge at 1.114 Hz over the background noise created by the chaotic oscillator, with a power of $`3\text{db}`$, relative to the maximum power of the ciphertext spectrum, while the power density of the ciphertext, at neighboring frequencies, falls below $`80\text{db}`$.
The chaotic receiver of was not used to recover the plaintext. Instead, the ciphertext was high-pass filtered to eliminate the chaotic masking component while retaining the plaintext information. The result is illustrated in Fig. 3. Comparing the result with the plaintext displayed in Fig. 1, it can be appreciated the good estimation of the plaintexts after an initial delay of approximately 29 seconds, due to the filter delay. The filter employed was a 2048 samples finite impulse response digital one, with a cut-off frequency of 1 Hz.
Note that this is the hardest case an attacker can face from the point of view of plaintext frequency, because for higher sound frequencies the spectrum of the background noise created by the chaotic oscillator is even lower.
This plaintext recovering method works equally well for different parameters values of the transmitter, because the maximum power components of its spectrum are concentrated in the frequency range between 0 and 0.3 Hz for all parameter values.
Thanks to the big separation between the plaintext frequency and the high amplitude components of the masking chaotic signal, our method works equally well with plaintext signals of much lower amplitudes than the plaintexts $`u_1(t)`$ and $`u_2(t)`$ of the examples described in . For instance, we present in Fig. 4 the retrieved text corresponding to the plaintext $`u_3(t)=0.0032(1+\mathrm{sin}(0.2t))\mathrm{cos}(7t)`$, that has a power level of -50db with respect to $`u_2(t)`$; but, as can be seen, the retrieved signal waveform is still perfectly preserved.
## 4 Plaintext retrieval by reduced order system synchronization
As mentioned in Sec. 1, many secure communication systems based on chaotic modulation and masking have been proposed in the past. In any of them the knowledge of the transmitter parameter values was mandatory to operate the receiver, since they played the role of system key. Some of them were compromised because it was possible to recover the system parameters carrying out an elaborate ciphertext signal analysis.
But the chaotic transmitter and the modulation and masking procedure described in was designed in such a way as to enable the plaintext retrieval with a reduced order receiver that did not even required the knowledge of any transmitter parameters. As a true encryption mechanism must necessarily make use of a key, this communication system may be only envisaged as an ordinary codification system, rather than a secure one, because the only required knowledge to recover the plaintext message is the receiver structure.
Moreover, it is possible to implement a whole family of alternative receivers to the one proposed in . Hence, a determined eavesdropper, still ignoring the precise structure of the transmitter nor its design parameters, may implement an alternative intruder receiver of its own design, also based on reduced order synchronization and feedback, capable of retrieving the ciphertext just as well as the authorized one.
To demonstrate this threat we have developed an extremely simple intruder receiver of order two, with linear feedback, constructed as follows:
$`\dot{x}_{1s}=`$ $`x_{2s},`$
$`\dot{x}_{2s}=`$ $`100\widehat{u}(t)x_{2s},`$
being $`\widehat{u}(t)=y_Tx_{1s}`$ the retrieved plaintext. The initial conditions were arbitrarily chosen as
$$x_{1m}(0)=x_{2m}(0)=x_{3m}(0)=x_{4m}(0)=0,x_{1s}(0)=0.1,x_{2s}(0)=1.$$
This simple receiver may decrypt the ciphertext as well as the receiver proposed in . Fig. 5 illustrates the perfect synchronism between the transmitter variable $`x_{1m}`$ and our intruder receiver variable $`x_{1s}`$, attained after a transient of 4 seconds, when no plaintext signal is present. The efficiency as intruder decoder is illustrated in Fig. 6, where the retrieved $`\widehat{u}_1(t)`$ and $`\widehat{u}_2(t)`$ texts corresponding to the plaintexts $`u_1(t)=\mathrm{cos}(7t)`$ and $`u_2(t)=(1+\mathrm{sin}(0.2t))\mathrm{cos}(7t)`$ are shown, comparing with the plaintexts illustrated in Fig. 1 it can be appreciated the perfect decoding after a short initial transient.
## 5 Conclusion
In summary, the chaotic masking cryptosystem proposed in is rather weak, since it can be broken in two different ways, without knowing the system parameters nor it detailed structure: by high pass filtering and by an intruder receiver based on reduced order synchronization and feedback. There is no mention about what the key is, nor what the key space is, a fundamental aspect in every secure communication system. The total lack of security discourages the application of this synchronization scheme to secure applications.
This work was supported by Ministerio de Ciencia y Tecnología of Spain, research grant SEG2004-02418. We thank the anonymous reviewer for his valuable suggestions.
## Figures
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# Where are ELKO Spinor Fields in Lounesto Spinor Field Classification?
## 1 Introduction
In order to find an adequate mathematical formalism for representing dark matter, Ahluwalia-Khalilova and Grumiller have recently introduced the *Eigenspinoren des Ladungskonjugationsoperators* (ELKO) spinor fields<sup>1</sup><sup>1</sup>1Dual-helicity eigenspinors of the charge conjugation operator. , which are shown to belong to a non-standard Wigner class, and to exhibit non-locality. They claim that an ELKO spinor field is a new fermion described by a spinor field that has not been identified, in the Physics literature, with any particle or more general physical entity yet. In the low-energy limit ELKO comports as a representation of the Lorentz group. However, mathematicians have already known since sometime ago that all spinors that are elements of the carrier spaces of the $`D^{(1/2,0)}D^{(0,1/2)}`$ or $`D^{(1/2,0)}`$, or $`D^{(0,1/2)}`$ representations of $`Sl(2,)`$ belong to one of the six classes found by Lounesto in his theory of the classification of spinor fields. Such an algebraic classification is based on the values assumed by their bilinear covariants, the Fierz identities, aggregates and boomerangs (see Eq.(7) below) . In this paper we prove that from the algebraic point of view ELKO spinor fields belong to Lounesto class 5 spinor fields, also called flagpole spinor fields due to the intrinsic flagpole structure they carry. It is a general property of class 5 spinor fields that they satisfy the Majorana condition, i.e., if $`\psi `$ denotes an ELKO spinor field, $`\lambda `$ a complex number of unitary modulus and $`𝒞`$ the charge conjugation operator then $`𝒞\psi =\lambda \psi `$. In Section 2, after presenting the bilinear covariants, that completely characterize a spinor field through Fierz identities and through Fierz aggregates (or boomerangs), Lounesto classification of spinor fields is reviewed. Section 3 recalls the definition of ELKO spinor fields and shows that ELKO is indeed a flagpole spinor field. We show moreover in Section 4 that algebraic constraints imply that any class 5 spinor field is such that their 2-component spinor fields have opposite helicities. Now, it is well known that Majorana spinor fields on Minkowski spacetime are defined as eigenvalues of the charge operator, and it is easy to verify that any spinor field that is eigenspinor of the charge operator necessarily belongs to Lounesto class 5. So, what differentiates ELKO from Majorana spinor fields? In authors quote that the difference is that according to Peskin and Schroeder and Marshak and Sudarshan it is imposed that the 2-component spinor fields of a Majorana spinor field have the same helicity. Is this reasonable? We briefly discuss such issue.
## 2 Bilinear Covariants
In this paper all spinor fields live in Minkowski spacetime $`(M,\eta ,D,\tau _\eta ,)`$. Here, the manifold $`M`$ $`^4`$, $`\eta `$ denotes a constant metric of signature $`(1,3)`$, $`D`$ denotes the Levi-Civita connection of $`\eta `$, $`M`$ is oriented by the 4-volume element $`\tau _\eta \mathrm{sec}{\displaystyle \stackrel{4}{}}T^{}M`$ and time-oriented by $``$. As usual $`T^{}M`$ denotes the cotangent bundle and $`TM`$ the tangent bundle over $`M`$. By a constant metric we mean the following: let $`\{x^\mu \}`$ be global coordinates in the Einstein-Lorentz gauge, naturally adapted to an inertial reference frame $`𝐞_0=/x^0`$. Let also $`𝐞_i=/x^i`$, $`i=1,2,3`$. Then, $`\eta (/x^\mu ,/x^\nu )=\eta _{\mu \nu }=\mathrm{diag}(1,1,1,1)`$. Also, $`\{𝐞_\mu \}`$ is a section of the frame bundle $`𝐏_{\mathrm{SO}_{1,3}^e}(M)`$ and $`\{𝐞^\mu \}`$ is its reciprocal frame satisfying $`\eta (𝐞^\mu ,𝐞_\nu ):=𝐞^\mu 𝐞_\nu =\delta _\nu ^\mu `$. Let $`\mathrm{\Xi }`$ be a section of the principal spin structure bundle $`𝐏_{\mathrm{Spin}_{1,3}^e}(M)`$ such that $`s(\mathrm{\Xi })=\{𝐞_\mu \}`$. Classical spinor fields<sup>2</sup><sup>2</sup>2Quantum spinor fields are operator valued distributions, as well known. It is not necessary to introduce quantum fields in order to know the algebraic classification of ELKO spinor fields. carrying a $`D^{(1/2,0)}D^{(0,1/2)}`$, or $`D^{(1/2,0)}`$, or $`D^{(0,1/2)}`$ representation of $`Sl(2,)\mathrm{Spin}_{1,3}^e`$ are sections of the vector bundle
$$𝐏_{\mathrm{Spin}_{1,3}^e}(M)\times _\rho ^4,$$
where $`\rho `$ stands for the $`D^{(1/2,0)}D^{(0,1/2)}`$ (or $`D^{(1/2,0)}`$ or $`D^{(0,1/2)}`$) representation of $`Sl(2,)\mathrm{Spin}_{1,3}^e`$ in $`^4`$. Other important spinor fields, like Weyl spinor fields are obtained by imposing some constraints on the sections of $`𝐏_{\mathrm{Spin}_{1,3}^e}(M)\times _\rho ^4`$. See, e.g., for details. Given a spinor field $`\psi `$ $`\mathrm{sec}𝐏_{\mathrm{Spin}_{1,3}^e}(M)\times _\rho ^4`$ the bilinear covariants are the following sections of $`{\displaystyle TM}={\displaystyle \underset{r=0}{\overset{4}{}}}`$ $`{\displaystyle \stackrel{r}{}}TMC\mathrm{}(M,\eta )`$, where $`TM`$ denotes the exterior algebra bundle of multivector fields and $`𝒞\mathrm{}(M,\eta )`$ denotes the Clifford bundle of multivector fields of Minkowski spacetime :
$`\sigma `$ $`=\psi ^{}\gamma _0\psi ,𝐉=J_\mu 𝐞^\mu =\psi ^{}\gamma _0\gamma _\mu \psi 𝐞^\mu ,𝐒=S_{\mu \nu }𝐞^{\mu \nu }={\displaystyle \frac{1}{2}}\psi ^{}\gamma _0i\gamma _{\mu \nu }\psi 𝐞^\mu 𝐞^\nu ,`$
$`𝐊`$ $`=\psi ^{}\gamma _0i\gamma _{0123}\gamma _\mu \psi 𝐞^\mu ,\omega =\psi ^{}\gamma _0\gamma _{0123}\psi ,`$ (1)
with $`\sigma ,\omega \mathrm{sec}{\displaystyle \stackrel{0}{}}TMC\mathrm{}(M,\eta )`$, $`𝐉,𝐊\mathrm{sec}{\displaystyle \stackrel{1}{}}TMC\mathrm{}(M,\eta )`$ and $`𝐒\mathrm{sec}{\displaystyle \stackrel{2}{}}TMC\mathrm{}(M,\eta )`$. In the formulas appearing in Eq.(1) the set $`\{\gamma _\mu \}`$ refers to the Dirac matrices in chiral representation (see Eq.(12)). Also,
$$\{1,𝐞^\mu ,𝐞^\mu 𝐞^\nu ,𝐞^\mu 𝐞^\nu 𝐞^\rho ,𝐞^0𝐞^1𝐞^2𝐞^3\},$$
where $`\mu ,\nu ,\rho =0,1,2,3`$, and $`\mu <\nu <\rho `$ is a basis for $`C\mathrm{}(M,\eta )`$, and
$$\{\mathrm{𝟏}_4,\gamma _\mu ,\gamma _\mu \gamma _\nu ,\gamma _\mu \gamma _\nu \gamma _\rho ,\gamma _0\gamma _1\gamma _2\gamma _3\}$$
is a basis for $`(4)`$. In addition, these bases satisfy the respective Clifford algebra relations
$`\gamma _\mu \gamma _\nu +\gamma _\nu \gamma _\mu `$ $`=2\eta _{\mu \nu }\mathrm{𝟏}_4\text{,}`$
$`𝐞^\mu 𝐞^\nu +𝐞^\nu 𝐞^\mu `$ $`=2\eta ^{\mu \nu },`$ (2)
where $`\mathrm{𝟏}_4`$ $`(4)`$ is the identity matrix, and $`\eta ^{\mu \nu }=\mathrm{diag}(1,1,1,1)`$. When there is no opportunity for confusion we shall omit the $`\mathrm{𝟏}_4`$ identity matrix in our formulas. We observe that the Clifford product of two Clifford fields is denoted by juxtaposition of symbols. For the orthonormal vector fields $`𝐞^\mu `$ and $`𝐞^\nu `$, $`\mu \nu ,`$ their Clifford product $`𝐞^\mu 𝐞^\nu `$ is equal to the exterior product of those vectors, i.e., $`𝐞^\mu 𝐞^\nu =𝐞^\mu 𝐞^\nu =𝐞^{\mu \nu }`$. Also, for $`\mu \nu \rho ,`$ $`𝐞^\mu {}_{}{}^{\nu }{}_{}{}^{\rho }=`$ $`𝐞^\mu 𝐞^\nu 𝐞^\rho `$, etc. More details on our notations, if needed can be found in .
Since we are interested only in the algebraic classification of spinor fields it is helpful, in order to consistently perform calculations with algebraic methods known to the majority of physicists, to introduce operator fields associated with the bilinear covariant fields. In a fixed spin frame these operator fields are, for each $`xM`$, mappings $`^4^4.`$ They will be represented by the same symbols, since from this usage (hopefully) no confusion will result. So, in what follows the bilinear covariants are considered as being the following operator fields:
$`\sigma `$ $`=\psi ^{}\gamma _0\psi ,𝐉=J_\mu \gamma ^\mu =\psi ^{}\gamma _0\gamma _\mu \psi \gamma ^\mu ,𝐒=S_{\mu \nu }\gamma ^{\mu \nu }={\displaystyle \frac{1}{2}}\psi ^{}\gamma _0i\gamma _{\mu \nu }\psi \gamma ^{\mu \nu },`$
$`𝐊`$ $`=\psi ^{}\gamma _0i\gamma _{0123}\gamma _\mu \psi \gamma ^\mu ,\omega =\psi ^{}\gamma _0\gamma _{0123}\psi .`$ (3)
In the case of the electron, described by Dirac spinor fields (classes 1, 2 and 3 below), $`𝐉`$ is a future-oriented timelike current vector which gives the current of probability. This means that the Clifford product of $`𝐉\mathrm{sec}{\displaystyle \stackrel{1}{}}TMC\mathrm{}(M,\eta )`$ with itself, i.e., $`𝐉^2`$ is such that
$$𝐉^2=J_\mu 𝐞^\mu J_\nu 𝐞^\nu =J_\mu J_\nu \frac{1}{2}(𝐞^\mu 𝐞^\nu +𝐞^\mu 𝐞^\nu )=\eta ^{\mu \nu }J_\mu J_\nu =J_\mu J^\mu >0.$$
(4)
Of course, if $`𝐉:Mx(4)`$ is interpreted as a vector (field) operator we have $`𝐉^2=J_\mu J^\mu \mathrm{𝟏}_4`$. In this case writing $`𝐉^2>0`$ means $`J_\mu J^\mu >0`$.
Moreover, the bivector $`𝐒`$ is associated with the distribution of intrinsic angular momentum, and the spacelike vector $`𝐊`$ is associated with the direction of the electron spin. For a detailed discussion concerning such entities, their relationships and physical interpretation, and generalizations, see, e.g., .
The bilinear covariants satisfy the Fierz identities
$$𝐉^2=\omega ^2+\sigma ^2,𝐊^2=𝐉^2,𝐉𝐊=0,𝐉𝐊=(\omega +\sigma \gamma _{0123})𝐒.$$
(5)
and also satisfy<sup>3</sup><sup>3</sup>3Note that S<sup>-1</sup> exists of course only if $`\omega `$ and $`\sigma `$ are not simultaneously null.:
$`𝐒\mathrm{}𝐉`$ $`=\omega 𝐊𝐒\mathrm{}𝐊=\omega 𝐉,(\gamma _{0123}𝐒)\mathrm{}𝐉=\sigma 𝐊,`$
$`(\gamma _{0123}𝐒)\mathrm{}𝐊`$ $`=\sigma 𝐉,𝐒\mathrm{}𝐒=\omega ^2+\sigma ^2,(\gamma _{0123}𝐒)\mathrm{}𝐒=2\omega \sigma ,`$
$`\mathrm{𝐉𝐒}`$ $`=(\omega +\sigma \gamma _{0123})𝐊,\mathrm{𝐊𝐒}=(\omega +\sigma \gamma _{0123})𝐉,\mathrm{𝐒𝐉}=(\omega \sigma \gamma _{0123})𝐊`$
$`\mathrm{𝐒𝐊}`$ $`=(\omega \sigma \gamma _{0123})𝐉,𝐒^2=(\omega \sigma \gamma _{0123})^2=\omega ^2\sigma ^22\omega \sigma \gamma _{0123},`$
$`𝐒^1`$ $`=𝐒{\displaystyle \frac{(\sigma \omega \gamma _{0123})^2}{(\omega ^2+\sigma ^2)^2}}={\displaystyle \frac{\mathrm{𝐊𝐒𝐊}}{(\sigma ^2+\omega ^2)^2}}.`$ (6)
In the formulas above $``$ denotes the scalar product and $``$ refers to the exterior product, while $`\mathrm{}`$ to the right contraction product of Clifford fields (or Clifford operators). For details, please consult, e.g., . Introduce the complex multivector field $`Z\mathrm{sec}\mathrm{}(M,\eta )`$ (where $`\mathrm{}(M,\eta )`$ denotes the complexified spacetime Clifford bundle, in which the typical fiber is $`_{1,3}_{4,1}`$ ) and the corresponding complex multivector operator (represented by the same letter):
$$Z=\sigma +𝐉+i𝐒+i𝐊\gamma _{0123}+\omega \gamma _{0123.}$$
(7)
When the multivector operators $`\sigma ,\omega ,𝐉,𝐒,𝐊`$ satisfy the Fierz identities, then the complex multivector operator $`Z`$ is denominated a *Fierz aggregate*, and, when $`\gamma _0Z^{}\gamma _0=Z`$, which means that $`Z`$ is a Dirac self-adjoint aggregate<sup>4</sup><sup>4</sup>4It is equivalent to say that $`\omega ,\sigma ,𝐉,𝐊,𝐒`$ are real multivectors., $`Z`$ is called a *boomerang*.
A spinor field such that *not both* $`\omega `$ and $`\sigma `$ are null is said to be regular. When $`\omega =0=\sigma `$, a spinor field is said to be singular. In this case the Fierz identities are in general replaced by the more general conditions (which obviously also holds for $`\omega ,\sigma 0`$). These conditions are:
$`Z^2`$ $`=4\sigma Z,Z\gamma _\mu Z=4J_\mu Z,Zi\gamma _{\mu \nu }Z=4S_{\mu \nu }Z,`$
$`Zi\gamma _{0123}\gamma _\mu Z`$ $`=4K_\mu Z,Z\gamma _{0123}Z=4\omega Z.`$ (8)
Now, any spinor field (regular or singular) can be reconstructed from its bilinear covariants as follows. Take an arbitrary spinor field $`\xi `$ satisfying $`\xi ^{}\gamma _0\psi 0.`$ Then the spinor field $`\psi `$ and the multivector field $`Z\xi `$, differ only by a phase. Indeed, it can be written as
$$\psi =\frac{1}{4N}e^{i\alpha }Z\xi ,$$
(9)
where $`N=\frac{1}{2}\sqrt{\xi ^{}\gamma _0Z\xi }`$ and $`e^{i\alpha }=\frac{1}{N}\xi ^{}\gamma _0\psi `$. For more details see, e.g., .
Lounesto spinor field classification is given by the following spinor field classes , where in the first three classes it is implicit that $`𝐉`$, $`𝐊`$, $`𝐒`$ $`0`$:
1. $`\sigma 0,\omega 0`$.
2. $`\sigma 0,\omega =0`$.
3. $`\sigma =0,\omega 0`$.
4. $`\sigma =0=\omega ,𝐊0,𝐒0`$.
5. $`\sigma =0=\omega ,𝐊=0,𝐒0`$.
6. $`\sigma =0=\omega ,𝐊0,𝐒=0`$.
The current density $`𝐉`$ is always non-zero. Type 1, 2 and 3 spinor fields are denominated Dirac spinor fields for spin-1/2 particles and type 4, 5, and 6 are respectively called flag-dipole, flagpole and Weyl spinor fields. Majorana spinor fields are a particular case of a type 5 spinor field. It is worthwhile to point out a peculiar feature of types 4, 5 and 6 spinor fields: although $`𝐉`$ is always non-zero, we have $`𝐉^2=𝐊^2=0`$. We shall see, below, that the bilinear covariants related to an ELKO spinor field, satisfy $`\sigma =0=\omega ,𝐊=0,𝐒0`$ and $`𝐉^2=0`$.
Lounesto proved that there are no other classes based on distinctions between bilinear covariants. So, ELKO spinor fields must belong to one of the six classes.
Before ending this section we remark that the sum of two spinor fields belonging to a given Lounesto class is not necessarily a spinor field of the same class, as it is easy to verify.
## 3 ELKO Spinor Fields
In this section we explore in details the algebraic properties of ELKO spinor fields as defined in .
A ELKO spinor field $`\mathrm{\Psi }`$ corresponding to a plane wave with momentum $`p=(p^0,𝐩)`$ can be written, without loss of generality, as $`\mathrm{\Psi }=\psi e^{ipx}`$ (or $`\mathrm{\Psi }=\psi e^{ipx}`$) with
$$\psi =\left(\genfrac{}{}{0pt}{}{i\mathrm{\Theta }\varphi _L^{}(𝐩)}{\varphi _L(𝐩)}\right),$$
(10)
where, given the rotation generators denoted by $`𝔍`$, the Wigner’s spin-1/2 time reversal operator $`\mathrm{\Theta }`$ satisfies $`\mathrm{\Theta }𝔍\mathrm{\Theta }^1=𝔍^{}`$. It is useful to choose $`i\mathrm{\Theta }=\sigma _2`$, as in , in such a way that it is possible to express
$$\psi =\left(\genfrac{}{}{0pt}{}{\sigma _2\varphi _L^{}(𝐩)}{\varphi _L(𝐩)}\right).$$
(11)
Here, as in , the Weyl representation of $`\gamma ^\mu `$ is used, i.e.,
$$\gamma _0=\gamma ^0=\left(\begin{array}{cc}0& \mathrm{𝟏}_2\\ \mathrm{𝟏}_2& 0\end{array}\right),\gamma _k=\gamma ^k=\left(\begin{array}{cc}0& \sigma _k\\ \sigma _k& 0\end{array}\right),$$
(12)
where
$$\sigma _1=\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right),\sigma _2=\left(\begin{array}{cc}0& i\\ i& 0\end{array}\right),\sigma _3=\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right)$$
(13)
are the Pauli matrices.
Omitting the subindex of the spinor $`\varphi _L(𝐩)`$, which is denoted heretofore by $`\varphi `$, the left-handed spinor field $`\varphi _L(𝐩)`$ can be represented by
$$\varphi =\left(\genfrac{}{}{0pt}{}{\alpha (𝐩)}{\beta (𝐩)}\right),\alpha (𝐩),\beta (𝐩).$$
(14)
Now using Eqs.(3) it is now possible to calculate explicitly the bilinear covariants for ELKO spinor fields:
$`\sigma `$ $`=\psi ^{}\gamma _0\psi =0,`$ (15)
$`\omega `$ $`=\psi ^{}\gamma _0\gamma _{0123}\psi =0`$ (16)
$`𝐉`$ $`=J_\mu \gamma ^\mu =\psi ^{}\gamma _0\gamma _\mu \psi \gamma ^\mu `$
$`=2(\alpha \beta ^{}+\alpha ^{}\beta )\gamma ^1+2i(\alpha ^{}\beta \alpha \beta ^{})\gamma ^2+2(\beta \beta ^{}\alpha \alpha ^{})\gamma ^3`$
$`+2(\alpha \alpha ^{}+\beta \beta ^{})\gamma ^0,`$ (17)
$`𝐊`$ $`=K_\mu \gamma ^\mu =\psi ^{}i\gamma _{123}\gamma _\mu \psi \gamma ^\mu =0,`$ (18)
$`𝐒`$ $`={\displaystyle \frac{1}{2}}S_{\mu \nu }\gamma ^{\mu \nu }={\displaystyle \frac{1}{2}}\psi ^{}\gamma _0i\gamma _{\mu \nu }\psi \gamma ^{\mu \nu }`$
$`={\displaystyle \frac{i}{2}}((\alpha ^{})^2+(\beta ^{})^2\beta ^2\alpha ^2)\gamma ^{02}+{\displaystyle \frac{1}{2}}((\alpha ^{})^2+(\beta ^{})^2+\beta ^2+\alpha ^2)\gamma ^{31}`$
$`+{\displaystyle \frac{1}{2}}((\beta ^{})^2+\beta ^2(\alpha ^{})^2\alpha ^2)\gamma ^{01}+{\displaystyle \frac{i}{2}}(\beta ^2\alpha ^2+(\alpha ^{})^2+(\beta ^{})^2)\gamma ^{02}`$
$`+(\alpha \beta +\alpha ^{}\beta ^{})\gamma ^{03}+{\displaystyle \frac{i}{2}}(\alpha \beta \alpha ^{}\beta ^{})\gamma ^{12}+{\displaystyle \frac{i}{2}}(\beta ^2\alpha ^2+(\alpha ^{})^2(\beta ^{})^2)\gamma ^{23}.`$ (19)
From the formulas in Eqs.(17, 18) it is trivially seen that that
$$𝐉𝐊=0.$$
(20)
Also, from Eq.(17) it follows that
$$𝐉^2=0,$$
and it is immediate that all Fierz identities introduced by the formulas in Eqs.(5) are trivially satisfied. It also follows directly from Eq.(19) (or easier yet, using the formula for $`𝐒^2`$ in Eq.(6) and Eq.(16)) that $`𝐒^2=0.`$
Now, any flagpole spinor field is an eigenspinor of the charge conjugation operator , here represented by $`𝒞\psi =\gamma ^2\psi ^{}`$. We must have:
$$\gamma ^2\psi ^{}=\lambda \psi ,|\lambda |\text{ }=1.$$
(21)
where $`\lambda `$ is a complex number of unitary modulus. Using Eq.(11) it follows that
$`\gamma ^2\psi ^{}`$ $`=\left(\begin{array}{cc}0& \sigma _2\\ \sigma _2& 0\end{array}\right)\left({\displaystyle \genfrac{}{}{0pt}{}{(\sigma _2\varphi ^{})^{}}{\varphi ^{}}}\right)`$
$`=\left({\displaystyle \genfrac{}{}{0pt}{}{\sigma _2\varphi ^{}}{\sigma _2\sigma _2^{}\varphi }}\right)`$
$`=\psi .`$ (22)
Now, recall that a Majorana spinor field (in Minkowski spacetime) is defined as an eigenvector of the charge operator, with $`\lambda =\pm 1`$. It follows, as can be easy verified that its structure implies immediately that it must be a flagpole spinor field, i.e., of Lounesto class 5 spinor fields. . So ELKO spinor field belongs to the same class as Majorana spinor fields, i.e., class 5 spinor fields, by Lounesto spinor field classification. So, the question arises: what is the difference between ELKO and Majorana spinor fields?
## 4 Helicities
Consider any class 5 spinor field $`\psi =\left(\genfrac{}{}{0pt}{}{\varphi _1}{\varphi _2}\right)`$. We already stated that any such $`\psi `$ satisfy the equation $`𝒞\psi =\lambda \psi `$ with $`\lambda \lambda ^{}=1`$. Such condition implies that $`\lambda \varphi _1=\sigma _2\varphi _2^{}`$.
Let as usual $`\sigma \widehat{𝐩}`$ be the helicity operator acting on 2-component spinor fields. Suppose now that $`\sigma \widehat{𝐩}\varphi _2=\varphi _2`$ (respectively, $`\sigma \widehat{𝐩}\varphi _1=\varphi _1`$). Then a trivial calculation shows that $`\sigma \widehat{𝐩}\varphi _1=\varphi _1`$ (respectively $`\sigma \widehat{𝐩}\varphi _2=\varphi _2`$), i.e., the 2-component spinor fields presented in the structure of a class 5 spinor field have necessarily opposite helicities.
Of course, this is also the case of an ELKO spinor field, since we have just proved that they belong to class 5 spinor fields. We are now prepared to give the answer to the question formulated at the end of the previous section, according to Ahluwalia-Khalilova and Grumiller . They asserted that the difference between ELKO and Majorana spinor fields resides in the fact that the 2-component spinor fields entering the structure of a Majorana spinor field have the same helicity. They attribute this statement, e.g., to Peskin and Schroeder and Marshak and Sudarshan . Of course, this assumption if used by or by any other author must be considered completely ad hoc from the algebraic point of view. There is no justification for it, except an eventual desire to give to Majorana particles a well-defined helicity, something that is not endorsed by the Mathematics of spinor fields. Reading carefully Peskin and Schroeder’s book we found that those authors propose as an exercise<sup>5</sup><sup>5</sup>5Exercise 3.4, page 73 of . the possibility of writing a field equation for a 2-component spinor field of definite helicity (with Grassmann algebra-valued entries) encoding the contents of a Majorana field. In part (e) of that exercise those authors call the 2-component spinor field (of definite helicity) a Majorana field. However, the true Majorana spinor field is a 4-component spinor field, eigenvector of the charge operator and thus, as already proved, it must be composed by two 2-component spinor fields of opposite helicities. At the heart of the issue it is a real confusion between the concepts of chirality and helicity for massive fermions. Indeed, we can find papers, e.g., one by Hannestad where it is stated that the Majorana quantum spinor field is indeed without definite chirality (i.e., “it is a linear combination of left handed and right handed parts”, these parts understood as chiral parts of a spinor field) but that the quantum state of a Majorana particle may be of definite helicity. Hannestad endorses his statement quoting the theory of Majorana particles as derived in the book by Mohapatra and Pal and also on the book by Kim and Pevsner . Also, Plaga stated that Majorana fermions may have states of definite helicity. However, adding power to the confusion he indeed exhibits a “Majorana field” with definite helicity, something that according to our view is equivocated. Plaga also stated that physical states cannot have definite chirality. A complete discussion of these issues will be postponed to another paper. Here we only quote that Ahluwalia-Khalilova and Grumiller showed that adhering to the correct mathematical result leads to interesting physical consequences, as, e.g., the issue of non-locality (see also ).
## 5 Concluding Remarks
We showed that ELKO spinor fields belong to the class of flagpole, class 5, spinor fields, according to Lounesto classification. We showed moreover that algebraic constraints imply that any class 5 spinor field is such that their 2-component spinor fields have opposite helicities. The statements that attributes to asserting that Majorana spinor fields, (a particular class 5 spinor field) are such that their 2-component spinor fields have the same helicity seems to be ad hoc. Moreover, it is easy to verify that<sup>6</sup><sup>6</sup>6As first observed by Ahluwalia-Khalilova and Grumiller in . while the *anticommutator* between the charge conjugation and parity operators acting on a Dirac spinor field is equal to zero, the *commutator* of those operators acting on an ELKO spinor field (which do not satisfy Dirac equation) is also zero. ELKO dynamics is to be analyzed in a forthcoming paper. Also, a relation between Lounesto’s classification and Wigner’s classification of spinor fields (which plays an important role in ) needs some further study and will be presented elsewhere.
Finally, we take the opportunity to call the reader’s attention to the fact that no use has been made until now (to the best of our knowledge) of class 4 spinor fields. Eventually they may be the important spinor fields to describe dark matter and/or dark energy. This possibility will be to explored elsewhere.
## Acknowledgements
The authors are grateful to Professor D. V. Ahluwalia-Khalilova for having asked to us to study the nature of ELKO spinor fields and to Dr. R. A. Mosna for his helpful suggestions. R. da Rocha is also grateful to CAPES for financial support.
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# Bubbling Orientifolds
## 1 Introduction
The recent discovery of an infinite set of $`\frac{1}{2}`$-BPS geometries of type IIB supergravity is of great interest. This discovery provides a correspondence between semiclassical states of the matrix harmonic oscillator and $`\frac{1}{2}`$-BPS geometries, in which the ground state of the matrix model (equivalent to a fermi fluid in a harmonic oscillator potential) corresponds to the well-known $`AdS_5\times S^5`$ geometry. Also, the $`pp`$-wave geometry emerges in the limit of relativistic fermions. A similar construction also exists for $`\frac{1}{2}`$-BPS backgrounds of 11-dimensional supergravity.
These geometries can have various interpretations, as D-branes and/or giant gravitons/dual giant gravitons, depending on the typical sizes of various regions. Their relation to a fermi fluid profile arises from the fact that one can pick an arbitrary shading of the complex plane into black and white regions (denoting occupied/unoccupied regions of fermion phase space) and input this data to construct a $`\frac{1}{2}`$-BPS supergravity solution with fluxes. The matrix oscillator in turn arises (see also Ref.) in the $`𝒩=4`$ supersymmetric Yang-Mills theory on a D3-brane wrapped on $`S^3`$, which has a coupling of the form $`\sqrt{g}R\varphi _i^2`$ for each complex scalar field $`\varphi _i`$ of the $`𝒩=4`$ supermultiplet.
This discovery provides a somewhat new insight into holography. Not only do $`\frac{1}{2}`$-BPS operators of the $`𝒩=4`$ SYM theory correspond to excited states of supergravity on $`AdS_5\times S^5`$, as was already known for a long time, but they can also be identified with entirely new $`\frac{1}{2}`$-BPS geometries<sup>3</sup><sup>3</sup>3 The LLM framework thus provides a tool for studying supergravity excitations of $`AdS_5\times S^5`$ when the backreaction is large. In other circumstances, when the backreaction is limited, other techniques may be more useful in the holographic context. One such example is the computations of anomalous dimensions in super Yang-Mills theory using spin chains, as pioneered by and reviewed for example in Refs.. . These geometries can possess quite different topologies and one can understand topology change as the process of fermi fluid profiles merging and separating. The phase space of the fermi fluid is realised on 2 of the 9 space dimensions of 10-dimensional type IIB supergravity. This implies that the gravity configuration space is partly noncommutative, explanations of which have been proposed in Refs. (related observations can be found in Ref.).
In the present work we would like to investigate the bubbling geometry construction in the presence of an intrinsically stringy effect, namely the orientifold plane. This ubiquitous object has led to important insights in the past. In gauge theory terms it can be used to modify the gauge group from $`SU(N)`$ to $`SO(2N),SO(2N+1)`$ or $`Sp(2N)`$ without breaking any supersymmetry. From the gravity point of view, parallel orientifold planes and D-branes break the same supersymmetries, so it is not surprising that we will find bubbling orientifolded geometries that are $`\frac{1}{2}`$-BPS. This adds a nontrivial class of new geometries to the ones proposed in Ref.. The “ground state” geometry in this case is $`AdS_5\times RP^5`$, an example studied in depth in Ref..
Besides providing more general examples of “bubbling geometries”, the introduction of orientifold planes in this context is likely to be of practical interest. For compactifications of string theory with space-filling branes and/or fluxes, it is well-known that the absence of tadpoles requires space-filling orientifolds. In Ref., it has been noted that the $`T^6/Z_2\times R^{3,1}`$ orientifold compactification has precisely $`AdS_5\times RP^5`$ as the local geometry around the space-filling D3-branes if they are brought near an orientifold. So besides representing a $`\frac{1}{2}`$-BPS geometry by itself, $`AdS_5\times RP^5`$ can also be thought of as a possible local geometry for a $`\frac{1}{2}`$-BPS compactification. This suggests a more general class of $`\frac{1}{2}`$-BPS compactifications (involving both branes and fluxes) for which the spacetimes discussed in the present paper could appear as local geometries.
In what follows we first briefly review bubbling geometries as well as orientifold 3-planes. Following this we describe the bubbling orientifold geometries and their realisation in terms of free fermions on a phase space that is the upper half-plane. We will see that the orientifold plane can be interpreted as the wall which truncates a regular harmonic oscillator potential to a “half-oscillator”, a familiar problem in quantum mechanics. We comment on the topology and other characteristics of the most general bubbling solutions, including discrete torsion and the possibility of two different types of topology change. Finally we briefly examine the analogue solutions in $`M`$-theory. In the Appendix, we collect some useful facts about random matrices in a harmonic oscillator potential, for the cases when the matrices are in the Lie algebra of $`SO(2N+1),Sp(2N),SO(2N)`$. These matrix models compute the $`\frac{1}{2}`$-BPS states of $`𝒩=4`$ super-Yang-Mills theory with the corresponding gauge groups.
## 2 Bubbling Geometries: A Brief Review
The type IIB $`\frac{1}{2}`$-BPS bubbling geometries of Ref. can be summarised in a family of classical solutions of type IIB supergravity given in terms of a single function $`z(y;x_1,x_2)`$. The metrics contain two 3-spheres $`S^3`$ and $`\stackrel{~}{S}^3`$ and are given by:
$`ds^2`$ $`=`$ $`\left[{\displaystyle \frac{y^2}{\frac{1}{4}z^2}}\right]^{\frac{1}{2}}(dt+V_idx_i)^2+\left[{\displaystyle \frac{\frac{1}{4}z^2}{y^2}}\right]^{\frac{1}{2}}(dy^2+dx_idx_i)`$ (1)
$`+\left[y^2{\displaystyle \frac{\frac{1}{2}+z}{\frac{1}{2}z}}\right]^{\frac{1}{2}}d\mathrm{\Omega }_3^2+\left[y^2{\displaystyle \frac{\frac{1}{2}z}{\frac{1}{2}+z}}\right]^{\frac{1}{2}}d\stackrel{~}{\mathrm{\Omega }}_3^2`$
Additionally the solution has suitable 5-form fluxes. The function $`z=z(y;x_1,x_2)`$ satisfies:
$$(_y^2+_i_i)z\frac{1}{y}_yz=0$$
(2)
and $`V_i`$ is a vector field determined in terms of $`z`$ via:
$$y_yV_i=ϵ_{ij}_jz,y_{[i}V_{j]}=ϵ_{ij}_yz$$
(3)
¿From the metric, it is clear that the range of the $`y`$ coordinate is restricted to $`y0`$. Also as $`y0`$, smoothness of the metric requires $`z\pm \frac{1}{2}`$ with $`y/\sqrt{\frac{1}{2}z}`$ fixed. Then if $`z+\frac{1}{2}`$, the 3-sphere $`\stackrel{~}{S}^3`$ with metric proportional to $`d\stackrel{~}{\mathrm{\Omega }}_3^2`$ shrinks to zero size, while the other 3-sphere $`S^3`$ with metric $`d\mathrm{\Omega }_3^2`$ remains finite. When $`z\frac{1}{2}`$ the reverse is true, and it is $`S^3`$ that shrinks to zero size. More details are given in Ref.. The geometries are parametrised by drawing smooth contours in the $`x_1,x_2`$ plane which demarcate regions of $`z=\frac{1}{2}`$ (black) from regions of $`z=+\frac{1}{2}`$ (white). In turn, these regions can be interpreted as semiclassical phase-space profiles of a fermi fluid where the black regions are occupied and the white regions are unoccupied. The constant phase-space densities of these droplets map to constant densities in the $`x_1,x_2`$ plane. This follows from the fact that the flux through a sphere formed by drawing a surface ending on a black/white region can be converted to an integral of a constant flux density over the region.
Fermi fluid profiles consisting of a black disc with a circular boundary correspond to the $`AdS_5\times S^5`$ solution. Deformations of this by adding/removing thin concentric shells correspond to giant gravitons or dual giant gravitons in this spacetime. But the most generic geometry can have very little to do with $`AdS_5\times S^5`$ and may contain, for example, an arbitrary number of $`S^5`$ factors.
The bubbling construction can be extended to $`\frac{1}{2}`$-BPS solutions of 11-dimensional supergravity. In this case one finds metrics parametrised in terms of a function $`D(y;x_1,x_2)`$. In what follows, we denote $`_yD`$ by $`D^{}`$. These metrics contain a five-sphere $`S^5`$ and a two-sphere $`\stackrel{~}{S}^2`$, and are given by:
$`ds^2`$ $`=`$ $`4\left[{\displaystyle \frac{y}{D^{}(1yD^{})^2}}\right]^{\frac{1}{3}}(dt+V_idx^i)^2+\left[{\displaystyle \frac{D_{}^{}{}_{}{}^{2}(1yD^{})}{y^2}}\right]^{\frac{1}{3}}\left(dy^2+e^D(dx_1^2+dx_2^2)\right)`$ (4)
$`+4\left[{\displaystyle \frac{y(1yD^{})}{D^{}}}\right]^{\frac{1}{3}}d\mathrm{\Omega }_5^2+\left[{\displaystyle \frac{y^2D^{}}{(1yD^{})}}\right]^{\frac{2}{3}}d\stackrel{~}{\mathrm{\Omega }}_2^2`$
along with a suitable 4-form flux. The function $`D(y;x_1,x_2)`$ satisfies the three-dimensional Toda equation:
$$_i_iD+_y^2e^D=0$$
(5)
and the vector field is determined in terms of $`D`$ by:
$$V_i=\frac{1}{2}ϵ_{ij}_jD$$
(6)
In this class of metrics too, the range of the coordinate $`y`$ is restricted to $`y0`$ and as $`y0`$, the function $`D`$ must obey one of two boundary conditions. One possibility is $`D^{}y`$ while the other is $`D^{}y^1`$. More precisely we have:
(i) $`Dy^2+g(x_1,x_2)`$
(ii) $`D\mathrm{log}y+h(x_1,x_2)`$ (7)
where $`g,h`$ are functions of $`(x_1,x_2)`$ satisfying the two-dimensional Liouville equation:
$$_i_ig(x_1,x_2)+e^{g(x_1,x_2)}=0$$
(8)
and similarly for $`h(x_1,x_2)`$.
In case (i), keeping $`D`$ finite in the limit, it is evident from the metric Eq. (4) that $`\stackrel{~}{S}^2`$ shrinks to zero size while $`S^5`$ remains finite. In case (ii) we keep $`D^{}y^1`$ finite in the limit and find that $`S^5`$ shrinks with $`\stackrel{~}{S}^2`$ remaining finite. Thus again the $`x_1,x_2`$ plane is divided into two types of regions, which can be associated to fermi fluid droplets. There is an important subtlety: the density of fermions in a droplet is not constant, unlike in the type IIB case. In fact the densities of the droplets are
(i) $`\rho (x_1,x_2)=2e^{g(x_1,x_2)}`$
(ii) $`\rho (x_1,x_2)=2e^{h(x_1,x_2)}.`$ (9)
However, as argued by the authors of Ref., two droplets related by a conformal mapping of the plane give rise to the same $`\frac{1}{2}`$-BPS solution. Therefore the topology of the bubbles in the $`(x_1,x_2)`$ plane is expected to be the same as that in the fermion phase space. In particular, they demonstrate that the M-theory pp-wave is given by one of the two boundary conditions for $`x_2<0`$ and the other one for $`x_2>0`$, just as one would expect from the correspondence of this system with the relativistic limit of free fermions.
## 3 The $`AdS_5\times RP^5`$ Orientifold
In Ref., $`\frac{1}{2}`$-BPS excitations of a ground state geometry corresponding to $`AdS_5\times S^5`$ in type IIB string theory were considered, and the full back reaction on the geometry was determined. We will instead be interested in an $`AdS_5\times RP^5`$ ground state, arising from a $`Z_2`$ orientifold projection of the $`S^5`$, as considered by Witten.
Introducing an orientifold plane changes the gauge group on a stack of D-branes to $`SO(2N)`$, $`SO(2N+1)`$ or $`Sp(2N)`$, instead of previously $`SU(N)`$. The orientifolded theory has no fixed points on the $`S^5`$, so there is no open string sector. Being an orientifold, traversing a non-contractible loop flips the orientation of the world-sheet (or any embedded orientable manifold). This means that there is also no “winding sector”. The spectrum therefore consists only of those $`AdS_5\times S^5`$ states which are invariant under the orientifold projection, at least for weak coupling<sup>4</sup><sup>4</sup>4 On the other hand, interactions will become more complicated, due to contributions from non-orientable world-sheets..
The topology of the 2-form fields $`B_{NSNS}`$ and $`B_{RR}`$ is in general non-trivial, in the sense that the field strength $`H=dB`$ belongs to some equivalence class of the cohomology group $`H^3(AdS_5\times RP^5,\stackrel{~}{Z})=H^3(RP^5,\stackrel{~}{Z})=Z_2`$, where $`\stackrel{~}{Z}`$ denotes twisted coefficients, i.e. coefficients flipping sign along a non-contractible loop. To preserve supersymmetry, it is necessary to choose $`H=0`$. We have already said that traversing a non-contractible loop flips the orientation of the world-sheet $`\mathrm{\Sigma }`$. Again using twisted coefficients, this means that the homology is $`\mathrm{\Sigma }H_2(RP^5,\stackrel{~}{Z})=Z_2`$, and we may represent the non-trivial element by $`RP^2`$. Consequently, for non-trivial choice of discrete torsion for the field strength, the contribution to the path integral is $`e^{i_\mathrm{\Sigma }B_{NSNS}}=(1)^s`$, where $`s`$ is the number of $`RP^2`$’s making up the world-sheet. Trivial choice of torsion always contributes a factor of unity.
There will be four distinct combinations of discrete torsion which can be labelled by
$`\theta _{NSNS}`$ $`=`$ $`{\displaystyle \frac{1}{2\pi }}{\displaystyle _\mathrm{\Sigma }}B_{NSNS}`$
$`\theta _{RR}`$ $`=`$ $`{\displaystyle \frac{1}{2\pi }}{\displaystyle _\mathrm{\Sigma }}B_{RR},`$ (10)
so that $`\theta _{NSNS}`$ and $`\theta _{RR}`$ can independently take the values $`0,\frac{1}{2}`$. The $`(0,0)`$ theory can be shown to be $`SL(2,Z)`$ invariant. The Montonen-Olive duality of such a theory identifies it as $`SO(2N)`$. Turning on $`\theta _{NSNS}=\frac{1}{2}`$ will change the gauge group to $`Sp(2N)`$, which can be seen in the following way. Feynman diagrams of $`Sp(2N)`$ are obtained from those of $`SO(2N)`$ by sending $`NN`$. Since the contribution of a generic diagram goes as $`N^{22gs}`$, where $`g`$ is the genus and $`s`$ is the number of glued copies of $`RP^2`$, we conclude that each $`RP^2`$ contributes a factor $`1`$. By definition of the discrete torsion, this is precisely the effect of turning on $`\theta _{NSNS}`$. This argument therefore indicates that turning on $`\theta _{NSNS}=\frac{1}{2}`$ (with $`\theta _{RR}=0`$ or $`\frac{1}{2}`$) is equivalent to changing the gauge group from $`SO(2N)`$ (or $`SO(2N+1)`$) to $`Sp(2N)`$. The only combination of discrete torsions which is left is $`(\theta _{NSNS},\theta _{RR})=(0,\frac{1}{2})`$, which can be identified with the “remaining” gauge group, $`SO(2N+1)`$.
Interestingly, the torsion need not be constant all over the manifold. To explain this, some background on 3-branes in this geometry is required. In addition to the ordinary $`D3`$ brane, the $`AdS_5\times RP^5`$ theory also contains 3-branes resulting from wrapping a 5-brane on an $`RP^2`$ submanifold of the $`RP^5`$. Any 3-brane is localised in one spatial direction on the $`AdS_5`$ space. The $`D3`$ brane is a source of 5-form flux, so the brane will act as a domain wall, separating $`SU(N)`$ and $`SU(N+1)`$ gauge theories on the boundary. On the orientifold, the flux $`N`$ through $`S^5`$, which covers the $`RP^2`$ twice, will instead shift by two units. Therefore, the gauge group will shift between $`SO(2N)`$ and $`SO(2N+2)`$ or between $`Sp(2N)`$ and $`Sp(2N+2)`$. But crossing a 3-brane made by wrapping a D5 or NS5-brane on $`RP^2`$ also makes the $`\theta `$ angle jump. This is because the wrapped 5-brane also acts as a source for the field $`H=dB`$. A wrapped $`D5`$ brane sources RR-form flux, and so makes $`\theta _{RR}`$ jump. Similarly, a wrapped $`NS5`$ brane has NS-flux on it, making $`\theta _{NSNS}`$ jump upon crossing it. We will discuss how these features relate to the LLM description in section 5.2.
In general, branes characterized by untwisted or twisted charges can only be wrapped on non-trivial cycles with untwisted or twisted coefficients, respectively. In addition, topological considerations may restrict the allowed values of the discrete torsions. In spite of these restrictions, the $`RP^5`$ geometry still introduces new types of objects, as compared to the unorientifolded theory, such as fat strings and Pfaffian particles (on the gauge theory side). They can be constructed by wrapping branes on non-trivial cycles of the $`RP^5`$ which lack any counterpart in the $`S^5`$ theory. We refer to Ref. for details.
Here, we will content ourselves by illustrating some of the general features involved by discussing some properties of the baryon vertex of $`SU(N)`$ gauge theory, which does exist both before and after the projection to the $`RP^5`$ geometry. In the unprojected case, the only non-trivial cycle is the full $`S^5`$ itself. The baryon vertex is constructed by wrapping a $`D5`$ brane on this $`S^5`$. Strings connect it to $`N`$ external quarks. Quarks are particles in the fundamental representation of the gauge group. The combination is fermionic, since it is made gauge invariant by contracting the colour wave functions using an antisymmetric tensor of order $`N`$.
That this makes sense from the gravity point of view can be seen as follows. There is a coupling $`\frac{1}{2\pi }_{S^5\times R}aF_5`$ on the world-volume of the $`D5`$ brane. As always, $`_{S^5}F_5=2\pi N`$, so the charge corresponding to the $`U(1)`$ gauge field $`a`$ gets a contribution of $`N`$ units due to the $`D5`$ brane. Similarly, each of the $`N`$ fundamental strings ending on the $`D5`$ contributes by $`1`$, making the total charge vanish, as required.
The other ends of the strings connect to quarks at the boundary, modelled by attaching their endpoints to a large $`D3`$ brane whose spatial world-volume is of the topology $`S^3\times P`$, where $`P`$ is a point on the $`S^5`$. Placing the wrapped $`D5`$ at the point $`Q`$ in $`AdS_5`$ and considering the “linking” numbers between the manifolds $`S^3\times P`$ and $`Q\times S^5`$, there will be a change of sign associated with interchanging two strings, confirming that the strings are fermionic.
In the projected theory, topological restrictions allow the $`D5`$ to wrap the $`RP^5`$ an even number of times only. Hence, the baryon vertex only couples to an even number of quarks on the orientifold.
## 4 Bubbling Orientifolds and Fermi Fluids
In the previous section we have seen that the spacetime $`AdS_5\times S^5`$ admits an involution that reflects all the directions of the 5-sphere and converts it into the real projective space $`RP^5`$. The dual gauge theory is the $`𝒩=4`$ super-Yang-Mills theory on a set of D3-branes parallel to an orientifold 3-plane. The gauge group can be $`SO(2N)`$, $`SO(2N+1)`$ or $`Sp(2N)`$ depending on the precise type of O3-plane. In keeping with the bubbling geometry idea, we expect that there should be an infinite set of $`\frac{1}{2}`$-BPS bubbling geometries that correspond to the excited states of this gauge theory.
To find these, we need to understand how the fermi fluid profile that corresponds to $`AdS_5\times S^5`$ is affected by the orientifolding. This profile is a circle of radius $`r_0=R_{AdS}^2`$ in the fermion phase space. In the geometry this space is realised as the $`x_1x_2`$ plane. Now $`AdS_5\times S^5`$ is parametrised as:
$$(ds)^2=r_0\left(\mathrm{cosh}^2\rho dt^2+d\rho ^2+\mathrm{sinh}^2\rho d\mathrm{\Omega }_3^2+d\theta ^2+\mathrm{cos}^2\theta d\stackrel{~}{\varphi }^2+\mathrm{sin}^2\theta d\stackrel{~}{\mathrm{\Omega }}_3^2\right)$$
(11)
with $`\theta [0,\frac{\pi }{2}],\stackrel{~}{\varphi }[0,2\pi ]`$. In terms of embedding coordinates in an $`\mathrm{IR}^6`$, one can write:
$`X_1`$ $`=`$ $`R\mathrm{cos}\gamma \mathrm{sin}\alpha \mathrm{sin}\theta `$
$`X_2`$ $`=`$ $`R\mathrm{sin}\gamma \mathrm{sin}\alpha \mathrm{sin}\theta `$
$`X_3`$ $`=`$ $`R\mathrm{cos}\beta \mathrm{cos}\alpha \mathrm{sin}\theta `$
$`X_4`$ $`=`$ $`R\mathrm{sin}\beta \mathrm{cos}\alpha \mathrm{sin}\theta `$
$`X_5`$ $`=`$ $`R\mathrm{cos}\stackrel{~}{\varphi }\mathrm{cos}\theta `$
$`X_6`$ $`=`$ $`R\mathrm{sin}\stackrel{~}{\varphi }\mathrm{cos}\theta `$ (12)
with $`\alpha ,\theta [0,\frac{\pi }{2}]`$ and $`\beta ,\gamma ,\stackrel{~}{\varphi }[0,2\pi ]`$. The embedding condition is $`_iX_i^2=R^2r_0`$.
The orientifold action on this $`\mathrm{IR}^6`$ is $`X_IX_I`$, $`I=1,2,\mathrm{},6`$. In terms of the angular variables this amounts to
$`\beta `$ $``$ $`\beta +\pi `$
$`\gamma `$ $``$ $`\gamma +\pi `$
$`\stackrel{~}{\varphi }`$ $``$ $`\stackrel{~}{\varphi }+\pi .`$ (13)
Going to the $`x_1x_2`$ plane of the bubbling geometry, we recall that it is described by polar coordinates $`(r,\varphi )`$ where
$`r`$ $`=`$ $`r_0\mathrm{cosh}\rho \mathrm{cos}\theta `$
$`\varphi `$ $`=`$ $`\stackrel{~}{\varphi }+t,`$ (14)
and $`\rho ,\stackrel{~}{\varphi }`$ are the coordinates appearing in Eq.(11) above. Therefore, under the orientifolding operation, the $`x_1x_2`$ plane undergoes the involution $`\varphi \varphi +\pi `$, which is the same as the reflection $`(x_1,x_2)(x_1,x_2)`$.
The precise picture is a little more complicated because at $`y=0`$, the full spatial geometry is not really 2-dimensional. In the regions where $`z=\pm \frac{1}{2}`$ (the white and black regions) the geometry is 5-dimensional, and consists of the $`(x_1,x_2)`$ plane together with one of the 3-spheres $`S^3`$ or $`\stackrel{~}{S}^3`$, parametrised respectively by $`d\mathrm{\Omega }_3`$ or $`d\stackrel{~}{\mathrm{\Omega }}_3`$, while the other 3-sphere has shrunk to zero size. The sphere $`\stackrel{~}{S}^3`$ lies inside $`S^5`$ (and is parametrised by the angles $`\alpha ,\beta ,\gamma `$ in Eq.(12)). Thus it is inverted by the orientifold action, while the other 3-sphere $`S^3`$ that lies in $`AdS_5`$ remains unaffected. Thus, at a generic point of the $`(x_1,x_2)`$ plane, the orientifold involution acts by reflecting $`(x_1,x_2)`$ and simultaneously inverting $`\stackrel{~}{S}^3`$. In the white regions, where $`z=\frac{1}{2}`$, the $`\stackrel{~}{S}^3`$ shrinks to zero size, while in the black regions, where $`z=\frac{1}{2}`$, it is $`S^3`$ that shrinks. It follows that in the white regions, the orientifolding operation acts solely by inverting the $`(x_1,x_2)`$ plane and turning it into $`\text{I}\mathrm{C}/\text{ZZ}_2`$. The same is true on boundaries between the black and white regions with $`z=\frac{1}{2},z=+\frac{1}{2}`$ respectively (where both $`S^3,\stackrel{~}{S}^3`$ shrink). In the black region, however, one has to keep in mind that the $`\stackrel{~}{S}^3`$ above a given point of the $`(x_1,x_2)`$ plane is identified with a reversed $`\stackrel{~}{S}^3`$ above the diametrically opposite point. Finally, at the origin $`x_1=x_2=0`$ which is the fixed point of $`\varphi \varphi +\pi `$, the $`\stackrel{~}{S}^3`$ gets an antipodal identification and becomes $`RP^3`$.
The above discussion was based on an involution that is a symmetry of the $`AdS_5\times S^5`$ background. Therefore strictly speaking it applies only to the simple fermi fluid profile consisting of a black disc of fixed radius centred at the origin. But now it is clear how to define orientifolding for the most general bubbling geometry: simply consider all fermi fluid profiles whose boundaries are well-defined on $`\text{I}\mathrm{C}/\text{ZZ}_2`$. Such boundaries have to be invariant under a rotation by $`\pi `$ in the $`(x_1,x_2)`$ plane. Alternatively they may be drawn on the fundamental domain of $`\text{I}\mathrm{C}/\text{ZZ}_2`$: the upper half-plane with the positive and negative halves of the $`x_1`$ axis identified with each other.
What is the fermi system dual to these spacetimes? For this system, the eigenvalue phase space should be an orbifold $`\text{I}\mathrm{C}/\text{ZZ}_2`$. This means the position $`x`$ of the fermion is strictly positive (recall that the identification between coordinate space and phase space is $`(x_1,x_2)(p,x)`$). The identification of the positive and negative momentum axes tells us that at $`x=0`$ the momentum is instantaneously reversed. This is consistent with the harmonic oscillator being truncated to the “half harmonic oscillator” with an infinite vertical wall at $`x=0`$, a well-known system in quantum mechanics. Thus we have a (Hermitian) matrix-valued particle, or equivalently fermionic eigenvalues, moving in this potential. The fermionic wave functions are made in the usual way out of one-particle wave functions, the latter now being the parity-odd solutions of the full harmonic oscillator.
¿From the correspondence between states of the matrix harmonic oscillator and $`\frac{1}{2}`$-BPS excitations of $`𝒩=4`$ SYM, we would expect the appropriate matrix model for an orientifolded geometry to be related to $`𝒩=4`$ SYM with gauge group $`SO(2N)`$<sup>5</sup><sup>5</sup>5Or $`SO(2N+1)`$ or $`Sp(2N)`$.. Then, following the arguments of Refs., one should get a usual, not semi-infinite, harmonic oscillator, the only change being that the Hermitian matrix (in the algebra of $`SU(N)`$) is replaced by an antisymmetric matrix (in the algebra of $`SO(2N)`$).
This appears to be a different theory, but in fact the two descriptions are equivalent, for the same reason that (in flat space) an orientifold plane projects $`SU(N)`$ Chan-Paton factors down to the $`SO`$ or $`Sp`$ subgroup. Consider the random matrix model for $`2N\times 2N`$ real antisymmetric matrices $`A`$ (for simplicity we choose constant matrices) with a potential $`V(A)`$:
$$Z=[dA]e^{\mathrm{𝐭𝐫}V(A)}$$
(15)
This is invariant under the orthogonal transformations, which can be used to reduce the matrix $`A`$ to the form
$`\mathrm{\Lambda }`$ $`=`$ $`\left(\begin{array}{cccc}\lambda _1\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right)& 0& \mathrm{}& 0\\ 0& \lambda _2\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right)& \mathrm{}& 0\\ \mathrm{}& \mathrm{}& & \mathrm{}\\ 0& 0& \mathrm{}& \lambda _N\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right)\end{array}\right)`$ (16)
$`=`$ $`i\sigma _2\mathrm{diag}(\lambda _1,\lambda _2,\mathrm{},\lambda _N).`$
Importantly, all $`\lambda _i`$ are real and can be chosen to be non-negative. For example, in each block, a negative $`\lambda _i`$ can be brought to a positive one by conjugating with $`\sigma _311`$.
The matrix measure
$$[dA]\underset{P<Q}{}dA_{PQ}$$
(17)
can then be shown to reduce to
$`[dA]`$ $`=`$ $`{\displaystyle \underset{j<k=1}{\overset{N}{}}}(\lambda _j^2\lambda _k^2)^2{\displaystyle \underset{i=1}{\overset{N}{}}}d\lambda _i`$ (18)
$`=`$ $`{\displaystyle \underset{j<k=1}{\overset{N}{}}}(\lambda _j\lambda _k)^2{\displaystyle \underset{j<k=1}{\overset{N}{}}}(\lambda _j+\lambda _k)^2{\displaystyle \underset{i=1}{\overset{N}{}}}d\lambda _i.`$
The nontrivial measure factor just corresponds to $`\sqrt{g}`$ for the metric $`g_{IJ}`$ on the space of deformations, in the coordinates made up of the $`\lambda _i`$ and the orthogonal transformations<sup>6</sup><sup>6</sup>6The formula above has appeared in, for example, Ref...
Once we know the metric, we easily derive the kinetic term in the Hamiltonian for this matrix model, which is just the Laplacian on the deformation space,
$`H`$ $`=`$ $`{\displaystyle \underset{i}{}}{\displaystyle \frac{1}{\sqrt{g}}}{\displaystyle \frac{}{\lambda _i}}\sqrt{g}{\displaystyle \frac{}{\lambda _i}}={\displaystyle \frac{1}{_{j<k}(\lambda _j^2\lambda _k^2)^2}}{\displaystyle \underset{i}{}}{\displaystyle \frac{}{\lambda _i}}{\displaystyle \underset{j<k}{}}(\lambda _j^2\lambda _k^2)^2{\displaystyle \frac{}{\lambda _i}}`$ (19)
$`=`$ $`{\displaystyle \frac{1}{_{j<k}(\lambda _j^2\lambda _k^2)}}{\displaystyle \underset{i}{}}{\displaystyle \frac{^2}{\lambda _i^2}}{\displaystyle \underset{j<k}{}}(\lambda _j^2\lambda _k^2).`$
The equality of the two lines above follows from the identity
$$\underset{i}{}\frac{^2}{\lambda _i^2}\left(\underset{j<k}{}(\lambda _j^2\lambda _k^2)\right)=0.$$
(20)
We see that with some changes, this works out much like the case of Hermitian matrix models. Absorbing the factor $`_{j<k}(\lambda _j^2\lambda _k^2)`$ into the wave function makes it fermionic. All this is in agreement with our argument that orientifolding confines the coordinates of the free fermions to a half-line. In fact it tells us something more precise – from Eq.(18), we see that the positive “eigenvalues”<sup>7</sup><sup>7</sup>7One should keep in mind that the actual eigenvalues are $`\pm i\lambda _i`$. $`\lambda _i`$ not only repel each other, but also repel their images $`\lambda _i`$. This is just what we would expect in the presence of an orientifold plane.
In the above we have mostly focused on the “standard” orientifold action that leads to the gauge group $`SO(2N)`$ on the gauge theory side. The above procedure admits a straightforward modification when the orientifold action is modified to produce $`SO(2N+1)`$ or $`Sp(2N)`$. For convenience, we provide a unified derivation of the matrix measures for all the cases $`SO(2N),Sp(2N),SO(2N+1)`$ in an Appendix. As reviewed in section 3, the latter groups arise when discrete $`B_{NSNS}`$ or $`B_{RR}`$ torsion is introduced in spacetime. We will return to this in section 5.2.
## 5 Properties of Orientifolded Fermi Fluids
### 5.1 General Properties
We now discuss some qualitative properties of the orientifolded fermi fluid and associated $`\frac{1}{2}`$-BPS geometries, using the standard orientifold projection of $`SO(2N)`$ type. As we have seen, the allowed fluid profiles have to be invariant under a rotation by $`\pi `$ about the origin of phase space, or equivalently under $`(x,p)(x,p)`$. The resulting space under this quotient is $`\text{I}\mathrm{C}/Z_2`$ and can be represented by the upper half-plane $`x>0`$ along with half of the $`x=0`$ axis, since points $`(0,p)`$ are identified with $`(0,p)`$. Fermi fluid profiles on this space can be represented either as bubbles in a fundamental region, such as the upper half-plane, or as $`Z_2`$-invariant configurations in the whole plane.
The simplest example is the semi-circle centred at the origin, which seeds the geometry $`AdS_5\times RP^5`$. As a fluid profile, it is the ground state of $`N`$ free fermions in the semi-infinite harmonic oscillator. We may note right away that the quotienting destroys translational invariance in the $`xp`$ plane and therefore the semi-circle has to be centred at the origin (Fig. 1). Here and in what follows, the fundamental part of the bubble is shaded while the $`Z_2`$ image is indicated with a dotted line to exhibit the fact that in the “upstairs” space this is a $`Z_2`$ invariant configuration.
We can also have circular configurations such as the one in Fig.2, where no part of the circle touches the horizontal axis. This arises as the near-horizon geometry of D3-branes parallel to, but far away from, an orientifold plane. This example shows that geometries which occur in the un-orientifolded case can also occur in the presence of the orientifold. All that is required is for the boundaries to be completely contained in the upper half-plane. The lack of translational invariance, alluded to above, means that the distance of the disc in this figure from the horizontal axis is physically meaningful, in contrast to the un-orientifolded case.
More general configurations are the union of bubbles of two basic types: those where the origin is included in the black region, and those where it is not. Examples of the two types are shown in Figs.3,4. From Fig.4 we notice that a boundary component of the latter type will be smooth only if, at the two points where it touches the real axis, the slope is the same.
As noted earlier, when the origin is in a black region there is an $`RP^3`$ in the spacetime geometry. This is the case for the bubble in Fig.4. For bubbles that do not include the origin, at first sight there appear to be two types: one illustrated in Fig.3 where the bubbles are completely contained in the upper half plane (along with their images in the lower half plane), and another type as in Fig.5 where the boundaries cross the horizontal axis at $`x=0`$.
The difference between these two examples is only superficial and can be eliminated by making a different choice of fundamental region. Choosing the right half plane as the fundamental region in Fig.5 puts the bubble completely inside the region, with its image on the other side. Thus we see that there are only two types of bubbles. Deforming these two types into each other leads to topology change via singular geometries, as we will see below.
In the models of Ref., particular types of bubbles were identified with “giant gravitons”. These are bubbles consisting of a black region with a small hole inside, the area of the hole being much smaller than the area of the black region surrounding it. The corresponding geometries were interpreted as containing giant gravitons made up of D3-branes wrapping a maximal $`\stackrel{~}{S}^3`$ in $`S^5`$. In the orientifolded case we may instead consider a hole inside the bubble of Fig.1. We see that there are two types of such holes, and correspondingly two types of giant gravitons (Fig.6). If the small hole is in a generic location then we have giant gravitons wrapping a 3-sphere $`\stackrel{~}{S}^3`$ in $`RP^5`$. On the other hand if the hole surrounds the origin, we have giant gravitons wrapping an $`RP^3`$ cycle of $`RP^5`$.
Interestingly, as the hole around the origin becomes large enough to interpret this as a new back-reacted geometry, we find that the $`RP^3`$ cycle has disappeared – for the reason, stated earlier, that the black region does not enclose the origin. Also, when the hole expands further so that the black region forms a thin semicircular ring (stuck to the horizontal axis), we can interpret the configuration as consisting of dual giant gravitons wrapping an $`S^3`$ of $`AdS_5`$, and uniformly distributed around an equator of $`RP^5`$.
### 5.2 Torsion Cycles
A unique feature of the $`AdS_5\times RP^5`$ theory as compared to the unorientifolded case is the appearance of discrete torsion, as explained by Witten and discussed in section 3. This means that there are topologically distinct configurations of the $`B`$-field, described by the theta angles $`\theta _{NSNS}`$ and $`\theta _{RR}`$. In this section, we wish to understand how this feature manifests itself in terms of the free fermion description on the distinguished $`(x_1,x_2)`$ geometry-seeding plane.
As a warm-up, consider first a $`D3`$ brane in $`AdS_5\times S^5`$. We may take it to be localized in the radial direction $`\rho `$ on $`AdS_5`$, as well as located at a point on the $`S^5`$. It then extends along an $`S^3`$ in the $`AdS_5`$. Suppose now that we choose a path $`T`$ whose endpoints are two points $`P`$ and $`Q`$ on opposite sides of the 3-brane, as shown in Fig.7. Consider the integral of $`dF_5`$ over the domain $`T\times S^5`$. Since the brane provides a source for the five-form flux $`F_5`$, and upon using Stokes’s theorem, we find
$`2\pi ={\displaystyle _{T\times S^5}}𝑑F_5={\displaystyle _{P\times S^5}}F_5{\displaystyle _{Q\times S^5}}F_5.`$ (21)
We can conclude that the flux through the $`S^5`$ changes by one upon crossing the brane, changing the the gauge group from $`SU(N)`$ to $`SU(N\pm 1)`$. On the orientifold, the gauge group instead changes from $`SO(2N)`$, $`SO(2N+1)`$ or $`Sp(2N)`$ to $`SO(2N\pm 2)`$, $`SO(2N+1\pm 2)`$ or $`Sp(2N\pm 2)`$ respectively, since $`RP^5`$ is covered twice by $`S^5`$.
¿From the outset, the orientifolded theory has no preferred set of discrete torsion, i.e. theta angles $`(\theta _{NSNS},\theta _{RR})`$. However, once a particular set of values has been assigned to some region, discrete torsions will automatically be induced on the entire space, depending on the distribution of magnetic sources. In the absence of magnetic sources for the $`B`$-fields, the torsions are constant over the entire space.
However, $`D5`$ or $`NS5`$ branes wrapped on $`RP^2`$ cycles, forming effective $`D3`$ branes on $`AdS_5`$, do provide such sources. Crossing such a $`D3`$ brane will, in addition to changing the rank of the gauge group as we just described, also cause the discrete torsion to change.
More concretely, suppose we are wrapping a $`D5`$ brane on an $`RP^2`$ cycle<sup>8</sup><sup>8</sup>8 The case of wrapping an $`NS5`$ brane on an $`RP^2`$ cycle is handled similarly, and will lead to a domain wall across which $`\theta _{NSNS}`$ jumps.. This will result in an effective $`D3`$ brane which additionally has an RR-form flux $`H_{RR}=dB_{RR}`$, which we choose to vanish to preserve supersymmetry. On the $`AdS_5`$ space, suppose the brane is again localized in the radial direction, but extended in the $`S^3`$ directions. On the $`RP^5`$, the brane is extended along an $`RP^2`$. This means that it must be localised on an $`RP^3`$, at least locally. The $`RP^2`$ and the $`RP^3`$ are generically intersecting at one point. Consider, then, the integral of $`dH`$ over $`T\times RP^3`$, where $`T`$ is the same path as before. Similar to the $`dF_5`$ integral dealt with previously, we now find
$`2\pi ={\displaystyle _{T\times RP^3}}𝑑H={\displaystyle _{P\times RP^3}}H{\displaystyle _{Q\times RP^3}}H.`$ (22)
Hence, the theta angle $`\theta _{RR}`$ jumps upon crossing the brane.
Our objective is now to understand what this picture corresponds to on the LLM plane $`(x_1,x_2)`$ of . An $`AdS_5\times RP^5`$ background geometry implies that we should start out with a “semicircular” black region centred at the origin. The quotation marks indicate that the correct designation of the geometry in the LLM plane requires taking the non-trivial identifications into account.
In general, an $`RP^i`$ submanifold is realized in $`RP^j`$, $`i<j`$, by the inclusion
$`(x_1,\mathrm{},x_{i+1})(x_1,\mathrm{},x_{i+1},0,\mathrm{},0).`$ (23)
To find an $`RP^2`$ inside $`RP^5`$, we must therefore choose three coordinates $`\{\stackrel{~}{X_1},\stackrel{~}{X_2},\stackrel{~}{X_3}\}`$ among the six embedding coordinates $`\{X_i\}`$, defined by Eq.(12). The remaining $`X_i`$’s should be set to zero.
For the $`S^3`$ of $`AdS_5`$ to separate two well-defined regions on the boundary plane, it should not shrink to zero size there. The magnetic $`D3`$ branes will then be associated with the white region. In this region, the radius of the $`\stackrel{~}{S}^3`$ embedded in $`S^5`$ shrinks to zero, corresponding to $`\mathrm{sin}\theta =0`$. Hence, $`X_1=X_2=X_3=X_4=0`$. We must therefore include at least one of the LLM coordinates in the generically non-vanishing set $`\{\stackrel{~}{X}_i\}`$, since otherwise the $`X_i`$’s cannot square to unity. Actually, $`X_5`$ and $`X_6`$ are only proportional to the LLM coordinates $`x_1`$ and $`x_2`$, differing by a $`\rho `$-dependent factor. Recalling Eq.(4), we can write
$`x_1`$ $`=\sqrt{r_0}X_5\mathrm{cosh}\rho `$ (24)
$`x_2`$ $`=\sqrt{r_0}X_6\mathrm{cosh}\rho .`$ (25)
This confirms that the brane will always be positioned outside of the black “semidisc” of radius $`r_0R^2`$.
Suppose first that precisely one of the two LLM coordinates belong to the generically non-vanishing set. All the other embedding coordinates must then vanish. Consequently, the non-vanishing coordinate, which is either $`X_5`$ or $`X_6`$, must be assigned the value $`R`$. Due to (25), this means that the brane appears as a point at radial position $`r_0\mathrm{cosh}\rho `$.
The other type of brane appears if we choose both of the LLM coordinates to belong to the generically non-vanishing set. Again, all other embedding coordinates vanish. The embedding condition then implies that $`X_5^2+X_6^2=R^2=r_0`$. Taking (25) into account, this leads us to conclude that in this case, the brane appears as a “semicircle” of radius $`r_0\mathrm{cosh}\rho `$. More precisely, the $`RP^2`$ collapses to an $`RP^1`$ in the LLM plane.
In conclusion, we have found two different types of magnetic $`D3`$ branes, appearing on the LLM plane as illustrated in Fig.8. One type appears as points, restricted to lie on one of the axes. The positions of two such branes are indicated by black dots in the figure. The other type appears as $`RP^1`$ cycles, one of which is shown. Every such cycle divides the LLM plane into two parts, allowing these D-branes to separate regions of unequal discrete torsion $`(\theta _{NSNS},\theta _{RR})`$. In both cases, the $`RP^2`$ is completely collapsed onto the LLM plane. In addition, the branes extend along the $`S^3`$ directions of the $`AdS_5`$.
We have illustrated the ideas involved assuming an $`AdS_5\times RP^5`$ background, but our conclusions should be equally valid also for the other $`\frac{1}{2}`$-BPS geometries. In particular, the topology on the LLM plane, including the distribution of domain walls, is expected to reflect the topology of the full bulk geometry.
### 5.3 Topology Change
In the absence of orientifolds, topology change in the family of $`\frac{1}{2}`$-BPS geometries of LLM type takes place at a fermi fluid boundary satisfying a local equation like $`x_1x_2=0`$. This can be thought of as a limit of the smooth boundary $`x_1x_2=\mu `$. Suppose this bounds a fermi fluid profile where for $`\mu >0`$ there is a single black region extending to infinity in, say, the upper left and lower right quadrants. Then for $`\mu <0`$ we have instead two separated black regions, one extending to infinity on the upper left and the other on the lower right. Clearly the connected black region splits into two disconnected ones as we pass through $`\mu =0`$ in the parameter space. It has been argued in Ref. that this is the “irreducible” transition responsible for topology change, and that more general topology changes can be decomposed into a sequence of such transitions. The 10-dimensional geometry changes by the appearance or disappearance of $`S^5`$ factors.
In the orientifolded theory there is another basic process of topology change that creates or destroys cycles of order two. The basic process of this type is the conversion of $`RP^5`$ into $`S^5`$ and vice versa. What happens is that an $`AdS_5\times RP^5`$ configuration consisting of a semi-disc anchored on the $`x_1`$ axis can be deformed until it develops a narrow neck connecting it to its image, as shown in Fig.9. At some stage the neck pinches off (Fig.10), and the bubble is no longer anchored to the horizontal axis.
Therefore this process represents a transition between $`RP^5`$ and $`S^5`$. We see that the origin was contained inside the bubble in the initial configuration, but is no longer contained at the end of the process. As we saw earlier, the former type of configuration has an $`RP^3`$ in the geometry while the latter does not. Thus this type of topology change causes us to lose or gain cycles of order two. At the transition, the (singular) configuration looks like a filled quadrant along with its mirror image, as shown in Fig.11. We expect that there will be topology-changing processes of “order $`n`$” built out of this basic process, as in Ref..
## 6 M-theory Orientifolds
M-theory contains orientifold 5-planes, so one can try to apply the construction above to this case. These orientifolds have been investigated further in Refs.. In particular one can take $`N`$ M5-branes parallel to an orientifold 5-plane, and the near-horizon geometry is $`AdS_7\times RP^4`$. Unlike the D3-O3 system in type IIB string theory, here the sphere factor becomes non-orientable after quotienting. This is because the space transverse to the original branes (and the orientifold plane) was $`R^5`$, whose orientation is reversed upon reflection.
What is the effect of orientifolding on the $`x_1x_2`$ plane of the LLM geometry? We cannot be as explicit as in the type IIB case. There we had a precise map between this plane and the phase space of free fermions. The orientifold plane was realised on the phase space as a wall, blocking off the region of $`x<0`$ (or more precisely identifying the lower half plane with a rotated copy of the upper half plane). Therefore this was also the case in the $`x_1x_2`$ plane and one could easily characterise functions $`z`$ which gave rise to the general orientifolded $`\frac{1}{2}`$-BPS geometry – as we did in the previous sections.
In the M-theory case, the free fermions are related to $`\frac{1}{2}`$-BPS excitations of the $`(2,0)`$ theory on M5-branes and arise from the transverse coordinates of the brane which are realised as scalar fields of the $`(2,0)`$ theory<sup>9</sup><sup>9</sup>9For a discussion of geometries in a pp-wave background, see Ref... Since all the transverse directions are being orientifolded, we again expect the fermions to move in a semi-infinite harmonic oscillator. Though we have seen in section 2 that the map from the fermion phase space to the $`x_1x_2`$ plane is not a simple identification, it was noted there that the topology of the two should be the same. Therefore we expect that the orientifold plane in the M-theory case is realised as the $`x_1`$ axis in the $`x_1x_2`$ plane of the geometry<sup>10</sup><sup>10</sup>10More precisely, there should be a choice of conformal transformation which maps the orientifold to the $`x_1`$ axis. Thereafter we will only be allowed to make conformal transformations that preserve the real axis, namely those which lie in $`SL(2,R)SL(2,C)`$. just as for type IIB strings. In that case most of the considerations in this paper will go through and we can generate a precise collection of $`\frac{1}{2}`$-BPS bubbling orientifolds in M-theory. We leave a more detailed analysis of this system for future work.
## 7 Conclusions
We have found an infinite class of new $`\frac{1}{2}`$-BPS geometries in type IIB string theory and, somewhat less explicitly, in M-theory. These have the same local geometry as the so-called “bubbling geometries” of Ref., but have global identifications. The geometries themselves have no orientifold plane and therefore the underlying string theory has no open-string sector, but the identifications can nevertheless be thought of as due to orientifold planes. The geometries have cycles of order 2, some of which support discrete torsion, in addition to the usual homology cycles. These backgrounds have the same amount of supersymmetry and, upto discrete factors, the same $`SO(4)\times SO(4)`$ symmetry as the original bubbling geometries. We saw that they are associated to free fermions in a half-oscillator potential, which in turn arise as the eigenvalues of matrices in the Lie algebra of $`SO(2N),Sp(2N)`$ and $`SO(2N+1)`$.
It would be interesting to extend these results to the $`\frac{1}{4}`$-BPS geometries associated to $`D3D7O3O7`$ systems. The introduction of $`D7`$-branes leads to a varying axion-dilaton background and the associated bubbling geometries are fairly simple and have already been found in Ref.. The introduction of orientifold 7-planes makes the system more interesting as it can then be related to “F-theory” compactifications. It has been argued that orientifold 7-planes are dynamical and behave in some regions of moduli space like non-perturbative D-branes (which at strong coupling exhibit remarkable effects including, for example, exceptional global symmetries). In this context, aspects of the AdS/CFT correspondence have been discussed in Refs. and it should be possible to find more general solutions of this kind using the bubbling prescription.
Finally, as we mentioned at the beginning, orientifolded bubbling geometries could be realised as local geometries of supersymmetric flux compactifications. Whenever the fermi fluid is localised in a bounded domain, the geometry is asymptotically $`AdS_5\times RP^5`$ (or $`AdS_5\times S^5`$), and therefore can be matched on to flat spacetime. But there might be a prescription as powerful as bubbling (with its connection to free fermions) that describes genuine compact geometries with orientifolds. This might make it much easier to classify supersymmetric flux vacua.
## Acknowledgements
We are grateful to Oleg Lunin, Martin Olsson, Kyriakos Papadodimas and Sandip Trivedi for helpful discussions.
## Appendix A Free fermions from matrix quantum mechanics
In this appendix, we show that free fermions emerge from all the traditional gauge groups. While some of this material is known or implicit in the matrix model literature, it is useful to compile all the needed results here.
The key observation is that
$`H`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{m}{}}}{\displaystyle \frac{1}{\sqrt{g}}}{\displaystyle \frac{}{\lambda _i}}\sqrt{g}{\displaystyle \frac{}{\lambda _i}}={\displaystyle \underset{i=1}{\overset{m}{}}}{\displaystyle \frac{1}{\mathrm{\Delta }^2}}{\displaystyle \frac{}{\lambda _i}}\mathrm{\Delta }^2{\displaystyle \frac{}{\lambda _i}}`$ (26)
$`=`$ $`{\displaystyle \frac{1}{\mathrm{\Delta }}}{\displaystyle \underset{i=1}{\overset{m}{}}}{\displaystyle \frac{^2}{\lambda _i^2}}\mathrm{\Delta },`$ (27)
where the measure factor $`\sqrt{g}=\mathrm{\Delta }^2`$ is given by
$`(\mathrm{\Delta }_{})^2`$ $`SU(m)`$ (28)
$`(\mathrm{\Delta }_{})^2(\mathrm{\Delta }_+)^2(\mathrm{\Delta }_0)^2`$ $`Sp(N),N=2m`$ (29)
$`(\mathrm{\Delta }_{})^2(\mathrm{\Delta }_+)^2`$ $`SO(N),N=2m`$ (30)
$`(\mathrm{\Delta }_{})^2(\mathrm{\Delta }_+)^2(\mathrm{\Delta }_0)^2`$ $`SO(N),N=2m+1,`$ (31)
where we have defined
$`\mathrm{\Delta }_{}`$ $`{\displaystyle \underset{i<j}{\overset{m}{}}}`$ $`(\lambda _i\lambda _j)`$ (32)
$`\mathrm{\Delta }_+`$ $`{\displaystyle \underset{i<j}{\overset{m}{}}}`$ $`(\lambda _i+\lambda _j)`$ (33)
$`\mathrm{\Delta }_0`$ $`{\displaystyle \underset{i}{\overset{m}{}}}`$ $`\lambda _i,`$ (34)
expressed in terms of the eigenvalues $`\{\lambda _i\}_{i=1,2,\mathrm{},m}`$. The last equality in (26) follows from the identity
$`{\displaystyle \underset{i=1}{\overset{m}{}}}{\displaystyle \frac{^2}{\lambda _i^2}}\mathrm{\Delta }=0.`$ (35)
The relation (26) means that a factor of $`\mathrm{\Delta }`$ can be absorbed by the wave function, giving rise to a system of free fermions.
It remains to establish (28). Let us begin by considering the unitary group, consisting of matrices $`U`$ satisfying $`UU^{}=𝕀`$. The Lie algebra $`su(2)`$ consists of anti-hermitian matrices $`A^{}=A`$. An infinitesimal variation $`\delta U=U^1dU`$ of a unitary matrix $`U`$ is an element of the Lie algebra, $`\delta U^{}=\delta U`$.
The hermitian matrix A can be diagonalized using unitary matrices,
$`U^1AU=\mathrm{\Lambda }=\text{diag}(\lambda _1,\lambda _2,\mathrm{},\lambda _m)`$ (36)
The significance is that we are trading degrees of freedom in the anti-hermitian matrix for those of the eigenvalues themselves and of the diagonalizing unitary transformation. Differentiating equation (36) and solving for $`dA`$ gives
$$\text{Tr }(dAdA^{})=\text{Tr }\left(2\delta U\mathrm{\Lambda }[\delta U,\mathrm{\Lambda }]+\left(d\mathrm{\Lambda }\right)^2\right),$$
(37)
which is manifestly invariant under unitary similarity transformations. Hence it defines a metric on the space of antihermitian matrices.
Denoting matrix elements of a matrix $`M`$ by $`M_{ij}`$, this becomes
$$\text{Tr }(dAdA^{})=2\underset{i<j}{\overset{m}{}}\delta u_{ij}\delta u_{ij}^{}(\lambda _i\lambda _j)^2+\underset{i=1}{\overset{m}{}}d\lambda _i^2,$$
(38)
We can read off the measure
$`\sqrt{g}=(\mathrm{\Delta }_{})^2,`$ (39)
(up to some numerical constant) in accordance with (28).
Next, we turn our attention to the symplectic gauge group $`\text{Sp}(2m)`$, consisting of matrices $`S`$ satisfying $`S^TJS=J`$, where
$`J`$ $`=`$ $`\left(\begin{array}{cc}0& 𝕀\\ 𝕀& 0\end{array}\right),`$ (40)
where $`𝕀`$ is the $`m\times m`$ identity matrix. The Lie algebra sp(2m) then consists of matrices $`A`$ such that $`A^T=JAJ`$. Infinitesimal variations $`\delta S=S^1dS`$ belong to the Lie algebra $`\text{sp}(2m)`$.
Again, a Lie algebra element can be diagonalized using symplectic matrices, $`S^1AS=\mathrm{\Lambda }`$. Differentiating $`S^1AS=\mathrm{\Lambda }`$ and solving for $`dA`$ gives
$$\text{Tr }(dAJdA^TJ)=\text{Tr }\left(2\delta S\mathrm{\Lambda }[\delta S,\mathrm{\Lambda }]+\left(d\mathrm{\Lambda }\right)^2\right).$$
(41)
Defining $`(e_{PQ})_{LM}=\delta _{PL}\delta _{QM}`$, a basis for $`\text{sp}(2m)`$ is
$`e_{j,k}^1`$ $`e_{j,k}e_{k+m,j+m}`$ (42)
$`e_{j,k}^2`$ $`e_{j,k+m}+e_{k,j+m}`$ (43)
$`e_{j,k}^3`$ $`e_{j,k}e_{k+m,j+m}.`$ (44)
Note that since $`0=det(A\lambda 𝕀)=det(A+\lambda 𝕀)`$, eigenvalues come in pairs, $`\{(\lambda _i,\lambda _i)\}_{i=1,2,\mathrm{},m}`$. Consequently, we can write $`\mathrm{\Lambda }`$ on the form
$`\mathrm{\Lambda }={\displaystyle \underset{i=1}{\overset{k}{}}}\mathrm{\Lambda }_ie_{ii}^1=\text{diag}(\lambda _1,\lambda _2,\mathrm{},\lambda _m,\lambda _1,\lambda _2,\mathrm{},\lambda _m).`$ (45)
Similarly, we write $`\delta S`$ as
$`\delta S`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{m}{}}}\delta S_i^de_{ii}^1`$ (46)
$`+`$ $`{\displaystyle \underset{j<k}{\overset{m}{}}}\left[\delta S_{jk}^i(e_{jk}^1+e_{kj}^1)+\delta S_{jk}^{ii}(e_{jk}^1e_{kj}^1)\right]`$ (47)
$`+`$ $`{\displaystyle \underset{jk}{\overset{m}{}}}\left[\delta S_{jk}^a(e_{jk^2+e_{jk}^3})+\delta S_{jk}^b(e_{jk}^2e_{jk}^3)\right].`$ (48)
The reason for this particular expansion is that it will make the metric diagonal. Indeed, differentiating $`S^1AS=\mathrm{\Lambda }`$ and solving for $`dA`$ gives
$`\text{Tr }(dA^TJdAJ)`$ $`=2`$ $`{\displaystyle \underset{i=1}{\overset{m}{}}}d\lambda _i^2`$ (49)
$`+8`$ $`{\displaystyle \underset{j<k}{}}(\lambda _j\lambda _k)^2\left[(\delta S_{jk}^{ii})^2(\delta S_{jk}^i)^2\right]`$ (50)
$`+8`$ $`{\displaystyle \underset{jk}{}}+(\lambda _j+\lambda _k)^2\left[(\delta S_{jk}^b)^2(\delta S_{jk}^a)^2\right].`$ (51)
Hence, the measure is
$`\sqrt{g}=(\mathrm{\Delta }_{})^2(\mathrm{\Delta }_+)^2(\mathrm{\Delta }_0)^2,`$ (52)
as anticipated in (28).
The orthogonal gauge group remains to be analysed. Orthogonal matrices $`O`$ satisfy $`OO^T=𝕀`$, and corresponding Lie algebra elements are antisymmetric, $`AA^T=A`$. The infinitesimal variation $`\delta O=O^1dO`$ of an orthogonal matrix $`O`$ is antisymmetric, $`\delta O^T=\delta O`$.
As in the symplectic case, the eigenvalues come in pairs, differing only by a sign. This means that $`SO(2m)`$ and $`SO(2m+1)`$ both have $`m`$ independent eigenvalues. Consider first the case $`SO(2m)`$. Using orthogonal similarity transformations, antisymmetric matrices can only be brought to the block-diagonal form
$`\mathrm{\Lambda }`$ $`=`$ $`\left(\begin{array}{cccc}\lambda _1\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right)& 0& \mathrm{}& 0\\ 0& \lambda _2\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right)& \mathrm{}& 0\\ \mathrm{}& \mathrm{}& & \mathrm{}\\ 0& 0& \mathrm{}& \lambda _N\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right)\end{array}\right)`$ (53)
$`=`$ $`i\sigma _2\mathrm{diag}(\lambda _1,\lambda _2,\mathrm{},\lambda _N).`$ (54)
Suppose that $`O`$ is the orthogonal matrix which block-diagonalizes the antisymmetric matrix $`A`$, i.e.
$$O^1AO=\mathrm{\Lambda }.$$
(55)
Differentiating equation (55) and solving for $`dA`$ gives
$$\text{Tr }(dAdA^T)=\text{Tr }\left(2\delta O\mathrm{\Lambda }[\delta O,\mathrm{\Lambda }]+\left(d\mathrm{\Lambda }\right)^2\right).$$
(56)
Defining
$`\sigma ^0`$ $`𝕀_{2\times 2}=\left(\begin{array}{cc}+1& 0\\ 0& +1\end{array}\right)\sigma ^2`$ $`i\sigma _2=\left(\begin{array}{cc}0& +1\\ 1& 0\end{array}\right)`$ (57)
$`\sigma ^1`$ $`\sigma _1=\left(\begin{array}{cc}0& +1\\ +1& 0\end{array}\right)\sigma ^3`$ $`\sigma _3=\left(\begin{array}{cc}+1& 0\\ 0& 1\end{array}\right)`$ (58)
as a basis of real $`2\times 2`$ matrices in terms of the identity matrix $`𝕀_{2\times 2}`$ and the Pauli matrices $`\{\sigma _i\}`$, we can write the variation $`\delta O`$ as
$`\delta O`$ $`=`$ $`\left(\begin{array}{cccc}\delta O_{11}\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right)& \delta O_{12}^A\sigma ^A& \mathrm{}& \delta O_{1N}^A\sigma ^A\mathrm{}\\ \delta O_{12}^A\sigma ^A& \delta O_{22}\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right)& \mathrm{}& \delta O_{2N}^A\sigma ^A\\ \mathrm{}& \mathrm{}& & \mathrm{}\\ \delta O_{1N}^A\sigma ^A& \delta O_{2N}^A\sigma ^A& \mathrm{}& \delta O_{NN}\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right)\end{array}\right)`$ (59)
$`=`$ $`\sigma ^A\delta O^A+i\sigma _2\mathrm{diag}(\delta O_{11},\delta O_{22},\mathrm{},\delta O_{NN}).`$ (60)
The antisymmetric matrix $`\delta O^A`$ is defined such that the entry at row $`i`$ and column $`j`$ is the coefficient $`\delta O_{ij}^A`$. We are using the convention that repeated indices $`A`$ are summed over.
Using the forms (53) and (59), the metric (56) becomes
$`\text{Tr }(dAdA^T)=`$ $``$ $`4{\displaystyle \underset{i<j}{}}{\displaystyle \underset{A=1,3}{}}\left(\delta O_{ij}^A\right)^2(\lambda _i+\lambda _j)^2`$ (61)
$`+`$ $`4{\displaystyle \underset{i<j}{}}{\displaystyle \underset{A=0,2}{}}(1)^{A/2+1}\left(\delta O_{ij}^A\right)^2(\lambda _i\lambda _j)^2+{\displaystyle \underset{i=1}{\overset{m}{}}}d\lambda _i^2`$ (62)
Consequently, the measure is
$`\sqrt{g}=(\mathrm{\Delta }_{})^2(\mathrm{\Delta }_+)^2,`$ (63)
as written in (28).
Consider now the gauge group $`SO(2m+1)`$. The analysis goes through in much the same way as for the $`SO(2m)`$ case, replacing
$`\delta O\delta O+{\displaystyle \underset{i=1}{\overset{m}{}}}{\displaystyle \underset{B=0,1}{}}\delta O_i^B\tau ^Be_{i,2k+1},`$ (64)
where $`\tau ^0\left(\begin{array}{c}1\\ 0\end{array}\right)`$, $`\tau ^1\left(\begin{array}{c}0\\ 1\end{array}\right)`$. Similarly, $`\mathrm{\Lambda }`$ is just extended with another row and column of zeros, since there are still only $`m`$ eigenvalues. With these modifications, the metric (61) gets an additional contribution,
$`\text{Tr }(dAdA^T)\text{Tr }(dAdA^T)+{\displaystyle \underset{i,B}{}}\lambda _i^2(\delta O_i^B)^2.`$ (65)
The measure (63) changes accordingly,
$`\sqrt{g}\sqrt{g}{\displaystyle \underset{i=1}{\overset{m}{}}}\lambda _i^2=(\mathrm{\Delta }_{})^2(\mathrm{\Delta }_+)^2(\mathrm{\Delta }_0)^2,`$ (66)
consistent with (28).
Note that the symplectic and orthogonal gauge groups share the following two features. The independent set of eigenvalues can always be chosen to be non-negative. In addition, absorbing a factor of $`\mathrm{\Delta }`$ into the wave function makes the eigenvalues $`\lambda _i`$ repel not only each other, but also their images $`\lambda _i`$. This is consistent with the orientifold interpretation, as explained in section 4.
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# Exploring Terrestrial Planet Formation in the TW Hydrae Association
## 1. Introduction
While calibrating the Infrared Astronomical Satellite (IRAS) Aumann et al. (1984) and Gillett (1986) discovered large far-infrared excesses around four nearby stars: Vega, Fomalhaut, $`\beta `$ Pictoris, and $`ϵ`$ Eridani. Release of the IRAS Point Source Catalog led investigators to many other young, nearby dwarf stars displaying excess emission at mid- and far-IR wavelengths (e.g., Walker & Wolstencroft, 1988; de la Reza et al., 1989; Gregorio-Hetem et al., 1992; Mannings & Barlow, 1998; Zuckerman & Becklin, 1993). As a result, numerous investigations of these objects from the ground have added critical insights concerning circumstellar material (e.g., Jayawardhana et al., 1998; Koerner et al., 1998; Kenyon & Bromley, 2002, 2004). The Hubble Space Telescope (HST) and the Infrared Space Observatory (ISO) (Habing et al., 2001; Spangler et al., 2001; Meyer & Beckwith, 2000, and references therein) have also given us key insights into specific objects, including images taken in the near-IR showing scattered light from thin disks of dust that are shown to absorb starlight and re-emit at much longer wavelengths (e.g., Schneider et al., 1999). Thus, in the last 20 years a body of observations, measurements, analysis and theory have accumulated that leave no doubt about the reality of terrestrial planetary material around young stars. Now, investigations of planetary debris systems (PDSs) are underway with the new tools provided by the Spitzer Space Telescope (Werner et al., 2004).
We have initiated a program using Spitzer to explore a unique and varied sample of late-type stars that are nearby, $``$40–100 pc, and show one or more indications of young age. Our total sample of over 120 stars has been prioritized into three groups: (1) the first 19 stellar systems identified by common space motion as members of the TW Hydrae association (TWA; Kastner et al., 1997; Webb et al., 1999), (2) $``$60 stars that show significant lithium absorption indicating young ages, noting that all of the TWA systems easily meet this requirement and therefore overlap this sample, (3) a group of $`60`$ stars with less certain ages based on either detection of X-rays or H$`\alpha `$ in emission.
A combination of special circumstances dictate that we start our exploration with the TW Hydrae association. First, the estimated age of the association ranges from 8–10 million years (Myr) and the ages of these stars are essentially equal (Stauffer, Hartmann, & Barrado y Navascues, 1995). Second, the TWA members are among the closest young stars available. Third, most association members lie in excellent, low-background regions of the sky for far-IR observations. IRAS has already shown that four of the brightest IR excesses known are found in the TWA. Additional important factors are that five of the TWA stars have Hipparcos distances averaging around 50 pc, their spectral types range from A0V to M5V with most objects being of M or late K spectral type, and almost all members show signs of T Tauri activity. TW Hydrae itself has been classified as a classic T Tauri system despite its apparent isolation (Herbig, G. H., 1978; Rucinski & Krautter, 1983). In contrast, and within the same association, the brightest and most dramatic PDS, HD 98800, shows no evidence for current accretion of primordial material onto the star and has served as the most direct proof of ongoing terrestrial planet formation (e.g., Low, Hines, & Schneider, 1999).
## 2. Observations
Observations of 24 members of the TW Hydrae association were made using Spitzer. These include 14 single stars and binaries that are unresolved at 24 $`\mu `$m and five binary systems at least partially resolved by Spitzer at 24 $`\mu `$m. The Multi-band Imaging Photometer for Spitzer (MIPS; Rieke et al., 2004) provided measurements or upper limits of the flux densities ($`F_\nu `$) at 24 $`\mu `$m and 70 $`\mu `$m for the entire sample. MIPS 160 $`\mu `$m photometry was also obtained for the four IR-excess systems discovered by IRAS (i.e., TW Hya, Hen 3-600, HD 98800B, and HR 4796A). The observations were made between 2004 January 30 and June 21, with each target observed in all of the intended wavelength bands during a single visit by Spitzer. Table 1 lists the stars in the association observed with MIPS along with basic data from previous measurements. Six other systems have been identified as members of the TWA (see e.g., Song et al., 2003; Webb, 2000a) since the final planning of the Spitzer/MIPS observations, but were not included in the target list because of time limitations.
The MIPS data are reduced using the data analysis tool (DAT) of the MIPS instrument team (Gordon et al., 2005). The images are corrected for distortion, registered, and mosaicked to compensate for the movements of the spacecraft and cryogenic scan mirror mechanism (CSMM) within MIPS that are employed to produce a series of spatially dithered images on the three detector arrays. Aperture photometry is then performed on the calibrated, mosaicked images. Table 2 summarizes the MIPS bandpasses, aperture photometry parameters, and the calibration factors adopted for the data reduction. The data have been color corrected using the relevant tables in the Spitzer Observer’s Manual.
The 24 $`\mu `$m data were collected using the standard MIPS point source observing template that obtains 14 images of the target at various standard locations on the 128$`\times `$128-pixel Si:As detector array during each cycle of the template. Either 3 s or 10 s integration times and 1–4 complete template cycles were used to ensure high signal-to-noise ratio (S/N) measurements of the targets. Total observation times were estimated for each target so that a S/N$`>`$10 measurement was achieved at the expected 24 $`\mu `$m flux level of the star’s photosphere, and all stars were detected to at least this S/N. Indeed, the median S/N is $``$160, though the internal photometric calibration is estimated to be 1–2% in this bandpass (e.g., Rieke et al., 2005).
For the aperture photometry of the 24 $`\mu `$m data, an aperture 15″ in radius ($`r`$) was typically used along with a sky annulus centered on the aperture and having an inner radius of $``$63″ and width of 25″. The large sky annulus was chosen to avoid as much flux as reasonable from the Spitzer plus MIPS point spread function (PSF) of the target. A larger aperture ($``$37$`\stackrel{}{\mathrm{.}}`$5 radius) was used for four binary TWA systems that are at least partially resolved at 24 $`\mu `$m (TWA 8, 9, 13, and 15). The relative brightnesses of individual components of the binaries were estimated by comparing the fluxes measured in $`r=2`$$`\stackrel{}{\mathrm{.}}`$5 apertures centered on the components. For TWA 9 and 13, a third source of 24 $`\mu `$m flux fell within the large photometric apertures. The flux from these sources were subtracted and are not included in the values for $`F_\nu `$ at 24 $`\mu `$m.
Only HD 98800 was bright enough that the core of the PSF saturated even in the minimum observation time allowed at 24 $`\mu `$m. In this case, the flux in the saturated core was estimated by measuring the flux in the wings of the PSF. This was done for the entire sample as well. The exercise on the non-saturated stars showed that the 24 $`\mu `$m flux in the saturated portion of the PSF for HD 98800 could be recovered to within 5%. After correcting for the saturated core, the 24 $`\mu `$m flux density was measured for HD 98800 in the usual manner.
Observations, data reduction, and aperture photometry for the 70 $`\mu `$m data are similar to those at 24 $`\mu `$m. However, the pixel scale of the 70 $`\mu `$m wide-field optical train ($``$10″ pixel<sup>-1</sup>) dictates that none of the binaries in the TWA can be resolved except for TWA 19. Therefore, we have used a uniform photometry aperture (3-pixel radius) and sky annulus (radius = 4–6 pixels) for the stars observed in the 70 $`\mu `$m wide-field mode. The chosen photometry aperture is very large compared to the $`1`$″ pointing accuracy of the telescope, eliminating the ambiguity in both the target position on the array and the determination of flux density upper limits for undetected sources.
TW Hya, HD 98800, and HR 4796A were observed in the 70 $`\mu `$m narrow-field mode to prevent saturation by these bright IR sources as this optical train halves the pixel scale. For these observations, an $`r=5`$$`\stackrel{}{\mathrm{.}}`$67 aperture was chosen. Instead of an annulus centered on the target to estimate the sky contribution within the photometric aperture, rectangular regions on either side of the target were utilized because of the narrow, 16-pixel width of the usable portion of the detector array.
Because of the degraded sensitivity and higher instrumental noise of the 70 $`\mu `$m Ge:Ga detector array relative to pre-launch estimates, it was not practical to obtain measurements accurate enough to detect objects at the expected levels of the TWA photospheres. As a result, the 70 $`\mu `$m observations were tailored to detect excesses at flux levels $``$10$`\times `$ the expected 70 $`\mu `$m photospheric flux density with a reasonably high degree of confidence ($`>`$3–4$`\sigma `$). Unfortunately, most of the TWA observations occurred before the bias level for the 70 $`\mu `$m detector was reduced. During the period before the bias change, the overall noise in the detector was greater and there were more pronounced changes seen in background levels across array modules. Because of these instrumental effects, a conservative criterion of 5$`\sigma `$ above background was adopted for a detection at 70 $`\mu `$m.
All four targets observed at 160 $`\mu `$m were detected at a S/N ranging from $``$6–160 using the standard MIPS photometry observing template. In all cases, an aperture of $`r=64`$″ was used for the photometry. The sky contribution within the photometric aperture was estimated using rectangular regions on either side of the target. As with the 24 $`\mu `$m and 70 $`\mu `$m data, the 160 $`\mu `$m photometry is corrected for the loss in flux caused by the use of apertures of finite size not encompassing the total flux from the PSF. The multiplicative aperture corrections to a hypothetical aperture of infinite extent are listed in Table 2.
The 160 $`\mu `$m data were also corrected for the blue filter leak discovered in this bandpass after the launch of Spitzer. This leak allows light at $``$$`\mu `$m in wavelength to reach the detector which is sensitive to near-IR radiation. Tests made during flight have shown that the contaminating flux is about 15$`\times `$ the flux received from a star’s photosphere at 160 $`\mu `$m assuming a Rayleigh-Jeans law between the near and far-IR. In all four cases in the TWA where 160 $`\mu `$m measurements were made, the expected photospheric 160 $`\mu `$m flux density is $``$1% of the measured signal. Given the bright far-IR excesses of these stars, the filter leak corrections are small.
In addition to broad-band photometry, observations of TW Hya and HD 98800 were made using MIPS in its SED mode. These observations were made on JD = 2453395, roughly a year after the MIPS photometry of TW Hya was obtained and about seven months after the photometry of HD 98800. In this observing mode, the CSMM brings light from the telescope to a reflective diffraction grating that also performs the functions of a slit and collimator. The collimated, dispersed light is then focused onto the 70 $`\mu `$m detector array. The result is a very low-resolution ($`R=\lambda /\mathrm{\Delta }\lambda 15`$-25) first-order spectrum covering a wavelength range of 52–97 $`\mu `$m. The slit/grating has a width of two pixels ($`20`$″) and the dispersion is 1.7 $`\mu `$m pixel<sup>-1</sup>. An unusable portion of the 70 $`\mu `$m array restricts the slit length to $`2`$$`\stackrel{}{\mathrm{.}}`$7. Also, a dead 4$`\times `$8-pixel readout at one end of the slit further reduces the effective length of the slit where the full spectral range can be sampled.
Observations in SED mode involve chopping the CSMM between the target and a region 1–3′ away to sample the background, and using small spacecraft moves to dither the target between two positions along the slit. Data reduction is again handled by the DAT in essentially the same manner as for 70 $`\mu `$m broad-band imaging. The result is co-added “on-” and “off-target” images. After the background image has been subtracted from the target, the spectrum is extracted using an extraction aperture of set width (8 pixels). Flux calibration of the TW Hya and HD 98800 spectra was accomplished using a MIPS SED-mode observation of $`\alpha `$ Boo made on JD = 2453196. The 50–100 $`\mu `$m continuum of $`\alpha `$ Boo is assumed to be $`\lambda ^2`$ and $`F_\nu =14.7`$ Jy at the effective wavelength of the 70 $`\mu `$m filter bandpass. The observations and spectral extractions of $`\alpha `$ Boo and the two TWA systems were identical and the same DAT calibration files were used for every target.
## 3. Results
Table 3 summarizes the MIPS photometric results for the TWA. Several stars in the association are known optical variables and we give the epoch of the start of the observations in the second column of the table. The 24 $`\mu `$m, 70 $`\mu `$m, and 160 $`\mu `$m flux densities are given in units of milli-Janskys (mJy; where 1 mJy = $`10^{26}`$ erg cm<sup>-2</sup> s<sup>-1</sup> Hz<sup>-1</sup>). Uncertainties in the flux densities are the square root of the quadratic sum of the measurement and calibration uncertainties, noting that most relative uncertainties are much lower. Also listed in Table 3 is the S/N of the detections in each MIPS band. Although some of the S/N values are extremely high, the uncertainties quoted for the photometry in all three bandpasses are dominated by uncertainties in the absolute flux calibration of the instrument. These are estimated to be $``$10% for the 24 $`\mu `$m band and $``$20% for the 70 $`\mu `$m and 160 $`\mu `$m bandpasses. Upper limits assigned in Table 3 to non-detections at 70 $`\mu `$m are at the level of 3$`\sigma `$ over the noise of the background.
Our primary goal for this observational program is to detect and measure IR excesses in the TWA. This requires that good estimates of $`F_\nu `$ be made for the TWA stellar photospheres at the effective wavelengths of the MIPS bandpasses. For all of stars in the sample, stellar photospheres were derived using a grid of Kurucz (Kurucz, 1979) and “NextGen” (Hauschildt, Allard, & Baron, 1999) photospheric models and fitting available optical photometry combined with near-IR flux densities compiled from the Two-Micron All Sky Survey (2MASS) Point Source Catalog. Temperature is a free parameter in the fitting, but in all cases solar metallicity and $`\mathrm{log}\mathrm{g}=4.5`$ are assumed. The NextGen models are known to yield better fits to late-type main sequence stars than Kurucz models, and were used for all TWA members except for HR 4796A and TWA 19A, which have much higher effective temperatures than the rest of the sample. The best-fitting model temperatures ($`T_{}`$) of the TWA are listed in Table 1.
### 3.1. 24 $`\mu `$m Excesses
Comparison of the 24 $`\mu `$m data for the TWA with the 2MASS $`K_s`$-band flux densities of these stars shows that the association can essentially be divided into two populations. Figure 1 plots the 24-to-2.17 $`\mu `$m flux density ratio against the flux density in the $`K_s`$-band. In all cases, the 2MASS photometry is consistent with the stellar photospheres emitting the near-IR light. The four stars with IR excesses discovered by IRAS are about a factor of 100 brighter at 24 $`\mu `$m relative to $`K_s`$ than the other stars in the sample. In fact, for the stars with flux ratios $`<`$0.02, the data are consistent with the 24 $`\mu `$m flux being solely from the photosphere. As a guide, a dashed line in Figure 1 marks the limit of $`F`$(24 $`\mu `$m)/$`F_{K_s}`$ if both bandpasses fall within the Rayleigh-Jeans regime of the stellar continuum. Since the TWA members are typically spectral type K and M, the $`K_s`$ band falls close to the emission peak of the photosphere. As a result, the stars typically lie above the Rayleigh-Jeans limit in Figure 1.
Figure 2 further explores the empirical relation between the 24 $`\mu `$m flux density with the near-IR. The stars with previously known IR excesses or $`T_{}>4500`$ K are omitted from Figure 2, and the flux ratio shown in Figure 1 is now plotted against the temperature of the model stellar photosphere. The majority of the 19 stars plotted form a locus that is offset from the $`F`$(24 $`\mu `$m)/$`F_{K_s}`$ ratios expected from blackbodies. Spectra derived from the NextGen models generally give a good approximation to the observed flux ratios, although the TWA members have not reached the main sequence. The unweighted average flux ratio $`F(24\mu \mathrm{m})/F_{K_s}=0.015`$ with an RMS of only 0.004. The concentration of the points close to the expected values from the NextGen models in Figure 2 suggests that we are seeing no IR excess at 24 $`\mu `$m for most of these late-type, young stars. The MIPS photometry is consistent with the 24 $`\mu `$m flux emanating from the photosphere for most of the members of the TWA. One example is TWA 6, where the MIPS 24 $`\mu `$m measurement supports the contention by Uchida et al. (2004) that the 5–20 $`\mu `$m spectrum of this star taken with Spitzer and the Infrared Spectrograph (IRS; Houck et al., 2004) is indeed the spectrum of the photosphere.
Although the NextGen models yield a much better brightness estimate of the photosphere at 24 $`\mu `$m than blackbodies in this temperature range, Figure 2 shows that the flux ratio is possibly more dependent on the model temperature than is actually observed. Note, TWA 12, 13B, and 17 fall well below the NextGen curve. These deviations and the intrinsic scatter of the data are intriguing, but the limited number of objects in the sample complicate a detailed investigation of their possible causes, especially considering that only temperature was varied in constructing the models and that the TWA stars are younger than the nominal zero-age main sequence stars described by the NextGen models. In the case of TWA 17, the 24 $`\mu `$m flux density is measured to be only $`1.5\pm 0.2`$ mJy and the source is detected at a $`S/N=13`$ with MIPS. Therefore, we do not confirm the possible excess at 12 and 18 $`\mu `$m reported from ground-based observations by Weinberger et al. (2004).
The measurement of $`F`$(24 $`\mu `$m)/$`F_{K_s}`$ for TWA 7 suggests that it possesses a 24 $`\mu `$m excess that contributes $``$40% to the total brightness of this system at this wavelength. Except for the four IRAS detections, TWA 7 has a higher flux ratio in Figures 1 and 2 than all of the other TWA members. The next highest ratio belongs to TWA 9B, but its deviation from the expected flux of the photosphere is less than half of that shown by TWA 7, and is no larger than the flux deficits at 24 $`\mu `$m measured for TWA 12, 13B, and 17. Since the 2MASS and MIPS measurements are not simultaneous, an alternate explanation for the apparent IR excess in TWA 7 could be the variability of the star which is not uncommon for T Tauri stars, although there are no published data suggesting that this star is variable. Also, the clustering of most measurements around $`F`$(24 $`\mu `$m)/$`F_{K_s}`$ is a clear demonstration that most stars in the TWA do not, in fact, exhibit significant variability.
### 3.2. 70 $`\mu `$m Excesses
Regardless of the strength of a possible 24 $`\mu `$m excess in the observed continuum of TWA 7, this object is detected at 70 $`\mu `$m at a flux level $``$40$`\times `$ brighter than the photospheric emission estimated for the star at this wavelength. Figure 3 presents the ratio between the observed and predicted 70 $`\mu `$m flux densities. Only six of the 20 systems (TWA 19 is resolvable at 70 $`\mu `$m) are detected with MIPS at 70 $`\mu `$m. TWA 13 joins TWA 7 and the four IRAS sources in having a detected 70 $`\mu `$m excess. TWA 13 cannot be resolved with Spitzer at 70 $`\mu `$m, so it is unknown which of the two components produces the excess, or if both components of the binary contribute to the observed flux.
Figure 3 shows the same basic division in the TWA sample as is found at 24 $`\mu `$m. The 70 $`\mu `$m excesses for TWA 7 and 13 are $``$10$`\times `$ fainter than the other detections. One-sigma upper limits are also shown in Figure 3 for the remainder of the sample. As discussed in §2, the sensitivity limits of the 70 $`\mu `$m array render detection of the photospheres of the TWA sample impractical. Excesses similar to those of TWA 7 and 13 cannot be ruled out for many members of the association, and it is clear from the data that it is not necessary for an object to show an excess at 24 $`\mu `$m to find one at 70 $`\mu `$m. For instance, there is no hint from the 24 $`\mu `$m measurement of TWA 13 of an excess at longer wavelengths. There are three stars where a 70 $`\mu `$m excess at $`>`$10$`\times `$ the flux density of the photosphere can be ruled out at the 3$`\sigma `$ confidence level: TWA 2, 5A, and 8A. Therefore, the minimum range in strength for 70 $`\mu `$m excesses in the TWA is greater than a factor of 300.
### 3.3. Spectral Energy Distributions
The SEDs of the six systems in the TWA with confirmed IR excesses are shown in Figure 4. It is readily apparent that this small sample of young stars exhibits a rich variety of IR properties. The T Tauri stars TW Hya and Hen 3-600 show no hint of their SEDs turning over even out to 160 $`\mu `$m. In fact, the MIPS 160 $`\mu `$m measurement for TW Hya reveals that the spectrum of the presumed optically thick accretion disk continues to rise into the infrared. Similar to TW Hya, the MIPS data for Hen 3-600 are consistent with the earlier IRAS results and show that the 70 and 160 $`\mu `$m flux densities are roughly equivalent.
In contrast, HD 98800A and B do not show the extreme activity associated with T Tauri stars. However, HD 98800B generates a huge IR excess that is well fit by a single blackbody with $`T=160`$-170 K (Low, Hines, & Schneider, 1999; Hines et al., 2004), whereas component A lacks an easily measured excess even though the total luminosities of the two companions are virtually identical. The MIPS photometry is consistent with earlier results. Likewise, the IR excess of HR 4796A produces an IR SED that can be fit by a single blackbody, but it is significantly cooler than found for HD 98800B ($`T=110`$ K; Jura et al., 1993). With the addition of the Spitzer observations, we fit this excess with a $`T=108`$ K blackbody. Given that HR 4796A is an A0V star, it is not surprising that it has an IR excess because of the finding that younger A stars have a higher probability of possessing a 24 $`\mu `$m excess and having stronger thermal excess emission than older A stars (Rieke et al., 2005). It is of note that HR 4796A has by far (by $``$10$`\times `$) the largest ratio of 24 $`\mu `$m excess emission-to-photospheric emission of the 265 A and B stars examined by Rieke et al. (2005).
Figure 5 presents the MIPS SED-mode observations of TW Hya and HD 98800B. These low-resolution spectra do not show any spectral features of high equivalent width in either object. Arcturus was used to calibrate the spectra and they are consistent with the 70 $`\mu `$m flux densities derived from the photometry. The IR spectral slopes of TW Hya and HD 98800B are quite different, with the continuum of TW Hya being much redder than that of HD 98800B. A power-law having a spectral index, $`\alpha `$, where $`F_\nu \nu ^\alpha `$, is adequate to describe these data. For TW Hya, $`\alpha =0.33\pm 0.04`$, and is consistent with the finding that the 160 $`\mu `$m flux density is higher than the 70 $`\mu `$m measurement. The slope of a single power-law fit to the continuum of HD 98800B between 52 and 97 $`\mu `$m yields $`\alpha =1.12\pm 0.04`$. The blackbody fit with $`T_{BB}=160`$ K to the PDS of HD 98800B is also shown in Figure 5, and its slope within this spectral region is in good agreement with the SED-mode data.
### 3.4. PDS Luminosity and Mass
Our primary reason for calculating the stellar and IR excess luminosity is to measure the fractional luminosities of those stars that show direct evidence of PDS formation and to provide well defined upper limits for those stars that fail to show an excess. Table 4 lists the stellar luminosity, $`L_{}`$, derived by numerical integration, and the fractional luminosity, $`L_{IR}/L_{}`$, where $`L_{IR}`$ is the luminosity of the infrared excess that is the result of numerically integrating the flux density in excess of the predicted level of the photosphere. Direct measurements of the distance to TWA members are limited to only five systems (Table 1). We have assigned a distance of 55 pc, the distance corresponding to the unweighted average of the parallaxes measured by Hipparcos excluding TWA 19, to the other stars in the sample. In general, $`L_{}`$ is higher than expected for the TWA given their spectral types because these young, rapidly evolving stars have higher luminosities, lower temperatures, and larger radii than when they reach the main sequence.
Conservative upper limits for $`L_{IR}/L_{}`$ have been assigned to objects listed in Table 4 that show no evidence for a 24 $`\mu `$m excess by assuming that the 1$`\sigma `$ upper limits estimated for the 70 $`\mu `$m flux densities represent the brightness of a possible IR excess plus photosphere at this wavelength. It is also assumed that any possible IR excess is well approximated by a blackbody with temperature, $`T_{PDS}`$. The upper limit on the temperature of a possible cold PDS is given in Table 4 and is set so that $`T_{PDS}`$ and the assumed value for $`F`$(70 $`\mu `$m) raise $`F`$(24 $`\mu `$m) by no more than 10%; the adopted uncertainty in the absolute flux calibration of the MIPS measurements at 24 $`\mu `$m. These upper limits on temperature do not preclude small amounts of warmer dust within systems that are still consistent with the 24 $`\mu `$m measurements, but they do constrain the properties of circumstellar material able to produce substantial excesses at 70 $`\mu `$m and be undetected by Spitzer.
The mass of circumstellar dust responsible for the luminosity of IR excesses observed in the TWA, or that could produce the upper limits found for $`T_{PDS}`$ and $`L_{IR}/L_{}`$, can be deduced following the inferences of Jura et al. (1995) and Chen & Jura (2001). That is, the minimum mass in dust can be estimated by
$$M_d(16/3)\pi (L_{\mathrm{IR}}/L_{})\rho R_{PDS}^2a$$
(Jura et al., 1995), where $`a`$ is the average radius of a dust grain, $`R_{PDS}`$ is the minimum distance from the star where the grains are in radiative equilibrium, and $`\rho `$ is the mass density of the grains. The equation above assumes that the dust is radiating from a thin shell at a distance $`R_{PDS}`$ from the illuminating star and that the debris disk is optically thin. Also, it is assumed that the dust grains are spherical and have cross sections equal to their geometric cross sections. In Table 4, we tabulate $`M_d`$ using the values for average grain size and density adopted by Chen & Jura (2001) ($`a=2.8`$ $`\mu `$m; $`\rho =2.5`$ g cm<sup>-3</sup>). The minimum dust mass responsible for detected IR excesses range from $`10^{25}`$ g for HR 4796A to $`10^{23}`$ g for TWA 13. For HD 98800B, the assumption that the debris disk is optically thin is not valid, accounting for the much smaller mass estimate than that determined by Low, Hines, & Schneider (1999).
The mass estimates in Table 4 for the TWA are conservative lower bounds for the total circumstellar dust mass in these systems since the MIPS observations are not sensitive to colder material located further from a star and emission from larger particles. For example, there is evidence for material around TWA 7 colder than the $``$80 K dust detected by Spitzer. Observations at 850 $`\mu `$m by Webb (2000a) detected TWA 7 with $`F_\nu 15`$ mJy. This flux density is about a factor of seven higher than the extrapolation to 850 $`\mu `$m of the $``$80 K blackbody determined by the 24 and 70 $`\mu `$m excesses. Webb (2000a) and Webb et al. (2000b) interpret the sub-millimeter emission as coming from 20 K dust with grain sizes on the order of a few hundred microns. As a consequence, the mass estimate for the dust disk around TWA 7 derived by the sub-millimeter measurement is $`10^3\times `$ greater than that estimated from the IR excess.
## 4. Discussion
### 4.1. Infrared Properties of the TWA
The main result of the Spitzer-MIPS observations of the TWA is to show that most stars in the association exhibit no evidence for circumstellar dust and that the 24 $`\mu `$m data place severe limits on the presence of warm ($`T100`$ K) dust within these systems. Although the TWA includes four remarkable IR-excess systems, there are no other objects that come close to their 24 $`\mu `$m output relative to the observed near-IR photospheric emission. The bimodal distribution of warm dust in the TWA, identified by Weinberger et al. (2004) at shorter wavelengths, is confirmed at 24 $`\mu `$m for a larger sample of objects (see Figure 1).
The bimodal nature of 70 $`\mu `$m excesses is not as pronounced as at shorter wavelengths, but there is still a gap of about a factor of 10 between the strengths of the 70 $`\mu `$m excess of Hen 3-600 and the Spitzer detections of TWA 7 and 13. There is also a range of at least $`300`$ between the brightest and faintest of the 70 $`\mu `$m excesses in the TWA relative to their photospheres at this wavelength. The luminosity of the observed IR excesses range from $`25`$% (TW Hya and HD 98800B) to $`0.1`$% (TWA 13) of $`L_{}`$.
It is difficult to fully characterize the IR excesses of TWA 7 and 13 given that they only begin to appear at $``$24 $`\mu `$m. In the case of TWA 13, there is also the complicating factor that the system is a binary and the distribution of the material responsible for the excess is unknown. However, calculating a color temperature for the two objects yields similar results: $`T_\mathrm{c}60`$–80 K. Assuming that the excesses are well fit by blackbodies, both systems have a fractional IR luminosity of $``$1% of that measured for HD 98800B; $`L_{IR}/L_{}=2.0\times 10^3`$ and $`1.7\times 10^3`$, for TWA 7 and 13, respectively. Although the lack of observed excess 24 $`\mu `$m emission constrains the temperatures of possible PDSs for most of the 70 $`\mu `$m non-detections to be $`<100`$ K, the MIPS observational upper limits cannot rule out IR excesses with fractional luminosity in the range of $`10^3`$$`10^4`$. Even at the estimated young age of the association it may already be more likely to detect cooler material, further away from the stars, than warmer material that results in an excess at 24 $`\mu `$m or shorter wavelengths. Clearly, at 70 $`\mu `$m wavelength, large amounts of additional Spitzer time will be needed to detect and measure the cool dust and debris around these stars.
The variety of the IR properties in the TWA is remarkable given the assumption that its members have similar ages. At 8–10 Myr after the formation of the TWA, there are still two systems (TW Hya and Hen 3-600) accreting presumably primordial dust and they have enough circumstellar material to reprocess a large fraction (10–30%) of the stellar luminosity into the infrared. Also in the association are two multiple-star systems, HD 98800 and HR 4796, that have a prominent PDS around only one of the stellar components and the debris is relatively warm ($`T_{PDS}=100`$–170 K). TWA 7 shows evidence for a weak 24 $`\mu `$m excess with the detected material being cooler ($`T_{PDS}80`$ K) and located in orbits $`7`$ AU from the star. The remaining 19 stars, including TWA 7 and 13, have apparently been cleansed to a high degree of dust with $`T>100`$ K. This variety of properties suggests that dust is swept from regions within several AU (see Table 4) of the star very early in the system’s evolution. The fact that stars in the association show no signature of warm dust while other members exhibit strong T Tauri activity implies that the transition period from the T Tauri stage to the more quiescent state observed for most of the TWA K- and M-type stars, is very short.
An efficient mechanism for the removal of dust is thought to be the formation of planets that sweep up dust and larger debris in the regions that they orbit. Weinberger et al. (2004) suggest that planetary formation was rapidly completed in the TWA at least in the region where terrestrial planets would be expected to form. This led to the rapid disappearance of dust in the inner several AU of these systems, leaving only the stellar photosphere to be detected at wavelengths of 24 $`\mu `$m and shorter. Strong stellar winds during the T Tauri stage and the Poynting-Robertson drag are also mechanisms that can decrease the IR emission from a system by ridding the environment of small dust grains. Of course, in most TWA systems, and within 10 Myr of their formation, dust destruction mechanisms must more than compensate for the dust production mechanism of collisions between larger bodies to clear out the terrestrial planet region. HD 98800B and HR 4796A are likely cases where their debris systems indicate that collisions between bodies are still important and produce dust close enough to the stars to result in large excesses at 24 $`\mu `$m.
Nearly all members of the observed sample are X-ray emitters, supporting the assertion that stellar winds in these young, low-mass stars may be an important dust destruction mechanism (see e.g., Hollenbach, Yorke, & Johnstone, 2000). We list the X-ray–to–stellar luminosity ratios, $`L_X/L_{}`$, for TWA stars in the last column of Table 4. The X-ray data are drawn from the ROSAT All Sky Survey Catalog and $`L_X`$ is calculated in the manner described by Sterzik et al. (1999). Chen et al. (2005) find an apparent correlation between IR and X-ray luminosity for their sample of F- and G-type stars with ages spanning 5–20 Myr in the Scorpius-Centaurus OB Association. This possible correlation is in the sense of brighter 24 $`\mu `$m excesses tending to be associated with fainter stars in X-ray flux, and Chen et al. (2005) suggest that the stronger stellar wind implied by the chromospheric activity may help explain the general lack of young systems with 24 $`\mu `$m excesses. Given that only three stars in the TWA sample have confirmed PDSs that produce measurable excesses at 24 $`\mu `$m, it is not possible to test the validity of the X-ray–IR flux correlation. However, HD 98800B, HR 4796A, and TWA 7, are all underluminous in X-rays compared to the average observed X-ray emission for the TWA.
### 4.2. Comparison of the TWA with the Young Solar System
We can compare the Spitzer observations of young stars with models for the early Solar System. Following the discussions in Gaidos (1999) and Jura (2004), we assume that the rate of dust production directly scales as the rate of lunar cratering during the Late Heavy Bombardment. The rate of dust production, $`\dot{M}_{dust}`$, is given by:
$$\dot{M}_{dust}=\dot{M}_0\left(1+\beta e^{\frac{t_9}{\tau }}\right),$$
where $`t_9`$ is the look-back time in 10<sup>9</sup> yr (Gyr), $`\dot{M}_0`$ is the current rate of dust production and $`\beta `$ and $`\tau `$ are fitting constants such that $`\beta =1.6\times 10^{10}`$ and $`\tau `$ = 0.144 (Chyba, 1991). In this model, the dust production rate has been approximately constant during the past 3.3 Gyr, but was as high as $`10^4`$ the current value when $`t_9=4.6`$. We set $`\dot{M}_0=3\times 10^6`$ g s<sup>-1</sup> based on the zodiacal light in the Solar System (Fixsen & Dwek, 2002), implying that the early Solar System may have had a dust production rate of $`3\times 10^{10}`$ g s<sup>-1</sup>. In a model where this quantity of dust is produced far from the star and then loses angular momentum under the action of the Poynting-Robertson effect, it is straightforward to show (e.g., Gaidos, 1999) that the luminosity of the dust, $`L_\nu `$, is given by the expression:
$$L_\nu =\frac{\dot{M}_{dust}c^2}{\nu }.$$
For the TWA, young M-type stars would have luminosities at 24 $`\mu `$m of $`2\times 10^{18}`$ erg s<sup>-1</sup> Hz<sup>-1</sup> and, for distances of $``$50 pc, we expect F<sub>ν</sub>(24 $`\mu `$m) $``$ 700 mJy, vastly greater than what is observed. This result suggests the absence of terrestrial planet-forming environments around these stars. However, in M-type stars where the luminosity is relatively low and the stellar wind rate is relatively large, dust grains mainly lose angular momentum by stellar wind drag rather than through Poynting-Robertson drag (see Jura, 2004; Plavchan et al., 2005). In this case, and if $`\dot{M}_{wind}`$ denotes the stellar wind loss rate, then:
$$L_\nu =\frac{\dot{M}_{dust}}{\dot{M}_{wind}}\frac{L_{}}{\nu }.$$
Although we do not know $`\dot{M}_{wind}`$ for young M-type stars, we can extrapolate from their X-ray emission since the winds are likely to be driven by the same hot corona which produces the X-rays. It is plausible that young M-type stars have mass loss rates 10<sup>2</sup> greater than the current Solar wind loss rate of $`2\times 10^{12}`$ g s<sup>-1</sup>, and we estimate that $`\dot{M}_{wind}=2\times 10^{14}`$ g s<sup>-1</sup>. For main sequence early M-type stars, we may adopt $`L_{}=0.05L_{}`$. Therefore, we expect that at 24 $`\mu `$m, $`L_\nu =2\times 10^{15}`$ erg s<sup>-1</sup> Hz<sup>-1</sup> and a predicted flux at 50 pc of only $``$0.7 mJy. A 24 $`\mu `$m excess this small is beyond the sensitivity of our current measurements, and it is possible that planet forming activity may still be occurring around the young M-type stars in the TWA even though we do not detect an infrared excess.
The TWA gives us an important example of the rapid evolution of the circumstellar material around generally low-mass stars with ages of $`10`$ Myr that is in line with time scales of terrestrial planet formation inferred for our solar system from radiochemistry of meteorites (Yin et al., 2002; Kleine et al., 2002). The diverse IR properties of the association also suggest that the evolutionary processes relevant to planetary formation and accretion disk dispersal may occur at different rates even for stars with similar spectral types. For most of the systems observed, the lack of emission from dust within a few AU of the stars leads to two possible conclusions. Either, (1) the early conditions around pre-main sequence M and late K stars prevent the formation of terrestrial planets in most cases, or that (2) terrestrial planet building has already progressed to the point where the mass is in the form of planetesimals.
This work is based on observations made with the Spitzer Space Telescope, which is operated by the Jet Propulsion Laboratory (JPL), California Institute of Technology (CIT), under National Aeronautics and Space Administration (NASA) contract 1407. We thank NASA, JPL, and the Spitzer Science Center for support through Spitzer, MIPS, and Science Working Group contracts 960785, 959969, and 1256424 to The University of Arizona. We thank M. Blaylock, C. Engelbracht, K. Gordon, K. Misselt, J. Muzerolle, G. Neugebauer, J. Stansberry, K. Stapelfeldt, K. Su, B. Zuckerman, and an anonymous referee for useful comments and discussions. V. Krause acknowledges a Summer Undergraduate Research Fellowship at JPL. This publication makes use of data products from the Two-Micron All Sky Survey, which is a joint project of the University of Massachusetts and the Infrared Processing and Analysis Center/CIT, funded by NASA and the National Science Foundation.
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# Regularity for solutions of the two-phase Stefan problem
## 1 Introduction
In this paper we will discuss the regularity of weak solutions to the two-phase Stefan problem
$$\frac{u}{t}=\mathrm{\Delta }\alpha (u)$$
(1.1)
in a domain $`\mathrm{\Omega }^n\times (0,T)`$, for some $`T>0`$. Here $`\alpha (u)=0`$ if $`1u1`$, $`\alpha (u)=u1`$ for $`u>1`$, and $`\alpha (u)=u+1`$ for $`u<1.`$
We will show that if $`uL_{loc}^2(\mathrm{\Omega })`$ is a solution in the sense of distributions of (1.1) (defined precisely below) then $`\alpha (u)`$ is continuous. In the case when $`u0`$, $`u`$ is a solution of the one-phase Stefan problem and Andreucci and Korten \[AnKo\] (see also Korten \[Ko\]) have shown that if $`uL_{loc}^1,`$ then $`\alpha (u)`$ is continuous. Although we believe this to be true in the two-phase case, we have not been able to obtain this generality and must assume $`uL_{loc}^2.`$
Under the assumption that $`u`$ is bounded and $`\alpha (u)L^2`$, Caffarelli and Evans \[CaE\] showed that $`\alpha (u)`$ is continuous. Similar results for more general singular parabolic equations were shown by Sacks \[S\], Ziemer \[Z\] and by DiBenedetto \[DiB\]. We will assume these results. We will show that a locally $`L^2`$ weak solution of (1.1) satisfies the hypotheses of these results (of any of these authors) to conclude the continuity of $`\alpha (u)`$.
Related to this equation is the porous medium equation $`u_t=\mathrm{\Delta }u^m`$, $`m>1`$. This has been studied extensively by many authors, but we mention in particular the regularity result of Dahlberg and Kenig \[DK\] who showed that a nonnegative $`L_{loc}^m`$ solution to the porous medium equation is a.e. equal to a continuous function. The methods in this present paper are descendants (via the work of Andreucci and Korten) of the methods of Dahlberg and Kenig found in \[DK\]. However, the fact that we are working with solutions which can be both positive and negative complicates matters. To achieve our results we will perform numerous integrations by parts and cannot determine the sign of the resulting boundary terms as in the one-phase case. Consequently we devise a different strategy and introduce new ideas and techniques.
Equation (1.1) is a formulation of the two-phase Stefan problem, describing the flow of heat within a substance which can be in a liquid phase or a solid phase, and for which there is a latent heat to initiate phase change. This allows for the presence of a “mushy zone”, that is, a region which is between the liquid and solid phases. In this model $`u`$ represents the enthalpy and $`\alpha (u)`$ the temperature. We have assumed that the thermal conductivity in both the solid and liquid phases is the same. These conductivities are determined by the slope of the function $`\alpha (u)`$ in the regions $`u1`$, and $`u1`$. The results below all continue to hold (with minor modifications) if the slope of $`\alpha (u)`$ differs in these regions.
We now state our main result. Suppose $`uL_{loc}^2(\mathrm{\Omega })`$ where $`\mathrm{\Omega }`$ is a domain contained in $`^n\times (0,T)`$. We consider distributional solutions of the equation $`u_t=\mathrm{\Delta }\alpha (u)`$, that is, $`u`$ which satisfy
$$_\mathrm{\Omega }\alpha (u)\mathrm{\Delta }\phi +u\phi _tdxdt=0$$
for every $`\phi C^{\mathrm{}}`$ with compact support in $`\mathrm{\Omega }.`$
###### Theorem 1.1.
Suppose $`uL_{loc}^2(\mathrm{\Omega })`$ is a solution of $`u_t=\mathrm{\Delta }\alpha (u).`$ Then $`\alpha (u)`$ is a.e. equal to a continuous function.
We do not expect, in general, such a result for $`u`$. As noted in Korten \[Ko1\], the solution to the Cauchy problem $`u_t=\mathrm{\Delta }\alpha (u)`$ on $`_+^{n+1}`$ with initial data $`0u_I(x)1`$ is just $`u(x,t)=u_I(x)`$. Thus, we cannot expect $`u(x,t)`$ to be any smoother than $`u_I(x).`$
The paper is structured as follows. In section 2 we prove energy estimates for weak solutions of the two phase problem. These show that $`\alpha (u)`$ and $`\alpha (u)_t`$ exist locally in $`L^2.`$ In section 3, we show that $`|\alpha (u)|`$ is subcaloric. An immediate consequence is that $`\alpha (u)`$ is locally bounded. This, combined with the energy estimates and previously mentioned theorem of DiBenedetto \[DiB\] (or others mentioned above) gives the continuity of $`\alpha (u)`$.
Throughout, the letter $`C`$ will denote a constant which may vary from line to line.
The work of the first author was partially supported by a Kansas EPSCoR grant under agreement NSF32169/KAN32170 and a Kansas State University mentoring grant. The second author would like to thank the National University of Ireland at Galway for their hospitality during part of this work.
## 2 Energy Estimates
We establish that $`\alpha (u)`$ has derivatives which are locally in $`L^2`$.
###### Theorem 2.1.
Suppose $`\mathrm{\Omega }_+^{n+1}`$ and $`uL_{\text{ loc }}^2(\mathrm{\Omega })`$ is a distributional solution of $`u_t=\mathrm{\Delta }\alpha (u)`$ on $`\mathrm{\Omega }`$. Suppose $`r<R`$, $`T_0<t_0<t_1<T_1`$, set $`\omega =(t_0,t_1)\times B(x_0,r)`$ and $`\stackrel{~}{\omega }=(T_0,T_1)\times B(x_0,R)`$, and suppose the closure of $`\stackrel{~}{\omega }`$ is contained in $`\mathrm{\Omega }`$. Then $`\alpha (u)`$, $`\alpha (u)_t`$ exist in $`L^2(\omega )`$ and there exists a constant $`C`$, depending only on $`\omega `$ and $`\stackrel{~}{\omega }`$ such that
$$_\omega |\alpha (u)|^2𝑑x𝑑tC_{\stackrel{~}{\omega }}u^2𝑑x𝑑t$$
(2.1)
and
$$_\omega \left|\frac{}{t}\alpha (u)\right|^2𝑑x𝑑tC_{\stackrel{~}{\omega }}u^2𝑑x𝑑t$$
(2.2)
###### Proof.
Let $`\phi _m(y,s)=\rho _m(y)\tau _m(s)`$, $`m=1,2,\mathrm{}`$ where $`\rho _m`$, $`\tau _m`$ are smooth mollifiers, radial, centered at $`0`$, compactly supported, and tending to $`\delta _0`$. For $`(x,t)\mathrm{\Omega }`$ and $`m`$ sufficiently large (depending on $`(x,t)`$), $`\phi _m(xy,ts)`$ is a test function supported in $`\mathrm{\Omega }`$ and thus
$$_\mathrm{\Omega }u(y,s)\frac{\phi _m}{t}(xy,ts)+\alpha (u(y,s))\mathrm{\Delta }\phi _m(xy,ts)dyds=0.$$
In the course of the proof, we will need to define three nested domains between $`\omega `$ and $`\stackrel{~}{\omega }`$. To simplify notation, set $`\omega _1=\omega `$, $`\omega _5=\stackrel{~}{\omega }`$ and we will define $`\omega _2`$, $`\omega _3`$ and $`\omega _4`$ with $`\omega _1\omega _2\omega _3\omega _4\omega _5.`$
For $`m=1,2,3,\mathrm{}`$ set
$$u_m(x,t)=_\mathrm{\Omega }u(y,s)\chi _{\omega _5}(y,s)\phi _m(xy,ts)𝑑y𝑑s.$$
Then for all $`(x,t)`$ and $`m`$, $`|u_m(x,t)|M(u\chi _{\omega _5}(x,t))`$ where $`M`$ denotes the Hardy-Littlewood maximal function (in both the variables $`(x,t)`$).
Choose $`a`$, $`\frac{3T_0+t_0}{4}<a<\frac{T_0+t_0}{2}`$ and $`b`$, $`\frac{t_1+T_1}{2}<b<\frac{t_1+3T_1}{4}`$ such that
$$_{B(x_0,R)}|M(u\chi _{\omega _5})(x,a)|^2𝑑xC_{\omega _5}|M(u\chi _{\omega _5})(x,t)|^2𝑑x𝑑t.$$
(2.3)
with a similar inequality for $`b`$.
In a similar fashion, set $`w_m=\alpha (u)\chi _{\omega _5}\phi _m`$. Define $`\omega _4=B(x_0,\frac{r+3R}{4})\times (\frac{3T_0+t_0}{4},\frac{t_1+3T_1}{4})`$. Then on $`\omega _4`$, $`\frac{}{t}u_m\mathrm{\Delta }w_m=0`$ for all $`m`$ sufficiently large. Using cylindrical coordinates we can choose an $`r_1`$, $`\frac{r+R}{2}<r_1<\frac{r+3R}{4}`$ so that
$$_{B(x_0,r_1)\times (T_0,T_1)}|M(\alpha (u)\chi _{\omega _5})|^2𝑑\sigma C_{\omega _5}|M(\alpha (u)\chi _{\omega _5})(x,t)|^2𝑑x𝑑t.$$
(2.4)
Then by (2.3), for all sufficiently large $`m`$,
$`{\displaystyle _{B(x_0,r_1)}}u_m(x,a)^2𝑑x`$ $`{\displaystyle _{B(x_0,R)}}|M(u\chi _{\omega _5})(x,a)|^2𝑑x`$ (2.5)
$`C{\displaystyle _{\omega _5}}|M(u\chi _{\omega _5})(x,t)|^2𝑑x𝑑t`$
$`C{\displaystyle _{\omega _5}}u^2𝑑x𝑑t`$
Likewise, by (2.4)
$`{\displaystyle _{B(x_0,r_1)\times (a,b)}}w_m^2𝑑\sigma `$ $`{\displaystyle _{B(x_0,r_1)\times (T_0,T_1)}}|M(\alpha (u)\chi _{\omega _5})|^2𝑑\sigma `$ (2.6)
$`C{\displaystyle _{\omega _5}}|M(\alpha (u)\chi _{\omega _5})|^2𝑑x𝑑t`$
$`C{\displaystyle _{\omega _5}}|\alpha (u)|^2𝑑x𝑑t.`$
Let $`\alpha _m(s)`$ be a smooth regularization of $`\alpha (s)`$ such that $`\alpha _m(s)=\alpha (s)`$ for $`|s|1+\frac{1}{m}`$, $`\alpha _m(s)`$ is strictly increasing and $`\alpha _m(s)0`$ except for $`s=0`$. Put $`\omega _3=B(x_0,r_1)\times (a,b).`$ Let $`v_m`$ be a solution to
$$\{\begin{array}{cc}v_t=\mathrm{\Delta }\alpha _m(v)\hfill & \text{ on }\omega _3\hfill \\ v(x,a)=u_m(x,a)\hfill & xB(x_0,r_1)\hfill \\ \alpha _m(v)=w_m\hfill & \text{ on }B(x_0,r_1)\times (a,b).\hfill \end{array}$$
Choose $`\varphi (x)`$ so that $`\varphi =0`$ on $`B(x_0,r_1)`$, $`\mathrm{\Delta }\varphi =1`$ on $`B(x_0,r_1)`$. Then $`\varphi <0`$ on $`B(x_0,r_1)`$ and $`\frac{\varphi }{n}=c_1>0`$ on $`B(x_0,r_1)`$, where $`n`$ is the outward normal and $`c_1`$ is a constant depending only on $`r_1`$ and the dimension.
By Green’s theorem we have
$`{\displaystyle _{B(x_0,r_1)}}(\alpha _m(v_m))^2\mathrm{\Delta }\varphi 𝑑x`$
$`={\displaystyle _{B(x_0,r_1)}}\mathrm{\Delta }(\alpha _m(v_m))^2\varphi 𝑑x+{\displaystyle _{B(x_0,r_1)}}(\alpha _m(u_m))^2{\displaystyle \frac{\varphi }{n}}𝑑\sigma {\displaystyle _{B(x_0,r_1)}}\varphi {\displaystyle \frac{}{n}}[\alpha _m(v_m)]^2𝑑\sigma `$
$`=2{\displaystyle _{B(x_0,r_1)}}\mathrm{\Delta }\alpha _m(v_m)\alpha _m(v_m)\varphi 𝑑x+2{\displaystyle _{B(x_0,r_1)}}|\alpha _m(v_m)|^2\varphi 𝑑x+c_1{\displaystyle _{B(x_0,r_1)}}\alpha _m(v_m)^2𝑑\sigma `$
$`2{\displaystyle _{B(x_0,r_1)}}v_{m_t}\alpha _m(v_m)\varphi 𝑑x+c_1{\displaystyle _{B(x_0,r_1)}}\alpha _m(v_m)^2𝑑\sigma `$
$`=2{\displaystyle \frac{d}{dt}}{\displaystyle _{B(x_0,r_1)}}A_m(v_m)\varphi 𝑑x+c_1{\displaystyle _{B(x_0,r_1)}}\alpha _m(v_m)^2𝑑\sigma `$
where $`A_m`$ is an antiderivative of $`\alpha _m`$. Integrate from $`a`$ to $`b`$ to obtain
$`{\displaystyle _a^b}{\displaystyle _{B(x_0,r_1)}}\alpha _m(v_m)^2\mathrm{\Delta }\varphi 𝑑x𝑑t2{\displaystyle _{B(x_0,r_1)}}`$ $`A_m(v_m(x,b))\varphi dx2{\displaystyle _{B(x_0,r_1)}}A_m(v_m(x,a))\varphi 𝑑x`$
$`+{\displaystyle _a^b}{\displaystyle _{B(x_0,r_1)}}\alpha _m(v_m)^2𝑑\sigma 𝑑t.`$
Now $`0A_m(x)x^2`$, $`\varphi <0`$, $`\mathrm{\Delta }\varphi =1`$ and recalling $`\omega _3=B(x_0,r_1)\times (a,b)`$, this yields:
$`{\displaystyle _{w_3}}\alpha _m(v_m)^2𝑑x𝑑t`$ $`C{\displaystyle _{B(x_0,r_1)}}v_m(x,a)^2𝑑x+{\displaystyle _a^b}{\displaystyle _{B(x_0,r_1)}}\alpha _m(v_m)^2𝑑\sigma 𝑑t`$ (2.7)
$`=C{\displaystyle _{B(x_0,r_1)}}u_m(x,a)^2𝑑x+{\displaystyle _a^b}{\displaystyle _{B(x_0,r_1)}}w_m^2𝑑\sigma 𝑑t`$
$`C{\displaystyle _{w_5}}u^2𝑑x𝑑t`$
where for the last inequality we have used (2.5) and (2.6). Let $`\psi (x)`$ be a nonnegative $`C_0^{\mathrm{}}(^n)`$ function such that $`\psi 1`$ on $`B(x_0,\frac{3r+R}{4})`$, $`\psi 0`$ outside $`B(x_0,\frac{r+R}{2})`$. To simplify notation set $`B(x_0,\frac{3r+R}{4})=B_1`$, $`B(x_0,\frac{r+R}{2})=B_2`$. Then
$`{\displaystyle _{B_2}}\psi \alpha _m(v_m)v_{m}^{}{}_{t}{}^{}𝑑x`$ $`={\displaystyle _{B_2}}\psi \alpha _m(v_m)\mathrm{\Delta }\alpha _m(v_m)𝑑x`$
$`={\displaystyle _{B_2}}\psi \alpha _m(v_m)\alpha _m(v_m)𝑑x{\displaystyle _{B_2}}\psi |\alpha _m(v_m)|^2𝑑x`$
$`={\displaystyle \frac{1}{2}}{\displaystyle _{B_2}}\mathrm{\Delta }\psi \alpha _m(v_m)^2𝑑x{\displaystyle _{B_2}}\psi |\alpha _m(v_m)|^2𝑑x.`$
Rearrange and integrate from $`a`$ to $`b`$ to obtain
$`{\displaystyle _a^b}{\displaystyle _{B_2}}\psi (\alpha _m(v_m))^2𝑑x𝑑t`$ $`={\displaystyle \frac{1}{2}}{\displaystyle _a^b}{\displaystyle _{B_2}}\mathrm{\Delta }\psi |\alpha _m(v_m)^2|𝑑x𝑑t{\displaystyle _a^b}{\displaystyle _{B_2}}\psi {\displaystyle \frac{d}{dt}}A_m(v_m)𝑑x𝑑t`$ (2.8)
$`C{\displaystyle _a^b}{\displaystyle _{B_2}}\alpha _m(v_m)^2𝑑x𝑑t+{\displaystyle _{B_2}}\psi (x)A_m(v_m(x,a))𝑑x`$
$`C{\displaystyle _a^b}{\displaystyle _{B_2}}\alpha _m(v_m)^2𝑑x𝑑t+{\displaystyle _{B_2}}v_m(x,a)^2𝑑x`$
$`C{\displaystyle _{w_3}}\alpha _m(v_m)^2𝑑x𝑑t+{\displaystyle _{B(x_0,r_1)}}v_m(x,a)^2𝑑x`$
$`C{\displaystyle _{w_5}}u^2𝑑x𝑑t`$
where we have used (2.7) and (2.5) and the definition of $`v_m`$ for the last inequality.
We now seek a similar estimate for the $`t`$ derivative. Let $`\eta (x)`$ be a nonnegative $`C_0^{\mathrm{}}(^n)`$ function such that $`\eta 1`$ on $`B(x_0,r)`$, $`\eta 0`$ outside $`B_1`$ and so that $`\frac{\eta }{\sqrt{\eta }}_{\mathrm{}}<\mathrm{}`$. Note that $`\alpha _m(v_m)_t=\alpha _m^{}(v_m)v_{m}^{}{}_{t}{}^{}`$ and $`0<\alpha _m^{}1`$ so that $`(\alpha _m(v_m)_t)^2\alpha _m(v_m)_tv_{m_t}.`$ Then
$`{\displaystyle _{B_1}}\eta (\alpha _m(v_m)_t)^2𝑑x`$ $`{\displaystyle _{B_1}}\eta \alpha _m(v_m)_tv_{m_t}𝑑x`$
$`={\displaystyle _{B_1}}\eta \alpha _m(v_m)_t\mathrm{\Delta }\alpha _m(v_m)𝑑x`$
$`={\displaystyle _{B_1}}\eta \alpha _m(v_m)\alpha _m(v_m)_t𝑑x{\displaystyle _{B_1}}\eta \alpha _m(v_m)_t\alpha _m(v_m)𝑑x`$
$`={\displaystyle _{B_1}}\eta \alpha _m(v_m)\alpha _m(v_m)_t𝑑x{\displaystyle \frac{1}{2}}{\displaystyle \frac{d}{dt}}{\displaystyle _{B_1}}\eta |\alpha _m(v_m)|^2𝑑x`$
Integrate from $`c`$ to $`d`$, where $`c`$ and $`d`$ are to be chosen momentarily. We obtain
$`{\displaystyle _c^d}{\displaystyle _{B_1}}`$ $`\eta (\alpha _m(v_m)_t)^2dxdt`$ (2.9)
$`{\displaystyle _c^d}\left|{\displaystyle _{B_1}}\sqrt{\eta }{\displaystyle \frac{\eta }{\sqrt{\eta }}}\alpha _m(v_m)\alpha _m(v_m)_t𝑑x\right|𝑑t+{\displaystyle \frac{1}{2}}{\displaystyle _{B_1}}\eta (x)|\alpha _m(v_m)(x,c)|^2𝑑x`$
$`{\displaystyle \frac{\eta }{\sqrt{\eta }}}_{\mathrm{}}\left({\displaystyle _c^d}{\displaystyle _{B_1}}|\alpha _m(v_m)|^2𝑑x𝑑t\right)^{\frac{1}{2}}\left({\displaystyle _c^d}{\displaystyle _{B_1}}\eta (\alpha _m(v_m)_t)^2𝑑x𝑑t\right)^{\frac{1}{2}}`$
$`+{\displaystyle \frac{1}{2}}{\displaystyle _{B_1}}\eta (x)|\alpha _m(v_m)(x,c)|^2𝑑x.`$
Choose $`c_m`$ (depending on $`m`$), $`\frac{T_0+3t_0}{4}<c_m<t_0`$, so that
$$_{B_1}\eta (x)|\alpha _m(v_m)(x,c_m)|^2𝑑xC_a^b_{B_2}\psi |\alpha _m(v_m)|^2𝑑x𝑑t.$$
(2.10)
Put $`d=t_1`$, $`c=c_m`$ in (2.9). Then recalling that $`\psi 1`$ on $`B_1`$, and using (2.10) and (2.8) we have
$`{\displaystyle _{c_m}^{t_1}}`$ $`{\displaystyle _{B_1}}\eta (\alpha _m(v_m)_t)^2𝑑x𝑑t`$
$`{\displaystyle \frac{\eta }{\sqrt{\eta }}}_{\mathrm{}}\left({\displaystyle _{c_m}^{t_1}}{\displaystyle _{B_2}}\psi |\alpha _m(v_m)|^2𝑑x𝑑t\right)^{\frac{1}{2}}\left({\displaystyle _{c_m}^{t_1}}{\displaystyle _{B_1}}\eta (\alpha _m(v_m)_t)^2𝑑x𝑑t\right)^{\frac{1}{2}}+C{\displaystyle _{\omega _5}}u^2𝑑x`$
$`{\displaystyle \frac{\eta }{\sqrt{\eta }}}_{\mathrm{}}\left({\displaystyle _{\omega _5}}u^2𝑑x𝑑t\right)^{\frac{1}{2}}\left({\displaystyle _{c_m}^{t_1}}{\displaystyle _{B_1}}\eta (\alpha _m(v_m)_t)^2𝑑x𝑑t\right)^{\frac{1}{2}}+C{\displaystyle _{\omega _5}}u^2𝑑x𝑑t`$
from which it follows that
$$_{c_m}^{t_1}_{B_1}\eta (\alpha _m(v_m)_t)^2𝑑x𝑑tC_{\omega _5}u^2𝑑x𝑑t$$
and consequently
$$_{\omega _1}(\alpha _m(v_m)_t)^2𝑑x𝑑tC_{\omega _5}u^2𝑑x𝑑t.$$
(2.11)
Thus, recalling $`\omega =\omega _1`$, $`\stackrel{~}{\omega }=\omega _5`$, (2.8) and (2.11) give
$$_\omega |\alpha _m(v_m)|^2𝑑x𝑑tC_{\stackrel{~}{\omega }}u^2𝑑x𝑑t\text{and}_\omega (\alpha _m(v_m)_t)^2𝑑x𝑑tC_{\stackrel{~}{\omega }}u^2𝑑x𝑑t.$$
(2.12)
To obtain (2.1) and (2.2) we will need to take limits. We first remark that with more care, similar estimates could be obtained with any compact set $`K\omega _3`$ replacing $`\omega =\omega _1`$ on the left hand side of the inequalities in (2.12); naturally, the constants on the right hand side depend on the position of $`K`$ within $`\omega _3.`$ Thus, from (2.7) and this observation, we have:
$`{\displaystyle _{w_3}}\alpha _m(v_m)^2𝑑x𝑑tC{\displaystyle _{\stackrel{~}{\omega }}}u^2𝑑x𝑑t,`$ $`{\displaystyle _K}|\alpha _m(v_m)|^2𝑑x𝑑tC(K){\displaystyle _{\stackrel{~}{\omega }}}u^2𝑑x𝑑t`$ (2.13)
$`\text{and }{\displaystyle _K}\left|{\displaystyle \frac{}{t}}\alpha _m(v_m)\right|^2`$ $`dxdtC(K){\displaystyle _{\stackrel{~}{\omega }}}u^2𝑑x𝑑t`$
for every compact $`K\omega _3.`$
By Rellich-Kondrachov there exists a subsequence $`\{\alpha _{m_k}(v_{m_k})\}`$ of $`\{\alpha _m(v_m)\}`$ (which we still write as $`\{\alpha _m(v_m)\}`$) and $`hL^2(\omega _3)`$ such that $`\alpha _m(v_m)h`$ in $`L^2(K)`$ for every compact set $`K\omega _3.`$ By taking subsequences, if necessary, we also may assume this convergence is a.e. By weak compactness, and again, by taking subsequences, we may assume that $`\alpha _m(v_m)h`$ weakly in $`L^2(\omega _3)`$. Equation (2.13) implies that the $`L^2(\omega _3)`$ norms of the $`v_m`$ are uniformly bounded, hence there exists a subsequence, (still denoted by $`v_m`$) such that $`v_mvL^2(\omega _3)`$ weakly.
We claim that $`\alpha (v)=h`$. First note that $`\alpha _m\alpha _{\mathrm{}}0`$, so that for a.e $`x\omega _3`$, $`\alpha (v_m)h`$. Consider the set where $`h>0`$. Then for a.e $`x`$ in this set, $`\alpha (v_m(x))h(x)>0`$, and hence $`v_m(x)h(x)+1`$. Thus, $`v(x)=h(x)+1`$ for a.e $`x`$ in the set where $`h(x)>0`$. Similarly, on $`h<1`$, $`v(x)=h(x)1`$ a.e. On the set $`h(x)=0`$ we must have $`1lim\; infv_m(x)lim\; supv_m(x)1`$ a.e. To see this consider an $`x`$ at which there exists a subsequence $`v_{m_k}(x)`$ which converges to $`y_0[1,1].`$ Then for this $`x`$, $`\alpha (v_{m_k}(x))\alpha (y_0)0`$ which implies $`\alpha (v_m(x))h(x)`$. Thus, $`1lim\; infv_m(x)lim\; supv_m(x)1`$ a.e. on $`h(x)=0`$, and hence $`1v(x)1`$ a.e. on $`h=0`$. We conclude that $`\alpha (v)=h`$ a.e.
Summarizing, we have $`v_mv`$ weakly in $`L^2(\omega _3)`$ and $`\alpha _m(v_m)\alpha (v)`$ weakly in $`L^2(\omega _3)`$, a.e. on $`\omega _3`$ and in $`L^2(K)`$ for every compact subset $`K`$ of $`\omega _3`$. To finish the proof we show that $`\alpha (u)=\alpha (v)`$ a.e. on $`\omega _3`$. Using integration by parts, and recalling that $`\alpha _m(v_m)=w_m`$ on $`B(x_0,r_1)\times (a,b)`$, we compute
$`{\displaystyle _{\omega _3}}(v_mu_m)`$ $`(\alpha _m(v_m)w_m)dxdt`$ (2.14)
$`={\displaystyle _a^b}{\displaystyle _{B(x_0,r_1)}}{\displaystyle _a^t}v_{m}^{}{}_{t}{}^{}(x,\tau )u_{m}^{}{}_{t}{}^{}(x,\tau )d\tau (\alpha _m(v_m(x,t))w_m(x,t))dxdt`$
$`={\displaystyle _a^b}{\displaystyle _{B(x_0,r_1)}}{\displaystyle _a^t}(\alpha _m(v_m)w_m)(x,\tau )𝑑\tau (\alpha _m(v_m)w_m)dxdt`$
$`={\displaystyle \frac{1}{2}}{\displaystyle _a^b}{\displaystyle _{B(x_0,r_1)}}{\displaystyle \frac{d}{dt}}\left|{\displaystyle _a^t}(\alpha _m(v_m)w_m)d\tau \right|^2𝑑x𝑑t`$
$`={\displaystyle \frac{1}{2}}{\displaystyle _{B(x_0,r_1)}}\left|{\displaystyle _a^b}(\alpha _m(v_m)w_m)d\tau \right|^2𝑑x0`$
We need to take limits as $`m\mathrm{}`$ in this inequality. Write
$`{\displaystyle _{\omega _3}}(v_mu_m)(\alpha _m(v_m)w_m)𝑑x𝑑t={\displaystyle _{\omega _3}}v_m\alpha _m(v_m)𝑑x𝑑t+{\displaystyle _{\omega _3}}v_m(w_m)𝑑x𝑑t`$
$`+{\displaystyle _{\omega _3}}(u_m)(\alpha _m(v_m))𝑑x𝑑t+{\displaystyle _{\omega _3}}u_mw_m𝑑x𝑑t=I+II+III+IV`$
Since $`u_mu`$ a.e. and in $`L^2(\omega _3)`$, $`w_m\alpha (u)`$ a.e. and in $`L^2(\omega _3)`$, $`v_mv`$ weakly, and $`\alpha _m(v_m)\alpha (v)`$ weakly, we conclude
$`II{\displaystyle _{\omega _3}}v(\alpha (u))`$ $`dxdt,III{\displaystyle _{\omega _3}}(u)(\alpha (v))𝑑x𝑑t,`$
$`\text{and }IV{\displaystyle _{\omega _3}}u\alpha (u)𝑑x𝑑t.`$
Expand out
$$_{\omega _3}(\alpha _m(v_m)\alpha _m(v))(v_mv)𝑑x0,$$
take $`m\mathrm{}`$ (make use of the fact that $`\alpha _m\alpha _{\mathrm{}}0`$) to conclude
$$\underset{m\mathrm{}}{lim\; inf}_{\omega _3}\alpha _m(v_m)v_m𝑑x_{\omega _3}\alpha (v)v𝑑x.$$
This combined with the estimates for II-IV and (2.14) yields
$$_{\omega _3}(vu)(\alpha (v)\alpha (u))𝑑x𝑑t0.$$
Since the integrand of this is nonnegative, we conclude $`\alpha (u)=\alpha (v)`$ a.e. on $`\omega _3`$. This completes the proof of the theorem. $`\mathrm{}`$
## 3 $`|\alpha (u)|`$ is subcaloric
###### Theorem 3.1.
$`|\alpha (u)|`$ is weakly subcaloric, that is, it satisfies
$$_\mathrm{\Omega }|\alpha (u)|\eta +|\alpha (u)|\eta _tdxdt0$$
for any nonnegative $`\eta W_0^{1,2}(\mathrm{\Omega }).`$
###### Proof.
Let $`0\eta W_0^{1,2}(\mathrm{\Omega }).`$ For $`h>0`$ set
$$\varphi _h(x)=\{\begin{array}{cc}1\hfill & \text{if }x>h\hfill \\ \frac{2}{h}x1\hfill & \frac{h}{2}x<h\hfill \\ 0\hfill & \text{if }|x|<\frac{h}{2}\hfill \\ \frac{2}{h}x+1\hfill & h<x\frac{h}{2}\hfill \\ 1\hfill & \text{if }x<h.\hfill \end{array}$$
Then $`\eta \varphi _h(\alpha (u))`$ is supported in $`\left\{|\alpha (u)|>\frac{h}{2}\right\}`$ and thus, $`(\alpha (u)u)[\eta \varphi _h(\alpha (u))]_t𝑑x𝑑t=0.`$
Then
$`0`$ $`={\displaystyle \alpha (u)[\eta \varphi _h(\alpha (u))]_t}\alpha (u)[\eta \varphi _h(\alpha (u))]dxdt`$ (3.1)
$`={\displaystyle \alpha (u)\eta _t\varphi _h(\alpha (u))𝑑x𝑑t}+{\displaystyle \alpha (u)\eta \varphi _h^{}(\alpha (u))\alpha (u)_t𝑑x𝑑t}`$
$`{\displaystyle \alpha (u)\eta \varphi _h(\alpha (u))𝑑x𝑑t}{\displaystyle \alpha (u)\eta \varphi _h^{}(\alpha (u))\alpha (u)𝑑x𝑑t}`$
$`=I+II+III+IV`$
We investigate each of these as $`h0`$. As $`h0`$, $`\varphi _h(\alpha (u))\text{sgn}(\alpha (u))`$ so that $`I|\alpha (u)|n_t𝑑x𝑑t`$. To estimate $`II`$, first note that
$$\varphi _h^{}(\alpha (u))=\frac{2}{h}\chi _{\left\{\frac{h}{2}<|\alpha |<h\right\}}.$$
Then
$$II=\alpha (u)\eta \frac{2}{h}\chi _{\{\frac{h}{2}<|\alpha (u)|<h\}}\frac{}{t}\alpha (u)𝑑x𝑑t$$
and consequently,
$$|II||\alpha (u)|\eta \frac{2}{h}\chi _{\{\frac{h}{2}<|\alpha (u)|<h\}}\left|\frac{}{t}\alpha (u)\right|𝑑x𝑑t2\eta \chi _{\{\frac{h}{2}<\alpha (u)<h\}}\left|\frac{}{t}(\alpha (u))\right|𝑑x𝑑t.$$
Since $`\frac{}{t}\alpha (u)L_{loc}^2(\mathrm{\Omega })`$, $`II0`$ as $`h0`$
To estimate $`III`$, note that when $`|\alpha (u)|>h`$, $`\alpha (u)\varphi _h(\alpha (u))=|\alpha (u)|`$. And when $`\frac{h}{2}<|\alpha (u)|<h`$,
$$|\alpha (u)\varphi _h(\alpha (u))||\alpha (u)|\left(\frac{2}{h}|\alpha (u)|+1\right)$$
Consequently,
$`\left|{\displaystyle _{\{\frac{h}{2}<|\alpha (u)|<h\}}}\alpha (u)\eta \varphi _h(\alpha (u))𝑑x𝑑t\right|`$ $`{\displaystyle _{\{\frac{h}{2}<|\alpha (u)|<h\}}}\left|\alpha (u)\right|\left[{\displaystyle \frac{2}{h}}|\alpha (u)|+1\right]|\eta |𝑑x𝑑t`$
$`0\text{ as }h0\text{ since }\alpha (u)L_{loc}^2(\mathrm{\Omega })).`$
Therefore, as $`h0`$, $`III|\alpha (u)|\eta dxdt.`$ Note that we can write $`IV`$ as
$$IV=|\alpha (u)|^2\eta \varphi _h^{}(\alpha (u))𝑑x𝑑t=|\alpha (u)|^2\eta \frac{2}{h}\chi _{\left\{\frac{h}{2}<|\alpha (u)|<h\right\}}𝑑x𝑑t$$
Thus, letting $`h0`$ in (3.1) yields:
$$0=|\alpha (u)|\eta _t𝑑x𝑑t|\alpha (u)|\eta dxdt\underset{h0}{lim}|\alpha (u)|^2\frac{2}{h}\chi _{\left\{\frac{h}{2}<|\alpha (u)|<h\right\}}\eta 𝑑x𝑑t$$
from which the theorem follows.
$`\mathrm{}`$
Remark. Suppose that instead of $`\varphi _h`$ as defined above, we defined
$$\varphi _h(x)=\{\begin{array}{cc}1\hfill & \text{if }x>h\hfill \\ \frac{2}{h}x1\hfill & \text{if }\frac{h}{2}<x<h\hfill \\ 0\hfill & \text{otherwise}.\hfill \end{array}$$
Then following the computations as in (3.1) we obtain (3.1) with this version of $`\varphi _h`$. In this case $`I\alpha (u)^+\eta _t𝑑x𝑑t`$ and as in the above case, $`II0`$, and $`III\alpha (u)^+\eta dxdt.`$ Now we may write
$$IV=|\alpha (u)^+|^2\eta \frac{2}{h}\chi _{\left\{\frac{h}{2}<\alpha (u)<h\right\}}𝑑x𝑑t.$$
We obtain:
$$0=\alpha (u)^+\eta _t𝑑x𝑑t\alpha (u)^+\eta dxdt\underset{h0}{lim}|\alpha (u)^+|^2\frac{2}{h}\chi _{\left\{\frac{h}{2}<\alpha (u)<h\right\}}\eta 𝑑x𝑑t$$
Thus, $`\alpha (u)^+`$ is subcaloric. In a similar fashion, we may use the function
$$\varphi _h(x)=\{\begin{array}{cc}\frac{2}{h}x+1\hfill & h<x<\frac{h}{2}\hfill \\ 1\hfill & \text{if }x<h\hfill \\ 0\hfill & \text{otherwise}\hfill \end{array}$$
and computations such as those above to obtain
$$0=\alpha (u)^{}\eta _t𝑑x𝑑t\alpha (u)^{}\eta dxdt+\underset{h0}{lim}|\alpha (u)^{}|^2\frac{2}{h}\chi _{\left\{h<\alpha (u)<\frac{h}{2}\right\}}\eta 𝑑x𝑑t$$
to conclude that $`\alpha (u)^{}`$ is supercaloric.
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# Second order formalism in Poincaré gauge theory
## 1 Introduction
In general relativity, the procedure to consider, during the variation of the Lagrangian, the connection and the metric as independent of each other is known as Palatini, or first order, formalism. Indeed the classical Hilbert-Einstein Lagrangian contains no second derivatives of those fields. On the other hand, in the conventional metric approach of general relativity, where the only independent field is the metric, the Lagrangian contains second derivatives of the metric, although incidentally, the terms containing them can be eliminated by subtracting a boundary term. This is known as second order formalism.
In Poincaré gauge theory (PGT), the usual procedure is to consider the connection $`\mathrm{\Gamma }_m^{ab}`$ and the tetrad $`e_m^a`$ as independent fields and consequently, the field equations are derived through variation with respect to those fields. For a detailed description of PGT, we refer to Refs. and . Usually, the Lagrangians are constructed from terms at most second order in curvature and torsion, and thus contain no second derivatives of the independent fields. We will therefore refer to this as the first order formalism.
However, in view of the relation $`\mathrm{\Gamma }_m^{ab}=\widehat{\mathrm{\Gamma }}_m^{ab}+K_m^{ab}`$, where $`\widehat{\mathrm{\Gamma }}_m^{ab}`$ is the Christoffel connection, which can be expressed as function of the tetrad (and its derivatives) only, and $`K_m^{ab}`$ is the contortion tensor, one might equally well consider $`K_m^{ab}`$ and $`e_m^a`$ as independent fields and carry out the variation with respect to those fields. Then, similar to the classical approach in general relativity (GR), the Christoffel part of the connection is considered as a function of the tetrad. We will refer to this as the second order formalism.
What could be the use of this formalism? Well, in conventional PGT, the most general Lagrangian is constructed under the requirement that it should not contain second and higher order derivatives of the independent fields $`e_m^a`$ and $`\mathrm{\Gamma }_m^{ab}`$ for the usual reasons (Cauchy problem). It is easy to see that, if we require instead the Lagrangian not to contain second and higher order derivatives of $`e_m^a`$ and $`K_m^{ab}`$, now seen as independent fields, then this will allow for other terms in the Lagrangian. Also, some terms previously allowed are now forbidden, as is the case, e.g., for terms of second order in the curvature, which will contain second derivatives of the tetrad field in this formalism. The Cauchy problem can now be stated with respect to $`e_m^a`$ and $`K_m^{ab}`$. Since the determination of the latter fields allows for the determination of $`\mathrm{\Gamma }_m^{ab}`$ too, no initial value problem will arise this way.
On the other hand, one could argue that, in view of the underlying gauge structure of the theory, it appears quite unnatural to consider as fundamental variables a set different to the gauge fields of the Poincaré group. It turns out, however, that both formalisms are completely equivalent. In other words, it is simply a matter of choice and convenience which formalism we use. This equivalence ultimately means that theories that contain higher derivatives in one formalism may be reformulated in the other formalism such that no higher derivatives occur. The change of the fundamental variables thus does not change anything concerning the viability or consistency of the theory, but it shows that certain theories that are apparently in conflict with the Cauchy initial value problem need not be so in reality.
Summarizing, the second order formalism allows us to investigate different Lagrangians than those usually considered in PGT. This, in turn, will allow us to construct a theory with a complete classical general relativity limit, and thus a complete agreement with the experimental situation, without any constraints on eventual parameters of the theory, but with dynamical torsion fields arising in the presence of spinning matter. As we have pointed out in Ref. , this is not possible in conventional PGT for the following reasons. Since in the usual first order approach, the independent fields are the tetrad $`e_m^a`$ and the Lorentz connection $`\mathrm{\Gamma }_m^{ab}`$, in order to get propagating modes for both fields, their first derivatives should appear quadratically in the Lagrangian. Therefore, apart from a term suitable of producing the general relativity limit, there has to be at least one term quadratic in the curvature tensor. In the presence of spinless matter, we wish the field equations to reduce to the Einstein equations of general relativity. If the Cartan equation (i.e., the connection equation) in this case leads to a vanishing torsion (Riemannian geometry), then the term quadratic in the curvature will necessarily contribute to the Einstein equation, which will thus not be of the GR form. (And therefore, in order to be in agreement with experiments, constraints will have to be imposed on the coupling constants of such theories.) The only exceptions to this are very artificial Lagrangians, with terms of the form $`R_{[ik]}R^{[ik]}`$ and $`R_{[ikm]}^lR_l^{[ikm]}`$, constructed from the antisymmetric part of the Ricci tensor and the completely antisymmetric (in the last three indices) part of the curvature tensor. Those terms will lead to contributions in the Einstein equation that vanish together with the torsion in the spinless case (in view of the Bianchi identities in the Riemannian limit). In the presence of spinning matter fields, those terms will eventually lead to propagating torsion, but in general, the field equations have very few and quite strange solutions, certainly not of the form expected for a Yang-Mills like theory. On the other hand, theories whose Cartan equation in the vanishing spin limit leads to a vanishing curvature (teleparallel geometry) and presenting a Yang-Mills like behavior for the Lorentz connection can easily be constructed using curvature squared terms, without problems with the general relativity limit. Those theories have been analyzed in detail in Ref. , where we have shown that they all present problems with an additional symmetry arising in the classical limit. Schematically, we can write the linear approximation of such theories in the following form of wave equations
$$\mathrm{}\mathrm{\Gamma }_m^{ab}=\sigma _m^{ab}\text{and}\mathrm{}g_{ik}=T_{ik}+𝒪(R),$$
(1)
which leads, in the spinless case, to $`\mathrm{\Gamma }_m^{ab}=0`$, and therefore the curvature corrections $`𝒪(R)`$ will vanish and we are left with the GR equation $`\mathrm{}g_{ik}=T_{ik}`$. The problem is that, even for $`\mathrm{\Gamma }_m^{ab}`$ fixed to zero, we can Lorentz rotate the tetrad field $`e_m^a\mathrm{\Lambda }_b^ae_m^b`$, without changing the metric $`g_{ik}`$. Thus, the tetrad field is not uniquely fixed by the field equations, and therefore, the behavior of spinning test matter, e.g., the spin precession and the trajectory of a Dirac particle, cannot be uniquely predicted . It is clear that this argumentation is valid independently of the formalism (first or second) one uses for the derivation of the field equations. It is a general problem of the set of equations $`\mathrm{\Gamma }_m^{ab}=0`$ and $`G_{ik}=T_{ik}`$, i.e., of the teleparallel form of GR.
On the other hand, a theory with an GR limit in a Riemannian geometry (vanishing torsion), should behave as
$$\mathrm{}K_m^{ab}=\sigma _m^{ab}\text{and}\mathrm{}g_{ik}=T_{ik}+𝒪(K),$$
(2)
where now the spinless case leads to a vanishing contortion $`K_m^{ab}=0`$, and thus to a vanishing torsion, and therefore, as before, the second equation reduces to the field equations of GR. In contrast to (1), the set of equations (2) fixes the geometry completely. Indeed, if we know the contortion and the metric, than the independent field tensors, torsion and curvature, are fixed, up to a Poincaré transformation of course, which is the underlying gauge symmetry of the theory .
However, as we have argued above, no Lagrangian of conventional PGT leads to equations of the form (2). In order to get such field equations using the conventional first order approach, one would have to use higher order derivatives of the fields $`(\mathrm{\Gamma }_m^{ab},e_m^a)`$ in the Lagrangian, which is (apparently) forbidden in view of the initial value problem.
We will solve this problem using the second order approach and construct an explicit example of a theory with equations of type (2), containing only first order derivatives of the independent fields $`(K_m^{ab},e_m^a)`$.
The article is structured as follows. In the next section, we will establish the equivalence of the first and second order formalism as far as the field equations are concerned. In section 3, we will compare the Noether currents and their conservation equations in the two approaches. Section 4 is devoted to the construction of an exemplifying theory that, using the second order formalism, allows for a full general relativity limit in the absence of spinning matter. Dynamical torsion fields will arise if the matter possesses a non-vanishing spin density. In section 5, extending slightly the Yasskin ansatz known from Einstein-Yang-Mills theory with internal symmetry group, we simplify the system of field equations for a special class of solutions, suitable for the description of spin polarized neutron stars. Finally, in section 6, explicit solutions are discussed and in section 7, the effects of the torsion fields on the spin precession of a test body are briefly analyzed in view of an eventual experimental detection of the non-Riemannian contributions.
## 2 Equivalence of first and second order formalisms
Let us first give a short review of the basic concepts of Riemann-Cartan geometry and fix our notations and conventions. For a complete introduction into the subject, the reader may consult Refs. and .
Latin letters from the beginning of the alphabet ($`a,b,c\mathrm{}`$) run from 0 to 3 and are (flat) tangent space indices. Especially, $`\eta _{ab}`$ is the Minkowski metric $`diag(1,1,1,1)`$ in tangent space. Latin letters from the middle of the alphabet ($`i,j,k\mathrm{}`$) are indices in a curved spacetime with metric $`g_{ik}`$ as before. We introduce the Poincaré gauge fields, the tetrad $`e_m^a`$ and the connection $`\mathrm{\Gamma }_m^{ab}`$ (antisymmetric in $`ab`$), as well as the corresponding field strengths, the curvature and torsion tensors
$`R_{lm}^{ab}`$ $`=`$ $`\mathrm{\Gamma }_{m,l}^{ab}\mathrm{\Gamma }_{l,m}^{ab}+\mathrm{\Gamma }_{cl}^a\mathrm{\Gamma }_m^{cb}\mathrm{\Gamma }_{cm}^a\mathrm{\Gamma }_l^{cb}`$ (3)
$`T_{lm}^a`$ $`=`$ $`e_{m,l}^ae_{l,m}^a+e_m^b\mathrm{\Gamma }_{bl}^ae_l^b\mathrm{\Gamma }_{bm}^a.`$ (4)
The spacetime connection $`\mathrm{\Gamma }_{ml}^i`$ and the spacetime metric $`g_{ik}`$ can now be defined through
$`e_{m,l}^a+\mathrm{\Gamma }_{bl}^ae_m^b`$ $`=`$ $`e_i^a\mathrm{\Gamma }_{ml}^i`$ (5)
$`e_i^ae_k^b\eta _{ab}`$ $`=`$ $`g_{ik}.`$ (6)
It is understood that there exists an inverse to the tetrad, such that $`e_i^ae_b^i=\delta _b^a`$. It can easily be shown that the connection splits into two parts,
$$\mathrm{\Gamma }_m^{ab}=\widehat{\mathrm{\Gamma }}_m^{ab}+K_m^{ab},$$
(7)
such that $`\widehat{\mathrm{\Gamma }}_m^{ab}`$ is torsion-free and is essentially a function of $`e_m^a`$. $`K_m^{ab}`$ is the contortion tensor (see below). Especially, the spacetime connection $`\widehat{\mathrm{\Gamma }}_{ml}^i`$ constructed from
$$e_{m,l}^a+\widehat{\mathrm{\Gamma }}_{bl}^ae_m^b=e_i^a\widehat{\mathrm{\Gamma }}_{ml}^i$$
(8)
is just the (symmetric) Christoffel connection of general relativity, a function of the metric only.
The gauge fields $`e_m^a`$ and $`\mathrm{\Gamma }_m^{ab}`$ are covector fields with respect to the spacetime index $`m`$. Under a local Lorentz transformation (more precisely, the Lorentz part of a Poincaré transformation, see Ref. ) in tangent space, $`\mathrm{\Lambda }_b^a(x^m)`$, they transform as
$$e_m^a\mathrm{\Lambda }_b^ae_m^b,\mathrm{\Gamma }_{bm}^a\mathrm{\Lambda }_c^a\mathrm{\Lambda }_b^d\mathrm{\Gamma }_{dm}^c\mathrm{\Lambda }_{c,m}^a\mathrm{\Lambda }_b^c.$$
(9)
The torsion and curvature are Lorentz tensors with respect to their tangent space indices as is easily shown. The contortion $`K_m^{ab}`$ is also a Lorentz tensor and is related to the torsion through $`K_{lm}^i=\frac{1}{2}(T_{lm}^i+T_{ml}^iT_{lm}^i)`$, with $`K_{lm}^i=e_a^ie_{lb}K_m^{ab}`$ and analogously for $`T_{lm}^i`$. The inverse relation is $`T_{lm}^i=2K_{[lm]}^i`$.
All quantities constructed from the torsion-free connection $`\widehat{\mathrm{\Gamma }}_m^{ab}`$ or $`\widehat{\mathrm{\Gamma }}_{lm}^i`$ will be denoted with a hat, for instance $`\widehat{R}_{km}^{il}=e_a^ie_b^l\widehat{R}_{km}^{ab}`$ is the usual Riemann curvature tensor.
We are now ready to compare the field equations arising in first and second order formalism and to establish their equivalence.
Let $`=_0+_m`$ be the Lagrangian density in the first order formalism and $`\stackrel{~}{}=\stackrel{~}{}_0+\stackrel{~}{}_m`$ the same Lagrangian in the second order formalism. The indices $`0`$ and $`m`$ refer to the gravitational (or free) part and to the matter part respectively. Thus, we can write
$$(e_m^a,\mathrm{\Gamma }_m^{ab})=\stackrel{~}{}(e_m^a,K_m^{ab}).$$
(10)
In concrete cases, the Lagrangians may differ by a total divergence, which does not influence the field equations and will therefore not be written out explicitely. For instance, in Ref. , we considered the teleparallel Lagrangian
$$e(R\frac{1}{4}T^{ikl}T_{ikl}\frac{1}{2}T^{ikl}T_{lki}+\frac{1}{2}T_{ik}^kT_m^{mi})=e\widehat{R},$$
(11)
where the left hand side is the first order form, while the right hand side is the usual Einstein-Hilbert Lagrangian expressed in terms of the tetrad only. In (11), a divergence term has been suppressed. (In other words, the left and right hand sides are not really equal, but can be rendered equal by adding a physically irrelevant divergence.)
The field equations in the conventional, first order formalism are well known,
$$\frac{\delta }{\delta \mathrm{\Gamma }_m^{ab}}=0\text{and}\frac{\delta }{\delta e_a^m}=0.$$
(12)
As usual, they will be referred to as Cartan and Einstein equations, respectively. Both formalisms are related by equation (7),
$$\mathrm{\Gamma }_m^{ab}=\widehat{\mathrm{\Gamma }}_m^{ab}(e_m^a)+K_m^{ab}.$$
For the Cartan equation in the second order formalism, we find
$$\frac{\delta \stackrel{~}{}}{\delta K_m^{ab}}=\frac{\delta }{\delta \mathrm{\Gamma }_l^{cd}}\frac{\delta \mathrm{\Gamma }_l^{cd}}{\delta K_m^{ab}}=\frac{\delta }{\delta \mathrm{\Gamma }_m^{ab}}=0.$$
(13)
Therefore, the Cartan equation has exactly the same form in both approaches. As to the Einstein equation, we easily find
$$\frac{\delta \stackrel{~}{}}{\delta e_a^m}=\frac{\delta }{\delta e_a^m}+\frac{\delta }{\delta \mathrm{\Gamma }_l^{cd}}\frac{\delta \mathrm{\Gamma }_l^{cd}}{\delta e_a^m}=\frac{\delta }{\delta e_a^m}+\frac{\delta }{\delta \mathrm{\Gamma }_l^{cd}}\frac{\delta \widehat{\mathrm{\Gamma }}_l^{cd}}{\delta e_a^m}=0.$$
(14)
This differs by the second term from the corresponding equation in (12). This can be seen explicitely by carrying out the variation of this term using the explicit form of $`\widehat{\mathrm{\Gamma }}_l^{cd}`$ from (8) and of the Christoffel symbols $`\widehat{\mathrm{\Gamma }}_{ik}^l`$.
However, we can instead use immediately the Cartan equation $`\delta /\delta \mathrm{\Gamma }_l^{cd}=0`$ (which we have shown to be equivalent to the Cartan equation in the second order formalism) to eliminate the second term in (14), and therefore we have the on-shell relation $`\delta \stackrel{~}{}/\delta e_a^m=\delta /\delta e_a^m`$.
To summarize, we have the following relations
$$\frac{\delta \stackrel{~}{}}{\delta e_a^m}=\frac{\delta }{\delta e_a^m}\text{and}\frac{\delta \stackrel{~}{}}{\delta K_m^{ab}}=\frac{\delta }{\delta \mathrm{\Gamma }_m^{ab}},$$
(15)
which shows clearly the equivalence of both sets of field equations. The above derivation also explains why some Lagrangians are allowed only in the first order formalism while others can only be considered in the second order formalism. The reason is that the terms resulting from the higher order derivatives (when the wrong formalism is applied to a certain Lagrangian), can actually be eliminated using the second field equation, thereby bringing the field equations into a form that would have resulted immediately if the right formalism had been used right from the start. (By wrong and right, we simply mean the formalism where we have higher order derivatives, or we have not, respectively. Both formalisms are correct, since they are equivalent.)
Note that in order to show the equivalence of both sets of field equations, we have to start with the full Lagrangian. We can of course not conclude that for given stress-energy tensor and spin density, the left hand sides of the field equations will have the same form in both procedures. The source tensors have to be modified at the same time.
Let us define the conventional canonical stress-energy and spin density tensors in the first order formalism
$$T_i^a=\frac{1}{2e}\frac{\delta _m}{\delta e_a^i}\text{and}\sigma _{ab}^m=\frac{1}{e}\frac{\delta _m}{\delta \mathrm{\Gamma }_m^{ab}},$$
(16)
as well as their second order analogues
$$\stackrel{~}{T}_i^a=\frac{1}{2e}\frac{\delta \stackrel{~}{}_m}{\delta e_a^i}\text{and}\stackrel{~}{\sigma }_{ab}^m=\frac{1}{e}\frac{\delta \stackrel{~}{}_m}{\delta K_m^{ab}}.$$
(17)
Just as before (see (13) and (14)), one shows that
$$\stackrel{~}{\sigma }_{ab}^m=\sigma _{ab}^m\text{and}\stackrel{~}{T}_i^a=T_i^a+\frac{1}{2}\sigma _{cd}^m\frac{\delta \widehat{\mathrm{\Gamma }}_m^{cd}}{\delta e_a^i}.$$
(18)
Now, use the relation (8), $`e_{i,k}^a+\widehat{\mathrm{\Gamma }}_{bk}^ae_i^b=e_l^a\widehat{\mathrm{\Gamma }}_{ik}^l`$, and the explicit form of the Christoffel symbols $`\widehat{\mathrm{\Gamma }}_{ik}^l`$ to derive
$$\widehat{\mathrm{\Gamma }}_m^{cd}=\frac{1}{2}\left[e^{ci}e_{i,m}^de^{di}e_{i,m}^c+e^{di}e_{m,i}^ce^{ci}e_{m,i}^d+e^{di}e^{ck}e_{bm}e_{k,i}^be^{ci}e^{dk}e_{bm}e_{k,i}^b\right].$$
We can now evaluate $`\stackrel{~}{T}_i^a`$ in (18). The result is
$$\stackrel{~}{T}^{lm}=T^{lm}+\frac{1}{2}(\sigma ^{lmk}\sigma ^{kml}\sigma ^{klm})_{;k},$$
(19)
where ; denotes the usual covariant derivative with the Christoffel connection.
Note that (19) looks basically like the Belinfante-Rosenfeldt relation between the symmetric Hilbert (or metric) stress-energy tensor $`\stackrel{~}{T}^{lm}`$ and the canonical (or Noether) tensor $`T^{lm}`$ known from purely Riemannian theory, see Refs. and , especially the discussion on relocalizations in Ref. ). We will see that $`\stackrel{~}{T}^{lm}`$ is not yet the symmetric tensor, but it is a tensor that reduces to the symmetric Hilbert tensor for vanishing torsion.
In order to prevent any kind of misunderstandings, let us emphasize once again that we do not modify the underlying gauge structure of the theory. Especially, the relation (10) makes clear the we do not change the coupling prescriptions in the matter Lagrangians. Thus, for instance, the Dirac particle still couples minimally to the gauge connection $`\mathrm{\Gamma }_i^{ab}`$ (which, from a mathematical point of view, is the fundamental field variable). The only difference is that we express this connection in terms of $`e_i^a`$ and $`K_i^{ab}`$. Any change in the coupling prescription would lead to a fundamentally different theory, especially concerning the gauge symmetry.
In the next section, we will derive the conservation equations for spin density and stress energy tensor in both approaches.
## 3 The Noether identities
### 3.1 First order formalism
Let us briefly review the derivation of the Noether identities related to local, tangent space Lorentz transformations and to general coordinate transformations. Under an infinitesimal Lorentz transformation (9), with $`\mathrm{\Lambda }_b^a=\delta _b^a+\epsilon _b^a`$ ($`\epsilon ^{ab}=\epsilon ^{ba}`$) the independent fields $`e_a^m`$ and $`\mathrm{\Gamma }_m^{ab}`$ undergo the following change:
$$\delta \mathrm{\Gamma }_m^{ab}=\epsilon _{,m}^{ab}\mathrm{\Gamma }_{cm}^a\epsilon ^{cb}\mathrm{\Gamma }_{cm}^b\epsilon ^{ac}\text{and}\delta e_a^m=\epsilon _a^ce_c^m.$$
(20)
The first equation can be written in the short form $`\delta \mathrm{\Gamma }_m^{ab}=\mathrm{D}_m\epsilon ^{ab}`$. The change in the matter Lagrangian therefore reads
$$\delta _m=\frac{\delta _m}{\delta e_a^m}\delta e_a^m+\frac{\delta _m}{\delta \mathrm{\Gamma }_m^{ab}}\delta \mathrm{\Gamma }_m^{ab}=e(2T^{[ac]}+\mathrm{D}_m\sigma ^{acm})\epsilon _{ac},$$
(21)
where we have omitted a total divergence. The covariant derivative operator $`\mathrm{D}_m`$ is defined to act with $`\mathrm{\Gamma }_m^{ab}`$ on tangent space indices and with $`\widehat{\mathrm{\Gamma }}_{ki}^l`$ (torsion free) on spacetime indices. The requirement of Lorentz gauge invariance therefore leads to
$$2T^{[ac]}+\mathrm{D}_m\sigma ^{acm}=0.$$
(22)
Slightly more complicated is the case of the coordinate invariance. Under an infinitesimal coordinate transformation,
$$\stackrel{~}{x}^i=x^i+\xi ^i,$$
(23)
the fields $`\mathrm{\Gamma }_m^{ab}`$ and $`e_m^a`$ transform as spacetime covectors (or one-forms), while the inverse tetrad $`e_a^m`$ is a contravariant vector with respect to the index $`m`$. Thus,
$$\stackrel{~}{e}_a^m(\stackrel{~}{x})=e_a^m(x)+\xi _{,k}^me_a^k\text{and}\stackrel{~}{\mathrm{\Gamma }}_m^{ab}(\stackrel{~}{x})=\mathrm{\Gamma }_m^{ab}(x)\xi _{,m}^k\mathrm{\Gamma }_k^{ab}.$$
Since we are interested in the change of the Lagrangian under an active transformation, we have to evaluate the change of the fields at the same point $`x`$, and thus have to express the transformed fields in the old coordinates, i.e.,
$$\stackrel{~}{e}_a^m(x)=\stackrel{~}{e}_a^m(\stackrel{~}{x})\stackrel{~}{e}_{a,k}^m(\stackrel{~}{x})\xi ^k,\stackrel{~}{\mathrm{\Gamma }}_m^{ab}(x)=\stackrel{~}{\mathrm{\Gamma }}_m^{ab}(\stackrel{~}{x})\xi ^k\stackrel{~}{\mathrm{\Gamma }}_{m,k}^{ab}(\stackrel{~}{x}).$$
In the $`\xi `$-terms, we can replace $`\stackrel{~}{e}_a^m(\stackrel{~}{x})`$ by $`e_a^m(x)`$ and $`\stackrel{~}{\mathrm{\Gamma }}_m^{ab}(\stackrel{~}{x})`$ by $`\mathrm{\Gamma }_m^{ab}(x)`$, since the difference will be of order $`\xi ^2`$. Finally, we find
$`\delta e_a^m`$ $`=`$ $`\stackrel{~}{e}_a^m(x)e_a^m(x)=\xi _{,k}^me_a^k\xi ^ke_{a,k}^m,`$ (24)
$`\delta \mathrm{\Gamma }_m^{ab}`$ $`=`$ $`\stackrel{~}{\mathrm{\Gamma }}_m^{ab}(x)\mathrm{\Gamma }_m^{ab}(x)=\xi _{,m}^k\mathrm{\Gamma }_k^{ab}\xi ^k\mathrm{\Gamma }_{m,k}^{ab}.`$ (25)
Let us note at this point that, if we define the parameters $`\epsilon ^a=e_m^a\xi ^m`$ and $`\epsilon ^{ab}=\xi ^k\mathrm{\Gamma }_k^{ab}`$, the above transformations can be written in the form
$`\delta e_m^a`$ $`=`$ $`\epsilon _c^ae_m^c\mathrm{D}_m\epsilon ^a\epsilon ^lT_{lm}^a`$
$`\delta \mathrm{\Gamma }_m^{ab}`$ $`=`$ $`\mathrm{D}_m\epsilon ^{ab}\epsilon ^lR_{lm}^{ab}.`$
Since we are also free to perform an additional, independent, Lorentz transformation (20), we may again consider $`\epsilon ^a`$ and $`\epsilon ^{ab}`$ as independent of each other. In this form, the transformations are sometimes referred to as generalized Poincaré gauge transformations, with $`\epsilon ^a`$ generating the translations and $`\epsilon ^{ab}`$ the Lorentz rotations . We will remain in the usual interpretation of active coordinate transformations and return to the form (24)-(25). The change of the Lagrangian reads
$`\delta _m`$ $`=`$ $`{\displaystyle \frac{\delta L_m}{\delta e_a^m}}\delta e_a^m+{\displaystyle \frac{\delta _m}{\delta \mathrm{\Gamma }_m^{ab}}}\delta \mathrm{\Gamma }_m^{ab}`$
$`=`$ $`2(eT_m^ae_a^k)_{,k}\xi ^m2eT_m^ae_{a,k}^m\xi ^k+(e\sigma _{ab}^m\mathrm{\Gamma }_k^{ab})_{,m}\xi ^ke\sigma _{ab}^m\mathrm{\Gamma }_{m,k}^{ab}\xi ^k.`$
Requiring $`\delta _m=0`$ and regrouping carefully the terms, we finally get
$$(\mathrm{D}_mT_b^m)e_k^b+T_b^mT_{mk}^b=\frac{1}{2}R_{mk}^{ab}\sigma _{ab}^m+\frac{1}{2}\mathrm{\Gamma }_k^{ab}(\mathrm{D}_m\sigma _{ab}^m+2T_{[ab]}).$$
(26)
The last term vanishes with (22) and the first two terms can be rewritten in a different form, using the relation $`T_{lm}^i=2K_{[lm]}^i`$ in the second term, and the explicit form of the covariant derivative $`\mathrm{D}_m`$ in the first term, to get
$$T_{k;m}^mK_k^{im}T_{mi}=\frac{1}{2}R_{mk}^{ab}\sigma _{ab}^m.$$
(27)
Recall that $`K_k^{im}`$ is antisymmetric in $`im`$, so that classical matter (with $`T_{[mi]}=0`$ and $`\sigma _{ab}^m=0`$) will satisfy the general relativistic conservation law $`T_{;m}^{km}=0`$.
### 3.2 Second order formalism
In the second order formalism, we have to evaluate the changes of the fields $`e_a^m`$ and $`K_m^{ab}`$ under Lorentz and coordinate transformations. The Lorentz transformation now reads
$$\delta e_a^m=\epsilon _a^ce_c^m,\delta K_m^{ab}=\epsilon _c^aK_m^{cb}+\epsilon _c^bK_m^{ac},$$
(28)
which shows clearly that $`K_m^{ab}`$ transforms as Lorentz tensor. Evaluating
$$\delta \stackrel{~}{}_m=\frac{\delta \stackrel{~}{}_m}{\delta e_a^m}\delta e_a^m+\frac{\delta \stackrel{~}{_m}}{\delta K_m^{ab}}\delta K_m^{ab},$$
and requiring $`\delta \stackrel{~}{}_m=0`$ we easily derive the following relation
$$0=2\stackrel{~}{T}^{[ac]}K_m^{bc}\sigma _b^{am}+K_m^{ba}\sigma _b^{cm}.$$
(29)
This does not look like a conservation equation for the spin density. We see that in the case where the torsion vanishes, the stress-energy tensor becomes symmetric. Therefore, we can see $`\stackrel{~}{T}^{im}`$ as a generalization of the symmetric Hilbert tensor. Let us note at this point that if we had considered as independent the fields $`K_l^{ik}`$ and $`e_m^a`$ (instead of $`K_l^{ab}`$ and $`e_m^a`$), we would have found a stress-energy tensor that is always symmetric. This, however, does not present any advantages (the variation of the gravitational Lagrangian would in most cases become more complicated) and so we preferred to stay a step closer to the conventional gauge approach, considering $`K_m^{ab}`$ with its natural tangent space indices. The advantage of this approach is the linear relation between $`\mathrm{\Gamma }_m^{ab}`$ and $`K_m^{ab}`$ and the resulting complete equivalence of the Cartan field equation in first and second order formalism. The alternative approach, with $`K_{lm}^i`$ and $`e_m^a`$, has been used in the framework of Einstein-Cartan-Dirac theory in the past . It leads of course to a symmetric Einstein equation.
Under coordinate transformations, the fields transform exactly as in the previous section (because just as $`\mathrm{\Gamma }_m^{ab}`$, $`K_m^{ab}`$ too is a spacetime covector), i.e.,
$`\delta e_a^m`$ $`=`$ $`\stackrel{~}{e}_a^m(x)e_a^m(x)=\xi _{,k}^me_a^k\xi ^ke_{a,k}^m,`$ (30)
$`\delta K_m^{ab}`$ $`=`$ $`\stackrel{~}{K}_m^{ab}(x)K_m^{ab}(x)=\xi _{,m}^kK_k^{ab}\xi ^kK_{m,k}^{ab}.`$ (31)
In the same way as in the previous section, we derive the conservation equation in the form
$`(\mathrm{D}_m\stackrel{~}{T}_a^m)e_k^a+\stackrel{~}{T}_a^mT_{mk}^a{\displaystyle \frac{1}{2}}\sigma _{li;m}^mK_k^{li}`$
$`={\displaystyle \frac{1}{2}}[K_{k,m}^{ab}K_{m,k}^{ab}+\widehat{\mathrm{\Gamma }}_{cm}^aK_k^{cb}+\widehat{\mathrm{\Gamma }}_{cm}^bK_k^{ac}`$
$`\widehat{\mathrm{\Gamma }}_{ck}^aK_m^{cb}\widehat{\mathrm{\Gamma }}_{ck}^bK_m^{ac}+K_{cm}^aK_k^{cb}+K_{cm}^bK_k^{ac}]\sigma _{ab}^m.`$
We used (29) to put the equation into this form. A shorter way of writing this equation is the following
$$(\mathrm{D}_m\stackrel{~}{T}_a^m)e_k^a+\stackrel{~}{T}_a^mT_{mk}^a\frac{1}{2}\sigma _{li;m}^mK_k^{li}=\frac{1}{2}(R_{mk}^{ab}\widehat{R}_{mk}^{ab})\sigma _{ab}^m.$$
(32)
### 3.3 Discussion
Comparing (22) and (29), we can derive the following relation between the antisymmetric parts of the stress-energy tensors:
$$2(\stackrel{~}{T}^{[lm]}T^{[lm]})=\sigma _{;k}^{lmk}.$$
(33)
This relation is easily verified using the explicit form of $`\stackrel{~}{T}^{lm}`$ in terms of $`T^{lm}`$, Eq. (19).
Also, from (26) and (32), we get a relation between the divergences of the stress-energy tensors of the form
$$(T_k^m\stackrel{~}{T}_k^m)_{;m}=\frac{1}{2}\widehat{R}_{mk}^{ab}\sigma _{ab}^m.$$
(34)
We used (33) to bring the relation into this form. Again, this relation can be derived directly from (19), using the symmetry properties of $`\sigma ^{lmk}`$ and of the Riemann tensor $`\widehat{R}_{lkm}^i`$.
Note that both sets of equations, (22,26) and (29,32), are valid, it is not the one or the other set. (The fact that the Lagrangian is invariant under Lorentz gauge and under general coordinate transformations does not depend on the choice of the independent field variables.) They are simply different equations for different quantities. The important thing is to use the right quantities as right hand side of the gravitational field equations (this depends on the choice of our formalism, first or second order).
Another point concerns the physical interpretation of $`T^{ik}`$ and $`\stackrel{~}{T}^{ik}`$. If we take the point of view that the measurable quantities are not the densities, but rather the integrated quantities, namely the momentum vector and the spin tensor , then it is easy to see that $`T^{ik}`$ and $`\stackrel{~}{T}^{ik}`$ lead, in general, to different definitions of the momentum. Apart from this ambiguity, there is the problem that even from a single stress-energy tensor, one can define many inequivalent momentum vectors, integrating $`T_i^a`$, $`T_a^i`$, $`T^{ik}`$ etc. (see Ref. for an interesting discussion on this point). Usually, one would like to define the momentum in a way that is conserved. This is of course not possible in our case, since whatever momentum vector one defines, it will always be subject to gravitational forces coming from curvature and torsion, and therefore it cannot be conserved. (The other way around, whatever momentum vector one uses, it will always be conserved in the limit of vanishing curvature and torsion.) As it seems, it is rather a matter of convention how one wishes to define spin and momentum, under the only restriction that for vanishing fields, they should be conserved.
In view of the rather unusual equation (29), which does not look like a conservation equation, one might be tempted to reject the set of current densities $`\stackrel{~}{T}_m^a,\sigma _{ab}^m`$ and to claim that the canonical currents $`T_m^a,\sigma _{ab}^m`$ as they arise in the conventional, first order formalism, should be considered as the true, physical currents. However, let us recall the fact that, for instance, the canonical stress-energy tensor $`T_m^a`$ of the Dirac Lagrangian contains non-axial torsion parts, which, as is well known, do not couple to the Dirac particle in the first place, and are neither contained in the Dirac equation, nor in the Lagrangian (where they cancel out with the hermitian conjugated term). We have brought this to attention in Ref. , to show that one has to be very careful on how to extract physical quantities from $`T_m^a`$. Thus, neither can the conventional, first order canonical tensors, be used without precautions to derive physically valid results.
In order to gain at least some insight into the contents of $`T^{ik}`$ and $`\stackrel{~}{T}^{ik}`$, let us take a look at the flat limit of our equations, i.e., assume a spacetime with $`g_{ik}=\eta _{ik}`$. In this case, we find from (19)
$$\stackrel{~}{T}^{lm}=T^{lm}+\frac{1}{2}(\sigma ^{lmk}\sigma ^{kml}\sigma ^{klm})_{,k}.$$
(35)
If we define the momentum vector as we do in special relativity,
$$P_k=T_k^idS_i,$$
(36)
then, for the alternative stress-energy tensor, we have
$$\stackrel{~}{P}_k=\stackrel{~}{T}_k^idS_i=P_k+\frac{1}{2}(\sigma _m^{il}\sigma _m^{li}\sigma _m^{li})_{,l}dS_i.$$
(37)
Clearly, the last term can be converted into a 4d volume integral over the divergence of the integrand, which vanishes however, due to the antisymmetry in $`il`$ of the expression in brackets. Thus, we have
$$\stackrel{~}{P}_k=P_k.$$
(38)
Of course, we also have (see (34) or (35))
$$\stackrel{~}{T}_{k,m}^m=T_{k,m}^m.$$
(39)
Usually, the relations (38) and (39) are said to show the equivalence of the two stress-energy tensors. (Same charges, same conservation law.) More generally, a change of the stress-energy tensor and the spin in the form
$`T^{lm}`$ $``$ $`T^{lm}+X_{,k}^{mlk},`$ (40)
$`\sigma _{km}^i`$ $``$ $`\sigma _{km}^i(X_{km}^iX_{mk}^i)`$ (41)
with $`X^{mkl}=X^{mlk}`$ is called relocalization of stress-energy and spin and is said not to affect physical quantities.
It is easily seen that (35) is a relocalisation of $`T^{lm}`$, without, however, the corresponding transformation of the spin density. Nevertheless, (26) and (29) coincide in flat space. The remaining relation is
$`\stackrel{~}{T}_{k,m}^m`$ $`=`$ $`T_{k,m}^m`$
$`=`$ $`{\displaystyle \frac{1}{2}}K_k^{im}(\sigma _{im,l}^lK_{il}^j\sigma _{jm}^lK_{ml}^j\sigma _{ji}^l)`$
$`+{\displaystyle \frac{1}{2}}(K_{k,m}^{ab}K_{m,k}^{ab}+K_{cm}^aK_k^{cb}+K_{cm}^aK_k^{ac})\sigma _{ab}^k.`$
Unfortunately, we cannot express the equation for the spin precession (22) or (29) in a form independent of the stress-energy tensor. (Using (35) in (33) leads to an identity.) However, since the spin density is the same in both approaches, it will be subject to the same equations of motion, expressed in one or the other way.
We therefore conclude that in the flat limit, $`g_{ik}=\eta _{ik}`$, the equations of motions for the current densities obtained in both pictures are identical as far as their physical contents is concerned.
In curved spacetime, this is not the case anymore. Especially, we get two different equations for the stress-energy tensor and we have to decide, based on physical arguments, how to define define the momentum vector.
Let us once again refer the reader to the extended discussion on canonical and metric stress-energy tensors and their relation via relocalizations in the general framework of metric affine theories in Hehl et al., .
## 4 Model construction
Comparing equations (27) and (32), we see that the following tensor
$`F_{mk}^{ab}`$ $`=`$ $`{\displaystyle \frac{1}{2}}(R_{mk}^{ab}\widehat{R}_{mk}^{ab})`$ (42)
$`=`$ $`K_{k,m}^{ab}K_{m,k}^{ab}+\widehat{\mathrm{\Gamma }}_{cm}^aK_k^{cb}+\widehat{\mathrm{\Gamma }}_{cm}^bK_k^{ac}`$
$`\widehat{\mathrm{\Gamma }}_{ck}^aK_m^{cb}\widehat{\mathrm{\Gamma }}_{ck}^bK_m^{ac}+K_{cm}^aK_k^{cb}+K_{cm}^bK_k^{ac}`$
plays, in the context of the second order formalism, a similar role as the curvature tensor $`R_{lm}^{ab}`$ in the conventional approach. Essentially, $`F_{mk}^{ab}`$ has the structure of a Lorentz Yang-Mills tensor, but with $`\mathrm{\Gamma }_m^{ab}`$ replaced by $`K_m^{ab}`$. The terms containing the Christoffel connection $`\widehat{\mathrm{\Gamma }}_m^{ab}`$ are needed to get the correct tensor behavior under a Lorentz transformation.
Apart from $`F_{mk}^{ab}`$, which we will use to produce the dynamical, Yang-Mills like, behavior of the torsion field, we will need a term that guarantees the correct general relativity limit of the theory. In the second order formalism, this is simply the Einstein-Hilbert Lagrangian, expressed in terms of the tetrad field, $`\widehat{R}`$. For reasons that will become clear later, we also include a mass term $`K_m^{ab}K_{ab}^m`$ into our Lagrangian.
Thus, we start with
$$_0=e(a\widehat{R}bF_{mk}^{ab}F_{ab}^{mk}cK_m^{ab}K_{ab}^m).$$
(43)
This Lagrangian contains no second order derivatives of the independent fields $`e_i^a`$ and $`K_i^{ab}`$, apart from the well known higher order terms contained in $`\widehat{R}`$, which can, however, be eliminated by subtracting a surface term. On the other hand, if $`_0`$ were to be treated in the first order formalism, there would arise higher order derivatives in the tetrad fields, coming from the term $`\widehat{R}_{lm}^{ab}\widehat{R}_{ab}^{lm}`$ contained in the second term of (43).
Including the matter Lagrangian $`_m`$, we are let to the following set of field equations
$`4b\mathrm{D}_kF_{ab}^{mk}+2cK_{ab}^m`$ $`=`$ $`\sigma _{ab}^m`$ (44)
$`a\widehat{G}_{ki}`$ $`=`$ $`\tau _{ki}^{(1)}+\tau _{ki}^{(2)}+\tau _{ki}^{(3)}+\stackrel{~}{T}_{ki},`$ (45)
where
$`\tau _{ki}^{(1)}`$ $`=`$ $`2b(F_{mi}^{cb}F_{cbk}^m{\displaystyle \frac{1}{4}}g_{ki}F_{ml}^{cb}F_{cb}^{ml})`$ (46)
$`\tau _{ki}^{(2)}`$ $`=`$ $`c(K_k^{cb}K_{cbi}{\displaystyle \frac{1}{2}}g_{ki}K_l^{cb}K_{cb}^l)`$ (47)
$`\tau _{ki}^{(3)}`$ $`=`$ $`2b[(F_{imk}^nK_n^{lm})_{;l}+(F_{in}^{lm}K_{km}^n)_{;l}(F_{km}^{ln}K_{in}^m)_{;l}`$ (48)
$`(F_{mkn}^lK_i^{mn})_{;l}+(F_{kmi}^nK_n^{lm})_{;l}(F_{mi}^{ln}K_{kn}^m)_{;l}].`$
We recognize in (44) a Yang-Mills type equation for the field $`K_m^{ab}`$, endowed with a mass term proportional to $`c`$. (Recall that we have defined $`\mathrm{D}`$ to act with the full Lorentz connection $`\mathrm{\Gamma }_m^{ab}`$ on tangent space indices, and with the Christoffel symbol on spacetime indices.) The corresponding symmetric, traceless stress-energy tensor is found in (46) and the stress-energy contribution from the mass term is given by (47).
The only non-symmetric contributions in (45) are contained in the unusual term $`\tau _{ki}^{(3)}`$, a term that results from the $`\widehat{\mathrm{\Gamma }}_m^{ab}`$ couplings in $`F_{mk}^{ab}`$. In this term, the Lorentz gauge structure of the theory is not obvious, since the spacetime and Lorentz indices of $`F_{mk}^{ab}`$ and $`K_m^{ab}`$ are completely mixed up. This is also the reason why we chose to write down Eq. (45) immediately in terms of spacetime tensors, as opposed to the original form $`\widehat{G}_i^a=\widehat{T}_i^a+\mathrm{}`$ that results from the variation with respect to $`e_a^i`$. In contrast to (44), which respects clearly the underlying gauge structure, nothing useful is revealed by keeping (45) in its mixed form, and it is convenient to pull down all quantities to physical spacetime such that one has to deal only with one kind of indices.
Concretely, we find for the antisymmetric part of (45)
$$\stackrel{~}{T}^{[ik]}=2b(F_n^{ilm}K_m^{kn}F_n^{klm}K_m^{in})_{;l}.$$
(49)
The calculations are rather lengthy, but one can directly verify that (49), together with (44), leads to the Noether identity (29).
In the classical, spinless limit, it is immediately clear that (44) leads to the groundstate solution $`K_m^{ab}=0`$, while (45) reduces to
$$\widehat{G}_{ik}=T_{ik},$$
which are the field equations of general relativity. (We took the opportunity to fix the parameter $`a`$ to $`a=1`$ to be consistent with the conventional choice of units.) Also, in view of (19), there is no need, in the classical limit, to distinguish between $`T^{ik}`$ and $`\stackrel{~}{T}^{ik}`$.
Thus, the theory described by (43) allows for a complete general relativity limit in the absence of spinning matter, while a source with intrinsic spin, e.g., an elementary particle or a spin polarized neutron star, will give rise to dynamical torsion fields via (44). This cannot be achieved within the conventional first order formalism if one is not willing to allow for higher order derivatives, and therefore, no such theory has yet been presented. In order to avoid misunderstandings, let us emphasize that this does not mean that there are no consistent Poincaré gauge theories based on the first order formalism without higher derivatives. What it means is that such theories do not have a full general relativity limit. Therefore, constraints on the free parameters of the theory will arise from the classical experiments, based on planetary motion and light propagation. This is rather disappointing, since it essentially means generalizing GR, but holding the changes small enough to ensure that they do not affect classical physics in a sensible way. On the other hand, in our model, classical measurements are not affected at all, independently of the coupling constants $`b,c`$, and therefore, the post-general relativistic effects will only arise if spinning matter is introduced. Thus, the generalization affects only spinning matter, in view of which the theory was generalized in the first place. In short: Additional fields (torsion) for additional degrees of freedom (spin).
Finally, let us make a remark on the mass term in (43). The reason for its introduction is easily seen in context of the classical limit. Suppose $`c=0`$ in (44)-(45). Then, every GR solution, satisfying in addition $`F_{mk}^{ab}=0`$, will be solution to the field equations with $`\sigma _{ab}^m=0`$. However, from $`F_{mk}^{ab}=0`$, we can not conclude to the unique solution $`K_m^{ab}=0`$ (we will see an explicit example in the next section). However, changing $`K_m^{ab}`$ by a term that leaves $`F_{mk}^{ab}`$ unaffected will have physical consequences (e.g., concerning the spin precession of a test particle) and therefore, this additional freedom cannot be accepted. This is not confined to the classical limit. From (43), it is obvious that, for $`c=0`$, the field equations will only determine $`F_{mk}^{ab}`$, but not $`K_m^{ab}`$, allowing therefore, in general, for a whole class of solutions for $`K_m^{ab}`$. This will become very clear in the context of the solutions we will analyze in the next section. The introduction of the mass term obviously resolves this problem.
## 5 Yasskin type reduction
Since the classical limit of our theory coincides completely with general relativity, the only way to test the new features related to the presence of torsion fields is to look at field configurations generated by a source with intrinsic spin. One way is for instance to study interactions between elementary particles. On the astrophysical level, one can consider massive, spin polarized bodies, especially neutron stars, and analyze, e.g., the evolution of a system of two such bodies or, which is simpler from the computational point of view, study the motion and spin precession of a test body (a smaller neutron star or elementary particles) in the vicinity of a central source.
In this section, we try to get solutions for the fields at the exterior of such a macroscopic spinning body. As discussed extensively in Ref. , such a matter distribution (located at the origin) is described by a spin density of the form
$$\sigma _{ab}^m=\sigma _{ab}u^m\rho (x),$$
(50)
where $`\rho (x)`$ is a suitably normalized density function and $`\sigma _{ab}`$ is the integrated spin tensor. This spin density also coincides with that describing the generalized Weyssenhoff fluid in Riemann-Cartan spacetimes . Then, following Yasskin , we look for solutions that have essentially the same structure as the current density, but in view of our rather special Lagrangian, we make the ansatz not for the Lorentz gauge potential $`\mathrm{\Gamma }_m^{ab}`$, but for the contortion $`K_m^{ab}`$, i.e., we set
$$K_m^{ab}=\sigma ^{ab}A_m,$$
(51)
where the field $`A_m`$ will have to be determined by the field equations. In Ref. , with a similar ansatz, the Einstein-Yang-Mills field equations were reduced essentially to an Einstein-Maxwell system of equations, under the assumption that the gauge charges (in our case, $`\sigma _{ab}`$) are constant. In order to preserve Lorentz invariance, one cannot require $`\sigma _{ab}`$ to be constant, but instead has to use the constraint
$$\widehat{\mathrm{D}}_i\sigma _{ab}=0,$$
(52)
which, with the help of (51), can also be written in the form
$$\mathrm{D}_i\sigma _{ab}=0.$$
(53)
In other words, the contortion fields, which in our second order approach play a role similar to the Yang-Mills potentials, drop out from the covariant derivative of $`\sigma _{ab}`$. This is similar to Ref. , where the gauge covariant derivatives coincide, for the specific ansatz, with the partial derivatives, which makes the requirement that the charges be constant a covariant one.
What does Eq. (52) mean physically? Well, in view of a physical interpretation of the spin tensor, one will have to use further constraints, the so-called spin supplementary condition (SSC), like $`\sigma _{ik}u^i=0`$ or similar. Roughly, there will be a local coordinate system where (52) reduces to $`\sigma _{ik}=const`$ and one might then choose the SSC $`\sigma _{i0}=0`$ and define the spin vector $`\sigma ^\mu =\epsilon ^{\mu \nu \lambda }\sigma _{\nu \lambda }`$ ($`\mu ,\lambda ,\nu =1,2,3`$). In other words, (52) is simply the covariant expression of the fact that the source is described by a constant spin vector. (Constant in the sense of a fixed vector, pointing in a fixed direction, as opposed to a radially orientated, spherically symmetric spin vector field, for instance.) This is just what we expect for a spin polarized body. Note, on the occasion, that $`\sigma _{ab}\sigma ^{ab}=2\stackrel{}{\sigma }^2`$. Therefore, let us define the spin $`s`$ through
$$\sigma _{ab}\sigma ^{ab}=2s^2.$$
(54)
In practice, we will apply (52) only to the contortion field (51), and not to (50), since we will only deal with exterior solutions here. The above arguments are therefore only a simplified version of what should be the result of an analysis of the interior dynamics of the neutron star, which is beyond the scope of this article. For our purposes, the star can equally well be considered to be pointlike (i.e., $`\rho (x)=\delta (x)`$ in (50)), as long as (52) is satisfied.
With the ansatz (50)-(52), the field tensor (42) reduces to
$$F_{ik}^{ab}=\sigma ^{ab}F_{ik},$$
(55)
where $`F_{ik}=A_{k,i}A_{i,k}`$, whereas the field equations (44) and (45) take the simple form
$`4bF_{;k}^{mk}`$ $`=`$ $`2cA^m+\rho (x)u^m`$ (56)
$`\widehat{G}_{ik}`$ $`=`$ $`4bs^2(F_{mi}F_k^m{\displaystyle \frac{1}{4}}g_{ik}F^{lm}F_{lm})`$ (57)
$`2cs^2(A_kA_i{\displaystyle \frac{1}{2}}g_{ik}A_lA^l)+\stackrel{~}{T}_{ik}.`$
Thus, with a slight extension of the method presented in Ref. , we have brought our equations into an Einstein-Proca system.
Several remarks are in order at this point. First, the ansatz (51) is supposed to hold at the exterior of the source. We do not know anything about the field configurations inside the source. Nevertheless, we have included the source terms in (56)-(57) in order to be able to fix eventual parameters of the solutions (like the Schwarzschild mass etc.). You may consider the source terms as boundary conditions, having in mind that the above form of the field equations is valid only at the exterior. Secondly, and not unrelated to the first remark, you may have noticed that equation (57) is symmetric, while one expects an asymmetric stress-energy tensor for a spinning body. This, of course, reflects again the fact that the equations do not hold inside the matter distribution. In order to ensure the continuity of the solutions at the boundary of the source, the stress-energy tensor should have a symmetric limit at that boundary. Recall however that the tensor $`\stackrel{~}{T}_{ik}`$, as we have outlined at the end of section 2, is not the canonical stress-energy tensor, but a tensor that is already quite close to the symmetric Hilbert tensor. More precisely, we have shown that its antisymmetric part is given by (see (29))
$$2\stackrel{~}{T}^{[ac]}=K_m^{bc}\sigma _b^{am}K_m^{ba}\sigma _b^{cm}.$$
Using the spin density (50) and assuming that $`K_m^{ab}`$ tends to the form (51) as we approach the boundary, we see that $`T^{[ac]}`$ vanishes without any restrictions on the stress-energy tensor. Therefore, the symmetry of (57) is not a constraint on the stress-energy tensor, but simply the result of the specific form (51) of the torsion.
On the other hand, if we assume that $`K_m^{ab}`$ has a similar form also at the interior of the source (this is to be understood as the mean value over a macroscopic region, not on an elementary particle level), then we can conclude that the stress-energy tensor is symmetric allover, and that it essentially coincides with the Hilbert tensor of general relativity. This is of course in contrast to conventional, first order, Poincaré gauge theory, where a symmetric stress-energy tensor is only possible for vanishing spin.
For the simplified equations (56)-(57), the reason for the necessity of the mass term (i.e., the $`c`$-terms) is especially clear. For $`c=0`$, we are dealing with an Einstein-Maxwell system. Thus, the torsion will only be determined up to a gauge transformation $`A_mA_m+f_{,m}`$. However, measurable results, like the evolution of spin and momentum of a test particle in the given field configuration , will depend on the full torsion tensor. Otherwise stated, the field equations for $`c=0`$ do not completely determine the physical fields.
Unfortunately, the mass term, necessary on theoretical grounds, is rather disturbing from a computational point of view, since the known solutions of the Einstein-Maxwell system do not easily generalize to the Proca field.
## 6 Discussion and solutions of the field equations
In the last section, we have reduced the field equations of our specific Poincaré gauge theory to a purely Riemannian problem: The description of a massive vector field in curved spacetime. The discussion of the system (56)-(57) is therefore subject to the framework of general relativity, and can be found in the corresponding literature. Nevertheless, for the sake of completeness, we will briefly sketch the main features of this system and refer to the original articles for details.
It is well known that in the cases of practical interest, no exact solution of the Einstein-Proca field equations has been found to date. Especially, no generalization of the spherically symmetric Reissner-Nordstrøm solution, nor of the axially symmetric Kerr-Newman solution are known for the case of a massive vector field. Moreover, it has been established that there is no black hole solution other then for $`A_i=0`$. This is a special case of the no-hair theorem for black holes. Here, we are interested in neutron stars. In view of the no-hair theorem, the vacuum solutions we eventually find for the exterior of the star will appear with a naked singularity at the origin. This, however, does not mean that they are unphysical. As vacuum solutions, they are restricted to the exterior of the star and will have to be completed by a suitable interior solution, thereby removing the singularity. On the other hand, the absence of black hole solutions indicates that there is a rather fundamental difference between the Maxwell and the Proca case. Especially, it seems that there are no solutions that tend to the Einstein-Maxwell black holes (Kerr-Newman or Reissner-Nordstrøm) for $`c0`$ ($`c`$ is the mass parameter of the vector field in (56)-(57)), as might have been expected.
A detailed analysis of the static, spherically symmetric case has been performed and numerical solutions have been derived in Refs. -. The authors of the three articles essentially agree in their conclusions, and there is no need to repeat the analysis here. The main feature of the solutions is that, for large distances from the source, the metric tends to the Schwarzschild solution. This is in agreement with the expected exponential decrease of the massive vector field known from the special relativistic limit. Indeed, for $`g_{ik}=\eta _{ik}`$, starting with the general spherically symmetric, static field $`A_i=(\phi (r),\psi (r)\frac{\stackrel{}{x}}{r})`$ in the comoving system of the source, $`u^i=\delta _0^i`$, one easily shows from (56) that $`\psi (r)=0`$, and that $`\phi (r)`$ is solution to the equation
$$4b\mathrm{\Delta }\phi +2c\phi =\rho ,$$
(58)
with solution
$$\phi \frac{e^{\sqrt{c/2b}r}}{r}.$$
In order to fix the proportionality factor, we have to take a look at equations (50) and (51). In contrast to the usual form of the Maxwell (or Proca) equations, the charge $`\sigma _{ab}`$ has essentially been removed from $`A_i`$ in the ansatz for $`K_m^{ab}`$, i.e., in the flat case, the density $`\rho (r)`$ is supposed to be normalized to $`\rho (r)\mathrm{d}^3x=1`$. It is then easy to see (consider the special case $`c=0`$ and $`\rho (r)=\delta (r)`$) that the complete solution reads
$$\phi (r)=\frac{1}{16\pi b}\frac{e^{\sqrt{c/2b}r}}{r}.$$
(59)
Since the systematic numerical analysis can be found in Refs. -, we will apply here a different approach to the system (56)-(57). Namely, we will proceed in the spirit of the WKB approach of quantum mechanics, which consists in expanding the system in terms of Planck’s constant. In our case, this means expanding in terms of the spin tensor $`\sigma _{ab}`$ or simply in terms of the parameter $`s`$ introduced in (54). To zeroth order (i.e., for $`s=0`$), we simply get the Schwarzschild solution with $`A_i=0`$. We are interested in the first order corrections to this solution. Thus, in (56)-(57), we simply neglect the terms in $`s^2`$. Again, form (57), we find the Schwarzschild solution. (Just as in Einstein-Maxwell theory, where the corrections to the metric are of second order in the electric charge, here, the corrections are of second order in the spin.) At this stage, we are dealing with a Proca field on a Schwarzschild background. Finally, in order to solve equation (56), we expand the metric in terms of the Schwarzschild mass $`m`$ and neglect the term of order $`\phi (r)\frac{m}{r}`$ and higher, leaving us with the solution (59).
Summarizing, at large distances, our solution looks like the Schwarzschild solution with $`A_i=0`$. Approaching the source, the first order corrections are given by the solution (59) for $`A_i=(\phi ,0,0,0)`$. The next order of approximation consists in interaction terms of the gravitational potential $`m/r`$ and the vector field $`A_i`$, leading to corrections in (59). Only at the next order ($`A^2`$) will the corrections to the Schwarzschild metric become apparent.
Note that the Schwarzschild metric is not only a very good approximation because of the exponential decrease of $`\phi `$, but also because of the absolute smallness of spin effects in general. In practice, rotational effects (leading to Kerr type corrections) will by far dominate the (intrinsic) spin effects. Not only will the rotational spin of a neutron star be much larger than its intrinsic spin, but moreover, the rotational effects lead to first order corrections in the metric, in contrast to the second order spin effects. Experimentally, this means that the spin effects cannot be observed by analyzing the geodesics, but rather by measuring directly the torsion field (through its interaction with spinning test bodies, see Ref. ).
Let us emphasize that, although the Schwarzschild metric is subject to higher order corrections, it will retain its spherical symmetry to each order. The same holds for $`A_i`$, but not for $`K_i^{ab}`$, which is given by (51).
Unfortunately, the analysis of Refs. - has not been extended to the axially symmetric case. Since this is the case of practical interest (neutron stars are supposed to be subject to a rotation of very high frequency), we will extend the simple approach of the above considerations also to this case. As before, we thus neglect terms of second order in the spin. This, clearly, leads to the Kerr metric. Then, we also neglect interactions of the gravitational potential with the vector field $`A_i`$.
First, let us recall the Kerr-Newman solution for the massless vector field, i.e., the axially symmetric solution of (56)-(57) for $`c=0`$. Using spherical coordinates $`x^i=(t,r,\vartheta ,\phi )`$, the Kerr-Newman metric reads
$`\mathrm{d}s^2={\displaystyle \frac{\rho ^22mr+q^2}{\rho ^2}}\mathrm{d}t^2{\displaystyle \frac{\rho ^2}{\mathrm{\Delta }}}\mathrm{d}r^2\rho ^2\mathrm{d}\vartheta ^2`$
$`{\displaystyle \frac{(2mrq^2)2a\mathrm{sin}^2\vartheta }{\rho ^2}}\mathrm{d}t\mathrm{d}\vartheta [a^2+r^2+a^2\mathrm{sin}^2\vartheta {\displaystyle \frac{2mrq^2}{\rho ^2}}]\mathrm{d}\phi ^2,`$ (60)
with $`\rho ^2=r^2+a^2\mathrm{cos}^2\vartheta ,\mathrm{\Delta }=r^22mr+q^2+a^2`$. Here, $`m,q,a`$ are constants of integration. As is easily seen taking suitable limits, $`m`$ is the Schwarzschild mass, $`a`$ is related to the angular momentum of the source (see any textbook on general relativity) and the charge parameter $`q`$ is proportional to the intrinsic spin $`s`$ (see below). The corresponding field $`A_i`$ has the form
$$A_i=\frac{q}{s\sqrt{2b}}(\frac{r}{\rho ^2},0,0,\frac{ar\mathrm{sin}^2\vartheta }{a^2\mathrm{cos}^2\vartheta +r^2}).$$
The constant $`q`$ can be evaluated by taking the Coulomb limit, i.e., take $`a=0`$ and compare with Eq. (59) for $`c=0`$. The result is
$$q=\frac{s\sqrt{2}}{16\pi \sqrt{b}},$$
(61)
and the solution takes the form
$$A_i=\frac{1}{16\pi b}(\frac{r}{\rho ^2},0,0,\frac{ar\mathrm{sin}^2\vartheta }{a^2\mathrm{cos}^2\vartheta +r^2}).$$
(62)
This is the exact solution for the massless vector field. A difference to the spherically symmetric case arises if we consider the flat limit of this solution. In the Reissner-Nordstrøm case, the solution $`A_i`$ is simply the Coulomb potential, which is a solution of $`F_{;i}^{ik}=0`$ in the flat case too. At first sight, it seems as if (62) is not a solution on the flat background $`\mathrm{d}s^2=\mathrm{d}t^2\mathrm{d}r^2r^2(\mathrm{d}\vartheta ^2+\mathrm{sin}^2\vartheta \mathrm{d}\phi ^2)`$. This, however, is the result of a false interpretation of the coordinate system. Indeed, the flat limit of (60) is found for $`q0,m0`$, i.e.,
$$\mathrm{d}s^2=\mathrm{d}t^2\frac{r^2+a^2\mathrm{cos}^2\vartheta }{r^2+a^2}\mathrm{d}r^2(r^2+a^2\mathrm{cos}^2\vartheta )\mathrm{d}\vartheta ^2(a^2+r^2)\mathrm{sin}^2\vartheta \mathrm{d}\phi ^2,$$
which is easily shown to be flat ($`\widehat{R}_{klm}^i=0`$) and related to the Minkowski metric $`\eta _{ik}=diag(1,1,1,1)`$ by the coordinate transformation
$`x`$ $`=`$ $`\sqrt{r^2+a^2}\mathrm{sin}\vartheta \mathrm{cos}\phi `$
$`y`$ $`=`$ $`\sqrt{r^2+a^2}\mathrm{cos}\vartheta \mathrm{cos}\phi `$
$`z`$ $`=`$ $`r\mathrm{cos}\vartheta .`$
We now find that (62) is indeed a solution in the flat background expressed in the non-spherical coordinates $`r,\vartheta ,\phi `$. Since it is rather difficult to generalize this to the massive case, we will neglect contributions that are of second order in the rotational momentum $`a`$. To that order, the coordinates used in (62) coincide with conventional spherical coordinates, and indeed the field
$$A_i=\frac{1}{16\pi b}(\frac{1}{r},0,0,\frac{a\mathrm{sin}^2\vartheta }{r}),$$
(63)
which differs from (62) by terms of order $`a^2`$ and higher, is a solution to the Maxwell equations on a flat background (in spherical coordinates). Physically, we see that (63) contains an electric monopole and a magnetic dipole contribution, while (62) contains, as is well known, the complete series of odd electric multipoles (monopole, quadrupole…) and even magnetic multipoles (dipole, octupole…). Thus, replacing (62) by (63), we essentially neglect quadrupole and higher moments. We then make the following ansatz for the massive vector field:
$$A_i=\frac{1}{16\pi b}(\frac{1}{r}f(r),0,0,\frac{a\mathrm{sin}^2\vartheta }{r}g(r)).$$
(64)
Putting this into (56) yields, on a flat background, the unique asymptotically flat solution
$`f(r)`$ $`=`$ $`e^{\sqrt{c/2b}r}`$
$`g(r)`$ $`=`$ $`(1+\sqrt{c/2b}r)e^{\sqrt{c/2b}r},`$ (65)
where the constants of integration have been fixed by requiring the limit (63) for $`c0`$.
Summarizing, the Kerr metric, as before the Schwarzschild metric, is a very good approximation for three reasons: (1) The spin is small compared to the angular momentum of the source. (2) Eventual corrections to the metric are of second order in $`A_i`$ and thus in the spin (in contrast to the angular momentum contributions which are first order). (3) The field $`A_i`$, and thus the metric corrections, decreases exponentially.
The first order approximation of $`A_i`$ is given by (64)-(65), neglecting interactions of the metric with $`A_i`$ (order $`A_im/r`$) and second order terms in the spin and in the rotational momentum.
Possibly, one can find exact solutions of (56) in Schwarzschild or Kerr backgrounds. This would yield the complete analogue of the WKB approximation, neglecting only terms of second and higher order in the spin, but including the spin-gravity couplings $`A_im/r`$. However, as we have argued before, it is not the scope of this article to analyze in detail the equations in the form (56)-(57), because this set of equations is well known from classical general relativity.
Finally, one could also consider the possibility of a very small $`c`$, too small for the exponential decrease to be sensible on astrophysical scales. (In analogy to theories with a small photon mass.) In this case, one can consider the term in $`c`$ as a gauge fixing term but irrelevant else. This means that we can simply take over the solutions from Einstein-Maxwell theory, i.e., the Reissner-Nordstrøm metric together with the Coulomb potential, or the Kerr-Newman metric with the field (62), but in the gauge that is consistent with $`c0`$, which is essentially $`A_{;i}^i=0`$, as can be easily derived taking the divergence of (56). Note that (62) and also the Coulomb solution in its usual form are already in this gauge. Thus, the term in $`c`$ forces us to choose the gauge $`A_{;i}^i=0`$, but the $`c`$-contributions will then be neglected in the field equations.
Does every solution of (56)-(57) tend to a solution of the Einstein-Maxwell equations in the limit $`c=0`$? Our approximate solutions seem to indicate that this is indeed the case. One should, however, be cautious, because there are nevertheless severe differences between the massive and the massless vector fields. One example is the previously mentioned absence of black hole solutions in the massive case. Thus, there might be surprises in the strong field region which have not been revealed by our approximations.
Let us remind that the whole section was based on the specific ansatz (51), which was adapted to the macroscopic spin distribution (50). This does not mean that, even if we retain (50), there are no other solutions. Especially, the existence of black holes cannot be excluded a priori. On the other hand, for a different spin density, for instance the totally antisymmetric tensor of the Dirac particle, a different ansatz will be necessary, leading to entirely different solutions.
## 7 Spin precession in torsion field
Finally, we address the question of the detectability of the torsion fields by analyzing the behavior of test bodies in a given field configuration. We confine ourselves to the concrete solutions of the previous section.
It is well known that the equations for the spin precession and for the momentum evolution in a Riemann-Cartan geometry depend on the nature of the test body in question. Especially, elementary particles with different spin couple in a different manner to the torsion. A complete review, as well as a simple method for the derivation of the equations of motion, has been presented in Ref. , where the references to the original articles can be found.
For simplicity, we consider again the case where the test particle is a macroscopic, spin polarized body. Thus, the system (central source + test body) we are interested in is essentially a neutron star binary, under the restriction, for computational simplicity, that the central star has a much larger mass and can be considered to be at rest.
The deviation from geodesic motion is given by a direct coupling of the test particles’ spin to the curvature tensor. Even for the classical spin (i.e., the angular momentum in the body’s rest frame), which presents a similar coupling to the curvature (Papapetrou equations in general relativity), these corrections have not yet been measured in experiment, and the contributions from the intrinsic spin and the non-Riemannian fields will be even smaller. Therefore, the only experimentally relevant equation is the spin precession equation, which reads
$$\mathrm{D}\sigma _{ik}^{(2)}+\widehat{\mathrm{D}}S_{ik}^{(2)}=P_iu_kP_ku_i,$$
(66)
where $`\mathrm{D}`$ denotes the covariant derivative with respect to proper time using the full connection $`\mathrm{\Gamma }_{lm}^i`$ and $`\widehat{\mathrm{D}}`$ the corresponding derivative with the Christoffel connection, i.e., $`\mathrm{D}\sigma _{ik}^{(2)}=\mathrm{d}\sigma _{ik}^{(2)}/\mathrm{d}\tau \mathrm{\Gamma }_{im}^l\sigma _{lk}^{(2)}u^m\mathrm{\Gamma }_{km}^l\sigma _{il}^{(2)}u^m`$ and $`\widehat{\mathrm{D}}S_{ik}^{(2)}=\mathrm{d}S_{ik}^{(2)}/\mathrm{d}\tau \widehat{\mathrm{\Gamma }}_{im}^lS_{lk}^{(2)}u^m\widehat{\mathrm{\Gamma }}_{km}^lS_{il}^{(2)}u^m`$. We denote by $`\sigma _{ik}^{(2)}`$ the (intrinsic) spin tensor of the test body and by $`S_{ik}^{(2)}`$ its angular momentum in the rest frame (due to rotation). The index $`(2)`$ is attached to avoid confusion with the corresponding quantities of the source. $`P_i`$ is the momentum vector of the test body, which is not necessarily parallel to the velocity $`u^i`$. However, the right hand side of (66) is of higher order and will be neglected in the following. Finally, we have to impose some conditions on the structure of the test body. It seems reasonable, for neutron stars, to make the ansatz (strong spin-rotation coupling)
$$S_{ik}^{(2)}=g\sigma _{ik}^{(2)},$$
(67)
which leaves us with
$$\widehat{\mathrm{D}}\sigma _{ik}^{(2)}+\frac{1}{1+g}K_{im}^l\sigma _{lk}^{(2)}u^m+\frac{1}{1+g}K_{km}^l\sigma _{il}^{(2)}u^m=0,$$
(68)
where we have separated into Riemannian and non-Riemannian contributions. Next, introduce the spin vector $`\sigma _{(2)}^i=\frac{1}{2}\eta ^{iklm}u_k\sigma _{lm}^{(2)}`$ ($`\eta ^{iklm}=|g|^{1/2}\epsilon ^{iklm}`$) and use the spin supplementary condition $`\sigma _{ik}^{(2)}u^k=0`$ to find
$$\widehat{\mathrm{D}}\sigma _{(2)}^i+\frac{1}{1+g}K_{lm}^i\sigma _{(2)}^lu^m=0,$$
(69)
where we have used the approximate validity of the geodesic equation $`\widehat{\mathrm{D}}u^i=0`$ and higher order terms have been omitted . This equation, together with the solutions of the previous section, can be used to evaluate the spin precession and especially to measure the influence of the torsion contributions.
Here, our scope is to illustrate the effects of the dynamical torsion fields. Therefore, we confine ourselves to the simple case of the non-rotating test body ($`S_{ik}^{(2)}=0`$, i.e., $`g=0`$) and use the spherically symmetric solution (no rotation of the source). This is not a very realistic case (as we have argued before, the rotational effects, when dealing with neutron stars, will in general be by far more important than the intrinsic spin effects), but it is a useful example to illustrate directly the post-general relativistic effects due to spin and torsion.
With $`K_l^{ab}`$ from (51), with $`A_m=(\phi ,0,0,0)`$ and introducing the spin vector of the source, $`\sigma _{(1)}^i=\frac{1}{2}\eta ^{iklm}u_k\sigma _{lm}^{(1)}`$, using only the Newtonian order of the Schwarzschild background (compare with Ref. ), we find to lowest order for the relevant spin precession equation
$$\frac{\mathrm{d}\stackrel{}{\sigma }^{(2)}}{\mathrm{d}t}=\frac{3}{2}\frac{m}{r^3}[\stackrel{}{L}\times \stackrel{}{\sigma }^{(2)}]+\phi (r)[\stackrel{}{\sigma }^{(1)}\times \stackrel{}{\sigma }^{(2)}],$$
(70)
where $`\phi (r)`$ is given by Eq. (59). The first term, with the orbital momentum $`\stackrel{}{L}`$, is identical to the precession of the rotational spin in general relativity (as derived from the lowest approximation of the Papapetrou equations). This is a general result of Poincaré gauge theory: In purely Riemannian spacetimes, no distinction can be made between intrinsic and classical (rotational) spin. The torsion, however, couples only to the intrinsic spin. Therefore, the other way around, only a particle with intrinsic spin can distinguish between Riemannian and non-Riemannian geometry.
The above result is easily generalized to the more realistic case of rotating neutron stars. One can also proceed to a more complete treatment considering the neutron star binary as a two body problem. On the other hand, Eq. (70) is already enough to discuss the possibility of an experimental detection. It is clear that the spin effects are very small, even for completely spin polarized stars. Moreover, the torsion decreases exponentially, and practical results will strongly depend on the coupling constants $`b`$ and $`c`$, especially on $`c`$. A possible, indirect, detection could be based on gravitational waves emitted by a merging binary system, since this will allow for conclusions on the evolution of the system in the strong field regime.
## 8 Conclusions
Considering, in the framework of Poincaré gauge theory, the Christoffel part of the Lorentz connection as a function of the tetrad field in a formalism similar to the second order formalism used in general relativity and in metrical theories in general, we were able to present a theory that allows for dynamically propagating torsion fields, preserving nevertheless a full general relativity limit in the absence of spinning matter. This is not possible in the framework of the conventional, first order formalism, without introducing higher order derivatives of the independent field variables and apparently running into trouble with the Cauchy initial value problem. The equivalence of the second order formalism with the conventional approach has been established and therefore, ultimately, we have shown that in the conventional formalism, certain Lagrangians lead to consistent theories despite the appearance of higher derivatives. A concrete model has been presented, and for a specific class of solutions, the field equations have been reduced to an Einstein-Proca system using a slight modification of the Yasskin method known from conventional Einstein-Yang-Mills theories. Approximate solutions have been discussed for concrete cases, corresponding to the exterior of spin polarized neutron stars. In view of an eventual experimental detection of the non-Riemannian effects, the spin precession equation of a test body in such a field configuration has been briefly analyzed.
## Acknowlegments
This work has been supported by EPEAEK II in the framework of “PYTHAGORAS II - SUPPORT OF RESEARCH GROUPS IN UNIVERSITIES” (funding: 75% ESF - 25% National Funds).
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# Variation in the number of points on elliptic curves and applications to excess rank
## 1. Introduction
Let $``$ be a one-parameter family of elliptic curves (equivalently, an elliptic surface) over $`(T)`$:
$$y^2+a_1(T)xy+a_3(T)y=x^3+a_2(T)x^2+a_4(T)x+a_6(T),a_i(T)[T].$$
(1.1)
For each integer $`t`$ we have an elliptic curve $`E_t`$ over $``$, with $`N_t(p)`$ the number of solutions modulo $`p`$. Set $`a_t(p)=pN_t(p)`$. If $`a_1(T)=a_3(T)=0`$ we have
$$a_t(p)=\underset{xmodp}{}\left(\frac{x^3+a_2(t)x^2+a_4(t)x+a_6(t)}{p}\right).$$
(1.2)
We are interested in evaluating the second moment for the family:
$`A_{2,}(p):={\displaystyle \underset{tmodp}{}}a_t^2(p).`$ (1.3)
For one-parameter families of elliptic curves with $`j(T)`$ non-constant, Michel \[Mic\] proves $`A_{2,}(p)=p^2+O(p^{3/2})`$ by using the Lefschetz-Groethendieck trace formula. We show his result is sharp by constructing a family where the second moment has a term of size $`p^{3/2}`$.
Moments of the number of solutions modulo $`p`$ provide enormous amounts of information about the family. Rosen and Silverman \[RoSi\] prove a conjecture of Nagao (unconditionally for rational elliptic surfaces, conditional on Tate’s conjecture in general) that the first moment is related to the rank of the family over $`(T)`$, denoted $`\mathrm{rank}((T))`$:
$`\underset{X\mathrm{}}{lim}{\displaystyle \frac{1}{X}}{\displaystyle \underset{pX}{}}{\displaystyle \frac{\mathrm{log}p}{p}}{\displaystyle \underset{tmodp}{}}a_t(p)=\mathrm{rank}((T)).`$ (1.4)
In \[Mil1, ALM\] it is shown how to construct families with moderate rank by choosing the $`a_i(T)`$ so that the first moment is computable and large.
Another application is in the connections between number theory and random matrix theory \[KaSa1, KaSa2\]. In showing the behavior of the low lying zeros (zeros near the central point) of $`L`$-functions of a family of elliptic curves agrees with that of eigenvalues near $`1`$ of orthogonal groups, the only needed inputs are the first and second moment of the number of solutions modulo $`p`$; to first order all families of elliptic curves with rank $`r`$ over $`(T)`$ agree with the same random matrix ensemble (see \[Mil2\], though a similar result with a global rather than local rescaling of the zeros is implicit in \[Sil2\]). An analysis of the lower order terms in the first and second moments leads to breaking this universality; i.e., seeing lower-order family dependent behavior in the low lying zeros. See \[Mil1, Yo1\] for more details on family dependent behavior. One application of these correction terms is a refinement on predicting the number of curves in a family with rank above the family rank. While these corrections vanish in the limit of large conductor, they lead to slight modifications of excess rank bounds for conductors in the range accessible by computers.
In §2 we determine the second moment of $`a_t(p)`$ for a specific family, showing Michel’s result is sharp. In §3 we analyze how the lower order terms in Michel’s theorem are related to bounds for the average rank of the family.
## 2. The Second Moment for $`y^2=x^3+Tx^2+1`$
We may expand the Legendre symbol $`\left(\frac{x}{p}\right)`$ in (1.2) by
$`\left({\displaystyle \frac{x}{p}}\right)={\displaystyle \frac{1}{G_p}}{\displaystyle \underset{c=1}{\overset{p1}{}}}\left({\displaystyle \frac{c}{p}}\right)𝐞_p\left(cx\right),𝐞_p\left(x\right)=e^{2\pi \mathrm{i}x/p}.`$ (2.5)
Here $`G_p=_{amodp}\left(\frac{a}{p}\right)𝐞_p\left(a\right)`$ is the Gauss sum, which equals $`\sqrt{p}`$ for $`p1mod4`$ and $`\mathrm{i}\sqrt{p}`$ for $`p3mod4`$. See, for example, \[BEW\]. Using Gauss sums to evaluate Legendre sums is a common technique; we sketch an alternate approach which avoids Gauss sums in Remark 2.2.
For the family $`:y^2=x^3+Tx^2+1`$, $`j(T)=\frac{256T^6}{4T^3+27}`$ and thus by Michel’s Theorem $`A_{2,}(p)=p^2+O(p^{3/2})`$. We determine an exact formula for $`A_{2,}(p)`$:
###### Theorem 2.1.
For the one-parameter family $`:y^2=x^3+Tx^2+1`$ over $`(T)`$, for $`p>2`$ the second moment of $`a_t(p)`$ is
$$A_{2,}(p)=\underset{tmodp}{}a_t(p)^2=p^2n_{3,2,p}p1+p\underset{xmodp}{}\left(\frac{4x^3+1}{p}\right),$$
(2.6)
where $`n_{3,2,p}`$ denotes the number of cube roots of $`2`$ modulo $`p`$. For any $`[a,b][2,2]`$ there are infinitely many primes $`p1mod3`$ such that
$$A_{2,}(p)\left(p^2n_{3,2,p}p1\right)[ap^{3/2},bp^{3/2}].$$
(2.7)
###### Proof.
Combining (1.2) and (2.5) yields
$`A_{2,}(p)={\displaystyle \underset{tmodp}{}}{\displaystyle \underset{xmodp}{}}{\displaystyle \underset{ymodp}{}}\left({\displaystyle \frac{x^3+1+x^2t}{p}}\right)\left({\displaystyle \frac{y^3+1+y^2t}{p}}\right)`$
$`={\displaystyle \underset{x,ymodp}{}}{\displaystyle \underset{c,d=1}{\overset{p1}{}}}{\displaystyle \frac{1}{p}}\left({\displaystyle \frac{cd}{p}}\right)𝐞_p\left(c(x^3+1)d(y^3+1)\right){\displaystyle \underset{tmodp}{}}𝐞_p\left((cx^2dy^2)t\right);`$ (2.8)
above we used the complex conjugate of (2.5) in expanding $`\left(\frac{y^3+1+y^2t}{p}\right)`$. The two Gauss sum expansions give $`\frac{1}{G_p\overline{G_p}}=\frac{1}{p}`$. It will be convenient to set $`g(x,y)=(xy)(x^2y^2(x+y))`$.
Note $`c`$ and $`d`$ are invertible modulo $`p`$ in (2). If the numerator in the $`t`$-exponential is non-zero, the $`t`$-sum vanishes. Thus it suffices to study (2) for $`cx^2dy^2modp`$. If exactly one of $`x`$ and $`y`$ vanishes, then $`cx^2dy^2modp`$. Hence $`x`$ and $`y`$ are both zero or non-zero. If both are zero the $`t`$-sum gives $`p`$, the $`c`$-sum gives $`G_p`$, the $`d`$-sum gives $`\overline{G}_p`$, for a total contribution of $`p`$. If $`x`$ and $`y`$ are non-zero then we must have $`dcx^2y^2modp`$. The $`t`$-sum gives $`p`$. Thus (2) is
$`A_{2,}(p)`$ $`=`$ $`p+{\displaystyle \underset{x,y=1}{\overset{p1}{}}}{\displaystyle \underset{c=1}{\overset{p1}{}}}{\displaystyle \frac{1}{p}}\left({\displaystyle \frac{x^2y^2}{p}}\right)𝐞_p\left(cy^2(x^3y^2+y^2x^2y^3x^2)\right)p`$ (2.9)
$`=`$ $`p+{\displaystyle \underset{x,y=1}{\overset{p1}{}}}{\displaystyle \underset{c=1}{\overset{p1}{}}}𝐞_p\left(cy^2(xy)(x^2y^2(x+y))\right)`$
$`=`$ $`p+{\displaystyle \underset{x,y=1}{\overset{p1}{}}}{\displaystyle \underset{c=0}{\overset{p1}{}}}𝐞_p\left(cy^2g(x,y)\right){\displaystyle \underset{x,y=1}{\overset{p1}{}}}1`$
$`=`$ $`p+{\displaystyle \underset{x,y=1}{\overset{p1}{}}}{\displaystyle \underset{c=0}{\overset{p1}{}}}𝐞_p\left(cy^2g(x,y)\right)(p1)^2`$
$`=`$ $`p\mathrm{\#}\{x,y0modp:g(x,y)0modp\}+p(p1)^2,`$
where the last equality follows from the fact that if $`g(x,y)0modp`$ then the $`c`$-sum is $`p`$ and otherwise it is $`0`$. We are left with counting how often $`g(x,y)0modp`$ for $`x`$, $`y`$ non-zero.
Whenever $`x=y`$ then $`g(x,y)0modp`$; therefore there are $`p1`$ solutions from $`x=y`$. Consider now $`x^2y^2x+ymodp`$, which we may rewrite as a quadratic in $`y`$: $`x^2y^2yx`$ $`0modp`$. By the Quadratic Formula modulo $`p`$ (recall $`p`$ is odd), if the discriminant $`4x^3+1`$ is a non-zero square modulo $`p`$ there are two distinct roots, if it is not a square modulo $`p`$ there are no roots, and if the discriminant vanishes there is one root. Equivalently, if $`\left(\frac{4x^3+1}{p}\right)=1`$ (respectively, $`1`$ or $`0`$) there are two (respectively, none or one) solutions to $`x^2y^2x+ymodp`$.
Recall neither $`x`$ nor $`y`$ is allowed to be zero. If $`y=0`$ then $`x^2y^2x+ymodp`$ reduces to $`x=0`$. Hence our solutions have $`x,y0modp`$, and for a non-zero $`x`$ the number of non-zero $`y`$ with $`x^2y^2x+ymodp`$ is $`1+\left(\frac{4x^3+1}{p}\right)`$. Hence the number of non-zero pairs with $`x^2y^2x+ymodp`$ is
$`{\displaystyle \underset{x=1}{\overset{p1}{}}}\left(1+\left({\displaystyle \frac{4x^3+1}{p}}\right)\right)=p1+{\displaystyle \underset{xmodp}{}}\left({\displaystyle \frac{4x^3+1}{p}}\right)1.`$ (2.10)
We must be careful about double counting solutions. If both $`xy0modp`$ and $`x^2y^2x+ymodp`$, then we find $`x^42xmodp`$. As $`x0modp`$, we obtain $`x^32modp`$. We have double counted all pairs $`(x,x)`$ with $`x`$ a cube root of $`2`$ modulo $`p`$. Let $`n_{3,2,p}`$ denote the number of cube roots of $`2`$ modulo $`p`$; $`|n_{3,2,p}|3`$. We have shown
$`\mathrm{\#}\{x,y\{1,\mathrm{},p1\}:g(x,y)0modp\}`$
$`=(p1)+\left(p1+{\displaystyle \underset{xmodp}{}}\left({\displaystyle \frac{4x^3+1}{p}}\right)1\right)n_{3,2,p}`$
$`=2p3n_{3,2,p}+{\displaystyle \underset{xmodp}{}}\left({\displaystyle \frac{4x^3+1}{p}}\right).`$ (2.11)
Thus (2.9) becomes
$`A_{2,}(p)`$ $`=`$ $`p\left(2p3n_{3,2,p}+{\displaystyle \underset{xmodp}{}}\left({\displaystyle \frac{4x^3+1}{p}}\right)\right)+p(p1)^2`$ (2.12)
$`=`$ $`p^2n_{3,2,p}p1+p{\displaystyle \underset{xmodp}{}}\left({\displaystyle \frac{4x^3+1}{p}}\right).`$
To complete the analysis, we need to determine the size of $`_{xmodp}\left(\frac{4x^3+1}{p}\right)`$. Note this is the number of solutions modulo $`p`$ to the elliptic curve $`y^2=4x^3+1`$, and this curve is equivalent to $`E:y^2=x^3+16`$. This curve has analytic rank $`0`$, as can be seen from $`L(E,1).5968`$. It has complex multiplication, and for $`p2mod3`$, $`a_E(p)=0`$. Write $`a_E(p)=2\sqrt{p}\mathrm{cos}\theta _{E,p}`$. As $`E`$ has complex multiplication, for $`p1mod3`$ the distribution of the angles $`\theta _{E,p}`$ is known; all we need is that for any $`[\mathrm{\Theta },\mathrm{\Theta }^{}][0,\pi ]`$, a positive percent of the time $`\theta _{E,p}[\mathrm{\Theta },\mathrm{\Theta }^{}]`$. This implies that a typical $`a_E(p)`$ is of size $`\sqrt{p}`$ if $`p1mod3`$, and hence for any $`[a,b][2,2]`$ we can find infinitely many primes $`p`$ with
$$A_{2,}(p)\left(p^2n_{3,2,p}p1\right)[ap^{3/2},bp^{3/2}].$$
(2.13)
Note that if $`p2mod3`$, as $`xx^3modp`$ is an automorphism then $`n_{3,2,p}=1`$ and $`a_E(p)=0`$. Thus, at least half the time, $`A_{2,}(p)=p^2p1`$.
###### Remark 2.2.
A few words should be said about how we cooked up this family. If instead of $`y^2=x^3+Tx^2+1`$ we had $`y^2=x^3+Tx+1`$, we would have found the condition $`dcxy^1modp`$. As we have $`\left(\frac{cd}{p}\right)`$ this would lead to $`\left(\frac{c^2}{p}\right)\left(\frac{xy}{p}\right)`$ times a similar $`c`$-exponential. It would not suffice to determine how often a similar $`g(x,y)`$ vanished; we would need to know the value of $`\left(\frac{xy}{p}\right)`$. Our analysis was greatly aided by the presence of $`\left(\frac{x^2y^2}{p}\right)`$. We also want to change the order of summation and do the $`t`$-sum first, which basically forces our family to be at most quadratic in $`t`$, and such that $`g(x,y)`$ factors easily. Instead of expanding by using Gauss sums, we could write the product of Legendre symbols (2) as the Legendre symbol of $`h(x,y,t)`$, where $`h`$ is quadratic in $`t`$ with leading term $`x^2y^2t^2`$:
$`\left({\displaystyle \frac{x^3+1+x^2t}{p}}\right)\left({\displaystyle \frac{y^3+1+y^2t}{p}}\right)`$
$`=\left({\displaystyle \frac{x^2y^2t^2+\left(y^2(x^3+1)+x^2(y^3+1)\right)t(x^3+1)(y^3+1)}{p}}\right).`$ (2.14)
We execute the $`t`$-sum first. Quadratic Legendre sums are easily determined; what matters is the discriminant modulo $`p`$. After some algebra we find the discriminant is $`g(x,y)^2`$ (with $`g(x,y)`$ as before), and then the argument proceeds identically. See \[Mil1, ALM\] for more on determining tractable families where the summation can be done in closed form. These families will be quadratic in $`t`$, although not necessarily in Weierstrass form.
## 3. Other Families and Applications to Excess Rank
We give some additional examples of families where the first and second moments can be determined exactly; see \[Mil1\] for the calculations (though we provide calculations of a representative set of these families in Appendices A through C).
Recall $`n_{3,2,p}`$ denotes the number of cube roots of $`2`$ modulo $`p`$, and set $`c_0(p)=\left(\frac{3}{p}\right)+\left(\frac{3}{p}\right)`$, $`c_1(p)=\left[_{xmodp}\left(\frac{x^3x}{p}\right)\right]^2`$ and $`c_{3/2}(p)=p_{xmodp}\left(\frac{4x^3+1}{p}\right)`$.
| Family | $`A_{1,}(p)`$ | $`A_{2,}(p)`$ |
| --- | --- | --- |
| $`y^2=x^3+Sx+T`$ | $`0`$ | $`p^3p^2`$ |
| $`y^2=x^3+2^4(3)^3(9T+1)^2`$ | $`0`$ | $`\{\genfrac{}{}{0pt}{}{2p^22pp2mod3}{0p1mod3}`$ |
| $`y^2=x^3\pm 4(4T+2)x`$ | $`0`$ | $`\{\genfrac{}{}{0pt}{}{2p^22pp1mod3}{0p3mod3}`$ |
| $`y^2=x^3+(T+1)x^2+Tx`$ | $`0`$ | $`p^22p1`$ |
| $`y^2=x^3+x^2+2T+1`$ | $`0`$ | $`p^22p\left(\frac{3}{p}\right)`$ |
| $`y^2=x^3+Tx^2+1`$ | $`p`$ | $`p^2n_{3,2,p}p1+c_{3/2}(p)`$ |
| $`y^2=x^3T^2x+T^2`$ | $`2p`$ | $`p^2pc_1(p)c_0(p)`$ |
| $`y^2=x^3T^2x+T^4`$ | $`2p`$ | $`p^2pc_1(p)c_0(p)`$ |
The first family is the family of all elliptic curves; it is a two parameter family and we expect the main term of its second moment to be $`p^3`$. Note that except for our family $`y^2=x^3+Tx^2+1`$, all the families $``$ have $`A_{2,}(p)=p^2h(p)p+O(1)`$, where $`h(p)`$ is non-negative. Further, many of the families have $`h(p)=m_{}>0`$. Note $`c_1(p)`$ is the square of the coefficients from an elliptic curve with complex multiplication. It is non-negative and of size $`p`$ for $`p3mod4`$, and zero for $`p1mod4`$ (send $`xxmodp`$ and note $`\left(\frac{1}{p}\right)=1`$). It is somewhat remarkable that all these families have a correction to the main term in Michel’s theorem in the same direction, and we analyze the consequence this has on the average rank. For our family which has a $`p^{3/2}`$ term, note that on average this term is zero and the $`p`$ term is negative.
Consider a one-parameter family of elliptic curves $``$ of rank $`r`$ over $`(T)`$. With our normalizations, under GRH the non-trivial zeros of $`E_t`$ are $`1+i\gamma _t`$, $`\gamma _t`$. We typically study $`t[N,2N]`$ with $`N\mathrm{}`$. Let $`C_t`$ be the conductor of the elliptic curve $`E_t`$, and let $`\mathrm{log}R=\frac{1}{N}_{t=N}^{2N}\mathrm{log}C_t`$ be the average log-conductor. For many families there is an integer $`a`$ such that $`\mathrm{log}C_t\mathrm{log}N^a`$ for most curves; this is true for the families listed above (see \[Mil1\] for the calculations). Assuming the Birch and Swinnerton-Dyer conjecture, by Silverman’s specialization theorem eventually all curves $`E_t`$ have rank at least $`r`$, and under natural standard conjectures (see \[He\]) a typical family will have equidistribution of signs of the functional equations. What is typically seen in studying the ranks of curves in a family is that roughly $`30\%`$ have rank $`r`$ and $`20\%`$ rank $`r+2`$, while about $`48\%`$ have rank $`r+1`$ and $`2\%`$ rank $`r+3`$. Random matrix theory predicts that in the limit $`50\%`$ should be rank $`r`$ and $`50\%`$ rank $`r+1`$ for an average rank of $`r+\frac{1}{2}`$, markedly different from the observed (approximately) $`r+\frac{1}{2}+.40`$. See \[Fe1, Fe2, Wa\] for numerical investigations and \[Br, H-B, FP, Mic, Sil2, Yo2\] for theoretical bounds of the average rank.
The excess rank question is whether this disagreement persists or is a result of small data. We often expect the rate of convergence for problems such as this to be like the logarithm of the conductors. As the conductors are often at most $`10^{12}`$, it is reasonable to believe the data is misleading (especially as random matrix theory predicts sub-families of higher rank of size $`N^{3/4}`$, and for small $`N`$ such families are a noticeable percentage; see for example \[CKRS, DFK, Go, GM, Mai, Ono, RoSi, ST, Yu\] for discussions of random matrix predictions and results from number theory).
For an even Schwartz test function $`\varphi `$ with $`\mathrm{supp}(\widehat{\varphi })(\sigma ,\sigma )`$, the $`1`$-level density (which is basically just the sum of the explicit formula for each curve) is defined by
$`{\displaystyle \frac{1}{N}}{\displaystyle \underset{t=N}{\overset{2N}{}}}{\displaystyle \underset{\gamma _t}{}}\varphi \left(\gamma _t{\displaystyle \frac{\mathrm{log}R}{2\pi }}\right)=\widehat{\varphi }(0)+\varphi (0){\displaystyle \frac{2}{N}}{\displaystyle \underset{t=N}{\overset{2N}{}}}{\displaystyle \underset{p}{}}{\displaystyle \frac{\mathrm{log}p}{\mathrm{log}R}}{\displaystyle \frac{1}{p}}\widehat{\varphi }\left({\displaystyle \frac{\mathrm{log}p}{\mathrm{log}R}}\right)a_t(p)`$
$`{\displaystyle \frac{2}{N}}{\displaystyle \underset{t=N}{\overset{2N}{}}}{\displaystyle \underset{p}{}}{\displaystyle \frac{\mathrm{log}p}{\mathrm{log}R}}{\displaystyle \frac{1}{p^2}}\widehat{\varphi }\left({\displaystyle \frac{2\mathrm{log}p}{\mathrm{log}R}}\right)a_t(p)^2+O\left({\displaystyle \frac{\mathrm{log}\mathrm{log}R}{\mathrm{log}R}}\right).`$ (3.15)
If $`\varphi `$ is non-negative, we obtain a bound for the average rank in the family by restricting the sum to be only over zeros at the central point. The error $`O\left(\frac{\mathrm{log}\mathrm{log}R}{\mathrm{log}R}\right)`$ comes from trivial estimation and ignores probable cancellation, and we expect $`O\left(\frac{1}{\mathrm{log}R}\right)`$ or smaller to be the correct magnitude. For most families $`\mathrm{log}R\mathrm{log}N^a`$ for some integer $`a`$.
The main term of the first and second moments of the $`a_t(p)`$ give $`r\varphi (0)`$ and $`\frac{1}{2}\varphi (0)`$, respectively, in (3). Assume the second moment of $`a_t(p)^2`$ is $`p^2m_{}p+O(1)`$, $`m_{}>0`$. We have already handled the contribution from $`p^2`$, and $`m_{}p`$ contributes
$`S_2`$ $``$ $`{\displaystyle \frac{2}{N}}{\displaystyle \underset{p}{}}{\displaystyle \frac{\mathrm{log}p}{\mathrm{log}R}}\widehat{\varphi }\left(2{\displaystyle \frac{\mathrm{log}p}{\mathrm{log}R}}\right){\displaystyle \frac{1}{p^2}}{\displaystyle \frac{N}{p}}(m_{}p)`$ (3.16)
$`=`$ $`{\displaystyle \frac{2m_{}}{\mathrm{log}R}}{\displaystyle \underset{p}{}}\widehat{\varphi }\left(2{\displaystyle \frac{\mathrm{log}p}{\mathrm{log}R}}\right){\displaystyle \frac{\mathrm{log}p}{p^2}}.`$
Thus there is a contribution of size $`\frac{1}{\mathrm{log}R}`$. A good choice of test functions (see Appendix A of \[ILS\]) is the Fourier pair
$`\varphi (x)={\displaystyle \frac{\mathrm{sin}^2(2\pi \frac{\sigma }{2}x)}{(2\pi x)^2}},\widehat{\varphi }(u)=\{\begin{array}{cc}\frac{\sigma |u|}{4}\hfill & \text{if }|u|\sigma \hfill \\ 0\hfill & \text{otherwise.}\hfill \end{array}`$ (3.17)
Note $`\varphi (0)=\frac{\sigma ^2}{4}`$, $`\widehat{\varphi }(0)=\frac{\sigma }{4}=\frac{\varphi (0)}{\sigma }`$, and evaluating the prime sum in (3.16) gives
$`S_2`$ $``$ $`\left({\displaystyle \frac{.986}{\sigma }}{\displaystyle \frac{2.966}{\sigma ^2\mathrm{log}R}}\right){\displaystyle \frac{m_{}}{\mathrm{log}R}}\varphi (0).`$ (3.18)
Let $`r_t`$ denote the number of zeros of $`E_t`$ at the central point (i.e., the analytic rank). Then up to our $`O\left(\frac{\mathrm{log}\mathrm{log}R}{\mathrm{log}R}\right)`$ errors (which we think should be smaller), we have
$`{\displaystyle \frac{1}{N}}{\displaystyle \underset{t=N}{\overset{2N}{}}}r_t\varphi (0)`$ $``$ $`{\displaystyle \frac{\varphi (0)}{\sigma }}+\left(r+{\displaystyle \frac{1}{2}}\right)\varphi (0)+\left({\displaystyle \frac{.986}{\sigma }}{\displaystyle \frac{2.966}{\sigma ^2\mathrm{log}R}}\right){\displaystyle \frac{m_{}}{\mathrm{log}R}}\varphi (0)`$
$`\text{Ave Rank}_{[N,2N]}()`$ $``$ $`{\displaystyle \frac{1}{\sigma }}+r+{\displaystyle \frac{1}{2}}+\left({\displaystyle \frac{.986}{\sigma }}{\displaystyle \frac{2.966}{\sigma ^2\mathrm{log}R}}\right){\displaystyle \frac{m_{}}{\mathrm{log}R}}.`$ (3.19)
###### Remark 3.1.
The Density Conjecture states that the $`1`$-level density (in the limit) should hold for all $`\sigma `$. In that case, the lower order terms from the second moment will *not* contribute to the bound for the average rank, as their contribution vanishes as $`\sigma \mathrm{}`$. Of course, the agreement with random matrix theory is a statement about the limit as $`N\mathrm{}`$; the correct finite-conductor model is still unknown
Let us examine the boost the $`m_{}p`$ term from the second moment gives to the upper bound for the average rank. As remarked, if our $`1`$-level density were true for all $`\sigma `$ then there would be no contribution from the correction term to the second sum, nor would the $`\frac{1}{\sigma }`$ term contribute, and we would obtain the average rank is bounded by $`r+\frac{1}{2}`$.
Let us assume we know the $`1`$-level density up to $`\sigma =1`$. (This is well beyond the range of current technology; the best result to date is for the family of all elliptic curves, where Young \[Yo2\] proves we may take any $`\sigma <\frac{7}{9}`$). Assume $`m_{}=1`$. The $`\frac{1}{\sigma }`$ term would contribute $`1`$, the lower correction would contribute $`.03`$ for conductors of size $`10^{12}`$, and thus the average rank is bounded by $`1+r+\frac{1}{2}+.03`$ $`=r+\frac{1}{2}+1.03`$. This is significantly higher than Fermigier’s observed $`r+\frac{1}{2}+.40`$.
If we were able to prove our $`1`$-level density for $`\sigma =2`$, then the $`\frac{1}{\sigma }`$ term would contribute $`\frac{1}{2}`$, and the lower order correction would contribute $`.02`$ for conductors of size $`10^{12}`$. Thus the average rank would be bounded by $`\frac{1}{2}+r+\frac{1}{2}+.02`$ $`=r+\frac{1}{2}+.52`$. While the main error contribution is from $`\frac{1}{\sigma }`$, there is still a noticeable effect from the lower order terms in $`A_{2,}(p)`$. Moreover, we are now in the ballpark of Fermigier’s bound; of course, we were already there without the potential correction term.
It seems hopeless to think about obtaining a $`1`$-level density for any family of elliptic curves with support $`\sigma =2`$ or more. Iwawniec, Luo and Sarnak \[ILS\] obtain such large support for families of weight $`k`$ cuspidal newforms of square-free level $`N`$, but only because of great averaging formulas (the Bessel-Kloosterman expansion in the Petersson formula) available for the family; the corresponding averaging formulas for elliptic curves are much weaker. We use the periodicity of $`a_t(p)`$ as a function of $`tmodp`$ to analyze complete sums of the moments for each prime; however, the error from the incomplete sum is bounded by Hasse and contributes a large error (this problem is avoided in \[ILS\] because the Petersson formula gives us sums over a basis of newforms, and there is no incomplete piece to be approximated). The random matrix models for the behavior of the zeros near the central point have been shown to hold as the conductors tend to infinity; in the small conductor ranges investigated, it is not surprising that there is disagreement. While it would be desirable to find a good model for small conductors (similar to Keating and Snaith’s \[KeSn1, KeSn2\] modeling zeros of $`\zeta (s)`$ at height $`T`$ by $`N\times N`$ matrices with $`N=\frac{\mathrm{log}T}{2\pi }`$), we can identify potential family dependent lower order terms in the $`1`$-level density arising from lower order terms in the second moment. For finite conductors these do lead to slightly larger predicted upper bounds for the average rank in a family.
The following appendices contain calculations for the second moment of a representative set of the other families mentioned in §3, and for completeness proofs of standard quadratic Legendre sums. For general one-parameter families of elliptic curves, there are not closed form expressions for the moments. Our hope is that these families may be useful for other investigations.
## Appendix A The Family $`y^2=x^3+T^2`$
###### Theorem A.1.
For the family $`:y^2=x^3+T^2`$,
$$A_{2,}(p)=\{\begin{array}{cc}2p^22p\hfill & \text{if }p1mod3\hfill \\ 0\hfill & \text{if }p2mod3\text{.}\hfill \end{array}$$
(A.20)
###### Proof.
If $`p2mod3`$ then $`a_t^2(p)=0`$ as $`xx^3modp`$ is an automorphism. Assume $`p1mod3`$.
$`A_{2,}(p)`$ $`=`$ $`{\displaystyle \underset{tmodp}{}}{\displaystyle \underset{xmodp}{}}{\displaystyle \underset{ymodp}{}}\left({\displaystyle \frac{x^3+t^2}{p}}\right)\left({\displaystyle \frac{y^3+t^2}{p}}\right)`$ (A.21)
$`=`$ $`{\displaystyle \underset{t=1}{\overset{p1}{}}}{\displaystyle \underset{xmodp}{}}{\displaystyle \underset{ymodp}{}}\left({\displaystyle \frac{x^3+t^2}{p}}\right)\left({\displaystyle \frac{y^3+t^2}{p}}\right)`$
$`=`$ $`{\displaystyle \underset{t=1}{\overset{p1}{}}}{\displaystyle \underset{xmodp}{}}{\displaystyle \underset{ymodp}{}}\left({\displaystyle \frac{t^4}{p}}\right)\left({\displaystyle \frac{tx^3+1}{p}}\right)\left({\displaystyle \frac{ty^3+1}{p}}\right)`$
$`=`$ $`{\displaystyle \underset{xmodp}{}}{\displaystyle \underset{ymodp}{}}{\displaystyle \underset{tmodp}{}}\left({\displaystyle \frac{tx^3+1}{p}}\right)\left({\displaystyle \frac{ty^3+1}{p}}\right)p^2.`$
We use inclusion / exclusion to reduce to $`xy0`$. If $`x=0`$, the $`t`$-sum vanishes unless $`y=0`$, in which case we get $`p`$. Similarly if $`y=0`$, the $`t`$ and $`x`$-sums give $`p`$. We subtract the doubly counted contribution from $`x=y=0`$, which gives $`p`$. Thus
$`A_{2,}(p)`$ $`=`$ $`{\displaystyle \underset{x=1}{\overset{p1}{}}}{\displaystyle \underset{y=1}{\overset{p1}{}}}{\displaystyle \underset{t(p)}{}}\left({\displaystyle \frac{tx^3+1}{p}}\right)\left({\displaystyle \frac{ty^3+1}{p}}\right)+2ppp^2.`$ (A.22)
By Lemma D.1, the $`t`$-sum is $`(p1)\left(\frac{x^3y^3}{p}\right)`$ if $`p|(x^3y^3)^2`$ and $`\left(\frac{x^3y^3}{p}\right)`$ otherwise. As $`p=6m+1`$, let $`g`$ be a generator of the multiplicative group $`/p`$. Solving $`g^{3a}g^{3b}modp`$ yields $`b=a`$, $`a+2m`$, or $`a+4m`$. Thus, $`x^3y^3modp`$ three times, and in each instance $`y`$ equals $`x`$ times a square ($`1`$, $`g^{2m}`$, $`g^{4m}`$).
$`A_{2,}(p)`$ $`=`$ $`{\displaystyle \underset{x=1}{\overset{p1}{}}}{\displaystyle \underset{\genfrac{}{}{0pt}{}{y=1}{y^3x^3}}{\overset{p1}{}}}p{\displaystyle \underset{x=1}{\overset{p1}{}}}{\displaystyle \underset{y=1}{\overset{p1}{}}}\left({\displaystyle \frac{x^3y^3}{p}}\right)+pp^2`$ (A.23)
$`=`$ $`(p1)3p+pp^2`$
$`=`$ $`2p^22p.`$
## Appendix B The Family $`y^2=x^3+x^2+T`$
###### Theorem B.1.
For the family $`:y^2=x^3+x^2+T`$ we have
$`A_{2,}(p)=p^22pp\left({\displaystyle \frac{3}{p}}\right).`$ (B.24)
###### Proof.
$`A_{2,}(p)`$ $`=`$ $`{\displaystyle \underset{t=0}{\overset{p1}{}}}{\displaystyle \underset{x=0}{\overset{p1}{}}}{\displaystyle \underset{y=0}{\overset{p1}{}}}\left({\displaystyle \frac{t+(x^3+x^2)}{p}}\right)\left({\displaystyle \frac{t+(y^3+y^2)}{p}}\right)`$
$`=`$ $`{\displaystyle \underset{t=0}{\overset{p1}{}}}{\displaystyle \underset{x=0}{\overset{p1}{}}}{\displaystyle \underset{y=0}{\overset{p1}{}}}\left({\displaystyle \frac{t^2+\left((x^3+x^2)+(y^3+y^2)\right)t+(x^3+x^2)(y^3+y^2)}{p}}\right).`$
Let $`\delta (x,y)=(x^3+x^2)(y^3+y^2)`$; note $`\delta (x,y)^2`$ is the discriminant of the quadratic regarded as a function of $`t`$. The $`t`$-sum is $`p1`$ if $`p|\delta (x,y)`$ and $`1`$ otherwise. Note
$$\delta (x,y)=(xy)(y^2+(x+1)y+(x^2+x)).$$
(B.26)
The first factor is congruent to zero when $`x=y`$; for fixed $`x`$, the discriminant of the second factor is $`(x+1)^24(x^2+x)=12x3x^2`$. Thus the number of solutions of the second factor, for fixed $`x`$, is $`1+\left(\frac{12x3x^2}{p}\right)`$. As the discriminant of $`12x3x^2`$ is $`16`$, summing over $`x`$ for $`p>2`$ yields $`p\left(\frac{3}{p}\right)`$ by Lemma D.2.
We must be careful about double counting. If both factors are congruent to zero, then $`3x^2+2x0`$, or $`x0,23^1`$. Hence we always double count two solutions.
$`A_{2,}(p)`$ $`=`$ $`\left[p+p\left({\displaystyle \frac{3}{p}}\right)2\right]p{\displaystyle \underset{x=0}{\overset{p1}{}}}{\displaystyle \underset{y=0}{\overset{p1}{}}}1`$ (B.27)
$`=`$ $`p^22pp\left({\displaystyle \frac{3}{p}}\right).`$
## Appendix C $`y^2=x^3T^2x+T^2`$
Consider the family $`:y^2=x^3T^2x+T^2`$. We calculate the first moment of $`a_t(p)`$, which shows the family is of rank $`2`$ over $`(T)`$, and then determine the second moment.
###### Theorem C.1.
For $`:y^2=x^3T^2x+T^2`$, $`A_{1,}(p)=2p`$. Thus by Rosen and Silverman the family has rank $`2`$ over $`(T)`$.
###### Proof.
$`A_{1,}(p)`$ $`=`$ $`{\displaystyle \underset{t(p)}{}}a_t(p)={\displaystyle \underset{t(p)}{}}{\displaystyle \underset{x(p)}{}}\left({\displaystyle \frac{x^3t^2x+t^2}{p}}\right)`$ (C.28)
$`=`$ $`{\displaystyle \underset{t=1}{\overset{p1}{}}}{\displaystyle \underset{x(p)}{}}\left({\displaystyle \frac{x^3t^2x+t^2}{p}}\right)={\displaystyle \underset{t=1}{\overset{p1}{}}}{\displaystyle \underset{x(p)}{}}\left({\displaystyle \frac{t^3x^3t^3x+t^2}{p}}\right)`$
$`=`$ $`{\displaystyle \underset{t=1}{\overset{p1}{}}}{\displaystyle \underset{x(p)}{}}\left({\displaystyle \frac{t^2}{p}}\right)\left({\displaystyle \frac{t(x^3x)+1}{p}}\right)`$
$`=`$ $`{\displaystyle \underset{t(p)}{}}{\displaystyle \underset{x(p)}{}}\left({\displaystyle \frac{t(x^3x)+1}{p}}\right){\displaystyle \underset{x(p)}{}}\left({\displaystyle \frac{1}{p}}\right)`$
$`=`$ $`{\displaystyle \underset{t(p)}{}}{\displaystyle \underset{x=0,\pm 1}{}}\left({\displaystyle \frac{t(x^3x)+1}{p}}\right)+{\displaystyle \underset{t(p)}{}}{\displaystyle \underset{\genfrac{}{}{0pt}{}{x(p)}{x0,\pm 1}}{}}\left({\displaystyle \frac{t(x^3x)+1}{p}}\right)p`$
$`=`$ $`{\displaystyle \underset{t(p)}{}}{\displaystyle \underset{x=0,\pm 1}{}}\left({\displaystyle \frac{1}{p}}\right)+{\displaystyle \underset{\genfrac{}{}{0pt}{}{x(p)}{x0,\pm 1}}{}}{\displaystyle \underset{t(p)}{}}\left({\displaystyle \frac{t+1}{p}}\right)p`$
$`=`$ $`3p+0p=2p.`$
###### Theorem C.2.
For $`:y^2=x^3T^2x+T^2`$,
$`A_{2,}(p)=p^2p\left[{\displaystyle \underset{x(p)}{}}\left({\displaystyle \frac{(x^3x)}{p}}\right)\right]^2\left({\displaystyle \frac{3}{p}}\right)\left({\displaystyle \frac{3}{p}}\right)=p^2+O(p).`$
###### Proof.
$`A_{2,}(p)`$ $`=`$ $`{\displaystyle \underset{t(p)}{}}a_t^2(p)`$ (C.30)
$`=`$ $`{\displaystyle \underset{t(p)}{}}{\displaystyle \underset{x,y(p)}{}}\left({\displaystyle \frac{x^3t^2x+t^2}{p}}\right)\left({\displaystyle \frac{y^3t^2y+t^2}{p}}\right)`$
$`=`$ $`{\displaystyle \underset{t=1}{\overset{p1}{}}}{\displaystyle \underset{x,y(p)}{}}\left({\displaystyle \frac{x^3t^2x+t^2}{p}}\right)\left({\displaystyle \frac{y^3t^2y+t^2}{p}}\right)`$
$`=`$ $`{\displaystyle \underset{t=1}{\overset{p1}{}}}{\displaystyle \underset{x,y(p)}{}}\left({\displaystyle \frac{t^3x^3t^3x+t^2}{p}}\right)\left({\displaystyle \frac{t^3y^3t^3y+t^2}{p}}\right)`$
$`=`$ $`{\displaystyle \underset{t=1}{\overset{p1}{}}}{\displaystyle \underset{x,y(p)}{}}\left({\displaystyle \frac{t^4}{p}}\right)\left({\displaystyle \frac{t(x^3x)+1}{p}}\right)\left({\displaystyle \frac{t(y^3y)+1}{p}}\right)`$
$`=`$ $`{\displaystyle \underset{t=0}{\overset{p1}{}}}{\displaystyle \underset{x,y(p)}{}}\left({\displaystyle \frac{t(x^3x)+1}{p}}\right)\left({\displaystyle \frac{t(y^3y)+1}{p}}\right){\displaystyle \underset{x,y(p)}{}}\left({\displaystyle \frac{1}{p}}\right)`$
$`=`$ $`{\displaystyle \underset{x,y(p)}{}}{\displaystyle \underset{t(p)}{}}\left({\displaystyle \frac{t(x^3x)+1}{p}}\right)\left({\displaystyle \frac{t(y^3y)+1}{p}}\right)p^2.`$
In Lemma D.2 we showed that, if $`a`$ and $`b`$ are not both zero,
$$\underset{t=0}{\overset{p1}{}}\left(\frac{at^2+bt+c}{p}\right)=\{\begin{array}{cc}(p1)\left(\frac{a}{p}\right)\hfill & \text{if }p|b^24ac\hfill \\ \left(\frac{a}{p}\right)\hfill & \text{otherwise.}\hfill \end{array}$$
(C.31)
In $`A_{2,}(p)`$ we have
$`a`$ $`=`$ $`(x^3x)(y^3y)=y(x^21)x(y^21)`$
$`b`$ $`=`$ $`(x^3x)+(y^3y)`$
$`c`$ $`=`$ $`1`$
$`\delta (x,y)`$ $`=`$ $`b^24ac=\left((x^3x)(y^3y)\right)^2.`$ (C.32)
We use inclusion / exclusion on $`x^3x`$ and $`y^3y`$ vanishing. Assume first that $`x^3x`$ equals zero (happens three ways: $`x=0,\pm 1`$). Then we have $`_t\left(\frac{t(y^3y)+1}{p}\right)`$, which is $`3p`$ from our $`A_{1,}(p)`$ computation, giving $`33p`$. Similarly we get $`33p`$ if $`y^3y`$ is zero. We subtract the doubly counted $`x^3xy^3y0`$ (nine ways), each of which gives $`_t\left(\frac{1}{p}\right)=p`$. Hence the contribution from at least one of $`x^3x`$ and $`y^3y`$ vanishing is $`9p`$.
Assume $`x,y\{0,\pm 1\}`$. When is $`\delta (x,y)=(x^3x)(y^3y)0(p)`$?
$`\delta (x,y)`$ $`=`$ $`(xy)(x^2+xy+y^21).`$ (C.33)
Therefore
$`A_{2,}(p)={\displaystyle \underset{\genfrac{}{}{0pt}{}{x,y0,\pm 1}{\delta (x,y)0}}{}}p\left({\displaystyle \frac{(x^3x)(y^3y)}{p}}\right){\displaystyle \underset{x,y0,\pm 1}{}}\left({\displaystyle \frac{(x^3x)(y^3y)}{p}}\right)+9pp^2.`$
Clearly, $`\delta (x,y)0(p)`$ if $`x=y`$, which happens $`p3`$ times. If $`x=y`$ then the second factor is $`3x^21`$, which is congruent to zero at most twice.
When is $`\delta _2(x,y)=x^2+xy+y^210`$? By the Quadratic Formula mod $`p`$,
$`y={\displaystyle \frac{x\pm \sqrt{43x^2}}{2}},`$ (C.35)
which reduces to finding when $`43x^2`$ is a square mod $`p`$. We get two values of $`y`$ if it is equivalent to a non-zero square, one value if it is equivalent to zero, and no values if it is not equivalent to a square. When solving $`\delta _2(x,y)0(p)`$, we make sure such $`y\{0,\pm 1\}`$. If $`y=0`$, $`x=\pm 1`$; $`y=1`$, $`x=0`$ or $`1`$; $`y=1`$, $`x=0`$ or $`1`$. Therefore, we don’t get an excluded $`y`$ (and similarly if we reverse the rolls of $`y`$ and $`x`$). Thus the number of solutions to $`\delta _2(x,y)0(p)`$ is
$`{\displaystyle \underset{x=2}{\overset{p2}{}}}\left[1+\left({\displaystyle \frac{43x^2}{p}}\right)\right]`$ $`=`$ $`p3+{\displaystyle \underset{x=2}{\overset{p2}{}}}\left({\displaystyle \frac{43x^2}{p}}\right)`$ (C.36)
$`=`$ $`p6+{\displaystyle \underset{x(p)}{}}\left({\displaystyle \frac{43x^2}{p}}\right).`$
We again use Lemma D.2. The discriminant now is $`0^24(3)4`$. For $`p5`$, $`p`$ does not divide the discriminant, hence this sum is $`\left(\frac{3}{p}\right)`$.
Thus, for $`x0,\pm 1`$, the number of solutions with $`x^2+xy+y^2`$ $`1`$ is $`p6\left(\frac{3}{p}\right)`$; the number with $`xy0`$ is $`p3`$. At most two of the pairs $`(x,y)`$ satisfying $`x^2+xy+y^210(p)`$ also satisfy $`x=y`$. These pairs satisfy $`3x^21`$, thus, if $`\left(\frac{3}{p}\right)=1`$ we have doubly counted two solutions; if it is $`1`$, there was no double counting. Thus, the number of doubly counted pairs is $`1+\left(\frac{3}{p}\right)`$, and the total number of pairs is
$`2p10\left({\displaystyle \frac{3}{p}}\right)\left({\displaystyle \frac{3}{p}}\right).`$ (C.37)
When $`x=y0,\pm 1`$, clearly $`\left(\frac{(x^3x)(y^3y)}{p}\right)=1`$. Hence these terms contribute $`1`$.
Consider $`xy`$ and $`x^2+xy+y^210`$. Thus $`x,y0,\pm 1`$. Then $`y^21x(x+y)`$ and $`x^21y(x+y)`$ and
$`\left({\displaystyle \frac{(x^3x)(y^3y)}{p}}\right)=\left({\displaystyle \frac{x(x^21)y(y^21)}{p}}\right)=\left({\displaystyle \frac{x^2y^2(x+y)^2}{p}}\right).`$ (C.38)
As long as $`xy`$, this is $`1`$. If $`x=y`$ then we would have $`x^2x^2+x^210`$. This implies $`x=\pm 1`$, which cannot happen as $`x,y0,\pm 1`$. Therefore all pairs have their Legendre factor $`+1`$, and we need only count how many such pairs there are. We’ve previously shown this to be $`p+O(1)`$, therefore
$`A_{2,}(p)`$ $`=`$ $`p\left[2p10\left({\displaystyle \frac{3}{p}}\right)\left({\displaystyle \frac{3}{p}}\right)\right]{\displaystyle \underset{x,y0,\pm 1}{}}\left({\displaystyle \frac{(x^3x)(y^3y)}{p}}\right)+9pp^2`$ (C.39)
$`=`$ $`p^2p\left[{\displaystyle \underset{x(p)}{}}\left({\displaystyle \frac{x^3x}{p}}\right)\right]^2\left({\displaystyle \frac{3}{p}}\right)\left({\displaystyle \frac{3}{p}}\right).`$
As $`x^3x`$ is a non-singular elliptic curve, by Hasse its sum above is bounded by $`4p`$. It has complex multiplication and analytic rank $`0`$. For $`p3`$ mod $`4`$ its $`a_E(p)=0`$ (change variables $`xx`$); for the remaining $`p`$, the angles of $`\frac{a_E(p)}{2\sqrt{p}}`$ are uniformly distributed. Hence $`A_{2,}(p)=p^2+O(p)`$. ∎
###### Remark C.3.
The reason this calculation succeeds is we have a very tractable expression for $`x(x^21)y(y^21)`$ when $`x^2+xy+y^210`$ mod $`p`$. It was non-trivial to find a family with high rank over $`(T)`$ and $`A_{2,}(p)`$ computable.
## Appendix D Quadratics Sums of Legendre Symbols
For completeness we include proofs of standard sums of Legendre symbols.
###### Lemma D.1.
For $`p>2`$
$$S(n)=\underset{x=0}{\overset{p1}{}}\left(\frac{n_1+x}{p}\right)\left(\frac{n_2+x}{p}\right)=\{\genfrac{}{}{0pt}{}{p1\mathrm{if}p|n_1n_2}{1\mathrm{otherwise}.}$$
(D.40)
###### Proof.
Shifting $`x`$ by $`n_2`$, we need only prove the lemma when $`n_2=0`$. Assume $`(n,p)=1`$ as otherwise the result is trivial. For $`(a,p)=1`$ we have
$`S(n)`$ $`=`$ $`{\displaystyle \underset{x=0}{\overset{p1}{}}}\left({\displaystyle \frac{n+x}{p}}\right)\left({\displaystyle \frac{x}{p}}\right)`$ (D.41)
$`=`$ $`{\displaystyle \underset{x=0}{\overset{p1}{}}}\left({\displaystyle \frac{n+a^1x}{p}}\right)\left({\displaystyle \frac{a^1x}{p}}\right)`$
$`=`$ $`{\displaystyle \underset{x=0}{\overset{p1}{}}}\left({\displaystyle \frac{an+x}{p}}\right)\left({\displaystyle \frac{x}{p}}\right)=S(an)`$
Hence
$`S(n)`$ $`=`$ $`{\displaystyle \frac{1}{p1}}{\displaystyle \underset{a=1}{\overset{p1}{}}}{\displaystyle \underset{x=0}{\overset{p1}{}}}\left({\displaystyle \frac{an+x}{p}}\right)\left({\displaystyle \frac{x}{p}}\right)`$ (D.42)
$`=`$ $`{\displaystyle \frac{1}{p1}}{\displaystyle \underset{a=0}{\overset{p1}{}}}{\displaystyle \underset{x=0}{\overset{p1}{}}}\left({\displaystyle \frac{an+x}{p}}\right)\left({\displaystyle \frac{x}{p}}\right){\displaystyle \frac{1}{p1}}{\displaystyle \underset{x=0}{\overset{p1}{}}}\left({\displaystyle \frac{x}{p}}\right)^2`$
$`=`$ $`{\displaystyle \frac{1}{p1}}{\displaystyle \underset{x=0}{\overset{p1}{}}}\left({\displaystyle \frac{x}{p}}\right){\displaystyle \underset{a=0}{\overset{p1}{}}}\left({\displaystyle \frac{an+x}{p}}\right)1`$
$`=`$ $`01=1`$
Where do we use $`p>2`$? We used $`_{a=0}^{p1}\left(\frac{an+x}{p}\right)=0`$ for $`(n,p)=1`$. This is true for all odd primes (as there are $`\frac{p1}{2}`$ quadratic residues, $`\frac{p1}{2}`$ non-residues, and $`0`$); for $`p=2`$, there is one quadratic residue, no non-residues, and $`0`$.∎
###### Lemma D.2 (Quadratic Legendre Sums).
Assume $`a`$ and $`b`$ are not both zero mod $`p`$ and $`p>2`$. Then
$$\underset{t=0}{\overset{p1}{}}\left(\frac{at^2+bt+c}{p}\right)=\{\begin{array}{cc}(p1)\left(\frac{a}{p}\right)\hfill & \text{if }p|b^24ac\hfill \\ \left(\frac{a}{p}\right)\hfill & \text{otherwise.}\hfill \end{array}$$
(D.43)
###### Proof.
Assume $`a0(p)`$ as otherwise the proof is trivial. Let $`\delta =4^1(b^24ac)`$. Then
$`{\displaystyle \underset{t=0}{\overset{p1}{}}}\left({\displaystyle \frac{at^2+bt+c}{p}}\right)`$ $`=`$ $`{\displaystyle \underset{t=0}{\overset{p1}{}}}\left({\displaystyle \frac{a^1}{p}}\right)\left({\displaystyle \frac{a^2t^2+bat+ac}{p}}\right)`$
$`=`$ $`{\displaystyle \underset{t=0}{\overset{p1}{}}}\left({\displaystyle \frac{a}{p}}\right)\left({\displaystyle \frac{t^2+bt+ac}{p}}\right)`$
$`=`$ $`{\displaystyle \underset{t=0}{\overset{p1}{}}}\left({\displaystyle \frac{a}{p}}\right)\left({\displaystyle \frac{t^2+bt+4^1b^2+ac4^1b^2}{p}}\right)`$
$`=`$ $`{\displaystyle \underset{t=0}{\overset{p1}{}}}\left({\displaystyle \frac{a}{p}}\right)\left({\displaystyle \frac{(t+2^1b)^24^1(b^24ac)}{p}}\right)`$
$`=`$ $`{\displaystyle \underset{t=0}{\overset{p1}{}}}\left({\displaystyle \frac{a}{p}}\right)\left({\displaystyle \frac{t^2\delta }{p}}\right)`$
$`=`$ $`\left({\displaystyle \frac{a}{p}}\right){\displaystyle \underset{t=0}{\overset{p1}{}}}\left({\displaystyle \frac{t^2\delta }{p}}\right)`$
If $`\delta 0modp`$ we get $`p1`$. If $`\delta =\eta ^2,\eta 0`$, then by the Lemma D.1
$$\underset{t=0}{\overset{p1}{}}\left(\frac{t^2\delta }{p}\right)=\underset{t=0}{\overset{p1}{}}\left(\frac{t\eta }{p}\right)\left(\frac{t+\eta }{p}\right)=1.$$
(D.45)
We note that $`_{t=0}^{p1}\left(\frac{t^2\delta }{p}\right)`$ is the same for all non-square $`\delta `$’s (let $`g`$ be a generator of the multiplicative group, $`\delta =g^{2k+1}`$, change variables by $`tg^kt`$). Denote this sum by $`S`$, the set of non-zero squares by $``$, and the non-squares by $`𝒩`$. Since $`_{\delta =0}^{p1}\left(\frac{t^2\delta }{p}\right)=0`$ we have
$`{\displaystyle \underset{\delta =0}{\overset{p1}{}}}{\displaystyle \underset{t=0}{\overset{p1}{}}}\left({\displaystyle \frac{t^2\delta }{p}}\right)`$ $`=`$ $`{\displaystyle \underset{t=0}{\overset{p1}{}}}\left({\displaystyle \frac{t^2}{p}}\right)+{\displaystyle \underset{\delta }{}}{\displaystyle \underset{t=0}{\overset{p1}{}}}\left({\displaystyle \frac{t^2\delta }{p}}\right)+{\displaystyle \underset{\delta 𝒩}{}}{\displaystyle \underset{t=0}{\overset{p1}{}}}\left({\displaystyle \frac{t^2\delta }{p}}\right)`$ (D.46)
$`=`$ $`(p1)+{\displaystyle \frac{p1}{2}}(1)+{\displaystyle \frac{p1}{2}}S=0`$
Hence $`S=1`$, proving the lemma. ∎
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# Hopping dynamics for localized Lyapunov vectors in many-hard-disk systems
## I Introduction
The dynamical instability and many-body nature play essential roles in the justification of a statistical treatment for deterministic dynamical systems. The dynamical instability is described as a rapid expansion of the difference between two nearby trajectories, namely the Lyapunov vector, and the system is called chaotic if at least one exponential rate (Lyapunov exponent) of divergence or contraction of the amplitude of Lyapunov vector is positive. The Lyapunov exponent $`\lambda `$ is defined for each independent direction of the phase space, so the chaotic properties of many-body systems can be characterized by an ordered set of Lyapunov exponents, the so-called Lyapunov spectrum $`\{\lambda ^{(1)},\lambda ^{(2)},\mathrm{}\}`$ where $`\lambda ^{(1)}\lambda ^{(2)}\mathrm{}`$. The Lyapunov spectrum is connected to the contraction rate of the phase space volume (roughly speaking, the dissipation rate) through the sum of all the Lyapunov exponents, allowing the calculation of transport coefficients from the Lyapunov spectra Eva90a ; Gas98 ; Dor99 . The conjugate pairing rule for Lyapunov spectra reduces the calculation of the sum of all the Lyapunov exponents to the sum of just one pair of Lyapunov exponents Dre88 ; Eva90b ; Det96a ; Tan02a . The set of all positive Lyapunov exponents specifies the natural invariant measure Eck85 ; Eck86 , which is used to calculate various quantities using periodic orbit theory Cvi05 and led to the first form of the fluctuation theorem Eva93 ; Gal95 ; Eva02 . On the other hand, recently, much attention has been paid to individual Lyapunov exponents and their Lyapunov vectors for many-body systems. As each Lyapunov exponent has the dimensions of inverse time, the Lyapunov spectrum can be regarded as a time-scale spectrum. From this point of view, Lyapunov exponents with small absolute values are connected to large and macroscopic time-scale behavior of many-body systems, and in this region the wave-like structure of Lyapunov vectors, known as the Lyapunov modes, is observed Pos00 ; Tan02c ; Wij04 ; Yan04 ; Tan03a . The Lyapunov mode is a reflection of a collective movement (phonon mode) of many-body systems, and comes from dynamical conservation laws and translational invariance Tan03a ; Tan04b ; Eck05 ; Tan05a . On the other hand, large Lyapunov exponents are dominated by small and microscopic time-scale movement, and in this region the spatially localized behavior of Lyapunov vector, the so-called Lyapunov localization, appears Man85 ; Kan86 ; Liv89 ; Fal91 ; Pik01 ; Mil02 ; Tan03b . For the largest Lyapunov exponent of many-body chaotic systems, analytical calculations have been attempted Cas96 ; Cas00 ; Zon98 ; Bei00 ; Zon02 . The time-scale separation in many-body systems is crucial to extract a macroscopic dynamics from microscopic many-body dynamics, and the Lyapunov spectrum allows us to discuss it dynamically.
As one of the features of Lyapunov vectors for many-body systems, the Lyapunov localization appears as a behavior in which the Lyapunov vector components for a few particles are significantly larger than the other components. Moreover, the localized region of a Lyapunov vector moves as a function of time. The magnitude of the localization of each Lyapunov vector can be measured quantitatively by the Lyapunov localization spectrum, which is defined as a set of exponential functions of entropy-like quantities for the normalized amplitudes of the Lyapunov vector components Tan03b ; Tan05b ; For05 . We have previously reported that the Lyapunov localization spectra show a bending behavior at low density, and its connection with kinetic theory properties (e.g. the Krylov relation for the largest Lyapunov exponent, and the mean free time being inversely-proportional to density, etc.) is discussed for many-hard-disk systems Tan03b ; Tan05b . However, the Lyapunov localization spectrum requires taking a time-average and characterizes only the static localization of the Lyapunov vectors. The physical meaning of the movement of the localized regions of Lyapunov vectors is not clearly understood.
The principal aim of this paper is to discuss dynamically the movement of the localized region of the Lyapunov vectors in many-hard-disk systems at low density. To discuss this problem we use the fact shown in Ref. Tan03b that at low density only two particle components of the normalized Lyapunov vector for large Lyapunov exponents have a non-zero value. This is due to the short-range of particle interactions in a many-hard-disk system. The movement of the localized region of the Lyapunov vector appears to be a series of jumps or hops, so we can introduce a hopping rate to describe the dynamics. To simplify the problem, in this paper we consider quasi-one-dimensional systems, in which the system width is so narrow that disks always remain in the same order Tan03a ; Tan05a ; Tan03b ; Tan05b . We show that this hopping rate depends on the hopping distance, and is a decreasing function of the hopping distance. This implies that there is a spatial correlation among localized regions of Lyapunov vectors.
We explain the hopping-distance dependence of the hopping rate in many-hard-disk systems in two ways. In the first approach we use a simple model expressed as an accumulation of bricks. Here, the hopping of the localized region of the Lyapunov vectors is expressed as a change in the position of the highest brick site. This model is a one-dimensional version of the so-called clock model, which has been used to calculate Lyapunov exponents for many-hard-disk systems Zon98 ; Bei00 ; Zon02 . In this paper we demonstrate that this model can reproduce the largest Lyapunov exponent for quasi-one-dimensional hard-disk systems. As the second approach to the hopping behavior of the localized Lyapunov vectors, we propose an analytical method to calculate the hopping rate from the sum of probability distributions for the brick height configurations between two separated highest brick sites. Using this analytical approach we can also discuss the relation between the hopping rate and the Lyapunov exponent. Hopping rates calculated by these two approaches are in good agreement with the ones for quasi-one-dimensional hard-disk systems.
The outline of this paper is as follows. In Sec. II we introduce the quasi-one-dimensional hard-disk system, and show the localized behavior of the Lyapunov vector corresponding to the largest Lyapunov exponent at low density. In Sec. III we introduce the hopping rate of the localized region of the Lyapunov vectors, and show the hopping-distance dependence for many-hard-disk systems. In Sec. IV we discuss the hopping rate using a brick accumulation model. In Sec. V we propose an analytical expression for the hopping rate. Section VI is our conclusion and some remarks. In Appendix A we discuss the technical details of the calculation of the hopping rate. In Appendix B we give a microscopic derivation of the brick accumulation model.
## II Quasi-One-Dimensional System and Localization of Lyapunov Vectors
The system which we consider in this paper is a quasi-one-dimensional system consisting of $`N`$ hard-disks in periodic boundary conditions. All of the particles are identical with radius $`R`$ and mass $`M`$, and the shape of the system is rectangular with the length $`L_x`$ and the width $`L_y`$ satisfying the inequality $`2R<L_y<4R`$. The schematic illustration of the system is given in Fig. 1, in which we number particles $`1,2,\mathrm{},N`$ from the left to right in this system. For the actual numerical results shown in this paper, we used: the radius of a particle $`R=1`$, the mass of a particle $`M=1`$, the total energy of the system $`E=N`$, the system width $`L_y=2R(1+10^6)`$, and the system length $`L_x=NL_y(1+d)`$ with the constant $`d`$ controlling the density $`\rho N\pi R^2/(L_xL_y)`$. In the quasi-one-dimensional system, the particle interactions are restricted to nearest-neighbor particles only, so particles remain in the same order. These features require less calculation effort and a simpler representation of results for quasi-one-dimensional system compared with fully two-dimensional systems. The quasi-one-dimensional system has already been used to investigate the localized behavior of Lyapunov vectors Tan03b ; Tan05b , the wave-like structure of Lyapunov vectors Tan03a ; Tan05a ; Tan04b and the transition between quasi-one-dimensional and fully two-dimensional systems For04 .
The dynamics of Lyapunov vectors in many-hard-disk systems is well established, and readers should refer to the references, for example, Ref. Del96 , for more detailed discussions. In many-hard-disk systems the dynamics is separated into a free-flight part and a collision part, and the free-flight part of the dynamics is integrable. This property allow us to express the dynamical evolution as a simple multiplication of time-evolutional matrices for the free-flight dynamics and the collision dynamics, leading to a fast and more accurate numerical simulation than for soft-core interaction models. For numerical calculations of the Lyapunov vectors shown in this paper, we use the algorithm developed by G. Benettin, et al Ben76 and I. Shimada and T. Nagashima Shi79 (also see Refs. Ben80a ; Ben80b ). This algorithm is characterized by intermittent (e.g. after every collision) re-orthogonalization and renormalization of Lyapunov vectors, preventing a divergence of the amplitude of the Lyapunov vectors.
In this paper we use the notation $`\delta 𝜞^{(n)}(t)(\delta 𝜞_1^{(n)}(t),\delta 𝜞_2^{(n)}(t),\mathrm{},\delta 𝜞_N^{(n)}(t))`$ for the Lyapunov vector corresponding to the $`n`$-th Lyapunov exponent $`\lambda ^{(n)}`$ at time $`t`$. Here, $`\delta 𝜞_j^{(n)}(t)`$ is the Lyapunov vector component contributed by the $`j`$-th particle in the $`n`$-th Lyapunov exponent at time $`t`$. To express the localized behavior of the Lyapunov vectors we introduce the quantity $`\gamma _j^{(n)}(t)`$ as
$`\gamma _j^{(n)}(t){\displaystyle \frac{\left|\delta 𝜞_j^{(n)}(t)\right|^2}{\underset{k=1}{\overset{N}{}}\left|\delta 𝜞_k^{(n)}(t)\right|^2}},`$ (1)
which is the normalized amplitude of the Lyapunov vector component for the $`j`$-th particle for the $`n`$-th Lyapunov exponent $`\lambda ^{(n)}`$ at time $`t`$. The localized behavior of Lyapunov vector, namely the Lyapunov localization, is the phenomenon where only a few of the $`\gamma _j^{(n)}(t)`$, $`j=1,2,\mathrm{},N`$ have a non-zero value at any time $`t`$.
Fig. 2 is an example of the Lyapunov localization. It is a graph of the normalized amplitude $`\gamma _j^{(1)}`$ of the Lyapunov vector components for the $`j`$-th particle corresponding to the largest Lyapunov exponent $`\lambda ^{(1)}`$ as a function of the particle index $`j`$ and the collision number $`n_t`$ ($`t/\tau `$ with time $`t`$ and the mean free time $`\tau `$). The system is a quasi-one-dimensional system consisting of $`50`$ hard-disks and with $`d=10^5`$ (the density $`\rho 7.85\times 10^6`$). From this figure we recognize that the non-zero components of the Lyapunov vector are concentrated at two nearest-neighbor particles. This characteristic is shown for hard-disk systems at low density and for large Lyapunov exponents in general, using the Lyapunov localization spectrum Tan03b ; Tan05b . Moreover, we observe that such localized regions of Lyapunov vector components move with time in discrete jumps or hops, in Fig. 2. This characteristic hopping has already been shown in Ref. Tan03b and that such spatial hopping of the non-zero components of $`\gamma _j^{(n)}(t)`$ are caused by some particle-particle collisions. However, not every collision causes a hopping of the localized Lyapunov vector. In this sense, particle collisions themselves are not sufficient to explain the hopping movement of the Lyapunov localization.
The Lyapunov localization appears in the Lyapunov vectors corresponding to large Lyapunov exponents in general, but in this paper for simplicity we consider only the Lyapunov vector $`\delta 𝜞^{(1)}`$ corresponding to the largest Lyapunov exponent $`\lambda ^{(1)}`$.
## III Hopping Rate of Localized Lyapunov Vectors
An advantage of the quasi-one-dimensional system in investigations of Lyapunov localization is that the movement of particles in the transverse direction are suppressed, and roughly speaking, the particle sequence corresponds to the particle’s position. Noting this feature, in this section we describe the hopping behavior of spatially localized Lyapunov vectors as the hopping of particle indices whose $`\gamma _j^{(n)}`$ defined by Eq. (1) have non-zero values.
As shown in Refs. Tan03b ; Tan05b the particle indices with non-zero $`\gamma _j^{(1)}`$ are a pair of nearest-neighbor particles in the low density limit, and we can introduce the hopping distance $`h`$ at each collision as the change of particle indices. In this definition the hopping distance $`h`$ is an integer number satisfying the inequality $`[N/2]h[N/2]`$, where $`[x]`$ is the integer part of the real number $`x`$. It should also be noted that we take the particle index $`j`$ as the one equivalent to the index $`j\pm N`$, because of periodic boundary conditions, so the hopping distance $`h\pm N`$ is equivalent to $`h`$. More technical details of the calculation of the hopping distance $`h`$ are given in Appendix A. Using this hopping distance $`h`$, in numerical simulations, we count the number $`N_T(h)`$ of hops with hopping distance $`h`$ in a time-interval $`T`$, and we define the hopping rate $`P_N(h)`$ as a function proportional to $`N_T(h)`$ as $`T\mathrm{}`$. In this paper we use the hopping rate $`P_N(h)/P_N(1)=lim_T\mathrm{}N_T(h)/N_T(1)`$ normalized by $`P_N(1)`$. From the reflection symmetry of the quasi-one-dimensional system in the longitudinal direction, the hopping rate $`P_N(h)`$ must be symmetric, namely $`P_N(h)=P_N(h)`$. We use this hopping rate to quantitatively discuss the hopping dynamics for localized Lyapunov vectors.
Figure 3 shows log-log plots of the hopping rates $`P_N(h)/P_N(1)`$ normalized by $`P_N(1)`$ as a function of $`|h|`$ in a quasi-one-dimensional hard-disk system with $`d=10^3`$ at different numbers of particles; $`N=25`$ (circles), $`50`$ (triangles), and $`75`$ (squares). Noting the symmetric property $`P_N(h)=P_N(h)`$, we use $`|P_N(h)P_N(h)|/P_N(1)`$ as error bars in Fig. 3. It is clear from Fig. 3 that the hopping rate $`P_N(|h|)`$ decreases as $`|h|`$ increases. This implies that there is a spatial correlation among the localized regions of a Lyapunov vector, rather than a random hopping.
The turn-up in the tail of the normalized hopping rate $`P_N(h)/P_N(1)`$ in Fig. 3 can be explained as an effect of periodic boundary conditions. Under periodic boundary conditions the hopping distance $`h+jN`$ for any integer $`j`$ is observed as the hopping distance $`h`$ in the quasi-one-dimensional system consisting of $`N`$ hard disks. Therefore, using the hopping rate $`P_{\mathrm{}}(h)`$ for the thermodynamics limit ($`N\mathrm{}`$ at a fixed density) for $`h=0,\pm 1,\pm 2,\mathrm{}\pm \mathrm{}`$, the hopping rate $`P_N(h)`$ for a finite system should be represented as
$`P_N(h)={\displaystyle \underset{j=\mathrm{}}{\overset{\mathrm{}}{}}}P_{\mathrm{}}(h+jN)`$ (2)
for $`h=[N/2],[N/2]+1,\mathrm{},[N/2]`$. The terms on the right-hand side of Eq. (2) with $`j0`$ cause the turn-up in the tail of the hopping rate $`P_N(h)`$ for a finite system. In this sense we can explain the $`N`$-dependence of the hopping rate $`P_N(h)`$ using an $`N`$-independent asymptotic distribution $`P_{\mathrm{}}(h)`$. To make this explanation convincing we fitted the numerical data for $`P_N(h)/P_N(1)`$ shown in Fig. 3 to the function $`f(h)\alpha [h^\beta +_{j=1}^2(|h+jN|^\beta +|hjN|^\beta )]`$ assuming that the decay of $`P_{\mathrm{}}(h)=\alpha |h|^\beta `$ with fitting parameters $`\alpha `$ and $`\beta `$, and neglecting the higher order small terms for $`|j|3`$. Here, the fitting lines are dotted for $`N=25`$, the broken for $`N=50`$ and solid for $`N=75`$.
Note that these sets $`(\alpha ,\beta )`$ of fitting parameter values, summarized in Table 1, are almost independent of the number of particles $`N`$. As shown in Fig. 3, these fitting lines nicely reproduce the values of the numerical data, especially the turn-ups.
We can also calculate the hopping rate $`P_N(h)`$ at different densities, as long as the density is low enough so that only a few of the normalized Lyapunov vector component amplitudes $`\{\gamma _j^{(n)}\}_j`$ have non-zero values. In Appendix A we show that the normalized hopping rate $`P_N(h)/P_N(1)`$ is almost density-independent at least for $`10^3<d<10^5`$, (namely in the density region $`7.85\times 10^6<\rho <7.85\times 10^4`$). However, a subtle increase of the normalized hopping rate as the density decreases is recognizable in this low density region.
## IV Brick Accumulation Model
As shown in Sec. III the hopping rate of the localized region of the Lyapunov vector is correlated spatially. In this section we explain this characteristic using a simple one-dimensional model, which we call the *brick accumulation model*.
A schematic illustration of the brick accumulation model, or simply the brick model, is given in Fig. 4. It is a one-dimensional model with $`N`$ sites in the horizontal direction. The bricks are dropped at random and occupy a pair of neighboring sites. Each brick has width two and the height $`1`$ and they accumulate at sites without overlap. The dynamics of the brick model is described using the brick height $`𝒦_j(n)`$ at the $`j`$-th site after $`n`$ bricks have been dropped. The total brick height $`𝒦_j(n)`$ takes on integer values and its dynamics is expressed as follows. If the $`n`$-th brick is dropped on sites $`j_n`$ and $`j_n+1`$ then
$`𝒦_l(n)=\{\begin{array}{c}\text{max}\{𝒦_{j_n}(n1),𝒦_{j_n+1}(n1)\}+1\hfill \\ \text{if}l\{j_n,j_n+1\}\hfill \\ \\ 𝒦_l(n1)\text{if}l/\{j_n,j_n+1\}\hfill \end{array}`$ (7)
(noting that the particle index $`N+1`$ is equivalent to $`1`$). Here the site number $`j_n`$ is chosen randomly on $`[1,N]`$ for each $`n`$.
One may notice that the brick accumulation model is a one-dimensional version of the so-called *clock model* Zon98 ; Bei00 ; Zon02 , which has been used to calculate the Lyapunov exponents for many-hard-disk systems. In the clock model, the brick height $`𝒦_j(n)`$ is usually called the clock value of the $`j`$-th particle after the $`n`$-th particle-particle collision, and the dynamics (7) is expressed as an adjustment of the clock values of colliding particles. However we use the name brick accumulation model for the one-dimensional version of the clock model, because in a simple image of brick accumulations, it is easily to visualize the configuration of brick heights which is essential for the analytical approach to the hopping rate discussed in the next section V, although this image is applicable only for the one-dimensional case. The image of brick accumulations also helps us to easily recognize the similarity of this model to ballistic aggregation models, whose scaling properties have been studied analytically and numerically Fra00 ; Rom02 .
In the brick model described by the brick height $`𝒦_j(n)`$, the site number $`j`$ corresponds to the particle index, and the number $`n`$ of dropped bricks corresponds to the collision number $`n_t`$ for the quasi-one-dimensional hard-disk system. Moreover, the brick height $`𝒦_j(n)`$ itself is connected to the Lyapunov vector component amplitude $`|\delta 𝜞_j^{(1)}(t)|`$ by
$`\left|\delta 𝜞_j^{(1)}(\tau n_t)\right|\left|\delta 𝜞_j^{(1)}(0)\right|\mathrm{exp}\left\{𝒦_j(n_t)\mathrm{ln}\rho \right\}`$ (8)
for the particle index $`j`$ and the collision number $`n_t`$ with the mean free time $`\tau `$ in the asymptotic limit of low density $`\rho 0`$. The derivation of the relation (8) from microscopic dynamics for hard-particle systems is given in Appendix B.
From the relation (8) between the brick height and the Lyapunov vector, the amplitude of the localized Lyapunov vector components with the highest brick height have much larger values than the other components because of the huge factor $`\mathrm{ln}\rho `$ appearing in the exponent of Eq. (8) at low density $`\rho <<1`$. Therefore, to a good approximation, the normalized amplitude $`\gamma _j^{(n)}`$ given by Eq. (1) has non-zero components only for particles corresponding to these largest amplitudes, as shown in Fig. 2. The combination of Eqs. (7) and (8) also explains why the amplitudes of the localized components of the Lyapunov vectors for nearest-neighbor particles take almost the same value, thus appearing as a flat top in the localized regions in Fig. 2. After all, the localized region of the Lyapunov vector represented in Eq. (8) is given as the site indices whose brick height $`𝒦_j(n_t)`$ take the maximum value after the $`n`$-th brick is dropped, and such sites with the highest brick height can be calculated using the brick model dynamics (7) only, without referring to the hard-disk dynamics. Based on these features we can calculate the hopping distance $`h`$ for the brick accumulation model (see the end of Appendix A for more detail of the hopping distance for the brick model), and therefore the hopping rate $`P_N(h)`$, similarly to the quasi-one-dimensional hard-disk system.
Figure 5 shows the normalized hopping rates $`P_N(h)/P_N(1)`$ for the brick model with $`50`$ sites (circles) and for a quasi-one-dimensional hard-disk system consisting of $`50`$ hard-disks and $`d=10^5`$ (squares). Agreement between the brick model and the quasi-one-dimensional hard-disk system for the normalized hopping rate $`P_N(h)/P_N(1)`$ for $`|h|1`$ is satisfactory. Numerical simulations of quasi-one-dimensional hard-disk systems show that the normalized hopping rate $`P_N(h)/P_N(1)`$ increases very slightly as the density decreases as shown in Appendix A, so the deviation of $`P_N(h)/P_N(1)`$ in tails in the brick model and the hard-disk system in Fig. 5 may be a finite density effect.
So far we have discussed the hopping rate $`P_N(h)`$ for non-zero hopping distances $`h0`$. Now we discuss the hopping rate $`P_N(0)`$ for zero hopping distance $`h=0`$. Table 2 shows the ratio $`P_N(0)/P_N(1)`$ for the hopping rates at $`h=0`$ and $`1`$ in the quasi-one-dimensional hard-disk system with $`d=10^3`$ and the brick accumulation model for $`N=25,50,75`$ and $`100`$. As shown in this table, the ratio $`P_N(0)/P_N(1)`$ depends on the number of particles $`N`$, and roughly speaking it is proportional to $`N`$. To explain this linear dependence of $`P_N(0)/P_N(1)`$ with respect to $`N`$, we note that in the brick model the hopping distance $`1`$ occurs only when a new brick is dropped at one (and not both) of the sites with the highest brick height, meaning that the hopping rate $`P_N(1)`$ should be approximately inversely proportional to $`N`$ . On the other hand, the hopping distance $`0`$ occurs when a new brick is dropped at a site which does not have the highest nor the second highest brick height, or is dropped at both the sites which are currently the site of the highest brick height, meaning that the hopping rate $`P_N(0)`$ should be, roughly speaking, independent of $`N`$. These considerations for $`P_N(0)`$ and $`P_N(1)`$ explain why the ratio $`P_N(0)/P_N(1)`$ is proportional to $`N`$. However we need a more detailed consideration to explain the difference between the coefficient $`\kappa `$ in the relation $`P_N(0)/P_N(1)\kappa N`$ between the quasi-one-dimensional hard-disk system and the brick accumulation model.
Before finishing this section, we discuss one more property of the brick accumulation model, which we will use in the next section. Fig. 6 is the sum $`_{j=1}^N𝒦_j(n)`$ of brick heights as a function of $`n`$. Note that this sum must increase at a speed of more than 2 per dropped brick, because accumulated bricks capture spaces below them which cannot be occupied by further dropped bricks. It is clear from this figure that this sum increases linearly, and a fit of the data to the function $`_{j=1}^N𝒦_j(n)=\alpha +\beta n`$ with fitting parameters $`\alpha `$ and $`\beta `$ leads to the values $`\alpha 0.00`$ and $`\beta 3.99`$ noteSum . (Note that data points in the main graph of Fig. 6 look exactly like this fit because of the large scale. In the inset to Fig. 6 we show the graph of $`_{j=1}^N𝒦_j(n)`$ as a function of $`n`$ on a much smaller scale, to show its fluctuating behavior.) Therefore, on average, each dropped brick adds a contribution of 4 to the sum $`_{j=1}^N𝒦_j(n)`$, meaning that each brick occupies not only its own $`2`$ spaces but also captures $`2`$ empty spaces below it.
It is important to note that from Eq. (8) the brick height $`𝒦_j(n)`$ dominates the exponential growth rate of the Lyapunov vector component amplitude. Using this feature of the brick heights, their sum $`_{j=1}^N𝒦_j(n)`$ is connected to the largest Lyapunov exponent $`\lambda ^{(1)}`$ as
$`\lambda ^{(1)}\underset{n+\mathrm{}}{lim}{\displaystyle \frac{1}{n\tau }}\left[{\displaystyle \frac{1}{N}}{\displaystyle \underset{j=1}{\overset{N}{}}}𝒦_j(n)\right]\mathrm{ln}\rho `$ (9)
with the mean free time $`\tau `$ and the density $`\rho `$ in the asymptotic limit of low density. \[More detailed discussion for the formula (9) is given in Appendix B.\] As a numerical check of the formula (9) we show in Fig. 7, the largest Lyapunov exponent $`\lambda ^{(1)}`$ as a function of the density $`\rho `$ in a quasi-one-dimensional hard-disk systems and in the brick accumulation model using the formula (9) with $`N=50`$. Here, to calculate the largest Lyapunov exponent from Eq. (9) we used the relation $`_{j=1}^N𝒦_j(n)4n`$ as shown in Fig. 6, and the values of the mean free time $`\tau `$ and the density $`\rho `$ of the quasi-one-dimensional system whose Lyapunov exponents are plotted in Fig. 7, and the data points are connected by a dashed line for ease of visibility. Figure 7 shows that Eq. (9) reproduces successfully the values of the largest Lyapunov exponent for the quasi-one-dimensional hard-disk systems, not only in the limit of low density but also at relatively high density such as $`\rho <0.3`$. It may be noted that a linear dependence of the sum $`_{j=1}^N𝒦_j(n)`$ of brick heights with respect to $`n`$ is necessary to get a finite value of the largest Lyapunov exponent $`\lambda ^{(1)}`$ by Eq. (9).
## V Analytical Expression for the Hopping Rate
In this section we discuss another approach to the hopping rate of the localized region of the Lyapunov vectors. This approach is inspired by some of the characteristics of the brick accumulation model discussed in Sec. IV, although we greatly simplified the brick model by omitting some other aspects of the model. The advantage of this approach is that we can get an analytical expression for the hopping rate, while to a good approximation it still reproduces the hopping rate for hard-disk systems. Another important point in this approach is that it connects the dynamics of the localized region of the Lyapunov vectors with a static property, the probability distribution of brick-height differences between nearest-neighbor sites.
In the brick accumulation model, a hop of the highest brick site occurs when two (non-nearest-neighbor) sites have the same brick height. Noting this characteristic we consider the probability distribution $`\stackrel{~}{P}_{\mathrm{}}(h)`$ under the constraint that the two sites $`\mu `$ and $`\mu +h`$ have the highest brick height. We require that there is no other highest site between the two highest sites $`\mu `$ and $`\mu +h`$ for $`|h|2`$;
$`𝒦_l(n)<𝒦_\mu (n)[=𝒦_{\mu +h}(n)],`$
$`\text{in}l=\{\begin{array}{c}\mu +1,\mu +2,\mathrm{},\mu +h1\hfill \\ \text{for}h2\hfill \\ \mu +h+1,\mu +h+3,\mathrm{},\mu 1\hfill \\ \text{for}h2,\hfill \end{array}`$ (14)
Using this probability distribution $`\stackrel{~}{P}_{\mathrm{}}(h)`$, we can estimate the hopping rate $`P_{\mathrm{}}(h)`$ in the thermodynamic limit as the one proportional to $`\stackrel{~}{P}_{\mathrm{}}(h)`$: $`P_{\mathrm{}}(h)\stackrel{~}{P}_{\mathrm{}}(h)`$. Then we can calculate the hopping rate $`P_N(h)`$ for a finite size system using Eq. (2), apart from a constant factor.
To calculate the probability distribution $`\stackrel{~}{P}_{\mathrm{}}(h)`$ we introduce the distribution $`\mathrm{\Lambda }(k)`$ of brick-height differences $`k`$ between nearest-neighbor sites. These two distribution functions are connected by
$`\stackrel{~}{P}_{\mathrm{}}(h)`$ $`=`$ $`{\displaystyle \underset{\begin{array}{c}k_1\\ (k_11)\end{array}}{}}{\displaystyle \underset{\begin{array}{c}k_2\\ (k_1+k_21)\end{array}}{}}\mathrm{}{\displaystyle \underset{\begin{array}{c}k_{|h|1}\\ (_{j=1}^{|h|1}k_j1)\end{array}}{}}`$
$`\times \mathrm{\Lambda }(k_1)\mathrm{\Lambda }(k_2)\mathrm{}\mathrm{\Lambda }(k_{|h|1})\mathrm{\Lambda }\left({\displaystyle \underset{j=1}{\overset{|h|1}{}}}k_j\right).`$
In Eq. (LABEL:passInte1), $`k_l`$ is a brick-height difference of nearest neighbor sites ($`l=1,2,\mathrm{},|h|1`$). The probability distribution for the specific brick-height configuration with $`k_l`$, $`l=1,2,\mathrm{},|h|1`$ is given by $`\mathrm{\Lambda }(k_1)\mathrm{\Lambda }(k_2)\mathrm{\Lambda }(k_3)\mathrm{}\mathrm{\Lambda }(k_{|h|1})\mathrm{\Lambda }(_{j=1}^{|h|1}k_j)`$, and the hopping rate $`\stackrel{~}{P}_{\mathrm{}}(h)`$ is calculated as the summation of this probability distribution over possible values of $`k_l`$, $`l=1,2,\mathrm{},h1`$, leading to Eq. (LABEL:passInte1). Figure 8 is a schematic illustration of a brick-height configuration connecting two highest sites and the brick-height differences $`k_j`$ of nearest neighbor sites in the case of $`h2`$. In Eq. (LABEL:passInte1), the lower bounds of the sums on the right-hand side of Eq. (LABEL:passInte1) come from the condition (14). Note that from Eq. (LABEL:passInte1) the distribution $`\stackrel{~}{P}_{\mathrm{}}(h)`$ is an even function of $`h`$, namely $`\stackrel{~}{P}_{\mathrm{}}(h)=\stackrel{~}{P}_{\mathrm{}}(h)`$.
Now, for simplicity, we assume that the distribution function $`\mathrm{\Lambda }(k)`$ can be expressed as an exponential function
$`\mathrm{\Lambda }(k)=𝒲\mathrm{exp}\{\eta |k|\}`$ (16)
with constants $`𝒲`$ and $`\eta (>0)`$. (We will discuss the validity of this assumption later.) Inserting Eq. (16) into Eq. (LABEL:passInte1), and replacing the sums over $`k_1`$ in Eq. (LABEL:passInte1) with the ones over $`k_1^{}k_1+1`$, we obtain
$`\stackrel{~}{P}_{\mathrm{}}(h)=\{\begin{array}{c}𝒲^2\mathrm{{\rm Y}}\underset{k=0}{\overset{+\mathrm{}}{}}\mathrm{{\rm Y}}^k\text{for}|h|=2\hfill \\ \\ 𝒲^{|h|}\mathrm{{\rm Y}}\underset{k_1=0}{\overset{+\mathrm{}}{}}\mathrm{{\rm Y}}^{k_1}\underset{k_2=k_1}{\overset{+\mathrm{}}{}}\mathrm{{\rm Y}}^{k_2\theta (k_2)}\underset{k_3=k_1k_2}{\overset{+\mathrm{}}{}}\mathrm{{\rm Y}}^{k_3\theta (k_3)}\hfill \\ \mathrm{}\underset{k_{|h|1}=\underset{j=1}{\overset{|h|2}{}}k_j}{\overset{+\mathrm{}}{}}\mathrm{{\rm Y}}^{k_{|h|1}\theta (k_{|h|1})}\hfill \\ \text{for}|h|3\hfill \end{array}`$ (22)
where $`\mathrm{{\rm Y}}`$ is defined by
$`\mathrm{{\rm Y}}\mathrm{exp}\{2\eta \}`$ (24)
and $`\theta (x)`$ is the Heaviside function taking the value $`1`$ for $`x>0`$ and the value $`0`$ for $`x0`$. The summations appearing in Eq. (LABEL:passInte2) can be carried out successively, and we obtain
$`\stackrel{~}{P}_{\mathrm{}}(h)={\displaystyle \frac{𝒲^{|h|}\mathrm{{\rm Y}}}{(1\mathrm{{\rm Y}})^{|h|1}}}\mathrm{\Omega }(|h|)`$ (25)
where the function $`\mathrm{\Omega }(k)`$ of $`k`$ is can be written as
$`\mathrm{\Omega }(2)`$ $`=`$ $`1,`$ (26)
$`\mathrm{\Omega }(3)`$ $`=`$ $`1+\mathrm{{\rm Y}},`$ (27)
$`\mathrm{\Omega }(4)`$ $`=`$ $`1+3\mathrm{{\rm Y}}+\mathrm{{\rm Y}}^2,`$ (28)
$`\mathrm{\Omega }(5)`$ $`=`$ $`(1+\mathrm{{\rm Y}})(1+5\mathrm{{\rm Y}}+\mathrm{{\rm Y}}^2),`$ (29)
$`\mathrm{\Omega }(6)`$ $`=`$ $`1+10\mathrm{{\rm Y}}+20\mathrm{{\rm Y}}^2+10\mathrm{{\rm Y}}^3+\mathrm{{\rm Y}}^4,`$ (30)
$`\mathrm{\Omega }(7)`$ $`=`$ $`(1+\mathrm{{\rm Y}})(1+14\mathrm{{\rm Y}}+36\mathrm{{\rm Y}}^2+14\mathrm{{\rm Y}}^3+\mathrm{{\rm Y}}^4),`$ (31)
and so on. Eq. (25), with Eqs. (26), (27), (28), (29), (30), (31), etc., gives an analytical expression for the hopping rate, for example, using the relation $`P_{\mathrm{}}(h)/P_{\mathrm{}}(1)=\stackrel{~}{P}_{\mathrm{}}(h)/\stackrel{~}{P}_{\mathrm{}}(1)`$.
The coefficients $`𝒲`$ and $`\eta `$ appearing in Eq. (16) can be determined from the two sum rules:
$`{\displaystyle \underset{k=\mathrm{}}{\overset{+\mathrm{}}{}}}\mathrm{\Lambda }(k)`$ $`=`$ $`1,`$ (32)
$`{\displaystyle \underset{k=\mathrm{}}{\overset{+\mathrm{}}{}}}|k|\mathrm{\Lambda }(k)`$ $`=`$ $`\mathrm{\Delta }𝒦,`$ (33)
where $`\mathrm{\Delta }𝒦`$ is the mean value of the absolute value of the brick height difference between nearest-neighbor sites. The first condition (32) is the normalization of the probability distribution $`\mathrm{\Lambda }(k)`$ which leads to
$`𝒲={\displaystyle \frac{\mathrm{exp}\{\eta \}1}{\mathrm{exp}\{\eta \}+1}}.`$ (34)
Using Eq. (34) the second condition (33) gives
$`\eta =\mathrm{ln}\left\{{\displaystyle \frac{1+\sqrt{1+\mathrm{\Delta }𝒦^2}}{\mathrm{\Delta }𝒦}}\right\},`$ (35)
satisfying the inequality $`\eta >0`$. Inserting Eq. (35) into Eqs. (24) and (34) we obtain
$`\mathrm{{\rm Y}}`$ $`=`$ $`\left({\displaystyle \frac{\mathrm{\Delta }𝒦}{1+\sqrt{1+\mathrm{\Delta }𝒦^2}}}\right)^2`$ (36)
$`𝒲`$ $`=`$ $`{\displaystyle \frac{1+\sqrt{1+\mathrm{\Delta }𝒦^2}\mathrm{\Delta }𝒦}{1+\sqrt{1+\mathrm{\Delta }𝒦^2}+\mathrm{\Delta }𝒦}}.`$ (37)
From Eqs. (36) and (37), there is only one parameter $`\mathrm{\Delta }𝒦`$ remaining to determine the hopping rate using Eq. (25).
To estimate the value of $`\mathrm{\Delta }𝒦`$, we use a property of the brick accumulation model discussed at the end of Sec. IV. Previously, we showed an approximate relation $`_{j=1}^N𝒦_j(n)4n`$ for the sum of brick heights, which means that in the brick accumulation model each dropped brick gives a mean contribution of 4 to this sum. This contribution consists of a contribution of 2 as the space occupied by a brick itself and another 2 as empty space below the brick which can now not be occupied by other bricks. This implies that the averaged brick-height difference between nearest-neighbor sites is about 2, so that
$`\mathrm{\Delta }𝒦2.`$ (38)
We use this value to calculate the hopping rate based on Eq. (25). One may notice that from the relation $`_{j=1}^N𝒦_j(n)(\mathrm{\Delta }𝒦+2)n`$ and the formula (9) we obtain
$`\mathrm{\Delta }𝒦{\displaystyle \frac{N\tau }{\mathrm{ln}\rho }}\lambda ^{(1)}2`$ (39)
in the limit of low density. From the relation (39) between the parameter $`\mathrm{\Delta }𝒦`$ specifying the hopping rate and the largest Lyapunov exponent $`\lambda ^{(1)}`$, the hopping rate of the localized region of the Lyapunov vectors is connected to the largest Lyapunov exponent.
Before comparing the hopping rate based on the analytical expression (25) with the ones for hard-disk systems and the brick accumulation model, we note that this analytical expression for the hopping rate may not be appropriate for small hopping distances $`|h|`$. In particular, it does not give the correct value of $`P_N(1)`$, because in the brick model the hopping distance $`h=\pm 1`$ does not appear from separated highest sites with the same brick height, the assumption used to derive Eq. (25). Therefore it is not appropriate to calculate the hopping rate normalized by $`P_N(1)`$ from this analytical expression and to compare it with the numerical results. The hopping rate $`P_N(\pm 2)`$ from Eq. (25) may also be problematic, because in the brick model the hopping distance $`h=\pm 2`$ occurs when 4 consecutive sites have the highest brick height, while to derive Eq. (25) we assumed that $`h=\pm 2`$ occurs when non-consecutive separate sites $`\mu `$ and $`\mu +h`$ have the same highest brick height. Based on these considerations we do not calculate the value $`P_N(\pm 1)`$ from the analytical approach in this section, and plot the hopping rate so that $`P_N(5)`$ by Eq. (25) gives the same value as that from the brick model.
In Fig. 5 for the normalized hopping rate $`P_N(h)/P_N(1)`$, we plotted the function $`\mathrm{\Psi }(h)\stackrel{~}{P}_N(h)P_N(5)/[\stackrel{~}{P}_N(5)P_N(1)]`$ using the value $`P_N(5)/P_N(1)`$ of the brick model with $`\stackrel{~}{P}_N(h)=\stackrel{~}{P}_{\mathrm{}}(h)+\stackrel{~}{P}_{\mathrm{}}(Nh)`$ using the analytical expression $`\stackrel{~}{P}_{\mathrm{}}(h)`$ given by Eq. (25) for $`|h|2`$. Note $`\mathrm{\Psi }(5)=P_N(5)/P_N(1)`$ so that $`\mathrm{\Psi }(h)`$ coincides with $`P_N(h)/P_N(1)`$ of the brick model at $`|h|=5`$. Here, the value of the hopping rate values are given at integer values of $`h`$, but we connect them with a broken line for ease of visibility. Figure 5 shows that the analytical expression (25) for the hopping rate reproduces the hopping rate for the brick accumulation model as well as the quasi-one-dimensional hard-disk system to a good approximation. It may be noted that for this plot we used the first two dominant terms on the right-hand side of Eq. (2), and part of the small deviation of the hopping rate in the tail between the brick model and the analytical expression should come from the omission of higher order terms in Eq. (2).
We notice that the approach in this section is simple enough to get an analytical expression for the hopping rate but it is not completely consistent with the brick accumulation model discussed in the previous section IV. Previously, we have mentioned the irrelevance of Eq. (25) as a description of a small hopping rate in the brick model. Actually, in the approach of this section we omitted the essential characteristic of the bricks as components of the brick accumulation model, except for the property (38), and treated the model components as blocks (or half bricks). As another example, we show in Fig. 9 the numerical result for the distribution $`\mathrm{\Lambda }(k)`$ of brick-height differences $`k`$ between nearest-neighbor sites in the brick accumulation model with $`N=50`$. Here, $`\mathrm{\Lambda }(k)`$ is normalized by $`_{k=2N}^{2N}\mathrm{\Lambda }(k)=1`$, instead of Eq. (32), because we cannot calculate $`\mathrm{\Lambda }(k)`$ in $`|k|+\mathrm{}`$ numerically. In this figure we added the exponential distribution (16) using Eqs. (36), (37) and (38). Figure 9 shows that the probability distribution $`\mathrm{\Lambda }(k)`$ does not coincide with an exponential distribution (16), although it may be justified as a first approximation. (On the other hand, the numerical evaluation of $`\mathrm{\Delta }𝒦`$ from the distribution $`\mathrm{\Lambda }(k)`$ as a numerical result in Fig. 9 gives the value $`1.99`$, then Eq. (38) is still justified.) On another point, in the brick model the highest brick sites appear as a pair of nearest-neighbor sites, but we did not take into account this characteristic in the analytical approach in this section. Despite the omission of characteristics of the brick accumulation model, the analytical approach in this section reproduces the hopping rate for many-hard-disk systems reasonably well, and it suggests that this approach still keeps enough of the essential characteristics that describe the dynamics of the Lyapunov localization.
## VI Conclusion and Remarks
In this paper we discussed the dynamics of the spatially localized region of the Lyapunov vector corresponding to the largest Lyapunov exponent in a quasi-one-dimensional hard-disk system. To discuss the dynamics of the localized region of the Lyapunov vector we introduced a hopping rate for the localized region, and showed that the hopping rate decreases as the absolute value of hopping distance increases. This hopping-distance dependence of the hopping rate was explained quantitatively in two ways: a brick accumulation model and an analytical approach. In the brick accumulation model, the hopping behavior of the localized Lyapunov vectors was explained as the movement of the highest position in the brick accumulations. It was shown that using this brick model we can calculate the largest Lyapunov exponent for quasi-one-dimensional hard-disk systems successfully. On the other hand, in the analytical approach the hopping rate was calculated from probability distributions for brick height differences of nearest neighbor sites via multiple summations over possible configurations that connect two separated highest sites. The result is related to the largest Lyapunov exponent. Both of the approaches successfully reproduced the hopping-distance dependence of the hopping rate for the localized Lyapunov vectors of quasi-one-dimensional hard-disk systems.
As a remark, it may be useful to mention a previous conjecture for the dynamics of the localized region of Lyapunov vectors corresponding to largest Lyapunov exponent. Before we started this work, there had been a view that the origin of the hopping behavior of localized Lyapunov vectors was:
The localized region of the Lyapunov vector hops to the position of a new grazing collision.
This conjecture was suggested from the fact that a change of Lyapunov vectors in particle collisions may be dominated by the factor $`1/(𝝈\mathrm{\Delta }𝒑)`$ in the collision dynamics for Lyapunov vector (see Appendix B), where $`𝝈`$ is the normalized collision vector and $`\mathrm{\Delta }𝒑`$ is the momentum difference of the colliding particles before the collision. In this argument, if two particles collide at a small angle (a grazing collision) with a small value of $`𝝈\mathrm{\Delta }𝒑`$, then the Lyapunov vector can change significantly and the position of the localized region may move. However, this scenario cannot be correct for the following reasons. First, this conjecture implies that as the position at which a grazing collision occurs is random, the hopping distance should also be random. This contradicts our numerical result that the hopping rate decreases with increasing hopping distance, as shown in Fig. 3. Second, we show in Fig. 10, the distribution $`D(𝝈\mathrm{\Delta }𝒑)`$ of $`𝝈\mathrm{\Delta }𝒑`$ in general collisions (dashed line), and the distribution of collisions causing hopping of the localized Lyapunov vector (solid line) noteB . It is clear from Fig. 10 that the quantity $`𝝈\mathrm{\Delta }𝒑`$ is not smaller in collisions which lead to jumps of the localized region of Lyapunov vector compared with the general case. It may also be noted that the hopping dynamics of the Lyapunov vector from the brick accumulation model described in Sec. IV is independent of collision parameters like collision angles, the momentum difference of colliding particles, and also the quantity $`𝝈\mathrm{\Delta }𝒑`$. This also suggests that the conjecture for the origin of the hopping of the localized region of the Lyapunov vector cannot be justified. On the other hand, this collision parameter independence of the hopping behavior in the brick model cannot explain why the distribution $`D(𝝈\mathrm{\Delta }𝒑)`$ is different from the general case and the hopping case in Fig. 10. This is an open problem.
As another open problem, the fits in Fig. 3 suggest that the hopping-distance dependence of the hopping rate in the thermodynamic limit seems to be a power law: $`P_{\mathrm{}}(h)h^\beta `$ with $`\beta 1.7`$. It remains to be determined whether this power behavior of the hopping rate can be justified analytically in the brick accumulation model or the analytical approach discussed in Sec. V.
In this paper we showed that the hopping rate of the localized region of a Lyapunov vector, and the largest Lyapunov exponent, are well described by the brick accumulation model. Then, one may ask a more direct numerical check of the justification of the brick accumulation model in quasi-one-dimensional hard-disk systems, for example, to check Eq. (8) or to observe numerically an actual brick configuration like that presented in Fig. 4. However, such a check of the brick model is not trivial for the following reasons. First, the brick accumulation model is only justified in the limit of low density, but numerical simulations have to be at some finite density. This effect appears, for example, as gradual changes of the Lyapunov vector component amplitudes, as shown in Fig. 11, which do not appear in the brick accumulation model. Second, even if we could simulate at an extremely low density in which such finite density effects can be neglected (although the case presented in this paper is not at such a low density), the factor $`\mathrm{ln}\rho `$ in Eq. (8) may be too large for an actual numerical calculation. Finally, to calculate Lyapunov vectors in this paper we used the algorithm developed by Benettin et al Ben80a ; Ben80b . This algorithm includes intermittent renormalizations of Lyapunov vectors, to prevent a divergence of the amplitudes of Lyapunov vectors in numerical calculations, but the brick accumulation model does not have such a normalization procedure in its dynamics. This difference makes a direct numerical check of Eq. (8) difficult in hard-disk systems. Different from Eq. (8) itself, the localized region of the Lyapunov vectors, which is required to calculate the hopping rate, is given simply by particle indices with the largest Lyapunov vector component amplitude, which is not influenced by such a difference of normalization procedure in calculations of Lyapunov vectors.
In this paper, in order to introduce the hopping rate we used the property of quasi-one-dimensional systems, that the order of particles is an invariant. In this sense, it is not trivial to generalize our argument to fully two- (or three-) dimensional systems. An effective way to describe the dynamics for the localized region of the Lyapunov vectors for higher spatial dimensions remains an important future problem. Related to this it should still be noted that it is known that the clock model version of the brick accumulation model itself can be easily generalized to any spatial dimensional case, although in higher dimension we do not have the concept of the accumulation of bricks, as we do for the quasi-one-dimensional case.
Finally, one should notice that the brick accumulation model (more generally the clock model) used in this paper has been justified for hard-disk (or hard-sphere) systems only (at least so far). On the other hand, the Lyapunov localization is observed not only in many-hard-disk systems but also in a wide variety of many-body chaotic systems, such as the Kuramoto-Sivashinsky model Man85 , a random matrix model Liv89 , map systems Kan86 ; Fal91 ; Gia91 ; Pik98 , coupled nonlinear oscillators Pik01 , etc. It should be an important future problem to develop approaches to the dynamics of Lyapunov localization in this wider class of systems.
## Acknowledgements
One of the authors (T.T) wishes to thank to R. van Zon, who indicated the possibility of describing the hopping dynamics for the localized Lyapunov vectors using the clock model, and supplied a computer program for the clock model. We also appreciate helpful discussions with E. Zabey and J. -P. Eckmann. In particular, they suggested the importance of the difference of brick heights for nearest-neighbor sites in the calculation of the hopping rate, and posed a question about the behavior of the sum of brick heights as shown in Fig. 6, and also indicated the analogy between the brick accumulation model and ballistic aggregation models. We also thank H. A. Posch for a comment on Lyapunov localization, which sparked our interest in this problem.
The authors appreciate the financial support of the Japan Society for the Promotion of Science.
## Appendix A Hopping Rate of a Localized Lyapunov Vector
In this appendix we give the detailed definition of the hopping rate for the localized region of the Lyapunov vectors, which is used in this paper.
As suggested in Ref. Tan03b ; Tan05b , the normalized Lyapunov vectors corresponding to the large Lyapunov exponents have non-zero components for only two particles in the low density limit, and changes of the non-zero components are caused by particle-particle collisions. Applying this characteristic of the Lyapunov vector to the quasi-one-dimensional system with periodic boundary conditions, we can introduce a hopping distance $`h^{[n]}`$ at the $`n`$-th particle-particle collision as
$`h^{[n]}=j_{n+1}j_nN\underset{¯}{\text{nint}}\left\{{\displaystyle \frac{j_{n+1}j_n}{N}}\right\}`$ (40)
in the low density limit. Here $`\{j_n,j_n+1\}`$ are the non-zero components before the $`n`$-th collision and $`\{j_{n+1},j_{n+1}+1\}`$ are the non-zero components just after the $`n`$-th collision. The function $`\underset{¯}{\text{nint}}\{x\}`$ is the closest integer to the real number $`x`$. We assume that changes in the position of the localized region of the Lyapunov vector are negligible during the free-flight interval, so $`\{j_{n+1},j_{n+1}+1\}`$ can also be interpreted as the set of the particle indices whose Lyapunov vector components take non-zero values just before the $`(n+1)`$-th collision. We count the number of times $`N_T(h)`$ that we see a hop of size $`h`$ in a time-interval $`T`$ where $`[N/2]h[N/2]`$. The normalized hopping rate $`P_N(h)`$ can be introduced as $`P_N(h)/P_N(1)=lim_T\mathrm{}N_T(h)/N_T(1)`$.
However, in actual numerical simulations, the particle density $`\rho `$ is always finite, and non-zero Lyapunov components of more than two particles can often be seen, at least down to a density $`\rho 10^5`$ which is the low density limit of the numerical simulations in this paper. This makes the above definition (40) of the hopping distance $`h^{[n]}`$ impractical. To explain this point concretely we show Fig. 11, which is a graph of the normalized Lyapunov vector component amplitude $`\gamma _j^{(1)}`$ defined by Eq. (1) as a function of the collision number $`n_t`$ and the particle index $`j`$ in a quasi-one-dimensional system with $`N=50`$ and $`d=10^4`$, (namely a density $`\rho 7.85\times 10^5`$). In this figure we can recognize three types of hops of the localized region of the Lyapunov vector. The first two hops keep the non-zero Lyapunov vector components of almost two particles sharp enough to apply the definition (40) of the hopping distance, but the third hop occurs gradually so that no clear hopping time can be determined. Note that the localized region of Lyapunov vector component amplitudes in the three dimensional plot 11 always has a flat top with a width of two-particles even for the third hop in Fig. (11).
In this paper, for concrete calculations, we define the localized region of the Lyapunov vector components as the particle indices $`j`$ for which $`\gamma _j^{(n)}>0.2`$. (Here, we use the similarity between particle positions and particle indices given in Fig. 1 for the quasi-one-dimensional system.) In Fig. 11, this localized region is approximately given by the region surrounded by the contour lines (level $`0.2`$) on the base of the graph. We introduce the quantity $`l_n`$ as the number of particles satisfying the inequality $`\gamma _j^{(n)}>0.2`$ just before the $`n`$-th collision. Note that $`0\gamma _j^{(n)}1`$ by definition (1) of $`\gamma _j^{(n)}`$, so $`l_n`$ cannot be larger than $`5`$. If $`l_n`$ is always 2 as in the low density limit, then we can use the hopping distance definition (40), but in numerical simulations at finite density $`l_n>2`$ can happen as shown in Fig. 11. The problem then is how do we define the hopping distance $`h^{[n]}`$ at the $`n`$-th collision, when $`l_n>2`$.
The definitions of the hopping distance for each case are categorized as follows.
Here only two particle indices are in the localized region before and after the collision, and they are always nearest-neighbors, so the hopping distance $`h^{[n]}`$ is given by Eq. (40).
Here we assume that the localized regions of the Lyapunov vector are given by $`\{j_n,j_n+1\}`$, $`\{j_{n+1},j_{n+1}+1,j_{n+1}+2\}`$ and $`j_{n+1}=j_n`$ or $`j_{n+1}=j_n1`$. We take the value of the hopping distance to be 1 (-1) for $`h^{[n]}`$ where $`j_{n+1}=j_n`$ ($`j_{n+1}=j_n1)`$.
In this case we assume that the localized regions of the Lyapunov vector are given by $`\{j_n,j_n+1,j_n+2\}`$, $`\{j_{n+1},j_{n+1}+1\}`$ and $`j_{n+1}=j_n`$ or $`j_{n+1}=j_n+1`$, and the hopping distance is $`0`$ in both cases.
In this case we assume that the localized regions of the Lyapunov vector are given by $`\{j_n,j_n+1,j_n+2\}`$, $`\{j_{n+1},j_{n+1}+1,j_{n+1}+2\}`$ and $`j_{n+1}=j_n+h^{[n]}`$ with the hopping distance $`h^{[n]}`$. We take into account the case $`h^{[n]}=1,0`$ or $`1`$ only.
In this case we assume that the localized regions of the Lyapunov vector are given by $`\{j_n,j_n+1\}`$, $`\{j_{n+1},j_{n+1}+1\}`$ and $`\{j_{n+1}^{},j_{n+1}^{}+1\}`$ satisfying $`j_{n+1}^{}=j_n`$ and $`|j_{n+1}^{}j_{n+1}|2`$. The hopping distance is defined by Eq. (40) using these $`j_{n+1}`$ and $`j_n`$.
In this case we assume that the localized regions of the Lyapunov vector are given by $`\{j_n,j_n+1\}`$, $`\{j_n^{},j_n^{}+1\}`$ and $`\{j_{n+1},j_{n+1}+1\}`$ satisfying $`|j_nj_n^{}|2`$ and $`|j_nj_{n+1}|1`$. The hopping distance $`h^{[n]}`$ is $`h^{[n]}=j_{n+1}j_n=1,0`$ or $`1`$.
In this case we consider only the case in which the localized regions of Lyapunov vector are given by $`\{j_n,j_n+1\}`$, $`\{j_n^{},j_n^{}+1\}`$, $`\{j_{n+1},j_{n+1}+1\}`$ and $`\{j_{n+1}^{},j_{n+1}^{}+1\}`$ satisfying $`j_{n+1}=j_n`$, $`j_{n+1}^{}=j_n^{}`$ and $`|j_nj_n^{}|2`$. The hopping distance takes the value $`0`$ in this case.
(For each case above, the corresponding schematic illustration is shown in Fig. 12.) Here, we use periodic boundary conditions for the particle index, so that the localized regions $`\{N,1\}`$, $`\{N,1,2\}`$ and $`\{N1,N,1\}`$ should be translated to $`\{0,1\}`$, $`\{0,1,2\}`$ and $`\{1,0,1\}`$, respectively, in the above definition of the hopping distance. In the examples shown in Fig. 11, the first two hops of the localized Lyapunov vector can be described as case (a), and the third hop is described as cases (c1) and (c2). Note that asymmetric definitions of hopping distances $`h^{[n]}`$ between cases (c1) and (c2) \[and similarly between cases (b1) and (b2)\] are required so that we can interpret the hopping distance of the third non-zero hop in Fig. 11 as only $`2`$ in spite of it involving both cases (c1) and (c2). Using the above hopping distance we count the number $`N_T(h)`$ of hops of distance $`h`$ in time interval $`T`$, and introduce the normalized hopping rate as $`P_N(h)/P_N(1)lim_T\mathrm{}N_T(h)/N_T(1)`$ for $`h=[N/2],[N/2]+1,\mathrm{},[N/2]`$. Notice that there are possibilities apart from those shown above, but it is observed that in numerical simulations the probabilities of these are extremely small (for example, more than 96 percent of all hops could be categorized this way).
In Fig. 13 we show the normalized hopping rate $`P_N(h)/P_N(1)`$ in a quasi-one-dimensional system of $`50`$ hard-disks for $`d=10^3`$ (circles), $`10^4`$ (triangles) and $`10^5`$ (squares). The hopping rate is almost density-independent in this low density range, although it may be very slightly larger at the lowest density for $`|h|=2,3,\mathrm{}`$.
We also calculate the hopping rate for the brick accumulation model explained in Sec. IV in a similar way. In the brick accumulation model we can introduce the localized region of the Lyapunov vectors as the site (particle) indices whose brick height are highest. For the brick model, cases (b1), (b2), and (b3) above cannot happen, and only cases (a), (c1), (c2) and (c3) above are taken into account in the numerical calculations. It may be noted that in the brick model, case (a) above can happen only when the hopping distance is $`1`$, $`0`$ or $`1`$, and Eq. (40) alone is not enough to calculate the hopping distance. This is another reason to take into account the case where $`l_n>2`$ in the calculation of the hopping distance.
## Appendix B Clock Model for Many-Hard-Particle Systems
Here we give an extension of the derivation of the clock model for the Lyapunov vector dynamics in many-hard-particle systems in the limit of low density. We also derive the formula (9) for the largest Lyapunov exponent from the clock model. The one-dimensional version of the clock model is the brick accumulation model used in this paper. In particular, we clarify the assumptions needed to justify the use of this model to discuss the Lyapunov localization. For the basic idea of the clock model, see, for example, Refs. Bei00 ; Zon02 . However, note that in this appendix we use notation and assumptions that are a little different from these references, so that the derived clock model is consistent with the discussions in this paper and able to be compared with the numerical results of Ref. Tan03b .
We consider a $`𝒟`$-dimensional system with $`N`$ hard-disks (or hard-spheres, etc.) with identical radius $`R`$ and mass $`M`$. We assume that there is no external field in the system so that the dynamics is simply free-flights, and collisions between two particles. We put $`\delta 𝒒_j`$ ($`𝒒_j`$) as the spatial part of the Lyapunov vector component (the spatial coordinate) of the $`j`$-th particle, and $`\delta 𝒑_j`$ ($`𝒑_j`$) as the momentum part of the Lyapunov vector component (the momentum) of the $`j`$-th particle.
We take $`t=t_n`$ to be the time of the $`n`$-th collision which involves particles $`j_n`$ and $`k_n`$, and define $`\tau _n`$ to be
$`\tau _n{\displaystyle \frac{t_nt_{n1}}{M}},`$ (41)
so that the $`n`$-th free flight time is given by $`\tau _nM`$. The free flight part of dynamics of the Lyapunov vector is represented as
$`\delta 𝒒_j(t_n^{})`$ $`=`$ $`\delta 𝒒_j(t_{n1}^+)+\tau _n\delta 𝒑_j(t_{n1}^+),`$ (42)
$`\delta 𝒑_j(t_n^{})`$ $`=`$ $`\delta 𝒑_j(t_{n1}^+)`$ (43)
Here, the argument $`t_n^\pm `$ refers to the limit of the quantity before the $`n`$-th collision ($``$) or after the $`n`$-th collision ($`+`$). On the other hand, the change in the Lyapunov vector in particle-particle collisions is represented as
$`\delta 𝒒_{j_n}(t_n^+)`$ $`=`$ $`\delta 𝒒_{j_n}(t_n^{})+\mathrm{\Theta }^{[n]}\delta 𝒒_{k_nj_n}(t_n^{}),`$ (44)
$`\delta 𝒒_{k_n}(t_n^+)`$ $`=`$ $`\delta 𝒒_{k_n}(t_n^{})\mathrm{\Theta }^{[n]}\delta 𝒒_{k_nj_n}(t_n^{}),`$ (45)
$`\delta 𝒒_l(t_n^+)`$ $`=`$ $`\delta 𝒒_l(t_n^{}),\text{for}l/\{j_n,k_n\},`$ (46)
$`\delta 𝒑_{j_n}(t_n^+)`$ $`=`$ $`\delta 𝒑_{j_n}(t_n^{})+\mathrm{\Theta }^{[n]}\delta 𝒑_{k_nj_n}(t_n^{})`$ (47)
$`+Q^{[n]}\delta 𝒒_{k_nj_n}(t_n^{}),`$
$`\delta 𝒑_{k_n}(t_n^+)`$ $`=`$ $`\delta 𝒑_{k_n}(t_n^{})\mathrm{\Theta }^{[n]}\delta 𝒑_{k_nj_n}(t_n^{})`$ (48)
$`Q^{[n]}\delta 𝒒_{k_nj_n}(t_n^{}),`$
$`\delta 𝒑_l(t_n^+)`$ $`=`$ $`\delta 𝒑_l(t_n^{}),\text{for}l/\{j_n,k_n\},`$ (49)
where $`\delta 𝒒_{k_nj_n}\delta 𝒒_{k_n}\delta 𝒒_{j_n}`$, $`\delta 𝒑_{k_nj_n}\delta 𝒑_{k_n}\delta 𝒑_{j_n},`$ and $`\mathrm{\Theta }^{[n]}`$ and $`Q^{[n]}`$ are defined by
$`\mathrm{\Theta }^{[n]}`$ $``$ $`𝝈^{[n]}𝝈^{[n]T},`$ (50)
$`Q^{[n]}`$ $``$ $`{\displaystyle \frac{𝝈^{[n]T}\mathrm{\Delta }𝒑^{[n]}}{2R}}\left(I+{\displaystyle \frac{𝝈^{[n]}\mathrm{\Delta }𝒑^{[n]T}}{𝝈^{[n]T}\mathrm{\Delta }𝒑^{[n]}}}\right)`$ (51)
$`\times \left(I{\displaystyle \frac{\mathrm{\Delta }𝒑^{[n]}𝝈^{[n]T}}{\mathrm{\Delta }𝒑^{[n]T}𝝈^{[n]}}}\right)`$
with
$`𝝈^{[n]}`$ $``$ $`{\displaystyle \frac{𝒒_{k_n}(t_n^{})𝒒_{j_n}(t_n^{})}{|𝒒_{k_n}(t_n^{})𝒒_{j_n}(t_n^{})|}},`$ (52)
$`\mathrm{\Delta }𝒑^{[n]}`$ $``$ $`𝒑_{k_n}(t_n^{})𝒑_{j_n}(t_n^{})`$ (53)
and the $`(𝒟N)\times (𝒟N)`$ identity matrix $`I`$. Note that in this appendix we introduce all vectors as column vectors, so for example, $`𝝈^{[n]T}\mathrm{\Delta }𝒑^{[n]}`$ is a scalar and $`𝝈^{[n]}\mathrm{\Delta }𝒑^{[n]T}`$ is a matrix where $`T`$ is the transpose. For later use we note that
$`\delta 𝒑_{j_n}(t_n^+)+\delta 𝒑_{k_n}(t_n^+)=\delta 𝒑_{j_n}(t_n^{})+\delta 𝒑_{k_n}(t_n^{})`$ (54)
which can be derived from Eqs. (47) and (48). For the derivation of Eqs. (44), (45), (46), (47), (48) and (49) for the Lyapunov vector dynamics, for example, see Ref. Del96 .
We consider a low density case, in which the free flight time $`\tau _nM`$ is large (as the free flight time is inversely proportional to the density). This justifies our first approximation for the Lyapunov vector:
$`\delta 𝒒_j(t_n^{})\tau _n\delta 𝒑_j(t_{n1}^+)`$ (55)
in the limit of low density, the term containing $`\tau _n`$ is much larger than the other term on the right-hand side of Eq. (42). This asymptotic relation (55) leads to
$`\delta 𝒒_{k_nj_n}(t_n^{})\tau _n\delta 𝒑_{k_nj_n}(t_{n1}^+)`$ (56)
from the definition of $`\delta 𝒒_{k_nj_n}(t_n^{})`$ and $`\delta 𝒑_{k_nj_n}(t_{n1}^+)`$. Using the relation (56), Eqs. (47), (48) and (49) can be rewritten as
$`\delta 𝒑_{j_n}(t_n^+)`$ $``$ $`\delta 𝒑_{j_n}(t_{n1}^+)+\mathrm{\Theta }^{[n]}\delta 𝒑_{k_nj_n}(t_{n1}^+)`$ (57)
$`+\tau _nQ^{[n]}\delta 𝒑_{k_nj_n}(t_{n1}^+),`$
$`\delta 𝒑_{k_n}(t_n^+)`$ $``$ $`\delta 𝒑_{k_n}(t_{n1}^+)\mathrm{\Theta }^{[n]}\delta 𝒑_{k_nj_n}(t_{n1}^+)`$ (58)
$`\tau _nQ^{[n]}\delta 𝒑_{k_nj_n}(t_{n1}^+)`$
$`\delta 𝒑_l(t_n^+)`$ $`=`$ $`\delta 𝒑_l(t_{n1}^+),\text{for}l/\{j_n,k_n\},`$ (59)
where we have used Eq. (43). Note that the spatial part of the Lyapunov vector does not appear in the dynamics described in (57) and (58), (59) anymore. The first and second terms on the right-hand side of (57) and (58) are negligible compared with the third term because of the large value of $`\tau _n`$, so we obtain
$`\delta 𝒑_{j_n}(t_n^+)`$ $``$ $`\tau _nQ^{[n]}\delta 𝒑_{k_nj_n}(t_{n1}^+),`$ (60)
$`\delta 𝒑_{k_n}(t_n^+)`$ $``$ $`\tau _nQ^{[n]}\delta 𝒑_{k_nj_n}(t_{n1}^+),`$ (61)
which lead to
$`\delta 𝒑_{j_n}(t_n^+)+\delta 𝒑_{k_n}(t_n^+)0.`$ (62)
On the other hand, for the dynamics given by (44), (45) and (46) for the spatial part of Lyapunov vector we obtain
$`\delta 𝒒_{j_n}(t_n^+)`$ $``$ $`\tau _n\left[\delta 𝒑_{j_n}(t_{n1}^+)+\mathrm{\Theta }^{[n]}\delta 𝒑_{k_nj_n}(t_{n1}^+)\right],`$ (63)
$`\delta 𝒒_{k_n}(t_n^+)`$ $``$ $`\tau _n\left[\delta 𝒑_{k_n}(t_{n1}^+)\mathrm{\Theta }^{[n]}\delta 𝒑_{k_nj_n}(t_{n1}^+)\right],`$ (64)
$`\delta 𝒒_l(t_n^+)`$ $``$ $`\tau _n\delta 𝒑_l(t_{n1}^+),\text{for}l/\{j_n,k_n\},`$ (65)
using Eqs. (55) and (56). Now we note
$`\delta 𝒑_{j_n}(t_{n1}^+)`$ $`=`$ $`{\displaystyle \frac{\delta 𝒑_{j_n}(t_{n1}^+)+\delta 𝒑_{k_n}(t_{n1}^+)}{2}}`$ (66)
$`+{\displaystyle \frac{\delta 𝒑_{j_n}(t_{n1}^+)\delta 𝒑_{k_n}(t_{n1}^+)}{2}}`$
$`=`$ $`{\displaystyle \frac{\delta 𝒑_{j_n}(t_n^{})+\delta 𝒑_{k_n}(t_n^{})}{2}}+{\displaystyle \frac{\delta 𝒑_{j_nk_n}(t_{n1}^+)}{2}}`$
$`=`$ $`{\displaystyle \frac{\delta 𝒑_{j_n}(t_n^+)+\delta 𝒑_{k_n}(t_n^+)}{2}}{\displaystyle \frac{\delta 𝒑_{k_nj_n}(t_{n1}^+)}{2}}`$
$``$ $`{\displaystyle \frac{\delta 𝒑_{k_nj_n}(t_{n1}^+)}{2}}`$
where we used Eqs. (43), (54) and (62). Similarly we have
$`\delta 𝒑_{k_n}(t_{n1}^+)`$ $``$ $`{\displaystyle \frac{\delta 𝒑_{k_nj_n}(t_{n1}^+)}{2}}.`$ (67)
Therefore Eqs. (63) and (64) can be rewritten as
$`\delta 𝒒_{j_n}(t_n^+)`$ $``$ $`\tau _n\mathrm{\Omega }^{[n]}\delta 𝒑_{k_nj_n}(t_{n1}^+),`$ (68)
$`\delta 𝒒_{k_n}(t_n^+)`$ $``$ $`\tau _n\mathrm{\Omega }^{[n]}\delta 𝒑_{k_nj_n}(t_{n1}^+)`$ (69)
where $`\mathrm{\Omega }^{[n]}`$ is defined by
$`\mathrm{\Omega }^{[n]}{\displaystyle \frac{1}{2}}I+\mathrm{\Theta }^{[n]}.`$ (70)
The asymptotic equations (59), (60), (61), (65), (68) and (69) give the Lyapunov vector dynamics in the limit of low density. It should also be noted that the assumptions used to derive this dynamics breaks some conservation laws in the original dynamics, for example, the quantity $`_{j=1}^N\delta 𝒑_j(t)`$ is conserved in the original dynamics (43), (47), (48) and (49) but this cannot be guaranteed exactly in the low density dynamics (59), (60) and (61).
Now we consider the Lyapunov vector $`\delta 𝜞`$ corresponding to a positive Lyapunov exponent. The positivity of the Lyapunov exponent means that the amplitude $`|\delta 𝜞|`$ of Lyapunov vector diverges exponentially in time. The dynamics given by (59), (60), (61), (65), (68) and (69) shows that the divergence of $`|\delta 𝜞|`$ must come from the Lyapunov vector components corresponding to colliding particles, as the other Lyapunov vector components diverge at most linearly. For this reason we neglect the change of the Lyapunov vector components for non-colliding particles. Under this assumption the Lyapunov vector dynamics for the Lyapunov vector component $`\delta 𝜞_j(\delta 𝒒_j,\delta 𝒑_j)^T`$ for the $`j`$-th particle is summarized as
$`\delta 𝜞_{j_n}(t_n^+)`$ $``$ $`\delta 𝜞_{k_n}(t_n^+)\tau _n\delta 𝚵^{[n1]}`$ (71)
$`\delta 𝜞_l(t_n^+)`$ $``$ $`\delta 𝜞_l(t_{n1}^+),\text{for}l/\{j_n,k_n\},`$ (72)
where $`\delta 𝚵^{[n1]}`$ is defined by
$`\delta 𝚵^{[n1]}\left(\begin{array}{c}\mathrm{\Omega }^{[n]}\delta 𝒑_{k_nj_n}(t_{n1}^+)\\ Q^{[n]}\delta 𝒑_{k_nj_n}(t_{n1}^+)\end{array}\right).`$ (75)
It is essential to note that the Lyapunov vector components $`\delta 𝜞_{j_n}(t_n^+)`$ and $`\delta 𝜞_{k_n}(t_n^+)`$ for the colliding particles have the same amplitude, and $`\delta 𝚵^{[n1]}`$ is independent of $`\tau _n`$.
We assume that the ratio
$`\mu {\displaystyle \frac{|\delta 𝒒_j|}{|\delta 𝒑_j|}}`$ (76)
between the amplitudes of the spatial and momentum parts of the Lyapunov vector for the $`j`$-th particle, are independent of the particle index $`j`$ noteA . Ref. Tan03b suggests that the ratio $`\mu `$ need not be of order $`1`$ in general. From Eq. (76) we have
$`|\delta 𝒑_j|{\displaystyle \frac{1}{\sqrt{1+\mu ^2}}}|\delta 𝜞_j|.`$ (77)
Using this assumption we estimate the magnitude of the vector $`\delta 𝚵^{[n1]}`$ as
$`\left|\delta 𝚵^{[n1]}\right|`$
$`=\sqrt{\left|\mathrm{\Omega }^{[n]}\delta 𝒑_{k_nj_n}(t_{n1}^+)\right|^2+\left|Q^{[n]}\delta 𝒑_{k_nj_n}(t_{n1}^+)\right|^2}`$
$`=\text{max}\{\left|\delta 𝒑_{j_n}(t_{n1}^+)\right|,\left|\delta 𝒑_{k_n}(t_{n1}^+)\right|\}`$
$`\times \sqrt{\left|\mathrm{\Omega }^{[n]}𝒆^{[n]}\right|^2+\left|Q^{[n]}𝒆^{[n]}\right|^2}`$
$`\text{max}\{\left|\delta 𝜞_{j_n}(t_{n1}^+)\right|,\left|\delta 𝜞_{k_n}(t_{n1}^+)\right|\}`$
$`\times \sqrt{{\displaystyle \frac{\left|\mathrm{\Omega }^{[n]}𝒆^{[n]}\right|^2+\left|Q^{[n]}𝒆^{[n]}\right|^2}{1+\mu ^2}}}`$ (78)
where $`𝒆^{[n]}`$ is defined by
$`𝒆^{[n]}{\displaystyle \frac{\delta 𝒑_{k_n}(t_{n1}^+)\delta 𝒑_{j_n}(t_{n1}^+)}{\text{max}\{\left|\delta 𝒑_{k_n}(t_{n1}^+)\right|,\left|\delta 𝒑_{j_n}(t_{n1}^+)\right|\}}},`$ (79)
and satisfies the inequality
$`\left|𝒆^{[n]}\right|{\displaystyle \frac{\left|\delta 𝒑_{k_n}(t_{n1}^+)\right|+\left|\delta 𝒑_{j_n}(t_{n1}^+)\right|}{\text{max}\{\left|\delta 𝒑_{k_n}(t_{n1}^+)\right|,\left|\delta 𝒑_{j_n}(t_{n1}^+)\right|\}}}2.`$ (80)
From Eq. (78) we can estimate the magnitudes $`|\delta 𝜞_{j_n}(t_n^+)|`$ and $`|\delta 𝜞_{k_n}(t_n^+)|`$ through $`|\delta 𝜞_{j_n}(t_n^+)|`$ $`=|\delta 𝜞_{k_n}(t_n^+)|=\tau _n|\delta 𝚵^{[n1]}|`$.
Now, we introduce the clock value $`𝒦_j(n)`$ of the $`j`$-th particle just after the $`n`$-th collision as
$`𝒦_j(n){\displaystyle \frac{1}{\mathrm{ln}\rho }}\mathrm{ln}{\displaystyle \frac{|\delta 𝜞_j(t_n^+)|}{|\delta 𝜞_j(0)|}},`$ (81)
or equivalently,
$`|\delta 𝜞_j(t_n^+)|=\left({\displaystyle \frac{1}{\rho }}\right)^{𝒦_j(n)}|\delta 𝜞_j(0)|`$ (82)
leading to Eq. (8). Here $`\rho `$ is the particle density whose value is $`0<\rho <1`$, and $`\delta 𝜞_j(0)`$ is the Lyapunov vector component for the $`j`$-th particle at the initial time. We choose the initial Lyapunov vector $`\delta 𝜞(0)`$ so that its components $`\delta 𝜞_1(0),\delta 𝜞_2(0),\mathrm{},\delta 𝜞_{N1}(0)`$ and $`\delta 𝜞_N(0)`$ have the same order of magnitude, and a larger clock value means larger magnitude of the Lyapunov vector component, $`|\delta 𝜞_j(t_n)|`$ .
We have already used the fact that the free flight time increases as the density $`\rho `$ decreases. To write down the clock value more meaningfully, we use the specific relation between the free flight time and the density, namely that the free-flight time $`\tau _nM`$ is approximately inversely proportional to the density $`\rho `$ at low density;
$`\tau _ns_n/\rho `$ (83)
where $`s_n`$ is independent of the density. Using Eq. (83), the expression (81) for the clock value $`𝒦_j(n)`$ can be rewritten as
$`𝒦_l(n)`$ $``$ $`\{\begin{array}{c}\text{max}\{𝒦_{j_n}(n1),𝒦_{k_n}(n1)\}+1+\mathrm{\Delta }\mathrm{\Phi }^{[n]}\hfill \\ \text{for}l=j_n\text{or}l=k_n\hfill \\ \\ 𝒦_l(n1)\text{for}l/\{j_n,k_n\},\hfill \end{array}`$ (88)
where $`\mathrm{\Delta }\mathrm{\Phi }^{[n]}`$ is defined by
$`\mathrm{\Delta }\mathrm{\Phi }^{[n]}{\displaystyle \frac{1}{\mathrm{ln}\rho }}\mathrm{ln}\left\{s_n\sqrt{{\displaystyle \frac{\left|\mathrm{\Omega }^{[n]}𝒆^{[n]}\right|^2+\left|Q^{[n]}𝒆^{[n]}\right|^2}{1+\mu ^2}}}\right\},`$
(89)
and where we have used (71), (78), (81) and $`|\delta 𝜞_{j_n}(0)||\delta 𝜞_{k_n}(0)|`$. Our final assumption to derive the clock model is that
$`\underset{\rho 0}{lim}\mathrm{\Delta }\mathrm{\Phi }^{[n]}=0.`$ (90)
To justify the assumption (90) note that the right-hand side of Eq. (89) for the quantity $`\mathrm{\Delta }\mathrm{\Phi }^{[n]}`$ has a factor $`1/\mathrm{ln}\rho `$ which goes to zero in the limit as $`\rho 0`$, and $`s_n`$, $`Q^{[n]}`$ and $`\mathrm{\Omega }^{[n]}`$ are almost independent of the density $`\rho `$, and the magnitude of the vector $`𝒆^{[n]}`$ is finite even in the limit $`\rho 0`$ from the inequality (80) noteA . Eqs. (88) and (90) lead to the clock dynamics
$`𝒦_l(n)`$ $``$ $`\{\begin{array}{cc}\text{max}\{𝒦_{j_n}(n1),𝒦_{k_n}(n1)\}+1\hfill & \\ \text{for}l=j_n\text{or}l=k_n\hfill & \\ & \\ 𝒦_l(n1)\text{for}l/\{j_n,k_n\},\hfill & \end{array}`$ (95)
in the low density limit, which is closed by the clock value $`𝒦_l(n)`$ only. From Eq. (95), the dynamics for the clock model is expressed as (i) the clock value is changed only when the corresponding particle collides, and (ii) the clock values of colliding particles are tuned to the same value given by $`1`$ plus the larger of the two clock values of the particles just before the collision.
In the quasi-one-dimensional system with periodic boundary conditions, particle indices of the colliding particles can be taken so that $`k_n=j_n+1`$ (note that index $`N+1`$ is equivalent to $`1`$). Therefore we obtain the dynamics (7) for the brick accumulation model explained in Sec. IV as the one-dimensional version of the clock model. Moreover, in the one-dimensional case, $`𝒦_j(n)`$ can be interpreted as the brick height of the $`j`$-site just after the $`n`$-th brick is dropped.
Finally we derive the formula (9) for the largest Lyapunov exponent from the clock value $`𝒦_j(n)`$. In the limit of low density $`\rho <<1`$, the amplitude $`|\delta 𝜞(t_n^+)|`$ of the Lyapunov vector can be approximated by
$`|\delta 𝜞(t_n^+)|`$ $`=`$ $`\sqrt{{\displaystyle \underset{j=1}{\overset{N}{}}}|\delta 𝜞_j(t_n^+)|^2}`$ (96)
$``$ $`\alpha _n\left({\displaystyle \frac{1}{\rho }}\right)^{\stackrel{~}{𝒦}^{[max]}(n)}|\delta 𝜞_j(0)|`$
noting that the sum $`_{j=1}^N|\delta 𝜞_j(t_n^+)|^2`$ is dominated by the Lyapunov vector component amplitude $`|\delta 𝜞_j(t_n^+)|`$ with the largest clock value $`\stackrel{~}{𝒦}^{[max]}(n)`$:
$`\stackrel{~}{𝒦}^{[max]}(n)\mathrm{max}\{𝒦_1(n),𝒦_2(n),\mathrm{},𝒦_N(n)\}`$ (97)
at $`n`$, because of the huge factor $`1/\rho `$. Here, $`\alpha _n`$ is the number of particle with the largest clock value $`\stackrel{~}{𝒦}^{[max]}(n)`$, and we have also used our assumption that $`|\delta 𝜞_j(0)|`$ is almost independent of the particle index $`j`$. Now we assume the approximate relation
$`{\displaystyle \frac{\stackrel{~}{𝒦}^{[max]}(n)}{n}}\stackrel{n\mathrm{}}{}{\displaystyle \frac{1}{nN}}{\displaystyle \underset{j=1}{\overset{N}{}}}𝒦_j(n)`$ (98)
in the limit of large $`n`$. (We have checked this relation numerically for the brick accumulation model.) Using the relations (96), (98) and $`tn\tau `$ we obtain
$`\lambda `$ $``$ $`\underset{t+\mathrm{}}{lim}{\displaystyle \frac{1}{t}}\mathrm{ln}|\delta 𝜞(t)|`$ (99)
$``$ $`\underset{n+\mathrm{}}{lim}{\displaystyle \frac{1}{n\tau }}\stackrel{~}{𝒦}^{[max]}(n)\mathrm{ln}\rho `$
$``$ $`\underset{n+\mathrm{}}{lim}{\displaystyle \frac{1}{n\tau N}}{\displaystyle \underset{j=1}{\overset{N}{}}}𝒦_j(n)\mathrm{ln}\rho `$
for the Lyapunov exponent $`\lambda `$ corresponding to the Lyapunov vector $`\delta 𝜞`$. Therefore we obtain the Eq. (9), which is independent of the system shape and the number of spatial dimensions.
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# The Intrinsic Properties of SMM J14011+0252
## 1. Introduction
SMM J14011+0252 (hereafter SMM J14011) is a bright submillimeter (submm) source identified in the field of the cluster lens, A 1835, by Ivison et al. (2000, I00) (see also Smail et al. 2002). The submm source has a radio counterpart, allowing the position of the submm emission to be pin-pointed to a close pair of blue galaxies. The redshift of the system was measured as $`z2.56`$ from a near-infrared spectrum by I00, and both components were confirmed to lie at the same redshift from optical spectroscopy by Barger et al. (1999). Molecular CO emission was subsequently identified at this redshift by Frayer et al. (1999). These observations appeared to show that the $`z=2.56`$ system is a massive, gas-rich far-infrared luminous source, only the second such system identified from submm surveys (e.g. Ivison et al. 1998).
SMM J14011 has become a archetype for submm galaxies, being studied extensively by a number of authors (Ivison et al. 2001, I01; Downes & Solomon 2003, DS03; Tecza et al. 2004, T04; Frayer et al. 2004; Swinbank et al. 2004; Motohara et al. 2005, M05; Carilli et al. 2005). From the outset it had been assumed that SMM J14011 was gravitationally amplified by the potential of the foreground cluster, by a factor of $`\mu 2.8`$$`3\times `$ (I00; Frayer et al. 1999). However, DS03 suggested that in fact the submm source is highly amplified by a superimposed, foreground galaxy – resulting in J1 and J2 being two images of a single background source, with an amplification of a factor of $`\mu 25`$–30 (see also M05). This scenario results in very different conclusions about the intrinsic properties of SMM J14011 and would mean that the source is not representative of the mJy-submm population, which are the subject of much research.
In this paper we reassess the available multiwavelength information on SMM J14011 to determine the likely influence of gravitational lensing on the measured properties of this galaxy. We assume a cosmology with $`\mathrm{\Omega }_m=0.27`$, $`\mathrm{\Omega }_\mathrm{\Lambda }=0.73`$ and $`H_o=71`$ km s<sup>-1</sup> Mpc<sup>-1</sup>, giving an angular scale of 3.88 kpc and 8.15 kpc per arcsec. at $`z=0.25`$ and $`z=2.56`$ respectively.
## 2. Analysis and Modelling
We show a variety of different views of SMM J14011 in Fig. 1 and mark on the main components following the naming scheme of M05. From ground-based optical imaging it is clear that the system is made up of two main components: J1 and J2 (I00). J1 and J2 are both relatively blue, compared to the bright early-type galaxies in the foreground cluster (I00). However, higher resolution optical imaging from Hubble Space Telescope (I01) identifies several knots within J1 – J1a/J1b/J1d in Fig. 1, which appear superimposed on the relatively smooth and regular underlying component J1c. While high-resolution ground-based near-infrared imaging highlights a low-surface brightness, very red extension to the north of J1, which terminates in a faint knot, J1n (I01; Frayer et al. 2004). The morphology of the redshifted H$`\alpha `$ emission from the $`z=2.56`$ galaxy has been mapped using both narrow-band imaging (Swinbank et al. 2004) and integral field spectroscopy (T04). These observations show that the H$`\alpha `$ emission is seen from the knots J1a/J1b/J1d and the $`2.5^{\prime \prime }`$-long, very red extension out to J1n (Fig. 1) – but not from J1c.
The exact position of the submm, molecular gas and radio emission relative to the optical imaging remains somewhat uncertain. Aligning the HST image to the radio reference frame using the five radio sources with bright optical counterparts (but not SMM J14011) yields a position for J1c of 14 01 04.96 +02 52 24.1 (J2000) with an uncertainty of $`0.3^{\prime \prime }`$ (consistent with the published position from I01). We plot the centroids and 1-$`\sigma `$ error-circles for the CO(4–3) and CO(7–6) emission from DS03 and our 1.4-GHz centroid on Fig. 1. Within the relative uncertainties the radio position could be associated with J1d or the near-infrared peak next to it, although the CO peaks are both somewhat south of this position.
### 2.1. The true nature of J1c
A detailed inspection of the optical spectrum of J1 from Barger et al. (1999), suggests that it comprises a mix of light from a UV-bright starburst at $`z=2.56`$ and a foreground galaxy as proposed by DS03. To isolate the light from the foreground galaxy we scale and subtract the spectrum of a high-redshift starburst (J2) from the spectrum of J1. The scaling factor is chosen so that the absorption features around $`4000`$Å are not highly negative in the resulting spectrum. The difference spectrum is shown in Fig. 2 and exhibits a number of absorption features which are identifiable as Balmer H$`\theta `$ \[4743.3Å observed wavelength\], Ca H \[4912.0Å\], Ca K \[4956.8Å\], H$`\delta `$ \[5126.7Å\], G-band \[5375.8Å\] and H$`\beta `$ \[6069.3Å\] at $`z=0.2489\pm 0.0004`$. Thus it appears that there is light from a foreground member galaxy of the A 1835 cluster mixed with that from the $`z=2.56`$ system in our spectroscopic slit. The identification of this foreground component is hampered by an unfortunate coincidence that some of the absorption lines in the $`z=0.25`$ galaxy can be interpreted as UV features at the redshift of the background source, e.g. Cii1335 corresponds to 4743.3Å and Siiv1394 falls close to 4956.8Å.
The lack of any detectable emission lines from the $`z=0.25`$ component in the J1 spectrum (e.g. \[Oii\] 3727, \[Oiii\] 5007), along with the apparent strength of the Balmer absorption lines (EW(H$`\delta )3`$Å) suggest that the superimposed galaxy could either be a passive galaxy or potentially a post-starburst system (Poggianti et al. 1999). However, we caution that the strength of these features is strongly dependent on the relative contributions from the foreground and background components in the spectrum, which is not well-determined.
Thus we confirm the suggestion by DS03 that one component of J1 is a foreground galaxy. Looking at the morphologies of the various components in Fig. 1, it seems most likely that J1c corresponds to the $`z=0.25`$ galaxy.
To gain a clearer view of the morphology and luminosity of J1c we take advantage of its apparent regularity and symmetry to remove the superimposed knots of emission. We rotate the HST image of SMM J14011 by 180 degrees and subtract it from itself – producing an image of just the knots of emission: J1a/J1b/J1d and J2. These can then in turn be subtracted from SMM J14011 to leave just the emission from J1c. This shows the restframe $`V`$-band structure of the galaxy – with a smooth, circular morphology indicative of a spheroid or face-on early-type disk. The measured ellipticity of J1c is $`ϵ=0.03\pm 0.01`$ (at PA=133 deg.) and its total magnitude (within a 6<sup>′′</sup>-diameter aperture) is $`R_{702}=21.27\pm 0.03`$ (our best estimate of the $`K`$-band magnitude of J1c is $`K=18.44\pm 0.10`$, giving a color of $`(RK)2.9`$, similar to that expected for a dwarf spheroid or early-type disk cluster member, Smith et al. 2002). Using gim-2d (Simard et al. 2002) we model the 2-D light distribution in J1c and obtain a best fit model ($`\chi ^2=1.3`$) with a half-light radius of $`r_{hl}=0.40\pm 0.05`$ or $`1.6`$ kpc and a bulge-to-disk ratio of $`0.30\pm 0.05`$, suggesting the galaxy is an early-type disk (Sab).
The apparent magnitude of J1c corresponds to an absolute magnitude of $`M_r19.2`$ at $`z=0.25`$, or roughly 0.1$`L^{}`$. Assuming the galaxy follows the Faber-Jackson relation measured in moderate redshift clusters (e.g. Zeigler et al. 2001), we would expect a velocity dispersion for the galaxy of $`<55\pm 15`$ km s<sup>-1</sup> given its absolute magnitude (the error is the 1-$`\sigma `$ scatter of galaxies around the Faber-Jackson relation). We quote this value as a limit due to the possibility that the mass-to-light ratio of J1c is higher than typical passive galaxies in the cluster if it has suffered recent star-formation. The low velocity dispersion we infer for this galaxy is consistent with the small half-light radius we measure if it follows the scaling ratios for early-type cluster galaxies at $`z0.2`$. Thus we conclude that J1c is a passive, early-type dwarf disk galaxy member of A 1835.
Fig. 2. — The spectrum of J1c obtained by subtracting a $`z=2.56`$ Lyman-break galaxy (a scaled version of the J2 spectrum) from the published J1 spectrum of Barger et al. (1999). The spectrum shown in bold is the smoothed difference spectrum, while the upper and lower traces show the spectra of component J1 and J2 respectively (offset for clarity). The difference spectrum represents the light from the $`z=0.25`$ galaxy, J1c. We mark some of the stronger spectral features visible in the foreground galaxy spectrum and we note that J1c contributes roughly half the light observed at 5000Å in the spectrograph slit.
We also use take advantage of the very different $`(JK)`$ colors of J1c and the red-extension, J1n (I01; Frayer et al. 2004) to scale the J-band emission of J1c and subtract it from the K-band image to leave a clear view of the K-band morphology of the high-redshift galaxy, Fig. 1. This unsurprisingly matches the H$`\alpha `$ morphology from T04, given the strong contribution from the H$`\alpha `$ line in the K-band (I00; T04; Swinbank et al. 2004; M05). However, we can use the good resolution of the Keck imaging (0.4<sup>′′</sup> FWHM) to show that this emission is resolved in both axes with a seeing-corrected size of $`1.8^{\prime \prime }\times 0.6^{\prime \prime }`$.
### 2.2. The strong lensing interpretation of SMM J14011
Two strong lensing models have been proposed for SMM J14011 based on the identification of various components of the system as multiple images of single background sources. DS03 discuss a toy model which seeks to explain J1a/J1b/J1n and J2 as four images of a single background source. Although they do not construct a rigorous model they suggest that a potential well associated with J1c with a velocity dispersion of $`200`$ km s<sup>-1</sup>, along with a contribution from the cluster mass distribution, could fit the four proposed images.
As shown by Swinbank et al. (2004) and T04, the spectroscopic information available for J1a/J1b/J1n and J2 indicates different redshifts for the two systems from their H$`\alpha `$ emission (e.g. Fig. 1). This rules out J2 as a counter image and hence allows us to discard the qualitative model suggested by DS03.
M05 also fit a strong-lensing model to components of J1 – this time J1a, J1b and J1d – using the gravlens model of Keeton (2001). Their model includes an external shear field arising from the cluster potential (constrained by the mass model of Schmidt et al. 2001) and is capable of reproducing the positions of the three proposed images using a highly-elliptical ($`ϵ=0.67\pm 0.09`$, PA $`=53\pm 10`$ deg.) potential well centered on J1c with a velocity dispersion of $`123_{14}^{+10}`$ km s<sup>-1</sup>.
The obvious problem with the M05’s lens model is that the parameters of the potential well of J1c differ substantially from those traced by the light distribution within this galaxy on similar scales. The ellipticity, position angle and likely velocity dispersion of the halo are all strongly at odds with those measured or expected for J1c ($`ϵ=0.03\pm 0.01`$, PA $`=133\pm 5`$ deg. and $`\sigma <55\pm 15`$ km s<sup>-1</sup>) if mass traces light on $`10`$ kpc scales within this galaxy in a similar manner to that found from modelling strong and weak lensing by individual galaxies (Treu & Koopman 2004; Kochanek et al. 2000; Natarajan et al. 1998). In particular, we note that the velocity dispersion required in this model would correspond to a galaxy with an absolute magnitude of $`M_r21`$, some $`6\times `$ brighter than J1c (assuming the $`z0.2`$ Faber-Jackson relation from Zeigler et al. 2001).
In addition, we find that the predicted flux ratios between J1a/J1b of 11:13 are at odds with the measured 1.0<sup>′′</sup> aperture magnitudes from the HST image of $`R_{702}=23.54`$ and 23.91 (with uncertainties of $`\pm 0.05`$, J1d has $`R_{702}24.2`$ in the same aperture). These yield a flux ratio between J1a and J1b of $`0.71\pm 0.07`$, compared to the predicted ratio of 1.18.
All of these discrepencies arise because of the incorrect identification of J1a/J1b/J1d as three images of a single background source. Looking at Fig. 1, it is clear that there is no compelling similarity between the morphologies of J1a, J1b and J1d: J1b is well resolved, with a FWHM of 0.57<sup>′′</sup>, while J1a is more compact (FWHM=0.30<sup>′′</sup>) and J1d is too faint to allow us to measure a reliable size.
### 2.3. Lens models
The observed properties of the various components of SMMJ14011 do not appear to require a strong-lensing interpretation. However gravitational magnification due to the known mass concentrations along our line-of-sight to SMMJ14011, i.e. J1c and A 1835, inevitably modify its observed properties. We quantify this effect, first considering the magnifying power of J1c alone. Adopting 55 km s<sup>-1</sup> as the line-of-sight stellar velocity dispersion of J1c, and a simple singular isothermal model ($`\sigma _{1\mathrm{D}}=\sqrt{4/3}\times 55`$ km s$`{}_{}{}^{1}=64`$ km s<sup>-1</sup>), the Einstein radius of J1c is just $`0.1^{\prime \prime }`$. In contrast, the observed separation of J1a/J1b from J1c is $`0.55^{\prime \prime }`$. To achieve an Einstein radius comparable with the observed separation, J1c would require a velocity dispersion of $`\sigma _{1\mathrm{D}}150`$ km s<sup>-1</sup> (as concluded by M05), implying an anomolously high mass-to-light ratio.
We now consider to what extent the contribution from the cluster potential can ameliorate this situation using the lenstool lens modelling code (Kneib et al. 1993, 1996; Smith et al. 2005, S05). The Schmidt et al. (2001) mass model for the cluster adopted by M05 uses ground-based optical photometry of selected arcs and the cluster’s X-ray emission to constrain a model of the cluster mass distribution. However, HST imaging now reveals one of the multiple-image systems (B and B) employed in their modelling by Schmidt et al. to be unreliable (S05). We therefore instead use the gravitational lens model developed using HST WFPC2 imaging by S05 to describe the mass distribution within the cluster. We refer the interested reader to S05 for full details of the model. The key features relevant to this study are (i) despite the presence of several gravitational arcs in A 1835, none of them have been spectroscopically identified to date (S05), hence the model is constrained by the weak shear signal detected in the outskirts of the WFPC2 frame, (ii) this weak shear signal was calibrated to 10% precision in projected mass using spectroscopically-confirmed strong-lensing clusters at the same redshift as A 1835, (iii) the family of acceptable models (defined by $`\mathrm{\Delta }\chi ^21`$) bracket Schmidt et al.’s models.
Taking our best-fit mass model for A 1835 we can now add a mass component at the position of J1c with ellipticity, position angle and velocity dispersion matching our estimated values (§2.2). The surface mass density of the combined A 1835+J1c lens is sub-critical at the position of J1a/J1b/J1d, i.e. the combined lens is not capable of multiple-imaging at the location of the submm emission. The same is true of models at the high and low-mass extremes allowed by the weak-shear data and the observational constraints on J1c. In summary, we estimate a conservative 1-$`\sigma `$ range for the amplification of the various components of: J1a $`\mu =2.8`$–4.9; J1b $`\mu =2.7`$–4.1; J1d $`\mu =2.7`$–4.0; J1n $`\mu =2.7`$–3.9; J2 $`\mu =3.0`$–5.0; the uncertainties are driven by the statistical error from the weak-shear constraints on A 1835. We therefore conclude that it is unlikely that SMM J14011 is strongly lensed, instead the different observed components are likely magnified by factors of $`\mu 3`$–5.
## 3. Discussion
Our analysis confirms the suggestion of DS03 that a component of the complex submm source, SMM J14011, is actually a foreground galaxy. However, we do not support their suggestion that the background far-infrared luminous source is multiply-imaged and hence highly-amplified by the foreground galaxy (in combination with the A 1835 cluster). There is clear spectroscopic evidence which rejects their proposed identification for the multiple images. Similarly, we have used the measured properties of the foreground galaxy, J1c, to reject a second lensing configuration suggested by M05. Both the observed ellipticity and estimated mass of the halo of J1c are very different from those required by M05’s model, while the flux ratios and morphologies of the proposed counter-images also do not agree with their predictions.
We suggest that J1c is probably a passive dwarf member of the A 1835 galaxy cluster. Our best estimates of the likely velocity dispersion of this galaxy, $`\sigma <55\pm 15`$ km s<sup>-1</sup>, indicate it is incapable of producing multiples images on the observed angular scales of J1a/J1b/J1d. Instead, we estimate a modest boosting of the amplification of these images, over that provided by the cluster alone. Assuming that the submm emission broadly traces the H$`\alpha `$ emission mapped by T04 gives a median amplification averaged over the source of $`\mu 3.5\pm 0.5`$, slightly higher than the original estimate assumed by I00 based purely on the lensing influence of the cluster potential. Thus SMM J14011 has an intrinsic submm flux of $`3.5\pm 0.5`$ mJy at 850 $`\mu `$m and a far-infrared luminosity of $`L_{FIR}4\times 10^{12}L_{}`$. SMM J14011 is an intrinsically luminous galaxy (Smail et al. 2002).
To study the structure of the background galaxy in more detail we can correct its observed shape (as displayed in the $`KJ`$ image in Fig. 1) for the distortion produced by the lens using our model. This indicates that the source is likely to have an intrinsic FWHM of $`0.5^{\prime \prime }`$ (4 kpc) and has a relatively circular morphology in the restframe optical (perhaps corresponding to a face-on orientation, which would help explain the relatively narrow CO line width for the system, DS03). The knots, J1a/J1b/J1d/J1n, visible in the HST imaging appear to be UV-bright clumps lying in a $`10`$ kpc region corresponding to the restframe optical extent of the galaxy and are most likely relatively unobscured star-forming regions within the galaxy. J2 represents a UV-bright companion with a separation of just $`20`$ kpc in projection. The close proximity of J2 to J1 suggests that their dynamical interaction may be responsible for the intense starburst currently underway in J1. The size and clumpy/multi-component morphology of SMM J14011 is thus very similar to that of typical submm galaxies with submm fluxes of $`5`$ mJy (Chapman et al. 2004; Smail et al. 2004).
We conclude that the intrinsic nature of SMM J14011 is much as originally stated by I00, at least regarding the long-wavelength properties (where J1c does not contribute). However, the identification of J1c as a foreground contaminant with only a weak lensing contribution does alter the conclusions of DS03 and M05 who both adopted strong-lensing models to correct the observed properties of SMM J14011. Using our prefered lens model, we reinterpret the apparent CO(7–6) size of SMM J14011 from DS03, who constrain the source to be $`2.3^{\prime \prime }\times <0.8^{\prime \prime }`$ (corrected for the beam), which corrected for lens amplification indicates an intrinsic CO(7–6) FWHM of $`6`$ kpc. This makes the CO emission in SMM J14011 comparable in size to the few SMGs studied at high-resolution with IRAM (Genzel et al. 2003; Tacconi et al. 2005). The discussion in DS03 then suggests that the source must plausibly consist of a series of compact knots distributed within the beam if the brightness temperature is not going to be too high.
In the case of M05, the main conclusion which changes as a result of our lens model is the estimated stellar mass of the system, which instead of being $`<10^9`$M (adopting an amplification of $`\mu 30`$), should be closer to $`<10^{10}`$M, comparable to the gas mass and dynamical mass limits derived from the CO line width (Frayer et al. 1999). Again this makes the restframe optically-derived stellar mass of SMM J14011 very similar to that of comparably luminous submm galaxies in the field (Smail et al. 2004), although we caution that these stellar mass estimates are highly uncertain due to the large extinction correction and sensitivity to the adopted ages and hence mass-to-light ratios of the stellar populations (Borys et al. 2005).
We note that the identification of J1c as a foreground galaxy will also alter the model of the near-infrared spectral energy distribution of this system derived from high-quality 2-dimensional spectroscopy by T04. They attempt to reproduce an apparent spectral break between the $`J`$\- and $`H`$-bands with the Balmer discontinuity in a young stellar population and derive a rough age of 200 Myrs and a reddening of $`A_V0.7`$ from this. From the observed $`K`$-band magnitude of J1c, we estimate that roughly half of the light in the $`K`$-band spectrum of T04 comes from the foreground galaxy and adopting reasonable near-infrared colors of $`(J_{AB}H_{AB})0.3`$ and $`(H_{AB}K_{AB})0.2`$ for this galaxy (assuming an unevolved early-type spectral energy distribution at $`z=0.25`$), that the true spectral break is likely to be stronger when the contribution from J1c is removed. This would likely decrease both the age and the reddening derived by T04 – bringing these closer to the values determined by M05. However, we stress that T04’s analysis of the baryonic mass of this system relies on the spectral rather than photometric properties of this system and hence is insensitive to the exact properties of J1c.
## 4. Conclusions
Our analysis of the spectral properties of J1c, a proposed component of the $`z=2.56`$ submm source SMM J14011, confirms that it is in fact a dwarf early-type member of the foreground cluster. However, our detailed modelling of the gravitational lensing contribution from J1c on the background submm source indicates that it is unlikely to significantly increase the amplification of the source over that from the cluster potential alone. We estimate conservative limits on the amplification of the various UV/near-infrared components of SMM J14011 in the range $`\mu 3`$–5. We conclude that SMM J14011 remains an intrinsically very luminous galaxy with properties that appear similar to comparably luminous systems now being studied in greater numbers at high redshifts.
We thank Dave Frayer, Jean-Paul Kneib, Nicole Nesvadba and Mark Swinbank for help and Andrew Baker, Dennis Downes and Kentaro Motoharo for useful discussions. We thank an anonymous referee for a report which helped improve the presentation of this paper. IRS acknowledges support from the Royal Society.
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# Study of the 𝐽/𝜓 decays to ΛΛ̄ and Σ⁰Σ̄⁰
## Abstract
With 58 million produced $`J/\psi `$ events collected by the BES-II detector at the BEPC, the decays $`J/\psi \mathrm{\Lambda }\overline{\mathrm{\Lambda }}`$ and $`\mathrm{\Sigma }^0\overline{\mathrm{\Sigma }}^0`$ are analysed. The branching ratios are measured to be $`Br(J/\psi \mathrm{\Lambda }\overline{\mathrm{\Lambda }})=(2.05\pm 0.03\pm 0.11)\times 10^3`$ and $`Br(J/\psi \mathrm{\Sigma }^0\overline{\mathrm{\Sigma }}^0)=(1.40\pm 0.03\pm 0.07)\times 10^3`$. The angular distribution is of the form $`\frac{dN}{dcos\theta }=N_0(1+\alpha cos^2\theta )`$, with $`\alpha =0.65\pm 0.12\pm 0.08`$ for $`J/\psi \mathrm{\Lambda }\overline{\mathrm{\Lambda }}`$ and $`\alpha =0.22\pm 0.17\pm 0.09`$ for $`J/\psi \mathrm{\Sigma }^0\overline{\mathrm{\Sigma }}^0,`$ respectively.
preprint: Draft-PRL
As is well-known from the helicity formalism, the angular distribution of $`B`$ in the decay of a neutral vector resonance $`V`$ into a baryon-antibaryon pair $`B\overline{B}`$ is given by 1
$$\frac{dN}{dcos\theta }1+\alpha cos^2\theta ,$$
where $`\theta `$ is the emission polar angle of $`B`$ in the $`V`$ rest frame. The first order calculations 2 3 of the perturbative QCD predict the theoretical value of $`\alpha `$ at $`J/\psi `$ energy. Table 1 summarizes the theoretical predictions of $`\alpha `$ for the decays $`J/\psi \mathrm{\Lambda }\overline{\mathrm{\Lambda }},\mathrm{\Sigma }^0\overline{\mathrm{\Sigma }}^0`$. Several experiments have measured $`\alpha `$ for $`J/\psi p\overline{p}`$ 4 5 6 7 , $`\mathrm{\Lambda }\overline{\mathrm{\Lambda }},\mathrm{\Sigma }^0\overline{\mathrm{\Sigma }}^0`$ 6 7 and $`\mathrm{\Xi }^{}\overline{\mathrm{\Xi }}^+`$ 6 .
In previous papers 8 9 , the analyses of $`J/\psi \mathrm{\Lambda }\overline{\mathrm{\Lambda }}`$ and $`\mathrm{\Sigma }^0\overline{\mathrm{\Sigma }}^0`$ using $`7.8\times 10^6J/\psi `$ events collected with the BES-I detector have been reported. The BES-I value of $`\mathrm{\Lambda }\overline{\mathrm{\Lambda }}`$ angular distribution coefficient $`\alpha `$ is in good agreement with the DM2 value 7 as well as the theoretical predictions in Ref.. However the BES-I value of $`\mathrm{\Sigma }^0\overline{\mathrm{\Sigma }}^0`$ angular distribution coefficient $`\alpha `$ with minus sign is obviously deviated from those of DM2 , MARKII as well as the theoretical prediction.
This letter presents the analyses of the decays $`J/\psi \mathrm{\Lambda }\overline{\mathrm{\Lambda }}`$ and $`\mathrm{\Sigma }^0\overline{\mathrm{\Sigma }}^0`$ based on BES-II $`58\times 10^6`$ $`J/\psi `$ events, with the purpose of improving the accuracy of branching ratio and $`\alpha `$ value measurements and clarifying the sign of the $`\alpha `$ for $`J/\psi \mathrm{\Sigma }^0\overline{\mathrm{\Sigma }}^0`$.
Table 1 theoretical expectations for $`\alpha `$ in $`J/\psi B\overline{B}.`$ Channel Ref. Ref. $`\mathrm{\Lambda }\overline{\mathrm{\Lambda }}`$ 0.32 0.51 $`\mathrm{\Sigma }^0\overline{\mathrm{\Sigma }}^0`$ 0.31 0.43
The BES-II detector has been described in detail elsewhere 10 .
Since the decays studied include $`\mathrm{\Lambda }\overline{\mathrm{\Lambda }}`$ pair in the final state, selection criteria are used to select $`J/\psi \mathrm{\Lambda }\overline{\mathrm{\Lambda }}+X`$ events, where X is a system of neutral particle(s) and/or undetected charged particle(s). The $`\mathrm{\Lambda }`$ is identified by its $`p\pi ^{}`$ decay mode. Candidates for $`J/\psi \mathrm{\Lambda }\overline{\mathrm{\Lambda }}+X`$ events are selected by requiring exactly four reconstructed charged tracks in the drift chamber with zero net charge. Tracks with $`cos\theta _{ch}<0.8`$ and transverse momemtum $`p_{xy}>0.07`$ GeV are accepted, where $`\theta _{ch}`$ is the polar angle with respect to the beam direction. Because the $`\mathrm{\Lambda }`$ is produced at the second vertex, no limits are required for the primary vertex position. The particles are identified by requiring that their combination weights of the time-of-flight (TOF) and the ionization energy loss (dE/dx) in the drift chamber be consistent with the corresponding particle hypothesis.
The events with four charged particles satisfying hypothesis of a pair of $`p\overline{p}`$ and a pair of $`\pi ^+\pi ^{}`$ are selected. Fig. 1 shows $`p\pi ^{}`$ invariant mass distribution for remaining $`J/\psi p\pi ^{}\overline{p}\pi ^++X`$ candidates (solid line). By fitting the $`M_{p\pi ^{}}`$ distribution to a gaussian distribution plus a quadratic polynomial, a 2.6 MeV mass resolution is estimated for the peak. Fig. 1 also shows $`p\pi ^{}`$ invariant mass for Monte Carlo simulation (dashed line) with the mass resolution of 2.6 MeV. To select $`\mathrm{\Lambda }`$ and $`\overline{\mathrm{\Lambda }}`$, the requirements of $`M_{p\pi ^{}}1.1156<0.008`$ GeV and $`M_{\overline{p}\pi ^+}1.1156<0.008`$ GeV are imposed.
The energy distribution of the $`\mathrm{\Lambda }\overline{\mathrm{\Lambda }}`$ pair, $`E_{\mathrm{\Lambda }\overline{\mathrm{\Lambda }}}`$ ,for $`J/\psi \mathrm{\Lambda }\overline{\mathrm{\Lambda }}+X`$ candidates is shown in Fig. 2, where $`E_{\mathrm{\Lambda }\overline{\mathrm{\Lambda }}}`$ is defined as the energy sum of $`\mathrm{\Lambda }`$ pair, data is represented by solid line and Monte Carlo simulation by dashed line.
The contamination from non-$`\mathrm{\Lambda }\overline{\mathrm{\Lambda }}+X`$ events in Fig. 2 can be estimated by the contribution from the sidebands, namely the off-$`\mathrm{\Lambda }`$ resonance background shown in Fig. 1. The events in the sidebands are defined as those in the squre area of $`M_{p\pi ^{}}1.1156<\sqrt{2}\times 3\sigma `$ and $`M_{\overline{p}\pi ^+}1.1156<\sqrt{2}\times 3\sigma `$ minus the squre area of $`M_{p\pi ^{}}1.1156<3\sigma `$ and $`M_{\overline{p}\pi ^+}1.1156<3\sigma `$ in two dimentional $`M_{p\pi ^{}}`$ vs. $`M_{\overline{p}\pi ^+}`$ plot, where $`\sigma =0.0026`$ GeV is the mass resolution of $`\mathrm{\Lambda }`$ $`(\overline{\mathrm{\Lambda }})`$. In Fig. 2 the sideband contribution is indicated by shaded area.
$`M_{p\pi ^{}}`$ (GeV)
Fig. 1. $`p\pi ^{}`$ invariant mass distribution in $`J/\psi p\pi ^{}\overline{p}\pi ^++X`$ candidates for data (solid line). The dashed line is Monte Carlo simulation for $`J/\psi \mathrm{\Lambda }\overline{\mathrm{\Lambda }}.`$
$`E_{\mathrm{\Lambda }\overline{\mathrm{\Lambda }}}`$ (GeV)
Fig. 2. Energy distribution of $`\mathrm{\Lambda }\overline{\mathrm{\Lambda }}`$ for $`J/\psi \mathrm{\Lambda }\overline{\mathrm{\Lambda }}+X`$ candidates. The shaded area is the sideband contamination described in the text. The dashed lines are from Monte Carlo simulation of $`J/\psi \mathrm{\Sigma }^0\overline{\mathrm{\Sigma }}^0`$ and $`J/\psi \mathrm{\Lambda }\overline{\mathrm{\Lambda }}`$.
A clear peak centered at the $`J/\psi `$ mass in Fig. 2 from the decay $`J/\psi \mathrm{\Lambda }\overline{\mathrm{\Lambda }}`$ is observed. The enhancement centered at 2.9 GeV is due to the decay $`J/\psi \mathrm{\Sigma }^0\overline{\mathrm{\Sigma }}^0`$,where $`\mathrm{\Sigma }^0\mathrm{\Lambda }\gamma `$. To obtain $`J/\psi \mathrm{\Lambda }\overline{\mathrm{\Lambda }}`$ from the selected $`J/\psi \mathrm{\Lambda }\overline{\mathrm{\Lambda }}+X`$ candidates, contamination from $`J/\psi \mathrm{\Sigma }^0\overline{\mathrm{\Sigma }}^0`$ decay is suppressed by imposing the total missing momentum cut $`p_{miss}<0.15`$ GeV and the $`\mathrm{\Lambda }`$ $`(\overline{\mathrm{\Lambda }})`$ momentum requirement $`0.98<p_\mathrm{\Lambda }`$ $`(p_{\overline{\mathrm{\Lambda }}})<1.2`$ GeV. Finally by requiring $`3.02E_{\mathrm{\Lambda }\overline{\mathrm{\Lambda }}}3.2`$ GeV a $`\mathrm{\Lambda }\overline{\mathrm{\Lambda }}`$ sample of 9462-503=8959 events is obtained, here 503 events are from the sidebands, which are obtained using the same criteria for $`\mathrm{\Lambda }\overline{\mathrm{\Lambda }}`$ selection to the aforementioned non-$`\mathrm{\Lambda }\overline{\mathrm{\Lambda }}+X`$ sideband events.
The detection efficiency for the events $`J/\psi \mathrm{\Lambda }\overline{\mathrm{\Lambda }}p\pi ^{}\overline{p}\pi ^+`$ depends on the $`\mathrm{\Lambda }`$ direction. A Monte Carlo simulation gives the detection efficiency $`ϵ(\theta _i)`$ in different $`\mathrm{\Lambda }`$ polar angle $`\theta _i`$ with the bin size $`\mathrm{\Delta }cos\theta _i=0.1`$. The efficiency corrected angular distribution of $`\mathrm{\Lambda }`$ for the decay $`J/\psi \mathrm{\Lambda }\overline{\mathrm{\Lambda }}`$ is shown as histogram line in Fig. 3. Fitting this angular distribution to the theoretical form
$$\frac{dN}{dcos\theta }=N_0(1+\alpha cos^2\theta ),$$
yields
$$N_0=(1990.9\pm 40.1)/0.1,$$
$$\alpha =0.65\pm 0.12,$$
where the error is statistical.
Using the $`N_0`$ and $`\alpha `$ values, the number of $`J/\psi \mathrm{\Lambda }\overline{\mathrm{\Lambda }}`$ events corrected for the efficiency $`ϵ_i`$ is deduced from the equation $`N_{\mathrm{\Lambda }\overline{\mathrm{\Lambda }}}=_1^1N_0(1+\alpha cos^2\theta )𝑑cos\theta ,`$ and the branching ratio is obtained by the equation
$$Br(J/\psi \mathrm{\Lambda }\overline{\mathrm{\Lambda }})=\frac{N_{\mathrm{\Lambda }\overline{\mathrm{\Lambda }}}}{Br^2(\mathrm{\Lambda }p\pi ^{})N_{J/\psi }},$$
where $`N_{J/\psi }=58\times 10^6`$ is the number of $`J/\psi `$ events with the error of $`4.72\%`$ 11 .
The remained contamination from $`J/\psi \mathrm{\Sigma }^0\overline{\mathrm{\Sigma }}^0`$ is $`0.8\%`$ using the branching ratio determined by this work. Subtracting this contamination the number of $`J/\psi \mathrm{\Lambda }\overline{\mathrm{\Lambda }}`$ signal events is 8887. The branching ratio and angular distribution coefficient are obtained to be
$$Br(J/\psi \mathrm{\Lambda }\overline{\mathrm{\Lambda }})=(2.05\pm 0.03\pm 0.11)\times 10^3,$$
$$\alpha =0.65\pm 0.12\pm 0.08,$$
where the first error is statistical and the second error is systematic. The systematic error includes the uncertainty of the detection efficiency due to imperfection of Monte Carlo simulation, the contamination from the indefinite branching ratio of background channel $`J/\psi \mathrm{\Sigma }^0\overline{\mathrm{\Lambda }}+c.c.`$ 12 and the uncertainty of the number of $`J/\psi `$ events.
$`cos\theta `$
Fig. 3. The $`\mathrm{\Lambda }`$ angular distribution for $`J/\psi \mathrm{\Lambda }\overline{\mathrm{\Lambda }}`$ events. The histogram is efficiency corrected data. The smooth dashed line is the fitting result.
The enhancement centered at 2.9 GeV in Fig. 2 is dominantly due to the decay $`J/\psi \mathrm{\Sigma }^0\overline{\mathrm{\Sigma }}^0`$,where $`\mathrm{\Sigma }^0\mathrm{\Lambda }\gamma `$. $`2.75<E_{\mathrm{\Lambda }\overline{\mathrm{\Lambda }}}<3.02`$ is required to select $`\mathrm{\Sigma }^0\overline{\mathrm{\Sigma }}^0`$ events. To remove the background from $`J/\psi \mathrm{\Lambda }\overline{\mathrm{\Lambda }}`$, $`\mathrm{\Lambda }\overline{\mathrm{\Lambda }}\pi ^0`$, $`\mathrm{\Xi }^0\overline{\mathrm{\Xi }}^0`$, $`\mathrm{\Lambda }\overline{\mathrm{\Lambda }}\gamma `$ and $`\mathrm{\Sigma }^0\overline{\mathrm{\Lambda }}+c.c.`$, the events are kinematically fitted to the $`J/\psi \mathrm{\Lambda }\overline{\mathrm{\Lambda }}\gamma \gamma `$ topology by imposing energy and momentum constraints (4C). The combination with the smallest $`\chi ^2`$ in the 4C fit is chosen to identify the radiative photons if there are more than two photon candidates in an event, here a photon is defined as a cluster with deposite energy larger than 30 MeV in the barrel shower counter, outside a $`25^{}`$ cone around $`\overline{p}`$ and outside a $`12^{}`$ cone around the other charged particles. Then a 6c fit ( two additional constraints $`M_{\mathrm{\Lambda }\gamma }=M_{\mathrm{\Sigma }^0}`$, $`M_{\overline{\mathrm{\Lambda }}\gamma }=M_{\overline{\mathrm{\Sigma }}^0}`$) is performed. It is required that $`\chi _{\mathrm{\Lambda }\overline{\mathrm{\Lambda }}\gamma \gamma }^2(4c)<\chi _{\mathrm{\Lambda }\overline{\mathrm{\Lambda }}\gamma }^2(4c).`$ The $`M_{\overline{\mathrm{\Lambda }}\gamma }`$ distribution is shown in Fig. 4, where the solid line is data and the dashed line is Monte Carlo simultaion $`J/\psi \mathrm{\Sigma }^0\overline{\mathrm{\Sigma }}^0.`$ The $`M_{\mathrm{\Lambda }\gamma }`$ distribution is similar to Fig. 4. Finally, after applying $`\chi _{\mathrm{\Lambda }\overline{\mathrm{\Lambda }}\gamma \gamma }^2(6c)<40`$, 2194-82=2112 events are remained, where 82 is from the sidebands which are obtained using the same criteria for $`\mathrm{\Sigma }^0\overline{\mathrm{\Sigma }}^0`$ selection to the aforementioned non-$`\mathrm{\Lambda }\overline{\mathrm{\Lambda }}+X`$ sideband events.
Table 2 The background for $`J/\psi \mathrm{\Sigma }^0\overline{\mathrm{\Sigma }}^0`$ Channel contamination ($`\%`$) $`\mathrm{\Lambda }\overline{\mathrm{\Lambda }}`$ 0.1 $`\mathrm{\Xi }^0\overline{\mathrm{\Xi }}^0`$ 0.3 $`\mathrm{\Lambda }\overline{\mathrm{\Lambda }}\pi ^0`$ 0.5 $`\mathrm{\Sigma }^0\overline{\mathrm{\Lambda }}+c.c.`$ $`<0.08`$ $`\mathrm{\Lambda }\overline{\mathrm{\Lambda }}\gamma `$ $`<0.08`$
The detection efficiency for the signal events depends on the $`\mathrm{\Sigma }^0`$ direction. The angular dependence of $`\mathrm{\Sigma }^0\overline{\mathrm{\Sigma }}^0`$ detection efficiency is obtained by selecting Monte Carlo events $`J/\psi \mathrm{\Sigma }^0\overline{\mathrm{\Sigma }}^0\mathrm{\Lambda }\gamma \overline{\mathrm{\Lambda }}\gamma p\pi ^{}\gamma \overline{p}\pi ^+\gamma `$ through the same selection criteria. The efficiency corrected angular distribution is shown in Fig. 5 as histogram. Fitting this angular distribution to the theoretical form
$$\frac{dN}{dcos\theta }=N_0(1+\alpha cos^2\theta ),$$
yields
$$N_0=(1786.9\pm 57.4)/0.1,$$
$$\alpha =0.22\pm 0.17,$$
where $`\theta `$ is the polar angle between the $`\mathrm{\Sigma }^0`$ and $`e^+`$ beam and the error is statistical. The number of $`J/\psi \mathrm{\Sigma }^0\overline{\mathrm{\Sigma }}^0`$ events corrected for the detection efficiency is $`N_{\mathrm{\Sigma }^0\overline{\mathrm{\Sigma }}^0}=_1^1N_0(1+\alpha cos^2\theta )𝑑cos\theta `$ and the branching ratio is calculated with
$$Br(J/\psi \mathrm{\Sigma }^0\overline{\mathrm{\Sigma }}^0)=\frac{N_{\mathrm{\Sigma }^0\overline{\mathrm{\Sigma }}^0}}{Br^2(\mathrm{\Lambda }p\pi ^{})N_{J/\psi }}.$$
Table 2 summarizes the remaining contamination to $`J/\psi \mathrm{\Sigma }^0\overline{\mathrm{\Sigma }}^0`$ candidates from background channels. The branching ratio of $`J/\psi \mathrm{\Lambda }\overline{\mathrm{\Lambda }}`$ is taken from this work, while the others are from the PDG. Subtracting all the contaminations of first three channels, the number of $`J/\psi \mathrm{\Sigma }^0\overline{\mathrm{\Sigma }}^0`$ signal events is 2093. The backgrounds of other two channels are considered in the systematic error. The results for the branching ratio of $`J/\psi \mathrm{\Sigma }^0\overline{\mathrm{\Sigma }}^0`$ and its angular distribution coefficeincy are
$$Br(J/\psi \mathrm{\Sigma }^0\overline{\mathrm{\Sigma }}^0)=(1.40\pm 0.03\pm 0.07)\times 10^3,$$
$$\alpha =0.22\pm 0.17\pm 0.09,$$
respectively, where the first error is statistical and the second systematic. The systematic error includes the uncertainty of detection efficiency due to imperfection of Monte Carlo simulation, the contaminations from the indefinite branching ratios of background channels $`J/\psi \mathrm{\Lambda }\overline{\mathrm{\Lambda }}\gamma `$ and $`J/\psi \mathrm{\Sigma }^0\overline{\mathrm{\Lambda }}+c.c.`$ and the uncertainty of the number of $`J/\psi `$ events.
Table 3 and Table 4 summarize the branching ratios and $`\alpha `$ values measured by this work for the decays $`J/\psi \mathrm{\Lambda }\overline{\mathrm{\Lambda }}`$ and $`J/\psi \mathrm{\Sigma }^0\overline{\mathrm{\Sigma }}^0`$ together with those previously reported by Mark II 6 and DM2 7 .
$`M_{\overline{\mathrm{\Lambda }}\gamma }`$
Fig. 4. The $`\overline{\mathrm{\Sigma }}^0`$ signal after 4c fit $`\chi _{\mathrm{\Lambda }\overline{\mathrm{\Lambda }}}^2(4c)<40`$. The solid line is data and the dashed line is Monte Carlo simultaion $`J/\psi \mathrm{\Sigma }^0\overline{\mathrm{\Sigma }}^0.`$
$`cos\theta `$
Fig. 5. The $`\mathrm{\Sigma }^0`$ angular distribution for $`J/\psi \mathrm{\Sigma }^0\overline{\mathrm{\Sigma }}^0`$. The histogram is efficiency corrected data. The smooth dashed line is the fitting result.
The statistics of the signal events in BES-II result is much better than those of Mark II, DM2 and BES-I. In this experiment the $`\alpha `$ value for $`J/\psi \mathrm{\Lambda }\overline{\mathrm{\Lambda }}`$ obtained by this experiment is consistent with the previous measurements, while the $`\alpha `$ value is negative for $`J/\psi \mathrm{\Sigma }^0\overline{\mathrm{\Sigma }}^0`$, which agrees with that of BES-I 9 while conflicts with the values of MARK-II and DM2 and theoretical expectation 2 3 .
In this analysis the 6c kinematic fit is used to select $`J/\psi \mathrm{\Sigma }^0\overline{\mathrm{\Sigma }}^0`$ and hence the contamination is small; while in DM2 7 and BES-I analyses 9 the $`\mathrm{\Sigma }^0`$ is not reconstructed and $`cos\theta `$ of $`\mathrm{\Sigma }^0`$ is replaced with that of $`\mathrm{\Lambda }`$. This makes $`\alpha `$ for $`J/\psi \mathrm{\Sigma }^0\overline{\mathrm{\Sigma }}^0`$ determined by this work rather believable.
It is worth noting that the central value of $`\alpha `$ reported by DM2 7 for $`J/\psi \mathrm{\Sigma }^0\overline{\mathrm{\Sigma }}^0`$ is positive, however, the angular distribution has a convex shape as a whole, with the only exception at $`\mathrm{cos}\theta =0.65`$.
Table 3
Experimental measurements for decay $`J/\psi \mathrm{\Lambda }\overline{\mathrm{\Lambda }}`$ Exp. $`\alpha `$ $`Br(\times 10^3)`$ Evts MARKII $`0.72\pm 0.36`$ $`1.58\pm 0.08\pm 0.19`$ 365 DM2 $`0.62\pm 0.22`$ $`1.38\pm 0.05\pm 0.20`$ 1847 BES-II $`0.65\pm 0.14`$ $`2.05\pm 0.03\pm 0.11`$ 8887
Table 4
Experimental measurements for $`J/\psi \mathrm{\Sigma }^0\overline{\mathrm{\Sigma }}^0`$ channel Exp. $`\alpha `$ $`Br(\times 10^3)`$ Evts MARK-II $`0.70\pm 1.10`$ $`1.58\pm 0.16\pm 0.25`$ 90 DM2 $`0.22\pm 0.31`$ $`1.06\pm 0.04\pm 0.23`$ 884 BES-II $`0.22\pm 0.19`$ $`1.40\pm 0.03\pm 0.07`$ 2093
The BES collaboration thanks the staff of BEPC for their hard efforts. This work is supported in part by the National Natural Science Foundation of China under contracts Nos. 10491300, 10225524, 10225525, 10425523, the Chinese Academy of Sciences under contract No. KJ 95T-03, the 100 Talents Program of CAS under Contract Nos. U-11, U-24, U-25, and the Knowledge Innovation Project of CAS under Contract Nos. U-602, U-34 (IHEP), the National Natural Science Foundation of China under Contract No. 10225522 (Tsinghua University), and the Department of Energy under Contract No.DE-FG02-04ER41291 (U Hawaii).
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# Revisiting the anomalous rf field penetration into a warm plasma
## I Introduction
A radio frequency electromagnetic wave does not penetrate into a plasma if the wave frequency $`\omega `$ is smaller than the electron plasma frequency $`\omega _p=\sqrt{4\pi e^2n_e/m}`$, where $`e`$ and $`m`$ are the electron charge and mass, respectively, and $`n_e`$ is the electron density. Electrons distribute their charge and current so as to shield out the electromagnetic wave. The shielding depends on the direction of the wave with regard to the plasma boundary. If the wave electric field is perpendicular to the plasma boundary, the rf field penetrates into the plasma only within a depth of the order of the Debye length $`v_T/\omega _p`$, where $`v_T=\sqrt{2T_e/m}`$ is the electron thermal velocity, determined by the electron temperature $`T_e`$, in eV. If the wave electric field is along the plasma boundary, the rf field penetrates into the plasma only within a depth of the order of the skin depth $`c/\omega _p`$, where $`c`$ is the speed of light in vacuum. Here, we consider a “collisionless” plasma, i.e. where the collision frequency is small compared to the wave frequency $`\nu \omega `$ and the electrons undergo rare collisions during the rf cycle; thus, collisions have little effect on wave screening by plasma.
Another important scale is the nonlocality or phase-mixing scale $`v_T/\omega `$, which determines the scale length of the electron current profile in the plasma. To demonstrate the concept of phase-mixing scale $`v_T/\omega `$ let us consider a simple model, where an electron acquires a prescribed velocity kick near the plasma boundary, in the direction perpendicular to the boundary
$$dv_x(t)=\mathrm{\Delta }V\mathrm{exp}(i\omega t).$$
(1)
The electron velocity at a distance $`x`$ from the boundary will be determined by the moment when velocity kick was acquired at the plasma boundary, i.e., by the time $`tx/v_x`$. The electron current in the plasma is given by an integration over all electrons with a velocity distribution function $`f(v_x)`$
$$j(x,t)=e\mathrm{\Delta }V_0^{\mathrm{}}f(v_x)\mathrm{exp}[i\omega (tx/v_x)]𝑑v_x.$$
(2)
Here, only electrons collided with the wall ($`v_x>0`$) have to be taken into account. For a Maxwellian distribution function $`f(v_x)=n_0e^{v_x^2/v_T^2}/v_T\sqrt{\pi }`$, the plasma current in Eq.(2) becomes
$$j(x,t)=\frac{j_0e^{i\omega t}}{\sqrt{\pi }}_0^{\mathrm{}}\mathrm{exp}\left(s^2+\frac{i\omega x}{v_Ts}\right)𝑑s,$$
(3)
where $`s=v_x/v_T`$ and $`j_0=en_0\mathrm{\Delta }V`$. The amplitude and phase of the current are shown in Fig.1. In the limit $`\omega x/v_T1`$, the integration in Eq. (3) can be performed analytically making use of the method of steepest descend Brillouin , see Appendix A for more details. This gives
$$j(x,t)\frac{j_0}{\sqrt{3}}\mathrm{exp}\left[i\omega t\frac{3}{4}\left(\frac{x}{\lambda _\omega }\right)^{2/3}+i\frac{3\sqrt{3}}{4}\left(\frac{x}{\lambda _\omega }\right)^{2/3}\right],$$
(4)
where $`\lambda _\omega =v_T/\sqrt{2}\omega `$ is the phase-mixing scale. Comparison of the asymptotic calculation result given by Eq. (4) with the exact result of numerical integration in Eq. (3) is shown in Fig. 1. From Fig. 1, it is evident that Eq. (4) approximates the exact result for any $`x`$ within a 15 percent error bar. The largest error occurs at $`x=0`$, where half of the electron population with velocity $`v_x>0`$ acquired the velocity kick, which gives rise to the electron current $`j(0)=j_0/2`$, whereas Eq. (4) predicts $`j(0)=j_0/\sqrt{3}`$, which corresponds to a 15 percent error.
Equation (4) describes the process of phase mixing - electrons with velocities different by $`\delta v_xv_T`$ have different phase lag of the order $`\omega x/v_T`$ at a distance $`x`$ from the plasma boundary. Therefore, at $`xv_T/\omega `$ the phase difference becomes considerable: contributions to the total current from electrons with different velocities $`v_x`$ cancel out each other, and the plasma current vanishes. Interestingly, the spatial profile of the current is not a simple exponential function, but an exponential function of $`\left(x/\lambda _\omega \right)^{2/3}`$. As it will be shown below this is typical for the spatial profiles of the electric field and electron current in warm plasmas due to nonlocal effects.
So far, we solved only test-particle problem and did not take into account the plasma polarization. The current in Eq. (3) is nonuniform; thus, there must be an electron density perturbation according to the continuity equation
$$e\frac{n_e}{t}=\frac{j}{x}.$$
(5)
The electron density perturbations polarize the plasma and generate an electric field, which in turn, affects the electron motion and the electron current profile. Thus, Eq.(3) has to be modified to include the self-consistent electric field. This requires solving the Vlasov equation together with the Poisson equation. In his famous 1946 paper, Landau obtained an analytic solution for the penetration of the longitudinal rf electric field into a plasma Landau . Note that he also described “Landau damping” in the same paper. We briefly review his solution for a small amplitude electric field in the linear approximation and discuss the more realistic case of a large amplitude electric field.
The structure of this review is as follows: In section II, the penetration of the longitudinal electric field into the plasma is described. This case corresponds to a capacitively coupled plasma. In section III, the penetration of the transverse electric field into the plasma is studied, which corresponds to an inductively coupled plasma. In subsection III.E, it is shown that anisotropy of the electron velocity distribution function can have a profound effect on the anomalous skin effect.
## II Penetration of the rf electric field directed perpendicular to the plasma boundary (capacitively-coupled plasma)
### II.1 Small-amplitude electric field
In the previous section, we considered a test particle current driven by artificially applied velocity modulations at the plasma boundary. Here, self-consistent penetration of a small amplitude rf electric field directed perpendicular to the plasma boundary is considered. Such a model provides some insight into the sheath structure of capacitively-coupled plasmas.
First, let’s consider a stationary negatively biased electrode. It is well-known that the externally applied electric field penetrates inside the plasma over distances of the order of the Debye length $`a=v_T/\sqrt{2}\omega _p=\sqrt{T_e/4\pi e^2n_0}`$. The plasma electrons are trapped by the plasma potential, $`\varphi (x)`$, in the potential well $`e\varphi (x)`$. The electron density obeys the Boltzmann distribution
$$n_e(x)=n_0\mathrm{exp}\left[e\varphi (x)/T_e\right].$$
(6)
The Poisson equation
$$\frac{d^2\varphi }{dx^2}=4\pi e(n_in_e)$$
(7)
can be simplified assuming small potential variations $`e\varphi (x)/T_e1`$ and a uniform background plasma with $`n_e=n_i=n_0`$. Thus, Eq. (7) becomes
$$\frac{d^2\varphi }{dx^2}=\frac{4\pi e^2n_0}{T_e}\varphi .$$
(8)
The solution of Eq. (8) is an exponentially decaying electric field $`E=d\varphi /dx`$
$$E=E_0\mathrm{exp}\left(\frac{x}{a}\right).$$
(9)
Here, $`E_0`$ is the value of the electric field at the plasma boundary. This is the solution for a steady state, time-independent sheath electric field. In the opposite case of the time-dependent electric field, the Boltzmann distribution given by Eq. (6) is no longer valid and the electron density has to be determined from the Vlasov equation. Landau solved the Vlasov equation coupled with the Poisson equation analytically in the linear approximation considering an electrostatic wave with small amplitude $`|e\varphi (x)|/T_e1`$ and small frequency $`\omega \omega _p`$ Landau . Details of the solution are described in Appendix B.
To summarize, the solution can be separated into three parts,
$$E_x(x,t)=\left[E_0\mathrm{exp}\left(\frac{x}{a}\right)+E_b+E_t(x)\right]e^{i\omega t}.$$
(10)
Here, $`E_0`$ is the amplitude of the electric field at the plasma boundary, $`E_b=E_0/\epsilon `$ is the electric field in the plasma bulk far away from the sheath region, $`\epsilon =1\omega _p^2/\omega ^2`$ is the dielectric constant of the cold plasma, and $`E_t(x)`$ is the electric field in a transient region with a spatial length of order $`v_T/\omega `$. The first term is the Debye screening of the external electric field. The second part describes a small, uniform electric field penetrating into the plasma far away from the boundary. The second and third terms are absent for a stationary applied electric field and appear only in the case of the rf electric field. The solution for the transient electric field $`E_t(x)`$ profile is derived in Appendix B and is given by
$$E_t(x)=\frac{2E_0}{\pi }_0^{\mathrm{}}\frac{1}{k}\frac{Im[\epsilon _{}(\omega ,k)]}{\epsilon _{}^{}(\omega ,k)\epsilon _{}(\omega ,k)}e^{ikx}𝑑k,$$
(11)
where $`\epsilon _{}(\omega ,k)`$ is the longitudinal plasma permittivity ($`𝐄𝐤`$),
$$\epsilon _{}(\omega ,k)1+\frac{2\omega _p^2}{k^2v_T^2}\left[1+\frac{\omega }{kV_T}Z\left(\frac{\omega }{kV_T}\right)\right],$$
(12)
and $`Z(\zeta )`$ is the plasma dispersion function Plasma Formulary
$$Z(\xi )=\pi ^{1/2}_{\mathrm{}}^{\mathrm{}}𝑑t\frac{\mathrm{exp}\left(t^2\right)}{t\xi },Im\xi >0.$$
(13)
In the limit $`x\lambda _\omega `$ only small $`k`$ contribute to the integral and $`\epsilon _{}(\omega ,k)`$ can be substituted by $`\epsilon _{}^{}(\omega ,0)\epsilon `$ in the denominator of Eq. (14), which gives
$$E_t(x)E_{appr}(x)=\frac{2E_0}{\pi \epsilon ^2}_0^{\mathrm{}}\frac{1}{k}Im[\epsilon _{}(\omega ,k)]e^{ikx}𝑑k,$$
(14)
Application of the method of steepest descend to Eq. (14) yields Landau
$$E_t(x)E_{std}(x)=\frac{2E_0}{\sqrt{3}\epsilon ^2}\frac{\omega _p^2}{\omega ^2}\left(\frac{x}{\lambda _\omega }\right)^{2/3}\mathrm{exp}\left[\left(\frac{3}{4}+i\frac{3\sqrt{3}}{4}\right)\left(\frac{x}{\lambda _\omega }\right)^{2/3}i\pi /3\right],$$
(15)
where $`\lambda _\omega =v_T/\sqrt{2}\omega `$ is the phase-mixing scale. The plots of the amplitude and phase of the electric field profile $`E_t(x)`$ given by Eq.(11) and the approximate $`E_{appr}(x)`$ given by Eq. (14), and asymptotic analytical result $`E_{std}(x)`$ given by Eq. (15) are shown in Fig. 2. Figure 2 shows that the steepest descend method given by Eq. (15) closely approximates Eq. (14) already for $`x>v_T/\omega `$. However, the both asymptotic solutions in Eq. (14) and Eq.(15) approximate the full solution in Eq. (11) only for very large $`x>40V_T/\omega `$. This is due to the made substitution $`\epsilon (\omega ,k)`$ by $`\epsilon (\omega ,0)`$, which results in a considerable error for $`k\omega /v_T`$ or $`xv_T/\omega `$.
*It follows from Eq. (15) that the electric field amplitude at $`x>v_T/\omega `$ is of order $`E_0/\epsilon `$, i.e., it is comparable with the electric field far away from the boundary (*$`E_tE_b`$*)*.
The origin of the electric field $`E_t(x)`$ can be explained by analyzing the individual electron dynamics. After passing through the region of the rf field, an electron acquires changes $`\mathrm{\Delta }\epsilon (v_x)`$ in energy and $`\mathrm{\Delta }u(v_x)`$ in velocity
$`\mathrm{\Delta }\epsilon (v_x)`$ $`=`$ $`{\displaystyle _{\mathrm{}}^{\mathrm{}}}v_xeE[x(t),t]𝑑t,`$ (16)
$`\mathrm{\Delta }u(v_x)`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Delta }\epsilon }{mv_x}}.`$
Here, the electron trajectory is $`x(t)=v_xt`$, $`v_x=|v_x|sgn(t)`$, and the electric field profile is given by Eq. (10). The total velocity kick is the summation over velocity kicks due to exponential, bulk and transitional electric fields
$$\mathrm{\Delta }u(v_x)=\mathrm{\Delta }u_0+\mathrm{\Delta }u_b+\mathrm{\Delta }u_t.$$
(17)
Substituting an exponential electric field into Eq.(16) gives the corresponding electron velocity kick
$$\mathrm{\Delta }u_0(v_x)\frac{2eE_0}{m\omega }\frac{\omega ^2a^2}{v_x^2+(\omega a)^2}.$$
(18)
Substituting the uniform electric field $`E_b`$ into Eq. (16), gives the electron velocity
$$\mathrm{\Delta }u_b(v_x)\frac{2eE_0}{im\omega \epsilon }.$$
(19)
This calculation can also be explained as follows: An electron has the oscillating velocity $`\mathrm{\Delta }u_s=eE_bi/m\omega `$ in a uniform rf electric field and a thermal velocity $`v_x`$. After a collision with the wall, an electron changes its velocity direction. If the initial average velocity was $`v_x<0`$, after the collision with the wall with specular reflection, the new average velocity $`v_x^{}>0`$ will change according to
$$v_x^{}+\mathrm{\Delta }u_s(t)=[v_x+\mathrm{\Delta }u_s(t)]$$
(20)
or the average velocity changes to
$$v_x^{}=v_x2\mathrm{\Delta }u_s(t),$$
(21)
which results in the effective velocity kick of Eq. (19).
The origin of the electric field in the transition region $`E_t(x)`$ is due to the plasma polarization. The velocity perturbations $`\mathrm{\Delta }u_s(v_x,t)`$ produce bunches in the electron density, which, in turn, generate the electric field $`E_t(x)`$. The decay of the electric field $`E_t(x)`$ is due to phase mixing similarly to the test-particle case in Eq. (4). Thus, generation of the transitional electric field $`E_t(x)`$ can be considered as a plasma self-consistency effect.
The electric field $`E_t(x)`$ generates a significant portion of the total velocity kick and thus noticeably influences the electron heating in the rf electric field. Figure 3 shows the amplitude of the electron velocity kick $`\mathrm{\Delta }u(v_x)`$ due to the interaction with the electric field given by Eq. (10). Electrons with small velocities $`v_x\omega a=v_T\omega /\omega _p`$ pick up a large velocity kick due to the exponential electric field $`E_0\mathrm{exp}(x/ai\omega t)`$, $`\mathrm{\Delta }u\mathrm{\Delta }u_02eE_0/m\omega `$. For very large electron velocities $`v_xv_T`$, the velocity kick given by Eq. (18) becomes small and the main contribution to the velocity kick comes from the uniform electric field $`E_b=E_0e^{i\omega t}/\epsilon `$ and the collision with the wall, $`\mathrm{\Delta }u\mathrm{\Delta }u_b2eE_0/m\omega \epsilon `$. In the intermediate range of velocities $`v_xv_T`$, the account of the electric field $`E_t(x)`$ is important, as in this case $`\mathrm{\Delta }u_t\mathrm{\Delta }u_b`$. As is evident from Fig. 3, taking this electric field $`E_t(x)`$ into account results in a considerable reduction of the electron velocity kick $`v_xv_T`$ for the bulk of the electron population compared with the case when this electric field is not taken into account. Note that most models neglect the electric field $`E_t(x)`$, see for example Libermann89 ,Lieberman& Godyak review .
### II.2 Large amplitude electric field
In many practical applications, the value of the external electric field is large: the potential drop in the sheath region $`V_{sh}`$ is typically of the order of hundreds of Volts and is much larger than the electron temperature $`T_e`$, which is of the order of a few Volts; consequently the electric field penetration has to be treated nonlinearly.
In the limit $`V_{sh}T_e`$, a wall is charged negatively all time with an alternating charge in a manner to conduct an ac current, driven by an external electric circuit. A negative charge pushes electrons away from the electrode up to a distance where the negative electric field is screened by a positive ion density. As $`V_{sh}T_e`$, the sheath width is much larger than the Debye length and the plasma sheath boundary can be considered as infinitely thin. The position of the boundary is determined by the condition that the external electric field is screened in the sheath regions when and where electrons are absent Libermann89 ; Me and Tsendin 1992 1 .
Electron interactions with the sheath electric field are traditionally treated as collisions with a moving potential barrier (wall). It is well known that multiple electron collisions with an oscillating wall result in electron heating, provided there is sufficient phase-space randomization in the plasma bulk. It is common to describe the sheath heating by considering electrons as test particles, and neglecting the plasma electric field Lieberman& Godyak review . As was pointed out in Refs. Libermann89 ; Me and Tsendin 1992 2 ; Aliev and me accounting for the electric field in the plasma reduces the electron sheath heating, and the electron sheath heating vanishes completely in the limit of uniform plasma density. Therefore, an accurate description of the rf fields in the bulk of the plasma is necessary for calculating the sheath heating. The electron velocity is oscillatory in the sheath, and as a result of these velocity modulations, the electron density bunches appear in the region adjacent to the sheath, similar to the previously described case of small-amplitude wave, see Fig. 4. These electron density perturbations decay due to phase mixing over a length of order $`v_T/\omega ,`$ where $`v_T`$ is the electron thermal velocity, and $`\omega `$ is the frequency of the electric field. The electron density perturbations polarize the plasma and produce an electric field in the plasma bulk. This electric field, in turn, changes the velocity modulations and correspondingly influences the electron density perturbations. Therefore, electron sheath heating has to be studied in a self-consistent nonlocal manner assuming a finite-temperature plasma.
Notwithstanding the fact that particle-in-cell simulations results have been widely available for the past decade Sommerer ; Surendra PRL , a basic understanding of the electron heating by the sheath electric field is being incomplete, because no one has studied the electric field in the plasma bulk using a kinetic approach, similar to the anomalous skin effect for the inductive electric field Lifshitz and Pitaevskii . In this regard, analytical models are of great importance because they shed light on the most complicated features of collisionless electron interactions with the sheath. In Ref.My PRL 2002 , an analytical model was developed to explore the effects associated with the self-consistent non-local nature of this phenomenon.
One of the approaches to study electron sheath heating is based on a fluid description of the electron dynamics. For the collisionless case, closure assumptions for the viscosity and heat fluxes are necessary. In most cases, the closure assumptions are made empirically or phenomenologically Surendra PRL , Gozadinos . The closure assumptions have to be justified by direct comparison with the results of kinetic calculations as is done, for example, in Refs. Hammett ; Furkal . Otherwise, inaccurate closure assumptions may lead to misleading results as discussed below.
To model the sheath-plasma interaction analytically, the following simplifying assumptions have been adopted in Ref. My PRL 2002 . The discharge frequency is assumed to be small compared with the electron plasma frequency. Therefore, most of the external electric field is screened in the sheath region by an ion space charge. The ion response time is typically larger than the inverse discharge frequency, and the ion density profile is quasi-stationary. There is an ion flow from the plasma bulk towards the electrodes. In the sheath region, ions are being accelerated towards the electrode by the large sheath electric field, and the ion density in the sheath region is small compared with the bulk ion density. In the present analytical treatment, the ion density profile is assumed fixed and is modelled in a two-step approximation: the ion density $`n_b`$ is uniform in the plasma bulk, and the ion density in the sheath $`n_{sh}<n_b`$ is also uniform (see Fig. 5). At the sheath-plasma boundary, there is a stationary potential barrier for the electrons ($`e\mathrm{\Phi }_{sh}`$), so that only the energetic electrons reach the sheath region. The potential barrier is determined by the quasineutrality condition, i.e., when the energetic electrons enter the sheath region, their instantaneous density is equal to the ion density \[$`n_e(\mathrm{\Phi }_{sh})=n_{sh}`$\].
The electron density profile is time-dependent in response to the time-varying sheath electric field. The large sheath electric field does not penetrate into the plasma bulk. Therefore, the quasineutrality condition holds in the plasma bulk, i.e., the electron density is equal to ion density, $`n_e=n_b.`$ In the sheath region, the electrons are reflected by the large sheath electric field. Therefore, $`n_e=n_{sh}`$ for $`x>x_{sh}(t)`$, and $`n_e=0`$ for $`x<x_{sh}(t)`$, where $`x_{sh}(t)`$ is the position of the plasma-sheath boundary Libermann89 . From Maxwell’s equations it follows that $`𝐉=0`$, where the total current $`𝐉`$ is the sum of the displacement current and the electron current. In the one-dimensional case, the condition $`𝐉=0`$ yields the conservation of the total current Landau ; Me and Tsendin 1992 1 :
$$en_eV_e+\frac{1}{4\pi }\frac{E_{sh}}{t}=j_0\mathrm{sin}(\omega t+\varphi ),$$
(22)
where $`j_0`$ is the amplitude of the rf current controlled by an external circuit and $`\varphi `$ is the initial phase. In the sheath, electrons are absent in the region of large electric field, and Eq.(22) can be integrated to give Me and Tsendin 1992 1
$$E_{sh}(x,t)=\frac{4\pi j_0}{\omega }[1\mathrm{cos}(\omega t+\varphi )]+4\pi |e|n_{sh}x,x<x_{sh}(t)$$
(23)
where Poisson’s equation has been used to determined the spatial dependence of the sheath electric field. The first term on the right-hand side of Eq. (23) describes the electric field at the electrode and the second term relates to the ion space charge screening of the sheath electric field. The position of the plasma-sheath boundary $`x_{sh}(t)`$ is determined by the zero of the sheath electric field, $`E_{sh}[x_{sh}(t),t]=0`$. From Eq. (23) it follows that
$$x_{sh}(t)=\frac{V_{sh0}}{\omega }[1+\mathrm{cos}(\omega t+\varphi )],$$
(24)
where $`V_{sh0}=j_0/(en_{sh})`$ is the amplitude of the plasma-sheath boundary velocity. The ion flux on the electrode is small compared with the electron thermal flux. Because electrons attach to the electrode, the electrode surface charges negatively, so that in a steady-state discharge, the electric field at the electrode is always negative, preventing an electron flux on the electrode. However, for a very short time ($`\omega t_n+\varphi \pi (1+2n)`$) the sheath electric field vanishes, allowing electrons to flow to the electrode for compensation of the ion flux. Note that there is a large difference between the sheath structure in the discharge and the sheath for obliquely incident waves interacting with a plasma slab without any bounding walls. Because electrodes are absent, electrons can move outside the plasma, and the electric field in the vacuum region, $`E_{sh}(x,t)=(4\pi j_0/\omega )\mathrm{cos}(\omega t+\varphi )`$, may have an alternating sign. Therefore, electrons may penetrate into the region of large electric field during the time when $`E_{sh}(x,t)>0`$ Brunel ; Yang . In the discharge, however, because the sheath electric field given by Eq. (23) always reflects electrons, the electrons *never* enter the region of the large sheath electric field, which is opposite to the case of obliquely incident waves.
The calculations based on the two-step ion density profile model are known to yield discharge characteristics in good agreement with experimental data and full-scale simulations Orlov .
For analytical calculation of the rf electric field inside the plasma, a linear approximation is used for the plasma conductivity. The validity of the linear approximation is based on the fact that the plasma-sheath boundary velocity and the mean electron flow velocity are small compared with the electron thermal velocity, $`V_{sh}v_T`$, Me and Tsendin 1992 1 ; Sommerer . The important spatial scale is the length scale for phase mixing, $`\lambda _\omega `$. The sheath width satisfies $`2V_{sh0}/\omega \lambda _\omega `$ because $`V_{sh}v_T`$. Therefore, the sheath width is neglected, and electron interactions with the sheath electric field are treated as a boundary condition. The collision frequency ($`\nu `$) is assumed to be small compared with the discharge frequency ($`\nu \omega `$), and correspondingly the mean free path is much larger than the length scale for phase mixing. Therefore, the electron dynamics is assumed to be collisionless. The discharge gap is considered to be sufficiently large compared with the electron mean free path, so that the influence of the opposite sheath is neglected. The effects of a finite gap width have been discussed in Refs. Me PRL 1999 ; Ulrich and me .
The electron interaction with the large electric field in the sheath is modelled as a collision with a moving oscillating rigid barrier with velocity $`V_{sh}(t)=dx_{sh}(t)/dt`$. After a collision with the plasma-sheath boundary - modelled as a rigid barrier moving with velocity $`V_{sh}(t)`$ \- an electron with initial velocity $`u`$ acquires a velocity $`u+2V_{sh}`$. Therefore, the power deposition density transfer from the oscillating plasma-sheath boundary is given by Libermann89
$$P_{sh}=\frac{m}{2}_{V_{sh}}^{\mathrm{}}𝑑u\left[u+V_{sh}(t)\right]\left[(2V_{sh}(t)+u)^2u^2\right]f_{sh}(u,t),$$
(25)
where $`m`$ is the electron mass, $`f_{sh}(u,t)`$ is the electron velocity distribution function in the sheath, and $`\mathrm{}`$ denotes a time average over the discharge period. Introducing a new velocity distribution function $`g(u^{},t)=f_{sh}[uV_{sh}(t),t]`$, Eq. (25) yields
$$P_{sh}=2mV_{sh}(t)_0^{\mathrm{}}u^2g(u^{},t)𝑑u^{},$$
(26)
where $`u^{}=uV_{sh}`$ is the electron velocity relative to the oscillating rigid barrier. From Eq.(26) it follows that, if the function $`g(u^{})`$ is stationary, then ($`P_{sh}=0`$) and there is no collisionless power deposition due to electron interaction with the sheath Libermann89 ; Gozadinos ; Raizer book . For example, in the limit of a uniform ion density profile $`n_{sh}=n_b`$, $`g(u^{})`$ is stationary (*in an oscillating reference frame of the plasma-sheath boundary*), and the electron heating vanishes Libermann89 , Me and Tsendin 1992 1 . Indeed, in the plasma bulk, the displacement current is small compared with the electron current, and from Eq. (22) it follows that the electron mean flow velocity in the plasma bulk, $`V_b(t)=j_0\mathrm{sin}(\omega t+\varphi )/en_b`$, is equal to the plasma-sheath velocity $`V_{sh}(t)`$, from Eq. (24). *Therefore, the electron motion in the plasma is strongly correlated with the plasma-sheath boundary motion.* From the electron momentum equation it follows that there is an electric field, $`E_b=m/edV_b(t)/dt`$, in the plasma bulk. In a frame of reference moving with the electron mean flow velocity, the sheath barrier is stationary, and there is no force acting on the electrons, because the electric field is compensated by the inertial force ($`eE_bmdV_b(t)/dt=0)`$. Therefore, electron interaction with the sheath electric field is totally compensated by the influence of the bulk electric field, and the collisionless heating vanishes Me and Tsendin 1992 2 . The example of a uniform density profile shows the importance of a self-consistent treatment of the collisionless heating in the plasma. If the function $`g(u^{},t)`$ is nonstationary, there is net power deposition. In Ref. My PRL 2002 , a kinetic calculation is performed to yield the correct electron velocity distribution function $`g(u^{},t)`$ and, correspondingly, the net power deposition.
The electron motion is different for low-energy electrons with an initial velocity in the plasma bulk $`|u|<u_{sh}`$, where $`u_{sh}^2=2e\mathrm{\Phi }_{sh}/m`$ and for energetic electrons with velocity $`|u|>u_{sh}`$. The low energy electrons with initial velocity $`u`$ in the plasma bulk are reflected from the stationary potential barrier $`e\mathrm{\Phi }_{sh}`$, and then return to the plasma bulk with velocity $`u`$. High energy electrons enter the sheath region with velocity $`u_1=(u^2u_{sh}^2)^{1/2}`$. They acquire a velocity $`u_2=2V_{sh}u_1`$ after collision with the moving rigid barrier, and then return to the plasma bulk with a velocity $`(u_2^2+u_{sh}^2)^{1/2}`$ multiple collisions .
As the electron velocity is modulated in time during reflections from the plasma-sheath boundary, so is the energetic electron density (by continuity of the electron flux). This phenomenon is identical to the mechanism of klystron operation klystron . The perturbations in the energetic electron density yield an electric field in the transition region adjusted to the sheath, see Fig.4.
The solution for the electric field $`E_t(x)`$ was obtained analytically in Ref.My PRL 2002 . Similar to the previous section, the solution is an expression for the inverse Fourier transform. It cannot be represented in an analytical form and has to be simulated numerically. This simulation has been performed for $`n_{sh}/n_b=1/3`$, $`\omega /\omega _p=1/100`$, and a Maxwellian electron distribution function. The electric field profile is close to $`E_t(x)E_{t0}\mathrm{exp}(x/\lambda _c)`$, where $`E_{t0}=0.72T_e/\lambda _\omega `$, and $`\lambda _c=(0.19+0.77i)\lambda _\omega `$ for $`x<6V_T/\omega `$. For $`x>6V_T/\omega `$, the electric field profile is no longer a simple exponential function, which is similar to the case considered in the previous section. The difference in phase of the currents of the energetic and low-energy electrons was observed in Ref.Surendra PRL , but it was misinterpreted as the generation of the electron acoustic waves. Electron acoustic waves can be excited if there is a complex value of $`k`$, with small damping $`\mathrm{𝐼𝑚}(k)\mathrm{𝑅𝑒}(k),`$ which is the root of the plasma dielectric function $`\epsilon (\omega ,k)=0`$ for a given $`\omega `$. For a Maxwellian electron distribution function, such root does not exist when $`\omega \omega _p`$. However, the electron acoustic waves can exist if the plasma contains two groups of electrons which have very different temperatures Mace . The wave phase velocity is $`\omega /k=\sqrt{n_c/n_h}\sqrt{T_h/m}`$ , where $`n_c`$ and $`n_h`$ are the electron densities of cold and hot electrons, respectively, and $`T_h`$ is the temperature of the hot electrons. The electron acoustic waves are strongly damped by the hot electrons, unless $`n_cn_h`$ and $`T_cT_h`$ , where $`T_c`$ is the electron temperature of the cold electrons Mace . In the opposite limit, $`n_c>4n_h`$, the electron acoustic waves do not exist Mace . In capacitively-coupled discharges, the electron population does stratify into two populations of cold and hot electrons, as has been observed in experiments Godyak 2EDF and simulation studies cold electron formation ; My CCP . Cold electrons trapped by the plasma potential in the discharge center do not interact with the large electric fields in the sheath region and have low temperature. Moreover, because of the nonlinear evolution of plasma profiles, the cold electron density is much larger than the hot electron density cold electron formation . Therefore, weakly-damped electron acoustic waves do not exist in the plasma of capacitively-coupled discharges. Reference Surendra PRL used the fluid equation and neglected the effect of collisionless dissipation, thus arriving at the incorrect conclusion about the existence of weakly-damped electron acoustic waves.
The power deposition is given by the sum of the power transferred to the electrons by the oscillating rigid barrier in the sheath region and by the electric field in the transition region,
$$P_{tot}=P_{sh}+P_{tr}.$$
(27)
Note that $`P_{tr}`$ can be negative. Calculations making use of the Vlasov equation yield My PRL 2002
$$P_{tot}=_0^{\mathrm{}}muD_u(u)\frac{df_0}{du}𝑑u,$$
(28)
where
$$D_u(u)=\frac{u|du|^2}{4}$$
(29)
is the diffusion coefficient in velocity space, and $`du`$ is the change in the electron velocity after passing through the transition and sheath regions,
$$du=2iV_b\left[\frac{u^{}}{u}\frac{n_b}{n_{sh}}\mathrm{\Theta }(|u|u_{sh})1\right]+\frac{eE_t(k=\omega /u)}{u},$$
(30)
where $`E_t(k)`$ is the Fourier transform of the electric field $`E_t(x)`$. First term describes the velocity acquired by fast electrons ($`|u|>u_{sh}`$) in collisions with the sheath; the second is due to the bulk electric field $`E_b`$ and collisions with either the potential barrier $`\mathrm{\Phi }_{sh}`$ or sheath; and the third is due the electric field in the transitional region $`E_t(x)`$. A plot of $`|du|^2/4`$ is shown in Fig. 6. Taking into account the electric field in the plasma (both $`E_b`$ and $`E_t`$) reduces $`|du|`$ for energetic electrons ($`u>u_{sh}`$) and increases $`|du|`$ for slow electrons ($`u<u_{sh}`$). *Therefore, the electric field in the plasma cools the energetic electrons and heats the low-energy electrons, respectively.* Similar observations were made in numerical simulations Surendra PRL .
Figure 7 shows the dimensionless power density as a function of $`n_b/n_{sh}`$. Taking into account the electric field in the plasma (both $`E_b`$ and $`E_1`$) reduces the total power deposited in the sheath region. Interestingly, taking into account only the uniform electric field $`E_b`$ gives a result close to the case when both $`E_b`$ and $`E_1`$ are accounted for. The electric field $`E_1`$ redistributes the power deposition from the energetic electrons to the low energy electrons, but does not change the total power deposition (compare lines (a) and (b) in Fig.6 and Fig. 7). Therefore, the total power deposition due to sheath heating can be calculated approximately from Eq. (28), taking into account only the electric field $`E_b`$. This gives
$$P_{tot}mV_b^2_0^{\mathrm{}}u^2\left[\frac{u^{}}{u}\frac{n_b}{n_{sh}}\mathrm{\Theta }(uu_{sh})1\right]^2\frac{df_0}{du}𝑑u.$$
(31)
The result of the self-consistent calculation of the power dissipation in Eq. (31) differs from the non-self-consistent estimate by the last term in Eq. (31), which contributes corrections of order $`n_{sh}/n_b`$ to the main term.
A future development should provide a self-consistent analysis of a more realistic, nonuniform, and self-consistent ion density profile $`n_i(x)`$. Such study has been currently performed for inductively coupled discharges only.
## III Penetration of the rf electric field into an inductively-coupled plasma
Low pressure inductively-coupled rf discharges are often operated in the non-propagating regime, when the driving rf field penetrates into plasma only within a skin layer of width $`\delta `$ near the antenna, i.e., exhibits a skin effect. Not only the rf field, but, in this case, also the resulting induced electric current is concentrated near the surface of the plasma. Depending on the local, or non-local nature of the relation between the electric current $`j`$ induced in plasma and the rf electric field $`E`$, the skin effect is called normal, if the dependence of the current on the electric field is local, or anomalous, if the dependence of the current on the electric field is nonlocal Kolobov review .
To differentiate between the two regimes of the skin effect, it is convenient to introduce the nonlocality parameter Weibel $`\mathrm{\Lambda }=(\lambda /\delta _0)^2`$, where $`\lambda v_T/(\omega ^2+\nu ^2)^{1/2}`$ is the effective electron mean free path and
$$\delta _0=\frac{c}{\omega _p(1+\nu ^2/\omega ^2)^{1/4}}$$
(32)
is the depth of the normal skin effect. The parameter $`\mathrm{\Lambda }`$
$$\mathrm{\Lambda }=\frac{v_T^2\omega _p^2}{\omega ^2c^2(1+\nu ^2/\omega ^2)^{1/2}}$$
(33)
is a fundamental measure of plasma current non-locality. In the local limit $`\mathrm{\Lambda }1`$, the effective mean free path is small compared with the skin depth $`\lambda \delta `$, and the current density at a particular point in space can be considered as a function of the electric field at the same point $`𝐣(𝐱)=\sigma (𝐱)𝐄(𝐱)`$ (Ohm’s law). In the opposite limit $`\lambda \delta `$, the mean free path exceeds the skin depth $`\lambda \delta `$, the relation between the current and the field $`𝐣(x)=\underset{¯}{\sigma }(𝐱,𝐱^{})𝐄(𝐱^{})𝑑𝐱^{}`$ is no longer local, because the conductivity $`\sigma (𝐱,𝐱^{})`$ has a spatial dispersion.
The penetration of the rf electric field into the plasma is described according to Faraday’s and Ampere’s laws
$$\times 𝐄=\frac{1}{c}\frac{𝐁}{t},$$
(34)
$$\times 𝐁=\frac{1}{c}\frac{𝐃}{t}+\frac{4\pi }{c}𝐣.$$
(35)
For a transverse harmonic wave in one-dimensional geometry $`E_y(x)e^{i\omega t}`$, the Faraday’s and Ampere’s laws give
$$\left(\frac{^2}{x^2}+\frac{\omega ^2}{c^2}\right)E_y=\frac{4\pi i\omega }{c^2}j_y,$$
(36)
where the current $`j`$ is the plasma electron current $`j_y=j_{ey}`$ (the ions are considered stationary), which has to be calculated making use of the electron kinetic equation, similar to the case of the penetration of the longitudinal wave into the plasma described in the previous section.
### III.1 Normal skin effect
In the limit of the normal skin effect ($`\mathrm{\Lambda }1`$), the electron thermal motion can be neglected. The electron flow velocity $`V_{ey}`$ may be obtained from Newton’s law taking into account the drag force due to the electron neutral collisions,
$$m\frac{}{t}V_{ey}=eE_y\nu V_{ey}.$$
(37)
This gives for the electron current ($`j_{ey}=en_eV_{ey})`$ the Ohm’s law relationship
$$𝐣_e(x)=\sigma _e𝐄(x),$$
(38)
where
$$\sigma _e=\frac{e^2n_e}{m(\nu i\omega )}.$$
(39)
The plasma current density is proportional to the electric field at the same point of space with a proportionality coefficient that is the complex conductivity of the cold plasma. Substituting Ohm’s law Eq. (38) with plasma conductivity from Eq. (39) into Eq. (36) gives the solution of the wave equation
$$E_y=E_{y0}e^{\alpha x},$$
(40)
where $`\alpha =\sqrt{4\pi i\omega \sigma _e/c^2}`$. Here, we neglected small terms associated with the displacement current in the limit $`\omega \omega _p`$, which is valid for the most plasma parameters in ICP discharges. The electric field can be equivalently expressed as
$$E_y(x,t)=E_{y0}e^{cos(ϵ/2)x/\delta _0}\mathrm{cos}[\omega t\mathrm{sin}(ϵ/2)x/\delta _0],$$
(41)
where $`\delta _0`$ is the normal skin depth in Eq.(32 ), and $`ϵ=\mathrm{arctan}(\nu /\omega )`$.
### III.2 Anomalous skin effect
The case of anomalous skin effect ($`\mathrm{\Lambda }1`$) for low-pressure inductively-coupled plasmas is more complicated comparing to the case of normal skin effect, and requires a more elaborate mathematical and numerical treatment to uncover its intrinsic complexity. In the limit $`\mathrm{\Lambda }1`$, the electron mean free path is large compared with the skin depth, and the electron current is determined not by the local rf electric field (Ohm’s law), but rather is a function of the whole profile of the rf electric field over distances of order $`\lambda `$. Therefore, a rather complicated nonlocal conductivity operator has to be determined for the calculation of the rf electric field penetration into the plasma.
In the case of a uniform plasma, the Vlasov and Maxwell equations can be solved by applying a Fourier transform Pippard . For a transverse harmonic wave in one-dimensional geometry $`E_y(x)e^{i\omega t},`$ a spatial Fourier harmonic of the current $`j_{yk}\mathrm{exp}(ikx)`$ simplifies to become Kolobov review ; me and Oleg
$$j_{yk}=\frac{e^2n}{imkV_T}Z\left(\frac{\omega }{kV_T}\right)E_{yk}.$$
(42)
Details of the solution are given in Appendix C. The electric field profile is given by the inverse Fourier transform of Eq.(36)
$$E_y(x)=\frac{2i\omega }{c^2}I_{\mathrm{}}^{\mathrm{}}\frac{e^{ikx}}{k^2\omega ^2\epsilon _t(\omega ,k)/c^2}𝑑k.$$
(43)
Here, $`I`$ is the surface current in the antenna and $`\epsilon _t(\omega ,k)`$ is transverse plasma permittivity, which for a Maxwellian EEDF is given by Lifshitz and Pitaevskii
$$\epsilon _t(\omega ,k)1+\frac{\omega _p^2}{\omega ^2}\frac{\omega }{v_T|k|}Z\left(\frac{\omega }{v_T|k|}\right).$$
(44)
Note the module sign as an argument of the plasma dispersion function. It reflects the proper symmetry of the continued electric field profile into semi-plane $`x<0`$ and also the proper pole position of the plasma dispersion function Landau ; Aliev and me . Neglecting the module sign results in erroneous results.
The solution for the electric field Eq.(43) has been described in many reviews and textbooks Lifshitz and Pitaevskii ; Aliev and me ; Lieberman& Godyak review ; Kolobov review . Here, we only focus on a property of the solution (43) not commonly acknowledged in the literature.
In the limit $`\mathrm{\Lambda }1`$ or $`\delta v_T/\omega `$, the plasma dielectric function can be substituted by its limiting value at small arguments $`Zi\sqrt{\pi }`$. Introducing the anomalous skin depth
$$\delta _a\frac{c}{\omega _p}\left(\frac{\omega _pv_T}{\omega c\sqrt{\pi }}\right)^{1/3},$$
(45)
and substituting $`Zi\sqrt{\pi }`$ into Eq. (44) and into Eq. (43) gives
$$E_y(x)=\frac{2i\omega }{c^2}I_{\mathrm{}}^{\mathrm{}}\frac{e^{ikx}}{k^2i/|k|\delta _a^3}𝑑k.$$
(46)
The integral in Eq. (46) cannot be calculated analytically, but it can be transformed into an integral in the complex $`k`$ plane by substituting $`|k|=\sqrt{k^2}`$. The contour of the integration should encompass branch point of the function $`\sqrt{k^2}`$ and has to come around the imaginary $`k`$axis. This gives Aliev and me
$`E_y(x)`$ $`=`$ $`E_0{\displaystyle \frac{(i\sqrt{3}+1)}{3\gamma _1}}\mathrm{exp}\left({\displaystyle \frac{x\gamma _2}{\delta _a}}\right)+{\displaystyle \frac{E_0}{3\gamma _1}}\mathrm{exp}\left({\displaystyle \frac{x}{\delta _a}}\right)`$ (48)
$`{\displaystyle \frac{2iE_0}{\pi \gamma _1}}P{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{\xi \mathrm{exp}\left(x\xi /\delta _a\right)}{1\xi ^6}}𝑑\xi .`$
where $`\gamma _1=2(\sqrt{3}+i)/3\sqrt{3},`$ $`\gamma _2=(1i\sqrt{3})/2`$ and $`E_0`$ is the electric field at the plasma boundary at $`x=0`$, $`P`$ stands for principal value of the integral. The last term represents the contribution of the integral around the imaginary $`k`$axis and the exponential terms originate from the poles. The electric field at $`x=0`$ can be calculated analytically
$$E_0=\frac{4i\omega I}{c^2}\frac{\pi (\sqrt{3}+i)\delta _a}{3^{3/2}}.$$
(49)
From Maxwell’s equations it follows that the magnetic field near the coil is $`B|_{0+}=2\pi I/c`$. Correspondingly, the derivative of the electric field at the plasma boundary is
$$\frac{dE_y}{dx}|_{x=0}=\frac{2\pi i\omega }{c^2}I.$$
(50)
The characteristic decay length of the electric field can be introduced as Kondratenko ; Rukhadze
$$\delta _s=\frac{E_0}{dE_y/dx}=\frac{2}{3}\left(1+i/\sqrt{3}\right)\delta _a.$$
(51)
The electric field profile from Eq. (48) is compared in Fig. 8 with the exponential profile
$$E_y(x)=E_0\mathrm{exp}\left[xRe(1/\delta _s)\right].$$
(52)
A more conventional plot of the amplitude and phase of the electric fields is shown in Fig. 9.
### III.3 Spatially averaged electric field, $`_0^{\mathrm{}}E_y𝑑x0`$ in the limit of a strong anomalous skin effect $`\mathrm{\Lambda }\mathrm{}`$.
The most apparent difference between the anomalous skin effect and the normal skin effect is that the amplitude of the rf filed is non-monotonic in the limit of anomalous skin effect and monotonic (exponential) for the normal skin effect. Moreover, in the case of the extremely anomalous skin effect, in the limit $`\mathrm{\Lambda }1`$, the spatially averaged rf electric field tends to zero Aliev and me
$$_0^{\mathrm{}}E_y𝑑x0,\mathrm{\Lambda }\mathrm{}.$$
(53)
In other words, the phase of the electric field changes by $`\pi `$ inside the skin layer, see Fig. 9(b). The spatially averaged electric field is given by the Fourier component at $`k=0`$, i.e.,
$$_0^{\mathrm{}}E_y𝑑x=\pi E(k=0).$$
(54)
Substituting the Fourier component of the electric field from Eq.(43) into Eq.(54) gives
$$_0^{\mathrm{}}E_y(x)𝑑x=\frac{2\pi i\omega I}{c^2}\frac{1}{\omega ^2\epsilon (\omega )/c^2},$$
(55)
and
$$\frac{_0^{\mathrm{}}E_y(x)𝑑x}{|E_0|\delta _a}=\frac{3^{3/2}\pi ^{1/3}}{\mathrm{\Lambda }^{1/3}}.$$
(56)
From Eq. (56) it is evident that as the nonlocality parameter tends to infinity, the averaged electric field tends to zero. This property of the electric field profile is consistent with nonlocality of the electron current. The electric field profile and the current profile are coupled to each other by Eq. (36). Therefore, the main part of the current and the electric field should decay on distances of order $`\delta _a`$, see Fig. 9. However, if the electric field profile has a non-zero average, the fast electrons will pick up a velocity kick from the skin layer and will transport the current over distances of order $`v_T/\omega \delta _a`$, where the electric field vanishes. This would contradict Maxwell’s equations. *Therefore, the zero average of the electric field is necessary and an important property of the electric field profile in the limit of the extreme anomalous skin effect* $`\mathrm{\Lambda }\mathrm{}`$.
The penetration length is defined in textbooks Kondratenko ; Rukhadze as
$$\lambda _E=\frac{_0^{\mathrm{}}E_y(x)𝑑x}{E_0}.$$
(57)
From the above discussion it follows that this definition is confusing, because in the limit of the anomalous skin effect the above defined penetration length is $`\lambda _E\delta _a`$ and is not a good measure of penetration length of the electric field. A better definition would be
$$\lambda _{|E|}=\frac{_0^{\mathrm{}}|E_y(x)|𝑑x}{|E_0|}.$$
(58)
In the limit of the strong anomalous skin effect, i.e. $`\mathrm{\Lambda }1`$, numerical calculation gives
$$\lambda _{|E|}=1.64\delta _a.$$
(59)
From Fig. 9 it is evident that in the region $`x2\delta _a`$ the amplitude of the electric field can be approximated by the exponential profile in Eq.(52) with the decay length
$$\delta _e=\frac{1}{Re(1/\delta _s)}=\frac{8}{9}\delta _a.$$
(60)
Note that the penetration length defined by Eq. (58), $`\lambda _{|E|}`$ is nearly twice as large as the initial decay length of the electric field amplitude near the plasma-wall boundary $`\delta _e`$. This is due to the pronounced long tail in the profile of the electric field.
Similarly, if we introduce the penetration length of the current
$$\lambda _{|j|}=\frac{_0^{\mathrm{}}|j_y(x)|𝑑x}{|j_0|},$$
(61)
numerical simulation gives
$$\lambda _{|j|}=1.87\delta _a\lambda _{|E|}.$$
(62)
This result contradicts to claim of Refs.Kondratenko ; Godyak review , that the magnetic field and current penetration lengths are much longer than the electric field penetration length. This claim is the result of an inaccurate definition of the penetration length.
In an attempt to reduce the phenomenon of the anomalous skin effect to the normal skin effect, many authors have substituted the correct profile of the electric field in Eq. (48) by an exponential profile $`E_0\mathrm{exp}(x/\delta _e)`$ with some fitting procedure for $`\delta _e`$ Vahedi ; Haas ; Tyshetskiy . By doing so, the property of the electric field in the limit of anomalous skin effect in Eq.(53) is violated. This leads to overestimation of the electron heating Aliev and me . Under the conditions of the anomalous skin effect $`v_T\delta _a\omega `$, electrons acquire a velocity kick
$$\mathrm{\Delta }v_y=\frac{2e}{mv_x}_0^{\mathrm{}}E_y(x)𝑑x.$$
(63)
If $`E_y(x)`$ satisfies the condition in Eq. (53), the electron velocity kick after passing through the skin layer is much smaller than in the case of an exponential electric field profile, which does not satisfy the property $`_0^{\mathrm{}}E_y𝑑x0`$, as $`\mathrm{\Lambda }\mathrm{}`$.
### III.4 Analytical separation of the electric field profile into an exponential part and a far tail.
Consider an exponential profile of the electric field in a plasma
$$E_y(x,t)=E_{y0}\mathrm{exp}\left(k_pxi\omega t\right),$$
(64)
where $`k_p`$ is a real positive number. The velocity perturbation in this electric field becomes
$$\mathrm{\Delta }v_y(x,t)=\frac{e}{m}_{\mathrm{}}^t𝑑\tau E_y[x(\tau ),\tau ].$$
(65)
The velocity kick $`\mathrm{\Delta }v_y`$ can be separated into a purely exponential part and a non-exponential part. Substituting the electron trajectory $`x(\tau )=xv_x(t\tau )`$ for $`v_x<0`$ gives
$$\mathrm{\Delta }v_y(x,t)=\frac{e}{m}\frac{Ey_0}{k_pv_xi\omega }\mathrm{exp}\left(k_pxi\omega t\right).$$
(66)
For $`v_x>0`$, the velocity acquired by an electron can be represented as the difference between the velocity kick acquired after a full pass through the skin layer and the contribution from the part of the skin layer $`[x;\mathrm{}]`$, i.e.,
$$\mathrm{\Delta }v_y(x,t)=\frac{e}{m}\left[_{\mathrm{}}^{\mathrm{}}_t^{\mathrm{}}\right]d\tau E_y[x(\tau ),\tau ].$$
(67)
The second part of the integral ($`\mathrm{\Delta }v_y^e`$) in Eq. (67) gives an exponential profile for the velocity kick, similar to Eq. (66)
$$\mathrm{\Delta }v_y^e(x,t)=\frac{e}{m}\frac{E_{y0}}{k_pv_xi\omega }\mathrm{exp}\left(k_pxi\omega t\right),v_x>0.$$
(68)
The first part of the integral ($`\mathrm{\Delta }v_y^{in}`$) in Eq. (67) gives
$$\mathrm{\Delta }v_y^{in}=\mathrm{\Delta }v_y^{\mathrm{}}e^{i\omega (tx/v_x)},$$
(69)
$$\mathrm{\Delta }v_y^{\mathrm{}}=\frac{e}{m}E_{y0}\left(\frac{1}{i\omega +k_pv_x}\frac{1}{i\omega k_pv_x}\right).$$
(70)
Here, $`\mathrm{\Delta }v_y^{\mathrm{}}`$ is the velocity kick acquired during the pass through the entire skin layer. The time $`tx/v_x`$ corresponds to the moment the electron collides with the wall.
Substituting $`\mathrm{\Delta }v_y^e(x,t)`$ from Eqs. (66) and (68 ) gives for the exponential part of the current
$$j_y^e=e\mathrm{\Delta }v_y\frac{f}{v_y}v_y𝑑𝐯,$$
(71)
$$j_y^e=\frac{e^2}{m}E_{y0}e^{i\omega tk_px}\frac{1}{k_pv_xi\omega }\frac{f}{v_y}v_y𝑑𝐯,$$
(72)
or, after integration, the exponential profile of the current becomes
$$j_y^e=\frac{e^2}{m}E_{y0}e^{i\omega tk_px}\frac{n}{k_pV_T}Z\left(\frac{i\omega }{k_pV_T}\right)^{}.$$
(73)
The asterisk denotes the complex conjugate. Note that Eq. (73) can be derived from Eq. (42) with the substitution $`k=ik_p`$ and by accounting for the following property of the dispersion function Plasma Formulary
$$Z(\xi ^{})=Z(\xi )^{}.$$
(74)
The exponential part of the profile should satisfy Maxwell’s equation (36). This gives an expression for $`k_p`$
$$k_p^2=\frac{\omega ^2}{c^2}+\frac{\omega _p^2}{c^2}\frac{i\omega }{k_pV_T}Z\left(\frac{i\omega }{k_pV_T}\right)^{}.$$
(75)
Note that because $`Z`$ in Eq. (75) has only purely imaginary and positive parts, $`k_p`$ is a real positive number, as it was assumed to be.
The non-exponential part of the electron velocity kick in Eq.(69) generates a non-exponential part of the current profile, which decays over a spatial scale of order $`V_T/\omega `$ due to the phase mixing, as the phase of the velocity kick $`\omega (tx/v_x)`$ in Eq. (69) is different for electrons with different $`v_x.`$ The current and electric field profiles are essentially non-exponential, similar to Eq.(4) for longitudinal velocity kicks, as discussed above.
Details of the exact analytical calculation of the electric field profile separation is given in Appendix C. Applying a procedure similar to that of Landau’s treatmentLandau for the longitudinal electric field, the integral in $`k`$space in Eq. (43) can be separated into an integral over an analytic function in the region $`k[\mathrm{},\mathrm{}]`$ and an integral over some non-analytic function in the region $`k[0,\mathrm{}]`$. To do so, the plasma permittivity has to be analytically continued from the real axis $`k<0`$, $`Imk=0`$, into the complex $`k`$ -plane, see Appendix C for details. The first integral can be readily calculated using the theory of residues. In the upper half-plane of the complex $`k`$, there exists only one pole of the analytically continued function of the plasma permittivity continued from $`k<0`$. The value of the pole is equal to $`ik_p`$, given by Eq. (75).
In the limit $`\omega k_pV_T`$, $`Z(\zeta )=1/\zeta `$, where $`\zeta =i\omega /k_pV_T`$. Substituting this value for the plasma dielectric function into Eq. (75) yields $`k_p=\omega _p/c`$, i.e., the normal skin layer length $`1/k_p=\delta _0`$ in Eq. (32) for $`\nu \omega `$ and $`\omega \omega _p.`$
Figure 10 shows the profile of the electric field for the same typical ICP parameters: plasma density $`n=10^{11}`$cm$`^3,`$ electron temperature $`T_e=3`$ eV, and discharge frequency $`f=13.56`$ MHz. Shown are the exact electric field profile $`E_y(x)`$ calculated according to Eq.(43), the exponential part of the electric field
$$E_{yp}(x)=E_0\mathrm{exp}(k_px)$$
(76)
with $`k_p`$ from Eq.(75), and the difference of the two
$$E_{yt}(x)=E_y(x)E_{yp}(x),$$
(77)
and the asymptotic calculation for $`E_{yt}(x)`$ in Eq.(151 ) $`E_{ystd}(x)`$. For these plasma parameters the skin effect is neither normal nor anomalous: $`\omega /k_pV_T=1.52`$. Notwithstanding the fact that the parameter $`\omega /k_pV_T`$ is of order unity, the main part of the electric field is close to the exponential profile in Eq. (76) with $`k_p`$ from Eq. (75), $`E_y(x)E_{yp}(x)`$. As evident from Fig. 10, the non-exponential part is small,$`E_{yt}(x)E_{yp}(x)`$, everywhere where the electric field is substantial, or up to distances five times of skin depth, for $`x<5/k_p=7.5V_T/\omega `$. The tail of the electric field profile for $`x>7V_T/\omega `$ is non-exponential and dominated by $`E_{yt}(x)`$.
In the limit of the anomalous skin effect $`\omega /k_pV_T1`$, $`Z(\zeta )=i\sqrt{\pi }`$, where $`\zeta =i\omega /k_pV_T`$. Substituting this value for the plasma dielectric function into Eq. (75) yields $`k_p=1/\delta _a,`$ which is very close to the skin impedance approximation in Eq. (52) which corresponds to $`k_p=9/8\delta _a`$ –a 12 % difference. As a result, the exponential profile in Eq. (76) approximates well the exact profile of the electric field over distances within a few skin depths even in the limit of the strong anomalous skin effect, as is evident in Fig. 11. However, the non-exponential part $`E_{yt}(x)`$ dominates $`E_{yp}(x)`$ at $`x>V_T/\omega `$ in accord with the requirement in Eq.(53).
### III.5 Surface impedance
An important plasma characteristic is the surface impedance, which is given by the ratio of the electric field to the rf magnetic field or the coil current at the plasma boundary Lifshitz and Pitaevskii
$$Z=\frac{E}{B}|_{x=0},$$
(78)
where
$$B|_{x=0}=\frac{2\pi }{c}I$$
(79)
is the magnetic field near the antenna. The total power $`P`$ deposited per unit area into the plasma is determined by the energy flux dissipated into the plasma or the time-averaged Poynting vector
$$P=<S_x>=\frac{1}{2}\frac{c}{4\pi }Re(EB^{}).$$
(80)
Substituting the electric field from Eq.(78) and the magnetic field Eq. (79) into Eq. (80) relates the power to the real part of the surface impedance
$$P=\frac{\pi }{2c}I^2ReZ.$$
(81)
The imaginary part of the surface impedance describes the plasma inductance.
The surface impedance can also be used to estimate the penetration length in the surface impedance approximation given by Eq. (51). Substituting the electric field from Eq. (78) and the magnetic field Eq. (79) into Eq.(51) relates the penetration depth and the surface impedance
$$\delta _s=\frac{cZ}{i\omega }.$$
(82)
The surface impedance can be calculated making use of Eq. (43) Lifshitz and Pitaevskii , i.e.,
$$Z=\frac{i\omega }{\pi c}_{\mathrm{}}^{\mathrm{}}\frac{1}{k^2\omega ^2\epsilon _t(\omega ,k)/c^2}𝑑k,$$
(83)
which requires numerical integration. On the other hand, we can use the results of the previous subsection that the main part of the electric field is an exponential function in Eq. (76) with $`k_p`$ given by Eq.(75). From Eq. (82), the imaginary part of the surface impedance can be obtained substituting $`\delta _s=1/k_p`$
$$Z_p=\frac{i\omega }{ck_p}.$$
(84)
A pure exponential profile yields only the imaginary part of the surface impedance. The real part of the impedance can be calculated by computing the power dissipated by electrons from the skin layer Vahedi
$$P=\frac{m}{4}v_xf|\mathrm{\Delta }v_y^{\mathrm{}}|^2\mathrm{𝐝𝐯},$$
(85)
where $`\mathrm{\Delta }v_y^{\mathrm{}}`$ is the velocity kick acquired by an electron after passing through the skin layer, which is given by Eq. (70). Here, $`m(\mathrm{\Delta }v_y^{\mathrm{}})^2/4`$ is the temporal average of the electron energy change in the skin layer and $`v_xf`$ is the electron flux on the wall. Equation (81) becomes
$$ReZ_p=\frac{2}{c}\omega _p^2|Z|^2fv_x\left(\frac{k_pv_x}{\omega ^2+(k_pv_x)^2}\right)^2𝑑v_x.$$
(86)
Because the imaginary part of impedance is large compared with its real part, only the imaginary part can be included on the right hand side in Eq. (86).
Figure 12 shows the real and imaginary parts of the surface impedance versus the discharge frequency calculated exactly, i.e., making use of Eq.(83), and approximately from Eqs. (84) and (86). Also shown at the top of this figure, is the ratio of the actual skin depth $`\delta =1/k_p`$ from Eq. (75) to the normal skin depth calculated in the cold plasma approximation $`\delta _0`$ given by Eq.(32). From Fig. 12 it is evident that within 50% accuracy, the impedance calculation can be based on the exponential profile in Eq. (76) for discharge frequencies higher than 1 MHz Haas . However, for lower frequencies, the assumption of purely exponential profile leads to overestimation of the electron heating and plasma resistivity up to a factor of 3 for $`f1`$ MHz, see Fig. 12. This is because the important property of the electric field profile under the conditions of strong anomalous skin effect in Eq. (53) is being violated. Note that at these low frequencies taking into account a small but finite collision frequency or nonlinear effects may be important.
### III.6 Anomalous skin effect for an anisotropic electron velocity distribution
The anomalous skin effect in a plasma with a highly anisotropic electron velocity distribution function (EVDF) is very different from the skin effect in a plasma with the isotropic EVDF. In Ref. anisotropic EEDF an analytical solution was obtained for the electric field penetrating into plasma with the EVDF described by a Maxwellian with two temperatures $`T_yT_x`$, where $`y`$ is the direction along the plasma boundary and $`x`$ is the direction perpendicular to the plasma boundary. Under the conditions
$$\frac{v_{Ty}}{\omega }\frac{c}{\omega _p};\omega _p\omega ,$$
(87)
the skin layer was found to consist of two distinct regions of width of order $`v_{Tx}/\omega `$ and $`v_{Ty}/\omega `$, where $`v_{Tx,y}=\sqrt{T_{x,y}/m}`$ are the thermal electron velocities in $`x`$ and $`y`$ directions, and $`\omega `$ is the incident wave frequency. The calculation is based on Eq.(43), where the dielectric permittivity has to be modified for an anisotropic EEDF to become
$$\epsilon _t(\omega ,k)=1\frac{\omega _p^2}{\omega ^2}\left\{1\frac{T_y}{T_x}\left[1+\frac{\omega }{\sqrt{2}v_{Tx}k}Z\left(\frac{\omega }{\sqrt{2}v_{Tx}k}\right)\right]\right\}.$$
(88)
In the case of anisotropic EEDF under conditions in Eq. (87), the integral in Eq.(43) has two poles and the integration over the branch point $`k=0`$ does not contribute. As a result, the profile of the electric field is a sum of the two complex exponents:
$$E(x)\frac{\omega }{\omega _p}B(0)\left[\frac{i\omega c}{\omega _pv_{Ty}}\mathrm{exp}(ik_{p1}x)+\sqrt{T_x/T_y}\mathrm{exp}(ik_{p_2}x)\right],$$
(89)
where $`k_{p1}`$ is given by
$$k_{p1}=i\frac{\omega }{v_{Ty}},$$
(90)
and $`k_{p2}`$ is given by
$$k_{p2}=\frac{\omega _p}{c}\sqrt{T_y/T_x}+i\frac{\sqrt{\pi }\omega }{2\sqrt{2}v_{Tx}}.$$
(91)
The profile of the electric field is shown in Fig. 13. The skin layer contains multiple oscillations of the electric field, in striking contrast to the case of isotropic EEDF.
## IV Conclusions
We showed that electrons can transport the plasma current away from the skin layer due to their thermal motion over distances of order $`v_T/\omega `$. As a result, the width of the skin layer increases when electron temperature effects are taken into account. The anomalous penetration of the rf electric field occurs not only for the wave transversely propagating to the plasma boundary (inductively coupled plasmas), but also for the wave propagating along the plasma boundary (capacitively coupled plasmas). It was shown that separating the electric field profile into exponential and nonexponential parts yields an efficient qualitative and quantitative description of the anomalous skin effect. Accounting for the non-exponential part of the profile is important for the calculation of the electron heating and the plasma resistivity. For example, the assumption of purely exponential profile leads to overestimation of up to a factor of 3 in the electron heating for $`f1`$ MHz, see Fig. 12.
Here, we considered only plasmas with a Maxwellian electron energy distribution function. However, in low pressure rf discharges, the EEDF is non-Maxwellian for plasma densities typically lower than $`10^{10}`$cm<sup>-3</sup> Godyak new exp . The nonlocal conductivity, and plasma density profiles and EEDF are all nonlinear and nonlocally coupled new Oleg . Hence, for accurate calculation of the discharge characteristics at low pressures, the EEDF needs to be computed self-consistently badri and me ; badri and me 2 ; Our article ; our1 . The effects of a nonMaxwellian EEDF, nonlinear phenomena, the effects of plasma non-uniformity and finite size,as well as influence of the external magnetic field on the anomalous skin effect will be reported in the second part of the review review part 2 .
Acknowledgments
This research was supported by the U.S. Department of Energy Office of Fusion Energy Sciences through a University Research Support Program and the University of Toledo. The authors gratefully acknowledge helpful discussions with R. Davidson, V. Godyak, Badri Ramamurthi, E. Startsev and L. D. Tsendin.
## Appendix A Analytical derivation of the current profile driven by velocity kicks near the plasma boundary
Consider that electrons acquire a velocity kick near the boundary, in the direction perpendicular to the boundary
$$dv_x=\mathrm{\Delta }V\mathrm{cos}(\omega t).$$
(92)
The electron velocity at a distance $`x`$ from the boundary will be determined by the exact moment of the collision with the boundary at a time $`tx/v_x`$. The electron current in the plasma is given by integration over all electrons with a distribution function $`f(v_x)`$
$$j(x,t)=e\mathrm{\Delta }V_0^{\mathrm{}}f(v_x)\mathrm{cos}(\omega t\omega x/v_x)𝑑v_x.$$
(93)
For a Maxwellian distribution function $`f(v_x)=n_0e^{v_x^2/v_T}/v_T\sqrt{\pi }`$ the current in Eq. (93) takes the form $`j(\xi ,t)=j_0A(\xi )\mathrm{cos}[\omega t\varphi (\xi )]`$ , where $`j_0=en_0\mathrm{\Delta }V`$ and $`A`$ and $`\varphi `$ are the amplitude and phase of the current, respectively, and $`\xi =\omega x/v_T`$. The functions $`A`$ and $`\varphi `$ are shown in Fig. 1. In the limit $`\xi 1`$, the integration in Eq. (93) can be performed analytically making use of the method of steepest descend Brillouin
$$j(\xi ,t)=\frac{j_0}{\sqrt{\pi }}Re\left(e^{i\omega t}_0^{\mathrm{}}e^{s^2+i\xi /s}𝑑s\right),$$
(94)
where and $`s=v_x/v_T`$. The integral in Eq.(94) can be calculated in the complex $`s`$ plane. The stationary phase point is given by $`d(s^2+i\xi /s)/ds=0`$ or $`s^3=i\xi /2`$. This gives the stationary point $`s_0=\left(i\xi /2\right)^{1/3}`$. In the neighborhood of this point, the function in the exponent can be expanded as a Taylor series, $`s^2+i\xi /s=s_0^2+i\xi /s_06(ss_0)^2/2=3s_0^23(ss_0)^2`$. Integration of the Gaussian gives $`_0^{\mathrm{}}e^{3(ss_0)^2}𝑑s=\sqrt{\pi /3}`$. Substituting this into the integration in Eq. (94) yields:
$$j(\xi ,t)=\frac{j_0}{\sqrt{3}}Re\left(\mathrm{exp}\left[i\omega t3\left(i\xi /2\right)^{2/3}\right]\right).$$
(95)
Substituting $`\left(i\right)^{2/3}=\left(e^{i\pi /2}\right)^{2/3}=e^{i\pi /3}=\mathrm{cos}(\pi /3)i\mathrm{sin}(\pi /3)=1/2\sqrt{3}i/2`$ into Eq. (95) gives
$$j(\xi ,t)=\frac{j_0}{\sqrt{3}}\mathrm{exp}\left(3\xi ^{2/3}/2^{5/3}\right)\mathrm{cos}\left(\omega t3\sqrt{3}\xi ^{2/3}/2^{5/3}\right).$$
(96)
## Appendix B Analytical derivation of the longitudinal rf electric field profile near the plasma boundary ($`𝐄𝐤`$)
The analytical solution for a longitudinal rf electric field involves solving the Vlasov equation for the electron velocity distribution function (EVDF) $`F`$
$$\frac{F}{t}+v_x\frac{F}{x}\frac{e}{m}E_x\frac{F}{v_x}=0,$$
(97)
together with the Poisson equation
$$\frac{dE}{dx}=4\pi e\left(n_i_{\mathrm{}}^{\mathrm{}}F𝑑v_x\right).$$
(98)
In the linear approximation, the EVDF can be split into two parts
$$F(t,x,v_x)=f_0(v_x)+f(t,x,v_x),$$
(99)
where $`f_0(v_x)`$ describes EVDF of a uniform plasma with uniform ion density $`n_e=n_i=n_0`$ and $`f(t,x,v_x)`$ is EVDF due a wave perturbation. Substituting Eq. (99) into Eqs. (97) and (98) yields the linearized Vlasov-Poisson system of equations
$$\frac{f}{t}+v_x\frac{f}{x}\frac{e}{m}E_x\frac{df_0}{dv_x}=\nu f,$$
(100)
$$\frac{dE_x}{dx}=4\pi e_{\mathrm{}}^{\mathrm{}}f(v_x)𝑑v_x.$$
(101)
In the first equation (100), the small collisional term with the collision frequency $`\nu \omega `$ is taken into account. In Ref. Landau Landau solved the linearized Vlasov-Poisson system making use of the Laplace transform for a semi-infinite plasma $`x>0`$. However, it is more convenient to apply a Fourier transform to an infinite plasma by artificially continuing the EVDF and the electric field in the semi-plane $`x<0`$ Aliev and me . Electrons moving with $`v_x<0`$ reflect from the boundary $`x=0`$ and change their velocity to $`v_x`$. This gives the boundary condition for the Vlasov equation in the semi-plane $`x>0`$
$$f(t,0,v_x)=f(t,0,v_x).$$
(102)
Instead of considering problem in the semi-plane $`x>0`$ with the boundary condition in Eq.(102), we can consider the entire plane $`x[\mathrm{},\mathrm{}]`$ by artificially continuing the electric field into the semi-plane $`x<0`$. The Vlasov equation is symmetric with respect to a change in variables according to the substitution
$$v_xv_x,xx,EE.$$
(103)
Therefore, electrons at $`x=0`$ with $`v_x>0`$, which are reflected from the wall can be represented as electrons which came from the semi-plane $`x<0`$ and interacted with the electric field
$$E_x(x<0)=E_x(x>0).$$
(104)
As a result, the electric field has to be continued anti-symmetrically into the semi-plane $`x<0`$.
Now we can apply the Fourier transform for the Vlasov-Poisson system of Eqs. (100) and (101). This gives for the components of the EVDF $`f_ke^{i\omega t+kx}`$ and the electric field $`E_ke^{i\omega t+kx}`$
$$i(\omega +i\nu v_xk)f_k\frac{e}{m}E_k\frac{df_0}{dv_x}=0,$$
(105)
$$ikE_k+2E_0=4\pi e_{\mathrm{}}^{\mathrm{}}f_k𝑑v_x.$$
(106)
Note that due the fact that the electric field is a discontinuous function, the Fourier transform of the derivative of the electric field $`dE/dx`$ is $`ikE+2E_0`$, where $`E_0=E(0)`$ is the electric field at the right side ($`x>0`$) of the plasma boundary. Substituting $`f_k`$ from Eq.(105) into (106) yields
$$E_k=\frac{2E_0}{ik}\frac{1}{\epsilon _{}(\omega ,k)},$$
(107)
where $`\epsilon _{}(\omega ,k)`$ is the longitudinal plasma permittivity
$$\epsilon _{}(\omega ,k)=1+\frac{\omega _p^2}{n_0k}_{\mathrm{}}^{\mathrm{}}\frac{1}{\omega +i\nu v_xk}\frac{df_0}{dv_x}𝑑v_x.$$
(108)
Substituting a Maxwellian EEDF
$$f_0=\frac{n_0}{\sqrt{\pi }v_T}\mathrm{exp}(v^2/v_T^2),$$
(109)
where $`v_T=\sqrt{2T/m}`$, into Eq.(108) and after some algebra Lifshitz and Pitaevskii , we obtain
$$\epsilon _{}(\omega ,k)1+\frac{2\omega _p^2}{k^2v_T^2}\left[1+\frac{1}{\sqrt{\pi }v_T}_{\mathrm{}}^{\mathrm{}}\frac{\mathrm{exp}(v^2/v_T^2)}{v_xk\omega i\nu }𝑑v_x\right].$$
(110)
The last term on the right hand side can be expressed in terms of the plasma dispersion function
$$Z(\zeta )=\frac{1}{\sqrt{\pi }}_{\mathrm{}}^{\mathrm{}}\frac{\mathrm{exp}(t^2)}{t\zeta }𝑑t,Im(\zeta )>0.$$
(111)
The dispersion function $`Z(\zeta )`$ in the form of Eq. (111) is only defined for $`Im(\zeta )>0`$ and is defined as an analytical continuation for $`Im(\zeta )<0`$. For $`k>0`$, in the limit $`\nu 0`$,
$$\frac{1}{\sqrt{\pi }v_T}_{\mathrm{}}^{\mathrm{}}\frac{\mathrm{exp}(v^2/v_T^2)}{v_xk\omega i\nu }𝑑v_x=\frac{1}{kV_T/\omega }Z(\omega /kV_T).$$
(112)
For $`k<0`$, the imaginary part of the $`(\omega +i\nu )/k`$ is negative and we have to transform the integral (110) so that the pole $`v_{xp}=(\omega +i\nu )/k`$ lies in the upper plane of the complex velocity. This can be achieved by substitution $`v_xv_x`$ , which gives for $`k<0`$
$$\frac{1}{\sqrt{\pi }v_T}_{\mathrm{}}^{\mathrm{}}\frac{\mathrm{exp}(v^2/v_T^2)}{v_x|k|\omega i\nu }𝑑v_x=\frac{1}{|k|v_T/\omega }Z(\omega /|k|v_T).$$
(113)
As a result,
$$\epsilon _{}(\omega ,k)1+\frac{2\omega _p^2}{k^2v_T^2}\left[1+\frac{1}{|kV_T/\omega |}Z(|\omega /kV_T|)\right].$$
(114)
Note that because the function $`f_0`$ is symmetric with respect to the substitution $`v_xv_x`$, $`\epsilon (\omega ,k)`$ is symmetric with respect to the substitution $`kk`$. Correspondingly the symmetry of the electric field in Eq.(104) is preserved.
The electric field profile is given by the inverse Fourier transform of Eq.( 107)
$$E_x(x)=\frac{1}{2\pi }_{\mathrm{}}^{\mathrm{}}\frac{2E_0}{ik}\frac{e^{ikx}}{\epsilon _{}(\omega ,k)}𝑑k.$$
(115)
In the limit $`x\mathrm{}`$, $`E(x)E_0/\epsilon `$, where $`\epsilon =\epsilon (\omega ,0)`$. This is in accord with the conservation of the total current in the one-dimensional geometry. The total current is the sum of the displacement current and the electron current,
$$\frac{1}{4\pi }\frac{E_x}{t}+j_e=I(t).$$
(116)
The total current conservation follows from the combination of the Poisson equation and the charge continuity equation. Indeed, taking the time derivative of the Poisson equation and making use of the charge continuity equation gives
$$\frac{}{t}E_x+4\pi j_e=0.$$
(117)
In one-dimensional geometry it can be integrated with a constant of space – the total current carrying through the plasma $`I(t)`$, which gives Eq.(116). For a harmonic electric field considered here, Eq.(116) gives
$$i\omega E_x4\pi i\omega \left(\frac{\epsilon 1}{4\pi }\right)E_x=i\omega E_0.$$
(118)
Here, we account for the relationship between the plasma conductivity ($`j_e=\sigma E`$) and the plasma dielectric function $`\epsilon =1+4\pi \sigma /(i\omega )`$. Eq.(118) gives
$$E_x(x\mathrm{})=E_0/\epsilon .$$
(119)
The same result can be obtained from Eq.(115) after substituting $`\epsilon (\omega ,k)\epsilon (\omega ,0)`$ and integrating. Thus, the electric field in the transition region is given by
$$E_x(x)E_0/\epsilon =\frac{1}{2\pi }_{\mathrm{}}^{\mathrm{}}\frac{2E_0}{ik}\left(\frac{1}{\epsilon _{}(\omega ,k)}\frac{1}{\epsilon _{}(\omega ,0)}\right)e^{ikx}𝑑k.$$
(120)
The dielectric function in the form given by Eq.(110) is not an analytic function of $`k`$. To apply the theory of residues, Landau proposed to split integral into two parts Landau according to
$`E_x(x)E_0/\epsilon `$ $`=`$ $`{\displaystyle \frac{1}{2\pi }}{\displaystyle _{\mathrm{}}^{\mathrm{}}}{\displaystyle \frac{2E_0}{ik}}\left({\displaystyle \frac{1}{\epsilon _1(\omega ,k)}}{\displaystyle \frac{1}{\epsilon (\omega ,0)}}\right)e^{ikx}𝑑k`$
$`+{\displaystyle \frac{1}{2\pi }}{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{2E_0}{ik}}\left({\displaystyle \frac{1}{\epsilon _1(\omega ,k)}}{\displaystyle \frac{1}{\epsilon _{}(\omega ,k)}}\right)e^{ikx}𝑑k,`$
where
$$\epsilon _1(\omega ,k)=1+\frac{2\omega _p^2}{\omega ^2k^2}\left[1\frac{1}{kv_T/\omega }Z(\omega /kv_T)\right].$$
(122)
The first integral can be calculated by moving the path of integration into the complex $`k`$plane and applying the theory of residues. For $`\omega \omega _p`$, $`\epsilon <0`$ and there is only one pole $`\epsilon _1(\omega ,k)=0`$ in the upper half-plane Landau . It corresponds to the usual screening with the Debye length. In the limit $`k\omega _p/v_T,`$ $`Z(|\omega /kv_T|)1`$ and $`Z(|\omega /kv_T|)/|kv_T/\omega |1`$, which gives
$$\epsilon _1(\omega ,k)1+\frac{2\omega _p^2}{k^2v_T^2}.$$
(123)
Calculation of the first term in Eq.(B) gives $`E_0\mathrm{exp}(x/a)`$, where $`a`$ is the Debye length $`a=v_T/\sqrt{2}\omega _p`$. Therefore,
$$E_x(x)=E_0/\epsilon +E_0\mathrm{exp}(x/a)+\frac{1}{2\pi }_0^{\mathrm{}}\frac{2E_0}{ik}\left(\frac{\epsilon _{}(\omega ,k)\epsilon _1(\omega ,k)}{\epsilon _1(\omega ,k)\epsilon _{}(\omega ,k)}\right)e^{ikx}𝑑k$$
(124)
For $`Im(k)=0`$, $`Z(\omega /kv_T)=Z(\omega /kv_T)^{}`$ Plasma Formulary and
$$\epsilon _1(\omega ,k)=1+\frac{2\omega _p^2}{v_T^2k^2}\left[1+\frac{1}{kv_T/\omega }Z(\omega /kv_T)^{}\right].$$
(125)
Substituting Eq.(125) into Eq.(124) gives for the last term $`E_t(x)`$
$$E_t(x)=\frac{4E_0}{\pi }\frac{\omega \omega _p^2}{v_T^3}_0^{\mathrm{}}\frac{1}{k^4}\frac{Im[Z(\omega /kv_T)]}{\epsilon _1(\omega ,k)\epsilon _{}(\omega ,k)}e^{ikx}𝑑k,$$
(126)
where Plasma Formulary
$$Im[Z(\zeta )]=\sqrt{\pi }\mathrm{exp}(\zeta ^2).$$
(127)
The last integral can be calculated analytically only in the limit $`xv_T/\omega `$ by applying the method of steepest descend. In this limit, $`kv_T/\omega ,`$ $`\epsilon _1(\omega ,k)\epsilon (\omega ,k)\epsilon `$ and
$$_0^{\mathrm{}}\frac{1}{k^4}\mathrm{exp}(ikx\omega ^2/k^2v_T^2)𝑑k\frac{\sqrt{2\pi }}{\sqrt{3}}(x\lambda _\omega )^{2/3}\lambda _\omega \mathrm{exp}\left[c\left(\frac{x}{\lambda _\omega }\right)^{2/3}i\pi /3\right],$$
(128)
where $`c=3(1+i\sqrt{3})/4`$, and $`\lambda _\omega =v_T/\sqrt{2}\omega `$ is the phase-mixing scale.
Substituting Eq.(127) into Eq.(126) and making use of Eq.(128) yields at $`x\lambda _\omega `$ Landau
$$E_t(x)\frac{2E_0}{\sqrt{3}\epsilon ^2}\frac{\omega _p^2}{\omega ^2}\left(\frac{x}{\lambda _\omega }\right)^{2/3}\mathrm{exp}\left[c\left(\frac{x}{\lambda _\omega }\right)^{2/3}i\pi /3\right].$$
(129)
The plots of amplitude and phase of the electric field profile $`E_t(x)`$ given by Eq.(126) and the approximate analytical result Eq. (129) are shown in Fig. 2.
## Appendix C Analytical derivation of the transverse rf electric field profile near the plasma boundary ($`𝐄𝐤`$)
The analytical solution involves solving the Vlasov equation for the electron velocity distribution function (EVDF) $`F`$
$$\frac{F}{t}+v_x\frac{F}{x}\frac{e}{m}(E_y+v_x\times B_z)\frac{F}{v_y}=0.$$
(130)
This equation has to be solved together with the Maxwell’s equation yielding
$$\left(\frac{d^2}{dx^2}+\frac{\omega ^2}{c^2}\right)E_y=\frac{4\pi i\omega }{c^2}\left[j+I\delta (x)\right],$$
(131)
where $`I`$ is the surface current. The plasma density is not perturbed in the transverse wave; therefore there is no need to solve the Poisson equation. In the linear approximation, the EVDF can be split into two parts
$$F(t,x,𝐯)=f_0(v)+f(t,x,𝐯),$$
(132)
where $`f_0(v)`$ describes EVDF of an isotropic, uniform plasma with uniform ion density $`n_e=n_i=n_0`$ and $`f(t,x,𝐯)`$ is the EVDF due a wave perturbation. Substituting Eq. (132) into Eqs. (130) yields the linearized Vlasov equation
$$\frac{f}{t}+v_x\frac{f}{x}\frac{e}{m}E_y\frac{f_0}{v_y}=\nu f.$$
(133)
In Eq. (133), the small collisional term with collision frequency $`\nu \omega `$ is taken into account. Similarly to the case of the longitudinal wave, we can consider the entire plane $`x[\mathrm{},\mathrm{}]`$ by artificially continuing the electric field in the semi-plane $`x<0`$. The Vlasov equation is symmetric relative to the change in variables according to the substitution
$$v_xv_x,xx,E_yE_y.$$
(134)
Therefore, electrons at $`x=0`$ with $`v_x>0`$ which are reflected from the wall can be represented as electrons which came from the semi-plane $`x<0`$ and interacted with the electric field
$$E_y(x<0)=E_y(x>0).$$
(135)
As a result, the electric field has to be continued symmetrically into the semi-plane $`x<0`$.
Now we can apply the Fourier transform for Eqs. (133) and (131). This gives for components of the EVDF $`f_ke^{i\omega t+ikx}`$ and the electric field $`E_{yk}e^{i\omega t+ikx}`$
$$i(\omega +i\nu v_xk)f_k\frac{e}{m}E_{yk}\frac{f_0}{v_y}=0,$$
(136)
$$\left(k^2+\frac{\omega ^2}{c^2}\right)E_{yk}=\frac{4\pi i\omega }{c^2}(j_k+I).$$
(137)
Substituting $`f_k`$ from Eq. (136) into (137) with the current $`j_k=ef_kv_y𝑑𝐯`$ yields
$$E_{yk}=\frac{4\pi i\omega }{c^2}I\frac{1}{k^2\frac{\omega ^2}{c^2}\epsilon _t(\omega ,k)},$$
(138)
where $`\epsilon _t(\omega ,k)`$ is the transverse plasma permittivity
$$\epsilon _t(\omega ,k)=1+\frac{\omega _p^2}{n_0\omega }_{\mathrm{}}^{\mathrm{}}\frac{v_y}{\omega +i\nu v_xk}\frac{f_0}{v_y}𝑑v_x.$$
(139)
Substituting a Maxwellian EEDF gives Lifshitz and Pitaevskii
$$\epsilon _t(\omega ,k)=1+\frac{\omega _p^2}{\omega ^2}\frac{\omega }{v_T|k|}Z\left(\frac{\omega }{v_T|k|}\right).$$
(140)
Note that because the function $`f_0`$ is symmetric relative to the substitution $`v_xv_x`$, $`\epsilon (\omega ,k)`$ is symmetric relative to the substitution $`kk`$. Correspondingly, the symmetry of the electric field in Eq. (135) is preserved.
The electric field profile is given by the inverse Fourier transform of Eq.(138)
$$E_y(x)=\frac{1}{2\pi }_{\mathrm{}}^{\mathrm{}}\frac{4\pi i\omega }{c^2}I\frac{e^{ikx}}{k^2\omega ^2\epsilon _t(\omega ,k)/c^2}𝑑k.$$
(141)
Similar to the analysis of the longitudinal wave, we split the integral in Eq. (138) into two parts
$$E_y(x)=E_{yp}(x)+E_{yt}(x),$$
(142)
where
$$E_{yp}(x)=\frac{2\pi i\omega }{c^2}I_{\mathrm{}}^{\mathrm{}}\frac{e^{ikx}}{k^2\omega ^2\epsilon _{t1}(\omega ,k)/c^2}𝑑k,$$
(143)
and
$$E_{yt}(x)=\frac{2i\omega }{c^2}I\frac{\omega ^2}{c^2}_0^{\mathrm{}}\frac{\left[\epsilon _t(\omega ,k)\epsilon _{t1}(\omega ,k)\right]e^{ikx}}{\left[k^2\omega ^2\epsilon _{t1}(\omega ,k)/c^2\right]\left[k^2\omega ^2\epsilon _t(\omega ,k)/c^2\right]}𝑑k,$$
(144)
$$\epsilon _{t1}(\omega ,k)=1\frac{\omega _p^2}{\omega ^2}\frac{\omega }{v_Tk}Z\left(\frac{\omega }{v_Tk}\right).$$
(145)
Note that $`\epsilon _{t1}(\omega ,k)=\epsilon _t(\omega ,k)`$ for $`k<0`$.
The first part $`E_{yp}(x)`$ of the electric field can be calculated by evaluating the integral in the complex $`k`$plane. A pole of $`E_{yp}(x)`$\- $`k_pi`$ lies on the imaginary axis of the $`k`$plane. The dielectric permittivity is real and negative on imaginary axis of the $`k`$plane
$$\epsilon _{t1}(\omega ,k_pi)=1\frac{\omega _p^2}{\omega ^2}\frac{\omega }{v_Tk_p}F\left(\frac{\omega }{v_Tk_p}\right),$$
(146)
where $`F(\zeta )=ImZ\left(i\zeta \right)=\sqrt{\pi }\mathrm{exp}(y^2)erfc(y)`$ Plasma Formulary . There is always a real value of $`k_p`$ as the root of
$$k_p^2=\omega ^2\epsilon _{t1}(\omega ,ik_p)/c^2.$$
(147)
Applying the theory of residues, the integral for $`E_{yp}(x)`$ gives
$$E_{yp}(x)=E_{op}e^{k_px},$$
(148)
where
$$E_{op}=\frac{2i\omega }{c^2}\frac{2\pi iI}{2k_pid\epsilon _{t1}(\omega ,k)/dk\omega ^2/c^2}.$$
(149)
In the limit $`x\delta `$, the last term $`E_{yt}(x)`$ can be calculated making use of the method of steepest descend. Substituting $`k^2\omega ^2\epsilon _{t1}(\omega ,k)/c^2`$ in the denominator of the expression for $`E_{yt}(x)`$ by its limit $`=\omega _p^2/c^2`$ at $`k0`$, gives
$$E_{yt}(x)=\frac{4\omega ^2\sqrt{\pi }}{\omega _p^2v_T}I\left[_0^{\mathrm{}}\frac{1}{k}\mathrm{exp}\left[\left(\frac{\omega }{v_Tk}\right)^2+ikx\right]𝑑k\right],$$
(150)
which yields
$$E_{ystd}(x)=\frac{4\omega ^2\pi }{\omega _p^2v_T}I\frac{\sqrt{2}}{\sqrt{3}}\left(\frac{x}{\lambda _\omega }\right)^{1/3}\mathrm{exp}\left[c\left(\frac{x}{\lambda _\omega }\right)^{2/3}i\pi /2\right],$$
(151)
where $`c=3(1+i\sqrt{3})/4`$ and $`\lambda _\omega =v_T/\sqrt{2}\omega `$.
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# Bogoliubov Excitations in a Kronig-Penney Potential
## I Introduction
In studies of Bose-Einstein condensation in ultracold atomic gases, elementary excitations play an important role in understanding various properties of the Bose-Einstein condensates, such as dynamics, thermodynamics and superfluidity. Elementary excitations have been experimentally investigated in trapped Bose-Einstein condensates, including the observation of collective modes rf:BEC ; rf:jila ; rf:mit and the measurement of the Bogoliubov excitation spectrum with use of Bragg spectroscopy rf:BEC ; rf:bogo .
An optical lattice is a periodic potential for atoms created by standing waves of laser beams. Experimental observations have also revealed properties of elementary excitations of Bose-Einstein condensates in an optical lattice. Stöferle et al. have studied excitation spectra of Bose-Einstein condensates in an optical lattice rf:sf\_mi . The Bose-Einstein condensates were excited by a modulation of the lattice depth, and a broad continuum of the excitation spectrum was observed; they suggested that the broad continuum characterizes the superfluid phase.
In the present paper, we study elementary excitations of Bose-Einstein condensates in an optical lattice with use of a Kronig-Penney potential. In previous theoretical papers rf:origin ; rf:dynam ; rf:tight ; rf:DNLS ; rf:comp ; rf:finite ; rf:sounv ; rf:phd ; rf:soft , elementary excitations in a sinusoidal lattice potential have been studied. Berg-Sørensen et al. numerically solved the Bogoliubov equations for a one-dimensional optical lattice, and they also calculated the low energy excitations analytically for shallow lattices within the Thomas-Fermi limit; they showed that the excitation spectrum is phonon-like at low energies rf:origin . The phonon dispersion of the excitation spectrum is directly connected to the superfluidity of a Bose-Einstein condensate in an optical lattice rf:land , and most of the papers support the phonon dispersion. In contrast, Ichioka et al. have reported on softening of the low energy excitations for deep lattices by numerically solving the Bogoliubov equations rf:soft . Such softening of the excitation spectrum, if it really exists, suggests an instability of superfluidity of a Bose-Einstein condensate in an optical lattice. One of our purposes is to investigate the possibility of the softening by analytically solving the Bogoliubov equations with a Kronig-Penney potential. We will show that the softening does not occur in the Kronig-Penney model.
In previous papers, analytical methods to calculate the excitation spectrum in the presence of a sinusoidal lattice potential are limited either to deep or shallow lattices. For deep lattices where the overlap of the condensate wave functions in neighboring sites is sufficiently small, a tight-binding approximation provides analytical expressions for the first band of the excitation spectrum rf:tight ; rf:DNLS ; rf:comp ; rf:finite . For shallow lattices, the perturbative treatment of the lattice potential is applicable rf:sounv ; rf:phd . An advantage of the Kronig-Penney model is that one can calculate the excitation spectrum exactly with analytic methods for arbitrary values of the lattice depth.
Another advantage of the Kronig-Penney model is that the excitation spectrum can be related to tunneling properties of the excitations through a single potential barrier. Kagan et al. studied the tunneling problem of excitations and predicted that a barrier is transparent for excitations within a limited range of low energies; they called such a behavior the anomalous tunneling rf:antun . We shall see, in fact, that the anomalous tunneling is crucial to the phonon-like form of the low energy excitations and that the softening of the excitation spectrum does not occur.
The outline of the present paper is as follows. In Sec. II, we introduce a formulation of the problem using the Bogoliubov theory and calculate the condensate wave function in a Kronig-Penney potential. In Sec. III, we analytically solve the Bogoliubov equations and obtain the band structure of the excitation spectrum in the Kronig-Penney potential. We derive the phonon-like form of the low energy excitation. We also calculate the phonon velocity and the band gap. Conclusions and open questions are discussed in the final section.
## II Condensate wave function in a Kronig-Penney potential
We consider a Bose-Einstein condensate confined in a combined potential of an axisymmetric harmonic potential and a one-dimensional periodic potential along the axial direction (the $`x`$ axis). As the periodic potential, we adopt a Kronig-Penney potential,
$`V(x)=V_0{\displaystyle \underset{n=\mathrm{}}{\overset{\mathrm{}}{}}}\delta (xna),`$ (1)
where $`a`$ is the lattice constant, and $`V_0`$ is the strength of the $`\delta `$-function potential barrier, expressing the lattice depth. It is assumed that the condensate is so elongated along the axial direction that the axial confinement can be neglected. We assume that the frequency of the harmonic potential of the radial direction is large enough compared to the excitation energy for the axial direction. In this situation, the one-dimensional treatment of the problem is justified.
Our formulation of the problem is based on the mean-field theory rf:BEC , which consists of the time-independent Gross-Pitaevskii equation and the Bogoliubov equations. They are
$`\left[{\displaystyle \frac{\mathrm{}^2}{2m}}{\displaystyle \frac{d^2}{dx^2}}+V(x)+g|\mathrm{\Psi }_0(x)|^2\right]\mathrm{\Psi }_0(x)=\mu \mathrm{\Psi }_0(x),`$ (2)
and
$`\begin{array}{cc}\left(\begin{array}{cc}H_0& g\mathrm{\Psi }_0(x)^2\\ g\mathrm{\Psi }_0(x)^2& H_0\end{array}\right)\left(\begin{array}{cc}u(x)& \\ v(x)& \end{array}\right)=\epsilon \left(\begin{array}{cc}u(x)& \\ v(x)& \end{array}\right),& \end{array}`$ (10)
$`H_0={\displaystyle \frac{\mathrm{}^2}{2m}}{\displaystyle \frac{d^2}{dx^2}}\mu +V\left(x\right)+2g|\mathrm{\Psi }_0(x)|^2.`$ (11)
Here $`\mu `$ is the chemical potential and $`m`$ is the mass of an atom. Since the radial confinement is harmonic and sufficiently tight, the coupling constant is affected by the harmonic oscillator length $`a_{}`$ of the radial confinement as $`g=\frac{2\mathrm{}^2a_s}{ma_{}^2}`$, where $`a_s`$ is the $`s`$-wave scattering length rf:1dime . The Gross-Pitaevskii equation determines the condensate wave function $`\mathrm{\Psi }_0(x)`$. The Bogoliubov equations determine the energy $`\epsilon `$ of the elementary excitations of the condensate and the wave functions of the excitation $`(u(x),v(x))^𝐭`$. The elementary excitations correspond to small fluctuations of the condensate wave function from the static equilibrium rf:BEC .
The periodic potential is sinusoidal in experiments of atomic gases trapped in an optical lattice. However, the Kronig-Penney model is useful to understand the problem of the periodic potential qualitatively, because it allows an analytical treatment of the problem. Schematic picture of the condensate in the Kronig-Penney potential is shown in Fig. 1.
The chemical potential is related to the number of condensate atoms $`N_0`$ in each well by the normalization condition:
$$_{na}^{(n+1)a}𝑑x|\mathrm{\Psi }_0(x)|^2=N_0.$$
(12)
At first, we shall solve Eq. (2) and obtain the condensate wave function in a Kronig-Penney potential. We assume that the condensate does not have supercurrent; therefore one can regard, without loss of generality, the condensate wave function as real. Multiplying Eq. (2) by $`\frac{d\mathrm{\Psi }_0}{dx}`$, one obtains the first integral of Eq. (2),
$`\left({\displaystyle \frac{d\mathrm{\Psi }_0}{dx}}\right)^2={\displaystyle \frac{\mu }{\xi ^2g}}(A^2{\displaystyle \frac{g}{\mu }}\mathrm{\Psi }_0^2)^2,`$ (13)
where $`\xi \frac{\mathrm{}}{\sqrt{m\mu }}`$ is the healing length, and $`A`$ is a constant corresponding to the value of the condensate wave function at the center of each lattice site. It is expressed as
$`A\sqrt{{\displaystyle \frac{g}{\mu }}}\mathrm{\Psi }_0((n+{\displaystyle \frac{1}{2}})a).`$ (14)
Integrating Eq. (13) again, one obtains
$`\mathrm{\Psi }_0(x)=\sqrt{{\displaystyle \frac{\mu }{g}}}A\mathrm{sn}(\sqrt{2A^2}(|xna|+x_0),{\displaystyle \frac{A}{\sqrt{2A^2}}}),`$
$`\left(n{\displaystyle \frac{1}{2}}\right)a<x<\left(n+{\displaystyle \frac{1}{2}}\right)a,`$ (15)
The boundary conditions at $`x=na`$ and $`x=(n+\frac{1}{2})a`$ determine the constants $`A`$ and $`x_0`$. This solution has been obtained in Refs. rf:kron ; rf:penn .
In order to calculate the excitation spectrum analytically, we assume that the lattice constant is sufficiently larger than the healing length. This assumption also allows us to relate the band structure of the excitation spectrum to the tunneling properties of the excitations. In this situation, the condensate wave function near the center of each lattice site is not affected by other potential barriers. Then, one approximately obtains the ground state solution of Eq. (2) as
$`\mathrm{\Psi }_0(x)=\sqrt{{\displaystyle \frac{\mu }{g}}}\mathrm{tanh}\left({\displaystyle \frac{|xna|+x_0}{\xi }}\right),`$
$`\left(n{\displaystyle \frac{1}{2}}\right)a<x<\left(n+{\displaystyle \frac{1}{2}}\right)a,`$ (16)
where $`x_0`$ is determined by the boundary condition at $`x=na`$
$`\mathrm{\Psi }_0(na+0)=\mathrm{\Psi }_0(na0),`$ (17)
$`{\displaystyle \frac{d\mathrm{\Psi }_0}{dx}}.|_{na+0}`$ $`=`$ $`{\displaystyle \frac{d\mathrm{\Psi }_0}{dx}}.|_{na0}+{\displaystyle \frac{2mV_0}{\mathrm{}^2}}\mathrm{\Psi }_0(na),`$ (18)
as
$`\mathrm{tanh}{\displaystyle \frac{x_0}{\xi }}={\displaystyle \frac{V_0+\sqrt{V_0^2+4(\mu \xi )^2}}{2\mu \xi }}.`$ (19)
This condensate wave function corresponds to the $`A1`$ limit of the Eq. (15).
Substituting Eq. (16) into the normalization condition of Eq. (12), we derive the relation between the chemical potential and $`N_0`$:
$$\mu \left(12\frac{\xi }{a}+2\frac{\xi }{a}\mathrm{tanh}\frac{x_0}{\xi }\right)=gn_0,$$
(20)
where $`n_0\frac{N_0}{a}`$ is the averaged density of the condensate. One can obtain approximate solutions of Eq. (20) in the limits of $`V_0gn_0\xi _0`$ and $`V_0gn_0\xi _0`$, where $`\xi _0\frac{\mathrm{}}{\sqrt{mgn_0}}`$. When $`V_0gn_0\xi _0`$, we expand Eq. (20) into power series of $`\frac{\xi _0}{a}`$ and $`\frac{gn_0\xi _0}{V_0}`$, and obtain
$`\mu gn_0(1+{\displaystyle \frac{2\xi _0}{a}}+{\displaystyle \frac{2\xi _0^2}{a^2}}{\displaystyle \frac{2gn_0\xi _0^2}{aV_0}}+{\displaystyle \frac{\xi _0^3}{a^3}}{\displaystyle \frac{6gn_0\xi _0^3}{a^2V_0}}).`$ (21)
In a similar way, when $`V_0gn_0\xi _0`$, we expand Eq. (20) into power series of $`\frac{\xi _0}{a}`$ and $`\frac{V_0}{gn_0\xi _0}`$, and obtain
$`\mu gn_0\left(1+{\displaystyle \frac{V_0}{gn_0a}}{\displaystyle \frac{V_0^2}{4a\xi _0(gn_0)^2}}\right).`$ (22)
In Eqs. (21) and (22), we express the expansions up to the third order of the small parameters. We show the chemical potential as a function of the potential strength $`V_0`$ in Fig. 2. The chemical potential increases monotonically as the potential strength increases, because the presence of the potential barriers makes the effect of repulsive interaction more significant.
## III Excitation spectrum
In this section, we shall solve the Bogoliubov equations with the condensate wave function of Eq. (16) and calculate the excitation spectrum of the condensate in a Kronig-Penney potential.
### III.1 Single barrier problem
A periodic potential can be viewed as a periodic array of potential barriers. The band structure of a single particle in a periodic potential can be expressed in terms of the tunneling properties of the particle in the presence of a single barrier potential rf:ashc . When one solves the Schrödinger equation with a single barrier, the wave function including the transmission amplitude is obtained. By imposing the Bloch’s theorem on the wave function, one obtains an equation to determine the band structure rf:ashc .
We shall determine the band structure of the excitation spectrum. In a similar way to the case of a single particle, one can relate the band structure of the excitation spectrum to the tunneling properties of the excitations in the presence of a single potential barrier, as discussed in the next subsection. In fact, once one solves the Bogoliubov equations with a single barrier and obtains the tunneling properties of the excitations, the band structure of the excitation spectrum can be easily obtained. Hence, we first focus our attention on the $`|x|<\frac{a}{2}`$ region, and solve the Bogoliubov equations analytically, with regard to the single barrier problem.
Kagan et al. have solved the Bogoliubov equations with a single rectangular barrier and numerically calculated the transmission amplitude rf:antun . Although they have also approximately obtained an analytical expression for the transmission amplitude, the approximation is not valid at very low energies. Since behaviors of the transmission amplitude around zero energy are crucial to the form of the excitation spectrum in an optical lattice, we need to obtain another analytical expression for the transmission amplitude. In the model of a $`\delta `$-function potential barrier which corresponds to the thin barrier limit of a rectangular barrier, one can calculate the transmission amplitude exactly.
There exist two independent solutions of the single barrier problem, corresponding to the types of scattering process. One solution $`\psi ^l(x)`$ describes the process where a Bogoliubov excitation comes from left, and the other solution $`\psi ^r(x)`$ describes the process where a Bogoliubov excitation comes from right. The schematic pictures of the solutions of the single barrier problem are shown in Fig. 3.
Substituting the condensate wave function of Eq. (16) into the Bogoliubov equations, let us obtain the solutions $`\psi ^l(x)`$ and $`\psi ^r(x)`$. On both sides of the barrier, one obtains four particular solutions analytically rf:antun under the boundary condition at $`|x|\xi `$:
$`(u_n(x),v_n(x))^𝐭e^{\frac{ip_nx}{\mathrm{}}}.`$ (23)
These solutions are
$`u_n(x)`$ $`=`$ $`\mathrm{\Lambda }_ne^{\frac{ip_nx}{\mathrm{}}}[.\mathrm{tanh}\left({\displaystyle \frac{|x|+x_0}{\xi }}\right)i\mathrm{sgn}(x){\displaystyle \frac{p_n\xi }{2\mathrm{}\epsilon }}`$
$`\times \left(\epsilon +\mu \mu \mathrm{tanh}^2\left({\displaystyle \frac{|x|+x_0}{\xi }}\right)\right)`$
$`+{\displaystyle \frac{E_{p_n}}{\epsilon }}\mathrm{tanh}\left({\displaystyle \frac{|x|+x_0}{\xi }}\right)i\mathrm{sgn}(x){\displaystyle \frac{E_{p_n}p_n\xi }{2\mathrm{}\epsilon }}.],`$
$`v_n(x)`$ $`=`$ $`\mathrm{\Lambda }_ne^{\frac{ip_nx}{\mathrm{}}}[.\mathrm{tanh}\left({\displaystyle \frac{|x|+x_0}{\xi }}\right)i\mathrm{sgn}(x){\displaystyle \frac{p_n\xi }{2\mathrm{}\epsilon }}`$
$`\times \left(\epsilon \mu +\mu \mathrm{tanh}^2\left({\displaystyle \frac{|x|+x_0}{\xi }}\right)\right)`$
$`{\displaystyle \frac{E_{p_n}}{\epsilon }}\mathrm{tanh}\left({\displaystyle \frac{|x|+x_0}{\xi }}\right)+i\mathrm{sgn}(x){\displaystyle \frac{E_{p_n}p_n\xi }{2\mathrm{}\epsilon }}.],`$
where
$`p_{1,2}`$ $`=`$ $`\pm p=\pm \sqrt{2m(\sqrt{\mu ^2+\epsilon ^2}\mu )},`$ (26)
$`p_{3,4}`$ $`=`$ $`i\gamma =i\sqrt{2m(\sqrt{\mu ^2+\epsilon ^2}+\mu )},`$ (27)
$`\mathrm{\Lambda }_n`$ $`=`$ $`\sqrt{{\displaystyle \frac{\mu ^2}{2g\epsilon }}}\times \{\begin{array}{cc}\frac{\sqrt{2\mu }+\mathrm{sgn}(x)i\sqrt{E_p}}{\sqrt{2\mu +E_p}},n=1,& \\ \frac{\sqrt{2\mu }\mathrm{sgn}(x)i\sqrt{E_p}}{\sqrt{2\mu +E_p}},n=2,& \\ 1,n=3,4& \end{array},`$ (31)
$`E_{p_n}`$ $`=`$ $`{\displaystyle \frac{p_n^2}{2m}}.`$ (32)
Wave functions $`(u_1(x),v_1(x))^𝐭`$ and $`(u_2(x),v_2(x))^𝐭`$ describe scattering components. Wave functions $`(u_3(x),v_3(x))^𝐭`$ at $`x<0`$ and $`(u_4(x),v_4(x))^𝐭`$ at $`x>0`$ describe the localized components around the potential barrier, because they decay exponentially at $`|x|\xi `$. Wave functions $`(u_3(x),v_3(x))^𝐭`$ at $`x>0`$ and $`(u_4(x),v_4(x))^𝐭`$ at $`x<0`$ diverge far from the potential barrier. It is noted that Eqs. (26) and (27) can be derived by solving
$`\epsilon =\sqrt{{\displaystyle \frac{p_n^2}{2m}}({\displaystyle \frac{p_n^2}{2m}}+2\mu )}.`$ (33)
Equation (33) expresses the Bogoliubov spectrum for a uniform system. The normalization constant $`\mathrm{\Lambda }_n`$ of the scattering components is determined to satisfy
$$u_n(x),v_n(x)=\sqrt{\frac{E_p+\mu \pm \epsilon }{2\epsilon }}e^{\frac{ip_nx}{\mathrm{}}}$$
(34)
at $`|x|\xi `$.
We omit the unphysical divergent components. The solutions $`\psi ^l(x)`$ and $`\psi ^r(x)`$ are expressed as superpositions of the remaining components. The solution $`\psi ^l(x)`$ is written as
$`\psi ^l(x)=(\begin{array}{cc}u^l& \\ v^l& \end{array})`$ $`=`$ $`\{\begin{array}{cc}(\begin{array}{cc}u_1& \\ v_1& \end{array})+r(\begin{array}{cc}u_2& \\ v_2& \end{array})\hfill & \\ +b(\begin{array}{cc}u_3& \\ v_3& \end{array}),\hfill & x<0,\hfill \\ t(\begin{array}{cc}u_1& \\ v_1& \end{array})+c(\begin{array}{cc}u_4& \\ v_4& \end{array}),\hfill & x>0,\hfill \end{array}`$ (50)
where the coefficients $`r`$, $`b`$, $`t`$, and $`c`$ are the amplitudes of the reflected, the left localized, the transmitted, and the right localized components, respectively. They are functions of the energy $`\epsilon `$ and the potential strength $`V_0`$. The boundary conditions at $`x=0`$ yield four equations to determine all the coefficients
$`\psi ^l(+0)=\psi ^l(0),`$ (51)
$`{\displaystyle \frac{d\psi ^l}{dx}}|_{+0}={\displaystyle \frac{d\psi ^l}{dx}}|_0+{\displaystyle \frac{2mV_0}{\mathrm{}^2}}\psi ^l(0).`$ (52)
These equations are linear simultaneous equations for the coefficients $`r`$, $`b`$, $`t`$ and $`c`$, and one can analytically solve them. Since the exact solutions of Eqs. (51) and (52) are unnecessarily complicated, we only write approximate forms of the coefficients. When $`\epsilon \mu `$ and $`V_0\mu \xi `$, one approximately obtains all the coefficients as
$`r`$ $`=`$ $`{\displaystyle \frac{\epsilon V_0\epsilon \mu \xi +i\frac{\epsilon ^2V_0}{\mu }}{Z}},`$ (53)
$`b`$ $`=`$ $`c={\displaystyle \frac{2\epsilon \mu \xi +\frac{4\epsilon \mu ^2\xi ^2}{V_0}+i\epsilon ^2\xi }{4Z}},`$ (54)
$`t`$ $`=`$ $`{\displaystyle \frac{\epsilon \mu \xi +i\mu ^2\xi }{Z}},`$ (55)
$`Z`$ $`=`$ $`\epsilon V_0\epsilon \mu \xi +i\mu ^2\xi .`$ (56)
Thus, we obtained $`\psi ^l(x)`$ analytically, and we can also calculate $`\psi ^r(x)`$ in the same way.
The transmission coefficient $`T|t|^2`$ and the phase shift $`\delta \mathrm{arg}(t)`$ are shown in Figs. 4 and 5, respectively, as a function of the energy. Expanding $`t`$ around $`\epsilon =0`$, one can analytically obtain approximate expressions of $`|t|`$ and $`\delta `$
$`|t|`$ $``$ $`1\alpha \left({\displaystyle \frac{\epsilon }{\mu }}\right)^2,`$ (57)
$`\delta `$ $``$ $`\beta {\displaystyle \frac{\epsilon }{\mu }}.`$ (58)
The coefficients $`\alpha `$ and $`\beta `$ are
$`\alpha =`$
$`{\displaystyle \frac{2(V_0\mu \xi )(V_0^3+\nu V_0^2+2\nu (\mu \xi )^24(\mu \xi )^3)+9(\mu \xi V_0)^2}{8(\mu \xi \nu )^2}},`$ (59)
$`\beta ={\displaystyle \frac{V_0^2+\nu V_03\mu \xi \nu +6(\mu \xi )^2}{2\mu \xi \nu }},`$ (60)
where
$`\nu =\sqrt{V_0^2+4(\mu \xi )^2}.`$ (61)
It is obvious from Figs. 4 and 5 and Eqs. (57) and (58) that the transmission coefficient $`T`$ approaches unity and the phase shift $`\delta `$ approaches zero as the energy is reduced to zero. This means that the potential barrier is transparent for low energy excitations. This behavior of the low energy excitations has been called the anomalous tunneling by Kagan $`\mathrm{𝑒𝑡}\mathrm{𝑎𝑙}.`$ rf:antun . Equations (55) and (56) shows that the peak of $`T`$ has a Lorentzian shape with half width $`\mathrm{\Delta }\epsilon V_0^1`$, and the anomalous tunneling becomes restricted to only the excitations with very low energies as strength of the potential barrier increases. This peculiar tunneling behavior of the excitations appears for arbitrary values of the potential strength.
### III.2 Band structure of excitation spectrum
In this subsection, we analytically calculate the band structure of the excitation spectrum with use of the solution of the single barrier problem obtained in the previous subsection. We show that the tunneling properties of the excitations determine the band structure.
Since the Bogoliubov equations are linear differential equations, a general solution of the equations can be described as a linear combination of independent solutions with the same energy. We can write a general solution of the Bogoliubov equations in the region $`|x|<\frac{a}{2}`$ as a linear combination of $`\psi ^l(x)`$ and $`\psi ^r(x)`$ rf:footnote :
$`\psi (x)=\chi \psi ^l(x)+\zeta \psi ^r(x),|x|<{\displaystyle \frac{a}{2}}.`$ (62)
where $`\chi `$ and $`\zeta `$ are arbitrary constants.
Now we extend this solution to all regions of $`x`$ by means of the Bloch’s theorem. The Bloch’s theorem asserts that $`\psi `$ satisfies
$`\psi (x+a)`$ $`=`$ $`e^{\frac{iqa}{\mathrm{}}}\psi (x),`$ (63)
$`{\displaystyle \frac{d\psi }{dx}}|_{x+a}`$ $`=`$ $`e^{\frac{iqa}{\mathrm{}}}{\displaystyle \frac{d\psi }{dx}}|_x,`$ (64)
where $`q`$ is the quasi-momentum. Equations (63) and (64) at $`x=\frac{a}{2}`$ yields an equation expressing the relation between the excitation energy $`\epsilon `$ and the quasi-momentum $`q`$:
$`\mathrm{cos}\left({\displaystyle \frac{qa}{\mathrm{}}}\right)={\displaystyle \frac{t^2r^2}{2t}}e^{\frac{ipa}{\mathrm{}}}+{\displaystyle \frac{1}{2t}}e^{\frac{ipa}{\mathrm{}}}.`$ (65)
Since we are assuming $`a\xi `$, the localized components around the potential barriers, which decay exponentially and vanish at $`|x|=\frac{a}{2}`$, do not appear explicitly in Eq. (65).
The Wronskian defined as
$`W(\psi ^j,\psi ^i)=u^j{\displaystyle \frac{d}{dx}}u^iu^i{\displaystyle \frac{d}{dx}}u^j+v^j{\displaystyle \frac{d}{dx}}v^iv^i{\displaystyle \frac{d}{dx}}v^j`$ (66)
yields relations between $`r`$ and $`t`$. We can simplify Eq. (65) by the relations. One can easily prove from the Bogoliubov equations that $`W`$ is independent of $`x`$ when $`\psi ^j`$ and $`\psi ^i`$ have the same energy. By evaluating $`W(\psi ^l,\psi ^l)`$, one obtains the conservation law of the energy flux rf:antun :
$`|t|^2+|r|^2=1.`$ (67)
Meanwhile, by evaluating $`W(\psi ^r,\psi ^l)`$, one obtains the relation
$`t=|t|e^{i\delta },r=\pm i|r|e^{i\delta }.`$ (68)
Substituting Eqs. (67) and (68) into Eq. (65), one obtains the simplified relation between the excitation energy and the quasi-momentum
$`{\displaystyle \frac{\mathrm{cos}\left(\frac{pa}{\mathrm{}}+\delta \right)}{|t|}}=\mathrm{cos}\left({\displaystyle \frac{qa}{\mathrm{}}}\right).`$ (69)
This relation is exactly the same form as that in the case of a single particle rf:ashc . We clearly see from Eq. (69) that the tunneling properties in the single barrier problem determine the relation between the excitation energy and the quasi-momentum, namely the band structure.
Solving Eq. (69) for the excitation energy $`\epsilon `$, we obtain the band structure of the excitation spectrum as shown in Fig. 6(a). Fig. 6(b) shows the left-hand side of Eq. (69) as a function of $`\epsilon `$. There exists no solution of Eq. (69) when the absolute value of the left-hand side exceeds unity, because the absolute value of the right-hand side is equal to or less than unity. Therefore, the energy regions where the absolute value of the left-hand side is greater than unity correspond to the forbidden regions of the excitation energy, namely the band gaps. They are expressed as the shaded regions in Fig. 6.
In the strong potential limit $`V_0\mathrm{}`$, the widths of all the energy bands become narrow. The energy bands approach the discrete eigenenergies of the Bogoliubov equations for a single potential well, because the condensate is perfectly divided into each well.
### III.3 Phonon dispersion
Fig. 6(a) shows that the first band of the excitation spectrum is phonon-like in the low energy regions. We shall verify that the form of the excitation spectrum at low energies is phonon-like for arbitrary values of the potential depth. The excitation spectrum in the low energy limit $`\epsilon 0`$ is calculated by expanding Eq. (69) around $`\epsilon =0`$. Substituting $`p\sqrt{\frac{m}{\mu }}\epsilon `$, Eqs. (57) and (58) into Eq. (69), one obtains the phonon dispersion of the excitation spectrum
$`\epsilon cq,`$ (70)
where the phonon velocity $`c`$ is
$`c=\sqrt{{\displaystyle \frac{\mu a^2}{m\left((a+\beta \xi )^22\xi ^2\alpha \right)}}}.`$ (71)
When $`V_0gn_0\xi _0`$, Eq. (71) can be approximated as
$`cc_0\sqrt{{\displaystyle \frac{gn_0a}{2V_0+gn_0a}}}\left(1+{\displaystyle \frac{5\xi _0}{2a}}\right),`$ (72)
where $`c_0\sqrt{\frac{gn_0}{m}}`$ is the phonon velocity in uniform systems. When $`V_0gn_0\xi _0`$, Eq. (71) can be approximated as
$`cc_0\left(1{\displaystyle \frac{3V_0^2}{16a\xi _0(gn_0)^2}}\right).`$ (73)
The phonon velocities for $`a=10\xi _0`$ and $`20\xi _0`$ as functions of $`V_0`$ are shown in Fig. 7. The phonon velocity decreases monotonically as the barrier strength increases, and this behavior is qualitatively consistent with the case of a sinusoidal potential rf:comp .
Thus the anomalous tunneling properties of Eqs. (57) and (58), namely the perfect transmission of very low energy excitations, leads to the phonon dispersion. It is pointed out in Ref. rf:soft that the excitation spectrum shows softening at large wavelengths for a deep sinusoidal lattice potential, relating to the instability of the superfluidity. In contrast, the excitation spectrum is always phonon-like in the Kronig-Penney potential because the anomalous tunneling occurs even for very strong potential barriers. The phonon dispersion of the excitation spectrum at low energies reflects the superfluidity of the condensate. In this sense, the anomalous tunneling is crucial to the stability of the superfluidity in the Kronig-Penney potential.
However, our results cannot eliminate the possibility of the softening in a sinusoidal periodic potential, because we have assumed in our calculation that the lattice constant is sufficiently larger than the healing length. In the Kronig-Penney model, this assumption assures that the size of the condensate wave function in each well hardly changes as the lattice potential becomes deep; consequently the size of the condensate is always sufficiently larger than the healing length. In contrast, even when the lattice constant is sufficiently larger than the healing length, the size of the condensate wave function in each well can be comparable to or smaller than the healing length in a deep sinusoidal potential. This is because the deeper the sinusoidal lattice potential is, the tighter the confinement of each well is. Thus, a possibility of the softening is remaining in the situation where the size of the condensate in each well is comparable to or smaller than the healing length.
### III.4 Band gap
The $`j`$th band gap $`\mathrm{\Delta }_{j\mathrm{th}}`$ is determined by the excitation energies of the top of lower band and the bottom of higher band, and they are obtained by solving Eq. (69) at $`q=\pm \frac{\mathrm{}\pi }{a}`$. Here we show the first band gap $`\mathrm{\Delta }_{1\mathrm{s}\mathrm{t}}`$ between the lowest and the second lowest band, as a function of the potential strength $`V_0`$ in Fig. 8.
When $`V_0gn_0a`$, the first band gap is much larger than the first band width. In this case, one can approximately obtain
$`\mathrm{\Delta }_{1\mathrm{s}\mathrm{t}}{\displaystyle \frac{\pi gn_0\xi _0}{a}}\left(1+{\displaystyle \frac{2\xi _0}{a}}\right)gn_0\sqrt{{\displaystyle \frac{gn_0\xi _0^2}{V_0a}}}(1+{\displaystyle \frac{5\xi _0}{2a}}).`$ (74)
When $`V_0gn_0\xi _0`$, one can obtain
$`\mathrm{\Delta }_{1\mathrm{s}\mathrm{t}}{\displaystyle \frac{\pi \xi _0V_0}{a^2}}.`$ (75)
Equation (75) and Fig. 8 show that at first the band gap increases linearly with the potential strength. The first term of the right-hand side of Eq. (74) corresponds to the lowest excitation energy of a condensate in a single well. Since the second term vanishes in the limit of $`V_0\mathrm{}`$, the band gap becomes the lowest excitation energy in a single well.
### III.5 Tight-binding limit
When the potential is sufficiently strong $`V_0gn_0a`$, one can approximately solve Eq. (69) and obtain the first band of the excitation spectrum. Since the excitation energy in the first band is much smaller than the chemical potential in this situation, the approximate expression of $`t`$ of Eqs. (55) and (56) is adequate. One can rewrite Eqs. (55) and (56) as
$`\mathrm{cos}\delta `$ $`=`$ $`|t|\left(1+{\displaystyle \frac{\epsilon ^2V_0}{\mu ^3\xi }}\right),`$ (76)
$`\mathrm{sin}\delta `$ $`=`$ $`|t|{\displaystyle \frac{\epsilon }{\mu }}\left({\displaystyle \frac{V_0}{\mu \xi }}2\right),`$ (77)
where
$`|t|={\displaystyle \frac{\mu ^2\xi }{\sqrt{\epsilon ^2V_0(V_02\mu \xi )+\mu ^4\xi ^2}}}.`$ (78)
Substituting Eq. (76), (77) and (78) into Eq. (69), one obtains
$$\epsilon gn_0\sqrt{\frac{2gn_0\xi _0^2}{aV_0}}\left(1+\frac{5\xi _0}{2a}\right)\left|\mathrm{sin}\left(\frac{qa}{2\mathrm{}}\right)\right|,$$
(79)
where $`\mu `$ is replaced with $`gn_0`$ by using Eq. (21).
Meanwhile, the first band of the excitation spectrum has been calculated using the tight-binding approximation in many theoretical papers rf:tight ; rf:DNLS ; rf:comp ; rf:finite . For a very deep lattice, the spectrum takes the form
$$\epsilon \sqrt{\frac{2\eta }{\kappa }}\left|\mathrm{sin}\left(\frac{qa}{2\mathrm{}}\right)\right|,$$
(80)
where $`\kappa `$ is the compressibility of the system and $`\eta `$ is the lowest energy bandwidth of stationary current-carrying states of the condensate rf:comp . The compressibility $`\kappa `$ is defined by the thermodynamic relation
$$\frac{1}{\kappa }=N_0\frac{\mu }{N_0}.$$
(81)
The parameter $`\eta `$ is related to the effective mass $`m^{}`$ through the relation
$$\eta =\frac{2\mathrm{}^2}{m^{}a^2}.$$
(82)
The definition of the effective mass is
$$\frac{1}{m^{}}=\frac{^2\stackrel{~}{ϵ}}{k^2}.|_{k=0}$$
(83)
where $`k`$ is the quasi-momentum of the current-carrying condensate and $`\stackrel{~}{ϵ}`$ is the energy per particle of the condensate. We should be careful not to confuse the energy dispersion of the Bogoliubov excitations $`\epsilon (q)`$, which is our major interest, and the energy per particle $`\stackrel{~}{ϵ}(k)`$ of the condensate itself.
Using Eqs. (81) and (82), one can obtain
$`{\displaystyle \frac{1}{\kappa }}`$ $``$ $`gn_0\left(1+{\displaystyle \frac{\xi _0}{a}}\right),`$ (84)
$`\eta `$ $``$ $`{\displaystyle \frac{(gn_0\xi _0)^2}{aV_0}}\left(1+{\displaystyle \frac{4\xi _0}{a}}\right).`$ (85)
Here we have used Eq. (21) to calculate the compressibility. Substituting Eqs. (84) and (85) into Eq. (80), we see that our calculation of the first band of the excitation spectrum is consistent with the result of the tight-binding approximation when the potential is sufficiently strong.
## IV CONCLUSION
In summary, we have investigated the band structure of the excitation spectrum of the condensate in a Kronig-Penney potential. Connecting the analytical solutions of Bogoliubov equations in the single barrier problem by means of the Bloch’s theorem, the relation between the excitation energy and the quasi-momentum have been obtained, which determines the band structure. We have shown that the excitation spectrum is gapless and linear in the low energy region due to the anomalous tunneling property of the excitations.
While we have considered current-free condensates in our calculations, the elementary excitations of condensates with large current in periodic potentials are known to exhibit the Landau and dynamical instabilities rf:dynam ; rf:DNLS ; rf:inst ; rf:dy\_inst . Since both instabilities are associated with the specific properties of the excitation spectra in the low energy region, the tunneling properties at low energy may affect the instability when the problem is considered with the Kronig-Penney model. Hence, it will be interesting to study the excitations of the current-carrying condensates in the Kronig-Penney potential and discuss the relation between the instability and the anomalous tunneling of the Bogoliubov excitations.
Note added. After the submission of this paper, we learned of a relevant paper by Kovrizhin rf:add . In Ref. rf:add , the tunneling problem of the Bogoliubov excitations through a $`delta`$-function potential barrier is investigated. Exact form of the wave function Eq. (50) is derived, and the anomalous behavior of the transmission coefficient is discussed.
###### Acknowledgements.
The authors would like to thank T. Kimura, N. Yokoshi, T. Nikuni and N. Hatano for fruitful comments and discussions. The work is partially supported by a Grant for the 21st Century COE Program (Holistic Research and Education center for Physics of Self-organization Systems) at Waseda University from the Ministry of Education, Sports, Culture, Science and Technology of Japan. I. D. is supported by Grant-in-Aid for JSPS fellows.
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# Fusion products, Kostka polynomials, and fermionic characters of (𝑠𝑢)̂(𝑟+1)_𝑘
## 1. Introduction
In an earlier paper (which we will refer to as I) we provided a physics application of the ideas of Feigin and Stoyanovsky (FS) , who showed that by counting the number of linearly independent symmetric polynomials that vanish when $`k+1`$ of their variables coincide one may obtain fermionic formulæ for the characters of integrable representations of the affine Lie algebra $`\widehat{su}(2)_k`$ . We also sketched how a similar count of the number of symmetric polynomials in two types of variables allows the computation of the character of the vacuum representation of $`\widehat{su}(3)_k`$. The reason for the restriction to the vacuum representation of $`\widehat{su}(3)_k`$ was that a naïve application of the FS strategy to most representations of higher rank groups does not yield the characters. Here we will explain both why the naïve method fails, and what must be done to correct it. To be more precise, we explain the ideas behind the corrected character formulæ in a form accessible to physicists. We will not give all the technical details of the proofs. The mathematical details and the proofs are given in .
The basic idea of is most simply explained in the language of conformal field theory: we compute matrix elements of ladder-operator currents between the highest weight state that defines the representation of interest and any other weight state in the representation. These matrix elements are rational functions (in the $`\widehat{su}(2)`$ case they are symmetric polynomials) in the co-ordinates of the current operators and have certain restrictions imposed on them by relations in the Lie algebra, and by the the highest weight condition. By counting the number of possible functions satisfying these constraints we are able to count the dimension of all weight subspaces.
For $`\widehat{su}(r+1)`$, and for a representation whose top graded component forms a finite-dimensional representation of $`su(r+1)`$ having only one non-zero Dynkin index (and hence corresponds to a rectangular Young diagram), it is not difficult to count the dimensions of the function space. For a general representation, it is hard. We find, however, that by introducing a fusion-product representation we can retain a function space whose dimensions we can count. The price to be paid is that the fusion-product representation is reducible. The generating function for the dimensions, although of fermionic form, is therefore a sum of characters of irreducible representations, moreover one with coefficients that are polynomials in the variable $`1/q`$. These $`q`$-dependent coefficients can be identified as being generalized Kostka polynomials . The decomposition may then be inverted to obtain the character of any desired irreducible representation in terms of the fermionic fusion-product characters.
Kostka polynomials and their generalizations are polynomials in $`q`$, with coefficients which are non-negative integers. The classical Kostka polynomials are known in the theory of symmetric functions as the transition coefficients between Hall-Littlewood and Schur polynomials. In physics, the Kostka polynomials made their first appearance in the study of the completeness of Bethe ansatz states in the Heisenberg spin chain . The methods introduced there, together with the theory of crystal bases, is the basis for the subsequent studies of .
In order to establish our notation, in section two we provide a brief review of affine Lie algebras. In section three we introduce the “principal subspace” of Feigin and Stoyanovsky, and the function spaces that are dual to them. In section four we show how the affine Weyl translations allow us, in the special case of rectangular representations, to use the character of the principal subspace to compute the character of the full space. This is the process that in I we called “filling the bose sea.” In section five we explain why the naïve extension of the method to non-rectangular representations fails, and exhibit the subtractions necessary to obtain correct formulæ for the characters. These formulæ were originally obtained numerically. In section six we show how these correct formulæ find their explanation in characters of reducible fusion-product representations, and also explain the origin of the $`q`$-dependent Kostka polynomial coefficients in the fusion-product decomposition. In section seven, we make the connection between these Kostka polynomials and the WZW conformal field theory. Finally, in section eight, we combine all the results to obtain a character formula for arbitrary highest-weight representations of $`\widehat{su}(r+1)_k`$, equation (8.2). Section nine is devoted to the conclusions and an outlook. Details concerning the characters of the principal subspaces can be found in appendix A, while appendix B contains details about the explicit formula for the (generalized) Kostka polynomials.
## 2. Notation
A simple Lie algebra $`𝔤`$ gives rise to the affine Lie algebra $`\widehat{𝔤}^{}=𝔤[t,t^1]\widehat{c}`$ with commutation relations
$$[at^m,bt^n]=[a,b]t^{n+m}+\widehat{c}na,b\delta _{m+n,0}.$$
(2.1)
Here $`a,b`$ denotes the Killing form, normalized so that a long root has length $`\sqrt{2}`$. The element $`\widehat{c}`$ is central, and so commutes with all $`a\widehat{𝔤}`$. It is convenient to adjoin to this algebra a grading operator $`\widehat{d}`$ that also commutes with $`\widehat{c}`$, and such that $`[\widehat{d},at^n]=n(at^n)`$. The algebra with the adjoined operator $`\widehat{d}`$ is denoted by $`\widehat{𝔤}`$. We will often write $`at^na[n]`$.
The affine algebra is defined for polynomials in $`t,t^1`$. We will later need to extend the definition to functions which are rational functions in $`t`$, or more generally, Laurent series in $`t^1`$. If $`f(t)`$ and $`g(t)`$ are such series, we can write the commutation relations as
$$[af,bg]=[a,b]fg+\frac{\widehat{c}a,b}{2\pi i}f^{}g𝑑t.$$
(2.2)
Here the integral treats the otherwise formal parameter $`t`$ as a complex variable taking values on a circle surrounding $`t=0`$. The result of the integration, $`\mathrm{Res}_{t=0}(f^{}gdt)`$, involves only a finite sum because of the polynomial condition on the negative powers of $`t`$ in $`f`$, $`g`$.
We will use the Chevalley basis for the generators of the simple algebra $`𝔤`$, in which each simple root $`\alpha _i`$ is associated with a step-up ladder operator $`e_{\alpha _i}E^{\alpha _i}`$, a step-down ladder operator $`f_{\alpha _i}E^{\alpha _i}`$, and the co-root element of the Cartan algebra $`h_{\alpha _i}\frac{2\alpha _iH}{\alpha _i^2}`$. The commutation relations of these generators are
$`[h_{\alpha _i},h_{\alpha _j}]`$ $`=0,`$ $`[h_{\alpha _i},e_{\alpha _j}]`$ $`=(𝐂_r)_{ji}e_{\alpha _j},`$
$`[h_{\alpha _i},f_{\alpha _j}]`$ $`=(𝐂_r)_{ji}f_{\alpha _j},`$ $`[e_{\alpha _i},f_{\alpha _j}]`$ $`=\delta _{ij}h_{\alpha _i}.`$ (2.3)
The remaining generators are obtained by repeated commutators of these, subject to the Serre relations
$$[\mathrm{ad}(e_{\alpha _i})]^{1(𝐂_r)_{ji}}e_{\alpha _j}=0,[\mathrm{ad}(f_{\alpha _i})]^{1(𝐂_r)_{ji}}f_{\alpha _j}=0.$$
(2.4)
In these expressions $`(𝐂_r)_{ij}`$ denotes the elements of the Cartan matrix of $`𝔤`$:
$$(𝐂_r)_{ij}\stackrel{\mathrm{def}}{=}\frac{2(\alpha _i\alpha _j)}{\alpha _j^2}.$$
(2.5)
In the case $`𝔤=su(r+1)`$ the elements of the Cartan matrix are given by $`(𝐂_r)_{ij}=2\delta _{i,j}\delta _{|ij|,1}`$. We will also have cause to use the inverse Cartan matrix whose elements (for $`𝔤=su(r+1)`$) are $`(𝐂_r^1)_{ij}=\mathrm{min}(i,j)\frac{ij}{r+1}`$.
Our interest is in integrable representations of $`\widehat{𝔤}`$, and in particular of $`\widehat{su}(r+1)`$. An integrable representation is one that, under restriction to any subgroup of $`\widehat{𝔤}`$ isomorphic to $`su(2)`$, decomposes into a set of finite-dimensional representations of this $`su(2)`$. From the work of Kac it is known that in any irreducible integrable representation the generator $`\widehat{c}`$ will act as a positive integer multiple of the identity, $`\widehat{c}k𝕀`$, and we will follow the physics convention of appending the integer $`k`$, the level of the representation, to the name of the group, and so write $`\widehat{𝔤}_k`$. These integrable representations are all highest-weight representations whose highest weight vector is annihilated by $`e_{\alpha _i}[n]`$, $`n0`$ and by $`h_{\alpha _i}[n]`$, $`f_{\alpha _i}[n]`$ with $`n>0`$. The top d-graded subspaces (where $`\widehat{d}`$ is taken to act as zero) form finite-dimensional representations of the simple algebra $`𝔤`$, but not all representations of $`𝔤`$ can form top components of integrable representations of $`\widehat{𝔤}_k`$. In the case of $`\widehat{su}(r+1)_k`$ the restriction on the representations of $`su(r+1)`$ that can appear as top components is that the number of columns in the Young diagram labeling the representation must be $`k`$.
We wish to obtain the dimensions $`\mathrm{mult}(\widehat{\mu })`$ of the weight spaces $`\widehat{\mu }(\mu ;k;d)`$ in an irreducible integrable level-$`k`$ representation $`H_{\widehat{\lambda }}`$ of $`\widehat{𝔤}_k`$, whose highest weight is $`\widehat{\lambda }=(\lambda ;k;0)`$. Here $`\lambda `$ and $`\mu `$ denote a highest weight and an arbitrary weight, respectively, of the associated finite-dimensional simple algebra $`𝔤`$, and $`d`$, a non-positive integer, is the eigenvalue of the grading operator $`\widehat{d}`$. The character of the representation is then defined to be
$$\mathrm{ch}H_{\widehat{\lambda }}(q,𝐱)=\underset{\widehat{\mu }}{}\mathrm{mult}(\widehat{\mu })x_1^{\mu _1}x_2^{\mu _2}\mathrm{}x_r^{\mu _r}q^d.$$
(2.6)
Here $`\mu _1,\mathrm{},\mu _r`$ are the components of the weight $`\mu `$ in the basis of the fundamental weights of the finite-dimensional algebra $`\omega _i=_j(𝐂_r^1)_{ji}\alpha _j`$, i.e. $`\mu =_{i=1}^r\mu _i\omega _i`$, and $`𝐱=(x_1,\mathrm{},x_r)`$. Because the eigenvalues $`d`$ are zero or negative, the character is a formal series in positive powers of $`q`$.
In this paper, we present fermionic formulas for the characters of all highest weight integrable representations of $`\widehat{su}(r+1)`$. These formulas are inspired by the construction of Feigin and Stoyanovsky, which we present below.
## 3. The principal subspace
In this section we briefly review the strategy of Feigin and Stoyanovsky , as we presented it in I. The reader should refer to I for more details.
Let $`|\widehat{\lambda }`$ be the highest-weight state of the irreducible representation $`H_{\widehat{\lambda }}`$. Recall that $`e_{\alpha _i}`$ and $`f_{\alpha _i}`$ are the step-up and step-down ladder operators corresponding to simple root $`\alpha _i`$ in the finite-dimensional algebra $`𝔤`$, and $`e_{\alpha _i}[n]`$ and $`f_{\alpha _i}[n]`$ are their $`\widehat{𝔤}`$ relatives. Since $`H_{\widehat{\lambda }}`$ is irreducible, by acting on $`|\widehat{\lambda }`$ with products of $`e_{\alpha _i}[n]`$, $`n<0`$, and $`f_{\alpha _i}[n]`$, $`n0`$, we generate the entire space $`H_{\widehat{\lambda }}`$. If we restrict ourselves to products containing only the $`f_{\alpha _i}[n]`$, $`n0`$, we obtain what is known as the principal subspace $`W_{\widehat{\lambda }}`$.
Note that in I we used a commuting set of operators applied to the highest weight state to generate a principal subspace (in the case of $`su(2)`$ and $`su(3)`$). This subspace had the advantage for physics applications that the resulting correlation functions were polynomials, rather than rational functions. For the purposes of this paper, however, it is necessary to use the Feigin-Stoyanovsky subspace, in order to obtain the simple closed-form formulæ for the principal subspace multiplicities for arbitrary representations that we report here.
Consider, for example, $`su(2)`$ where there is only one simple root $`\alpha _1`$ and so the root label is redundant. Here, since each application of an $`f[n]`$ moves us one step to the left and $`n`$ steps down in the weight diagram, we cannot reach any weight to the right of the column lying below $`\lambda `$. (See figure 2 where the upper part of the weight diagram of the vacuum representation of $`\widehat{su}(2)_2`$ is displayed.) Furthermore, we can reach weights close to this column via fewer distinct products $`f[n_1]f[n_2]\mathrm{}f[n_m]`$ than the dimension of the weight space. Because of this paucity of paths, even if all the $`f[n_1]f[n_2]\mathrm{}f[n_m]|\widehat{\lambda }`$ were linearly independent, we would be able to obtain only a subspace in each weight space. If, however, we look at weights in columns far to the left in the weight diagram, the number of distinct paths leading to a given weight grows rapidly, whilst the dimensions of the weight spaces near the head of any such column remain small. It is plausible, and indeed true, that we can obtain all states in such a weight space by applying suitable products of $`f[n]`$’s to $`|\widehat{\lambda }`$. Thus, by counting the number of linearly independent paths, we can obtain the multiplicity of these distant weights. Because the weight diagram is invariant under the action of affine Weyl translations, as we explain in section 4, this information provides the dimension of all the weight spaces, and hence the complete character.
We could in principal use Lie algebra relations to determine the number of independent paths between the weights. It is easier, however, to obtain the dimensions of the weights $`\widehat{\mu }`$ appearing in the principal subspace by counting the number of linearly independent functions $`F`$ that arise as matrix elements
$$F(\{z^{(1)}\},\{z^{(2)}\},\mathrm{})=v|R\{f_{\alpha _1}(z_1^{(1)})\mathrm{}f_{\alpha _1}(z_{m^{(1)}}^{(1)})f_{\alpha _2}(z_1^{(2)})\mathrm{}f_{\alpha _2}(z_{m^{(2)}}^{(2)})\mathrm{}\}|\widehat{\lambda },$$
(3.1)
where $`|v`$ is an element of a weight space $`\widehat{\mu }`$.
In addition, $`f_{\alpha _i}(z)`$ denotes the current operator
$$f_{\alpha _i}(z)=\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}f_{\alpha _i}[n]z^{n1},$$
(3.2)
and the $`R`$ denotes “radial ordering”, meaning that operators $`f(z)`$ with larger $`|z|`$ are placed to the left of those with smaller $`|z|`$’s<sup>1</sup><sup>1</sup>1Strictly speaking, we regard $`z_i^{(\alpha )}`$’s as formal variables and consider matrix elements of all possible orderings of the roots $`\alpha _i`$. We then count the resulting number of linearly independent functions in the formal variables $`z_i^{(\alpha )}`$.. By duality, the number of linearly independent functions $`F`$ with $`m^{(i)}`$ currents $`f_{\alpha _i}(z)`$ and fixed total degree $`N`$ is equal to the dimension, in the principal subspace $`W_{\widehat{\lambda }}`$ of the weight
$$\widehat{\mu }=(\lambda \underset{i}{}m^{(i)}\alpha _i;k;N\underset{i}{}m^{(i)}).$$
(3.3)
Thus, the numbers $`m^{(i)}`$ correspond to the number of roots $`\alpha _i`$ which need to be subtracted from $`\lambda `$ to obtain $`\mu `$, hence, in terms of the fundamental weights, we have $`\mu =_{i=1}^r(l_i_j(𝐂_r)_{ij}m^{(j)})\omega _i`$, where the $`l_i`$ are the components of $`\lambda `$, i.e. $`\lambda =_il_i\omega _i`$.
The linear dependencies between products of $`f_{\alpha _i}[n]`$ due to relations in the Lie algebra, the integrability, and the properties of the highest weight vector $`|\widehat{\lambda }`$ translate, through duality, into conditions satisfied by the function $`F`$. For example, operators $`f_{\alpha _i}(z)`$ corresponding to the same simple root commute with each other. The function $`F`$ is therefore a symmetric function in the $`m^{(i)}`$ variables $`z_j^{(i)}`$ for fixed $`i`$. It must also possess certain poles and zeros whose exact form we will specify below. Once we know all these properties, we can set out to count the number of independent functions. It is not, however, easy to be sure that we have obtained a complete set of constraints on $`F`$. If we have have missed some, we will over-count the multiplicities.
In the case of $`\widehat{su}(r+1)_k`$ representations with only one non-zero Dynkin index, i.e. $`\lambda =l\omega _p`$, it is not too hard to find all the constraints on functions $`F(\{z_j^{(i)}\})`$. We already observed that $`F(\{z_j^{(i)}\})`$ is symmetric under the exchange of the variables $`z_j^{(i)}z_j^{}^{(i)}`$. Next we observe that the commutator
$$[f_{\alpha _i},f_{\alpha _{i+1}}]=f_{\alpha _i+\alpha _{i+1}}$$
(3.4)
implies the operator product
$$f_{\alpha _i}(z)f_{\alpha _{i+1}}(w)=\frac{f_{\alpha _i+\alpha _{i+1}}(w)}{(zw)}+\text{regular terms},$$
(3.5)
and this in turn implies that $`F(\{z_j^{(i)}\})`$ can have a pole of at most order one when the coordinates of two currents corresponding to adjacent simple roots coincide, $`z_j^{(i)}=z_j^{}^{(i+1)}`$. (There is no pole for non-adjacent indices because $`\alpha _i+\alpha _j`$ is not a root unless $`j=i\pm 1`$.) We will refer to the index $`i`$ as the color of the variables. In the representation $`\lambda =l\omega _p`$, and for $`l>0`$, we have that $`f_{\alpha _p}[0]|\widehat{\lambda }0`$, and since $`f_{\alpha _p}[0]`$ comes with coefficient $`z^1`$, the function $`F(\{z_j^{(i)}\})`$ may have a pole of order one at $`z_j^{(p)}=0`$. (Note that $`f_{\alpha _i}[0]|\widehat{\lambda }=0`$ for $`ip`$.) These are the only possible poles, and we make them explicit by writing
$$F(\{z_j^{(i)}\})=\frac{f(\{z_j^{(i)}\})}{_j(z_j^{(p)})_{i=1}^{r1}_{j,j^{}}(z_j^{(i)}z_j^{}^{(i+1)})},$$
(3.6)
where $`f(\{z_j^{(i)}\})`$ is a now a polynomial symmetric under the exchange of variables of the same color: $`z_j^{(i)}z_j^{}^{(i)}`$. Because of relations in the algebra and properties of the representation $`f(\{z_j^{(i)}\})`$ is not an arbitrary symmetric polynomial, but must possess certain zeros. These we now describe:
* The integrability condition requires that $`[f_{\alpha _i}(z)]^{k+1}`$ annihilates any vector in the representation. This tells us that $`f(\{z_j^{(i)}\})=0`$ when $`z_1^{(i)}=z_2^{(i)}=\mathrm{}=z_{k+1}^{(i)}`$, for any color $`i`$.
* The integrability properties of the top component of the representation with highest weight $`\lambda =l\omega _p`$ tell us that $`f_{\alpha _p}^{l+1}[0]|\widehat{\lambda }=0`$. The function $`f(\{z_j^{(i)}\})`$ must therefore have a zero when $`l+1`$ of the $`z_j^{(p)}`$ become zero.
* The Serre relations (2.4) are
$$[f_{\alpha _i},[f_{\alpha _i},f_{\alpha _{i+1}}]]=[f_{\alpha _{i+1}},[f_{\alpha _{i+1}},f_{\alpha _i}]]=0,$$
(3.7)
and indicate that $`F`$ should have no pole if two currents of color $`i`$ are made to coincide with a current of adjacent color. This requires that $`f(\{z_j^{(i)}\})`$ have a pole-canceling zero when $`z_1^{(i)}=z_2^{(i)}=z_1^{(i+1)}`$ or when $`z_1^{(i)}=z_1^{(i+1)}=z_2^{(i+1)}`$, for $`i=1,\mathrm{},r1`$.
To summarize: we find that the polynomial $`f(\{z_j^{(i)}\})`$ is symmetric in the variables corresponding to the same color, and vanishes when any of the following conditions holds
$`z_1^{(i)}`$ $`=\mathrm{}=z_{k+1}^{(i)},i`$ (3.8)
$`z_1^{(i)}`$ $`=z_2^{(i)}=z_1^{(i+1)},z_1^{(i)}=z_1^{(i+1)}=z_2^{(i+1)},i=1,\mathrm{},r1`$ (3.9)
$`z_1^{(p)}`$ $`=\mathrm{}=z_{l+1}^{(p)}=0.`$ (3.10)
We will denote the space of rational functions $`F(\{z_j^{(i)}\})`$ (3.6) such that $`f(\{z_j^{(i)}\})`$ satisfies (3.8), (3.9) and (3.10) by $`_{l\omega _p;k}`$. We now define the character of this space to be
$$\mathrm{ch}_{l\omega _p;k}(q,𝐱)\stackrel{\mathrm{def}}{=}x_p^l\underset{\{F(z)\}}{}q^{\mathrm{deg}(F)+_{i=1}^rm^{(i)}}\left(\underset{i=1}{\overset{r}{}}x_i^{_j(𝐂_r)_{ji}m^{(j)}}\right),$$
(3.11)
where the sum is over all the functions $`F(\{z_j^{(i)}\})`$ in the space $`_{l\omega _p;k}`$. The powers of $`q`$ and $`x_i`$ are motivated by the form of the weights (3.3), which we rewrite here in terms of the fundamental weights $`\omega _i`$
$$\widehat{\mu }=(\lambda \underset{i,j}{}m^{(j)}(𝐂_r)_{ji}\omega _i;k;N\underset{i}{}m^{(i)}).$$
(3.12)
As a reminder, $`m^{(i)}`$ is the number of variables of color $`(i)`$, which corresponds to the number of roots $`\alpha _i`$ subtracted from the highest weight $`\lambda =l\omega _p`$. Note that the exponent of $`q`$ is not just the total degree of $`F(z)`$, but is shifted by $`_im^{(i)}`$ (compare with (3.3)). The reason is that the currents $`f_{\alpha _i}(z)`$ are defined by eq. (3.2), in which the power of $`z`$ is shifted by the scaling dimension.
Counting the dimensions of these function spaces is an intricate but tractable problem. We will explain the structure of this character in the appendix A, and refer to for details and the proof. The result is the function-space character
$$\mathrm{ch}_{l\omega _p;k}(q,𝐱)=x_p^l\underset{\begin{array}{c}i=1,\mathrm{},r\\ a=1,\mathrm{},k\\ m_a^{(i)}0\end{array}}{}\left(\underset{i=1}{\overset{r}{}}x_i^{_j(𝐂_r)_{ji}m^{(j)}}\right)\frac{q^{\frac{1}{2}m_a^{(i)}(𝐂_r)_{i,j}𝐀_{a,b}m_b^{(j)}𝐀_{a,l}m_a^{(p)}}}{_{i=1}^r_{a=1}^k(q)_{m_a^{(i)}}},$$
(3.13)
where it is understood that repeated indices are summed over. Here various symbols need to be defined: i) the matrix $`𝐀`$ has the entries $`𝐀_{a,b}=\mathrm{min}(a,b)`$; ii) for any integer $`m>0`$, we define $`(q)_m=_{i=1}^m(1q^i)`$ and $`(q)_0=1`$; iii) The sum over integers $`m_a^{(i)}`$ is to be understood as a sum over partitions of $`m^{(i)}`$, where $`m_a^{(i)}`$ denotes the number of rows of length $`a`$ in partition $`(i)`$. That is $`_{a=1}^kam_a^{(i)}=m^{(i)}`$, see figure 1 for an example.
The space of functions $`_{l\omega _p;k}`$ is dual to the principal subspace $`W_{l\omega _p;k}`$. It follows that the character (3.13) is the character of the principal subspace
$$\mathrm{ch}W_{l\omega _p;k}(q,𝐱)=\mathrm{ch}_{l\omega _p;k}(q,𝐱).$$
(3.14)
In the next section, we will explain how we can use the affine Weyl translations to obtain characters for the full representation from the characters of the principal subspaces.
We would like to note that it is straightforward to generalize the function space $`_{l\omega _p;k}`$, by changing the constraint (3.10) to $`z_1^{(i)}=\mathrm{}=z_{l_i+1}^{(i)}=0`$ for $`i=1,\mathrm{},r`$ and taking the product over all $`p`$ in the denominator of (3.6). We will denote this function space by $`_{𝐥;k}`$, where $`𝐥=(l_1,\mathrm{},l_r)^T`$. The character of this function space is
$$\mathrm{ch}_{𝐥;k}(q,𝐱)=\underset{i=1}{\overset{r}{}}x_i^{l_i}\underset{\begin{array}{c}i=1,\mathrm{},r\\ a=1,\mathrm{},k\\ m_a^{(i)}0\end{array}}{}\left(\underset{i=1}{\overset{r}{}}x_i^{_j(𝐂_r)_{ji}m^{(j)}}\right)\frac{q^{\frac{1}{2}m_a^{(i)}(𝐂_r)_{i,j}𝐀_{a,b}m_b^{(j)}𝐀_{a,l_i}m_a^{(i)}}}{_{i=1}^r_{a=1}^k(q)_{m_a^{(i)}}}.$$
(3.15)
We would like to stress that this character not the character of the principal subspace $`W_{\lambda ;k}`$, where $`\lambda =_il_i\omega _i`$. However, it can be interpreted as a character of a somewhat larger space, as we will see below.
## 4. Affine Weyl translations
In this section, we will exploit the symmetry of the representation spaces under the affine Weyl group. Elements of the affine Weyl group map weights of a highest weight representation to other weights in the same representation, in such a way that the weight-space multiplicities are preserved.
Elements of the affine Weyl group can be thought of as a product of a finite Weyl reflection and an affine Weyl translation. We will focus on the abelian subgroup generated by the affine Weyl translations, because they will enable us to obtain the character of the full integrable representation from the character of the principal subspace.
The affine Weyl translation $`T_{\alpha _i}^{N_i}`$ acts on a weight $`\widehat{\lambda }=(\lambda ;k;d)`$, where $`\lambda =_il_i\omega _i`$, by ‘translating’ it to $`\lambda +N_i\alpha _i`$ (no summation implied) and shifting the value of $`d`$ (see, for instance, , equation (6.5.2))
$$T_{\alpha _i}^{N_i}(\lambda ;k;d)=(\lambda +kN_i\alpha _i;k;dN_il_iN_i^2k).$$
(4.1)
More generally, we have
$$\underset{i=1}{\overset{r}{}}T_{\alpha _i}^{N_i}(\lambda ;k;d)=(\lambda +k\underset{i=1}{\overset{r}{}}N_i\alpha _i;k;d\underset{i=1}{\overset{r}{}}N_il_i\frac{k}{2}𝐍^T𝐂_r𝐍),$$
(4.2)
where $`𝐍=(N_1,\mathrm{},N_r)^T`$, and $`𝐂_r`$ is the Cartan matrix of $`su(r+1)`$.
We can use these affine Weyl translations to obtain the characters of the full representation. To illustrate how this works, we will use an explicit example, namely the vacuum representation of $`\widehat{su}(2)_2`$. The top part of the weight diagram of this representation is given in figure 2.
From this figure, we see that we can obtain the whole representation by acting with an arbitrary combination of the operators $`e[i]`$ and $`f[i]`$, with $`i<0`$ on the highest weight state $`|\widehat{\lambda }|v_0`$ (we dropped the subscript $`\alpha _1`$). Under the affine Weyl translation $`T^1`$, this weight is mapped to the weight $`|v_1`$. We can also generate the whole representation by acting with operators $`e[i]`$ and $`f[j]`$ on this state. In this case, we need to act with an arbitrary combination of operators taken from the sets $`\{e[i]|i<2\}`$ and $`\{f[j]|j<2\}`$. In general, the whole representation can be obtained by acting on the state $`|v_N=T^N|v_0`$ with the operators $`\{e[i]|i<2N\}`$ and $`\{f[j]|j<2N\}`$. We can now take the limit $`N\mathrm{}`$, and obtain that, we can generate the whole representation by acting on $`|v_{\mathrm{}}`$ with only the step down operators $`\{f[j]|j\}`$, see for the proof of this statement.
Let us come back to the principal subspace. As a reminder, the principal subspace $`W_{\widehat{\lambda }}`$ is the space generated by acting with the operators $`f[j]`$, with $`j<0`$ on $`|v_0`$ (again, we are focusing on the vacuum representation of $`\widehat{su}(2)_2`$). We can define a sequence of subspaces, by acting with the affine Weyl translation, i.e. $`W^{(N)}=T^NW_{\widehat{\lambda }}`$. The subspace $`W^{(N)}`$ is obtained by acting with the operators $`\{f[j]|j<2N\}`$ on the state $`|v_N`$. So, in the limit $`N\mathrm{}`$, we find that $`W^{(\mathrm{})}`$ is obtained by acting with $`\{f[j]|j\}`$, on the state $`|v_{\mathrm{}}`$. Comparing this subspace with the description of the full representation of the previous paragraph, we find that they are in fact the same.
Using this result, we can obtain the characters of the integrable representation by acting with the affine Weyl translation $`T^N`$ on the character of the principal subspace, and taking the limit $`N\mathrm{}`$.
In our paper , we showed that the effect of acting with the affine Weyl translation and taking the limit of $`N\mathrm{}`$ only alters the character of the principal subspaces of $`\widehat{su}(2)_k`$ in the following two ways. First of all, the summation over the variable $`m_k^{(1)}0`$ is extended to the negative integers. Secondly, the factor $`\frac{1}{(q)_{m_k^{(1)}}}`$ is replaced by $`\frac{1}{(q)_{\mathrm{}}}`$. In physical terms (see ), this corresponds to ‘filling the Bose sea’, and considering the excitations on top of a ‘large droplet’. In , we showed that this procedure also works in the case $`\widehat{su}(r+1)_k`$. Using this result, we obtain the characters for rectangular highest-weight representations
$$\mathrm{ch}H_{l\omega _p;k}(q,𝐱)=\frac{1}{(q)_{\mathrm{}}^r}x_p^l\underset{\begin{array}{c}m_k^{(i)}\\ m_{a<k}^{(i)}_0\\ i=1,\mathrm{},r\end{array}}{}\left(\underset{i=1}{\overset{r}{}}x_i^{_j(𝐂_r)_{ji}m^{(j)}}\right)\frac{q^{\frac{1}{2}m_a^{(i)}(𝐂_r)_{i,j}𝐀_{a,b}m_b^{(j)}𝐀_{a,l}m_a^{(p)}}}{_{i=1}^r_{a=1}^{k1}(q)_{m_a^{(i)}}}.$$
(4.3)
The character (4.3) can be written in a form which allows all $`l_i`$ to be non-zero, namely, in the same way as the character $`\mathrm{ch}_{𝐥;k}`$, equation (3.15)
$$\mathrm{ch}_{𝐥;k}^{\mathrm{}}(q,𝐱)\stackrel{\mathrm{def}}{=}\frac{1}{(q)_{\mathrm{}}^r}\left(\underset{i=1}{\overset{r}{}}x_i^{l_i}\right)\underset{\begin{array}{c}m_k^{(i)}\\ m_{a<k}^{(i)}_0\\ i=1,\mathrm{},r\end{array}}{}\left(\underset{i=1}{\overset{r}{}}x_i^{_j(𝐂_r)_{ji}m^{(j)}}\right)\frac{q^{\frac{1}{2}m_a^{(i)}(𝐂_r)_{i,j}𝐀_{a,b}m_b^{(j)}𝐀_{a,l_i}m_a^{(i)}}}{_{i=1}^r_{a=1}^{k1}(q)_{m_a^{(i)}}}.$$
(4.4)
We cannot however interpret this character as the character of the highest-weight representation $`H_{\widehat{\lambda }}`$, where the finite part of $`\widehat{\lambda }`$ is given by $`\lambda =_il_i\omega _i`$. This is because we did not take all the constraints into account for arbitrary representations. For example, for general highest weight, our restrictions on the space do not ensure that $`f_\alpha [1]|\widehat{\lambda }=0`$ whenever $`\alpha `$ is not a simple positive root. However, as we will show in the next section, the characters $`\mathrm{ch}_{𝐥;k}^{\mathrm{}}`$ are an essential ingredient of the characters $`\mathrm{ch}H_{\widehat{\lambda };k}`$ for arbitrary highest weights.
## 5. Arbitrary highest-weight representations
As mentioned above, we do not expect the function-space characters $`\mathrm{ch}_{𝐥;k}^{\mathrm{}}`$ defined in eq. (4.4) to be the characters of arbitrary highest-weight representations $`H_{\widehat{\lambda }}`$. This is because we have not yet imposed all the conditions on the function spaces. We have not imposed the remaining conditions, because the complexity of the resulting space makes counting its dimensions intractable.
We anticipate, therefore, that the function-space character (4.4) over-counts the dimensions of the weight spaces whenever $`\lambda `$ is not a rectangular representation. Nevertheless it is reasonable to conjecture some sort of relation between $`\mathrm{ch}H_{\widehat{\lambda }}`$ and $`\mathrm{ch}_{𝐥;k}^{\mathrm{}}`$.
To seek such a relation we used $`Mathematica^\text{®}`$ to compare the multiplicities given by $`\mathrm{ch}_{𝐥;k}^{\mathrm{}}`$ with those in the representation $`H_{\widehat{\lambda }}`$, obtained by using the affine version of Freudenthal’s recursion formula (see for instance, , page 578). We found that whenever more than one Dynkin index is non-zero, the character $`\mathrm{ch}_{𝐥;k}^{\mathrm{}}`$ does indeed over-count, and subtractions are necessary. It turns out that these subtractions can be written as a sum over $`\mathrm{ch}_{𝐥^{};k}^{\mathrm{}}`$, where the coefficients are polynomials in $`q^1`$.
To illustrate this we give a table expressing the characters of integrable level-$`4`$ representations of $`\widehat{su}(4)`$ in terms of the characters $`\mathrm{ch}_{𝐥^{};k}^{\mathrm{}}`$.
$`\mathrm{ch}H_{1,1,0;4}`$ $`=\mathrm{ch}_{1,1,0;4}^{\mathrm{}}{\displaystyle \frac{1}{q}}\mathrm{ch}_{0,0,1;4}^{\mathrm{}}`$
$`\mathrm{ch}H_{2,1,0;4}`$ $`=\mathrm{ch}_{2,1,0;4}^{\mathrm{}}{\displaystyle \frac{1}{q}}\mathrm{ch}_{1,0,1;4}^{\mathrm{}}+{\displaystyle \frac{1}{q^2}}\mathrm{ch}_{0,0,0;4}^{\mathrm{}}`$
$`\mathrm{ch}H_{1,2,0;4}`$ $`=\mathrm{ch}_{1,2,0;4}^{\mathrm{}}{\displaystyle \frac{1}{q}}\mathrm{ch}_{0,1,1;4}^{\mathrm{}}+{\displaystyle \frac{1}{q^2}}\mathrm{ch}_{1,0,0;4}^{\mathrm{}}`$
$`\mathrm{ch}H_{1,1,1;4}`$ $`=\mathrm{ch}_{1,1,1;4}^{\mathrm{}}\left({\displaystyle \frac{1}{q}}+{\displaystyle \frac{1}{q^2}}\right)\mathrm{ch}_{0,1,0;4}^{\mathrm{}}{\displaystyle \frac{1}{q}}\mathrm{ch}_{2,0,0;4}^{\mathrm{}}{\displaystyle \frac{1}{q}}\mathrm{ch}_{0,0,2;4}^{\mathrm{}}`$
$`\mathrm{ch}H_{3,1,0;4}`$ $`=\mathrm{ch}_{3,1,0;4}^{\mathrm{}}{\displaystyle \frac{1}{q}}\mathrm{ch}_{2,0,1;4}^{\mathrm{}}+{\displaystyle \frac{1}{q^2}}\mathrm{ch}_{1,0,0;4}^{\mathrm{}}`$
$`\mathrm{ch}H_{2,2,0;4}`$ $`=\mathrm{ch}_{2,2,0;4}^{\mathrm{}}{\displaystyle \frac{1}{q}}\mathrm{ch}_{1,1,1;4}^{\mathrm{}}+\left({\displaystyle \frac{1}{q^2}}+{\displaystyle \frac{1}{q^3}}\right)\mathrm{ch}_{0,1,0;4}^{\mathrm{}}+{\displaystyle \frac{1}{q^2}}\mathrm{ch}_{2,0,0;4}^{\mathrm{}}`$
$`\mathrm{ch}H_{2,1,1;4}`$ $`=\mathrm{ch}_{2,1,1;4}^{\mathrm{}}\left({\displaystyle \frac{1}{q}}+{\displaystyle \frac{1}{q^2}}\right)\mathrm{ch}_{1,1,0;4}^{\mathrm{}}{\displaystyle \frac{1}{q}}\mathrm{ch}_{3,0,0;4}^{\mathrm{}}{\displaystyle \frac{1}{q}}\mathrm{ch}_{1,0,2;4}^{\mathrm{}}+`$
$`\left({\displaystyle \frac{1}{q^2}}+{\displaystyle \frac{1}{q^3}}\right)\mathrm{ch}_{0,0,1;4}^{\mathrm{}}`$
$`\mathrm{ch}H_{1,3,0;4}`$ $`=\mathrm{ch}_{1,3,0;4}^{\mathrm{}}{\displaystyle \frac{1}{q}}\mathrm{ch}_{0,2,1;4}^{\mathrm{}}+{\displaystyle \frac{1}{q^2}}\mathrm{ch}_{1,1,0;4}^{\mathrm{}}{\displaystyle \frac{1}{q^3}}\mathrm{ch}_{0,0,1;4}^{\mathrm{}}`$
$`\mathrm{ch}H_{1,2,1;4}`$ $`=\mathrm{ch}_{1,2,1;4}^{\mathrm{}}\left({\displaystyle \frac{1}{q}}+{\displaystyle \frac{1}{q^2}}\right)\mathrm{ch}_{0,2,0;4}^{\mathrm{}}{\displaystyle \frac{1}{q}}\mathrm{ch}_{2,1,0;4}^{\mathrm{}}{\displaystyle \frac{1}{q}}\mathrm{ch}_{0,1,2;4}^{\mathrm{}}+{\displaystyle \frac{1}{q^2}}\mathrm{ch}_{1,0,1;4}^{\mathrm{}}`$
$`{\displaystyle \frac{1}{q^3}}\mathrm{ch}_{0,0,0;4}^{\mathrm{}}.`$ (5.1)
In addition, we found evidence for the relation
$$\mathrm{ch}H_{l_1,0,l_3;k}=\mathrm{ch}_{l_1,0,l_3;k}^{\mathrm{}}\frac{1}{q}\mathrm{ch}_{l_11,0,l_31;k}^{\mathrm{}},$$
(5.2)
for $`l_1,l_3>0`$ and $`k4`$.
It is interesting that the characters $`\mathrm{ch}H_{\lambda ;k}`$ can still be expressed in terms of $`\mathrm{ch}_{𝐥;k}^{\mathrm{}}`$, but the pattern of subtractions is at first sight obscure. However, inverting the relations so as to express the characters $`\mathrm{ch}_{𝐥;k}^{\mathrm{}}`$ in terms of the characters of the full representations gives a clue as to what is happening. We find
$`\mathrm{ch}_{1,1,0;4}^{\mathrm{}}`$ $`=\mathrm{ch}H_{1,1,0;4}+{\displaystyle \frac{1}{q}}\mathrm{ch}H_{0,0,1;4}`$
$`\mathrm{ch}_{1,1,1;4}^{\mathrm{}}`$ $`=\mathrm{ch}H_{1,1,1;4}+\left({\displaystyle \frac{1}{q}}+{\displaystyle \frac{1}{q^2}}\right)\mathrm{ch}H_{0,1,0;4}+{\displaystyle \frac{1}{q}}\mathrm{ch}H_{2,0,0;4}+{\displaystyle \frac{1}{q}}\mathrm{ch}H_{0,0,2;4}`$
$`\mathrm{ch}_{2,1,0;4}^{\mathrm{}}`$ $`=\mathrm{ch}H_{2,1,0;4}+{\displaystyle \frac{1}{q}}\mathrm{ch}H_{1,0,1;4}`$
$`\mathrm{ch}_{1,2,0;4}^{\mathrm{}}`$ $`=\mathrm{ch}H_{1,2,0;4}+{\displaystyle \frac{1}{q}}\mathrm{ch}H_{0,1,1;4}`$
$`\mathrm{ch}_{3,1,0;4}^{\mathrm{}}`$ $`=\mathrm{ch}H_{3,1,0;4}+{\displaystyle \frac{1}{q}}\mathrm{ch}H_{2,0,1;4}`$
$`\mathrm{ch}_{1,3,0;4}^{\mathrm{}}`$ $`=\mathrm{ch}H_{1,3,0;4}+{\displaystyle \frac{1}{q}}\mathrm{ch}H_{0,2,1;4}`$
$`\mathrm{ch}_{2,2,0;4}^{\mathrm{}}`$ $`=\mathrm{ch}H_{2,2,0;4}+{\displaystyle \frac{1}{q}}\mathrm{ch}H_{1,1,1;4}+{\displaystyle \frac{1}{q^2}}\mathrm{ch}H_{0,0,2;4}`$
$`\mathrm{ch}_{2,1,1;4}^{\mathrm{}}`$ $`=\mathrm{ch}H_{2,1,1;4}+\left({\displaystyle \frac{1}{q}}+{\displaystyle \frac{1}{q^2}}\right)\mathrm{ch}H_{1,1,0;4}+{\displaystyle \frac{1}{q}}\mathrm{ch}H_{3,0,0;4}+{\displaystyle \frac{1}{q}}\mathrm{ch}H_{1,0,2;4}+{\displaystyle \frac{1}{q^2}}\mathrm{ch}H_{0,0,1;4}`$
$`\mathrm{ch}_{1,2,1;4}^{\mathrm{}}`$ $`=\mathrm{ch}H_{1,2,1;4}+\left({\displaystyle \frac{1}{q}}+{\displaystyle \frac{1}{q^2}}\right)\mathrm{ch}H_{0,2,0;4}+{\displaystyle \frac{1}{q}}\mathrm{ch}H_{2,1,0;4}+{\displaystyle \frac{1}{q}}\mathrm{ch}H_{0,1,2;4}+{\displaystyle \frac{1}{q^2}}\mathrm{ch}H_{1,0,1;4}.`$ (5.3)
In addition, we have
$$\mathrm{ch}_{l_1,0,l_3;k}^{\mathrm{}}=\underset{j=0}{\overset{\mathrm{min}(l_1,l_3)}{}}\frac{1}{q^j}\mathrm{ch}H_{l_1j,0,l_3j;k}.$$
(5.4)
All the signs on the right hand side of the above decompositions are positive. This suggests that the $`\mathrm{ch}_{𝐥;k}^{\mathrm{}}`$ are indeed characters of $`\widehat{su}(r+1)`$ representations, but these representations are reducible. We need to understand which representations are occurring, and why the coefficients of their characters are polynomials in $`q^1`$. This we will do in the next section.
## 6. Kostka polynomials
In this section we will introduce Feigin and Loktev’s $`q`$-refinement of the Littlewood-Richardson coefficients that occur in the decomposition of tensor products of the finite-dimensional $`su(r+1)`$ representations .
Let $`\omega _i`$ denote the fundamental weights of the finite-dimensional $`su(r+1)`$ algebra. A finite-dimensional irreducible representation labeled by the Dynkin indices $`l_i_0`$ has highest weight $`\lambda =_{i=1}^rl_i\omega _i`$ and corresponds to a Young diagram with $`l_1`$ columns containing one box, $`l_2`$ columns with two boxes, and so on:
.
The Littlewood-Richardson rules provide a product that allows us to make the space of Young diagrams into an associative algebra. The multiplication operation in this algebra mirrors the multiplication and decomposition of the Schur functions corresponding to the Young diagrams, and, these being the characters of the associated $`su(r+1)`$ representations, reflect the decomposition of $`su(r+1)`$ tensor product representations into their irreducible components. For example, using the notation $`V_{l_1,l_2,l_3}`$ for the $`su(4)`$ representation with Dynkin indices $`l_1,l_2,l_3`$, the Young diagram manipulation
corresponds to the decomposition
$$V_{1,1,0}V_{2,0,0}=V_{3,1,0}V_{1,2,0}V_{2,0,1}V_{0,1,1}.$$
(6.1)
The coefficients are not always unity. For example
,
or
$$V_{0,0,1}V_{0,1,0}V_{1,0,0}=V_{1,1,1}V_{0,0,2}V_{2,0,0}2V_{0,1,0}.$$
(6.2)
This last expression is to be compared with the decomposition of our $`\widehat{su}(4)_{k=4}`$ function space
$$\mathrm{ch}_{1,1,1;4}^{\mathrm{}}=\mathrm{ch}H_{1,1,1;4}+\frac{1}{q}\mathrm{ch}H_{0,0,2;4}+\frac{1}{q}\mathrm{ch}H_{2,0,0;4}+\left(\frac{1}{q}+\frac{1}{q^2}\right)\mathrm{ch}H_{0,1,0;4}.$$
(6.3)
If we set $`q=1`$ in the coefficients in this decomposition, we recover the integers appearing in (6.2). All representations appearing in (5.3) are similarly accounted for: The $`H_{\widehat{\lambda }_i}`$ appearing in the decomposition of $`\mathrm{ch}_{l_1,l_2,l_3;4}^{\mathrm{}}`$ are precisely the $`\widehat{su}(4)`$ representations whose top grades are the $`V_{\lambda _i}`$ in the decomposition of the product $`V_{l_1,0,0}V_{0,l_2,0}V_{0,0,l_3}`$, and if $`q1`$ their $`q`$-coefficients reduce to the multiplicity of these $`V_{\lambda _i}`$.
We need to understand why the coefficients are $`q`$-dependent. Following Feigin and Loktev , we now show that these $`q`$-polynomial coefficients can be introduced even for finite-dimensional $`su(r+1)`$ representations. Let $`V_{\lambda _i}`$ denote a collection of irreducible finite-dimensional highest-weight representations of a Lie algebra $`𝔤`$, and consider the tensor product
$$V_{\lambda _1}\mathrm{}V_{\lambda _N}.$$
(6.4)
To each $`V_{\lambda _i}`$ we associate a distinct complex number $`\zeta _i`$ (to be thought of as “where” the representation is located). We then introduce the algebra $`𝔤[t]𝔤[t]`$ of $`𝔤`$-valued polynomials in $`t`$, with Lie bracket
$$[at^m,bt^n]=[a,b]t^{m+n}.$$
(6.5)
(Since only non-negative powers of $`t`$ are being allowed, there is no non-trivial central extension, so this is not the affine Lie algebra associated with $`𝔤`$.)
The representation $`V_{\lambda _i}`$ gives rise to the evaluation representation of $`𝔤[t]`$ by setting
$$(at^m)|u=\zeta _i^ma|u,\mathrm{for}|uV_{\lambda _i}.$$
(6.6)
We similarly define the action of $`𝔤[t]`$ on the tensor product space. If $`a𝔤`$, and $`|u_iV_{\lambda _i}`$, we set
$$(at^m)(|u_1\mathrm{}|u_N)=\underset{i}{}\zeta _i^m(|u_1\mathrm{}a|u_i\mathrm{}|u_N).$$
(6.7)
We will call this action the geometric, or fusion, co-product. For $`m=0`$ it reduces to the usual co-product action of $`𝔤`$ on the product space
$$a\mathrm{\Delta }(a)=𝕀\mathrm{}a\mathrm{}𝕀.$$
(6.8)
If $`|\lambda _i`$ are the highest weight (and hence cyclic) vectors for the representations $`V_{\lambda _i}`$, then, with the usual co-product action of $`𝔤`$, the vector $`|v|\lambda _1\mathrm{}|\lambda _N`$ is the highest weight vector for only one of the irreducible representations occurring in the decomposition of the tensor product of the $`V_{\lambda _i}`$. Under the action of $`𝔤[t]`$, however, and provided that the $`\zeta _i`$ are all distinct, $`|v`$ is a cyclic vector for the entire tensor product space. This is because the matrix $`\zeta _i^j`$ is invertible, and so any vector $`|u_1\mathrm{}a|u_j\mathrm{}|u_N`$ is obtainable as a linear combination of $`(at^m)(|u_1\mathrm{}|u_N)`$ for different $`m`$, and any vector $`|u_i`$ in each of the $`V_{\lambda _i}`$ is obtainable by the action of suitable $`a`$ on $`|\lambda _i`$.
The Lie algebra $`𝔤[t]`$ is graded by the degree of the polynomial in $`t`$. This grading extends to the universal enveloping algebra $`U(𝔤[t])`$ and gives rise to a filtration — a nested set of vector spaces —
$$F^0\mathrm{}F^iF^{i+1}\mathrm{},$$
(6.9)
where
$$F^i=U^i(𝔤[t])|v.$$
(6.10)
Here, $`U^i(𝔤[t])`$ denotes the elements of the universal enveloping algebra $`U(𝔤[t])`$, which have degree in $`t`$ less than or equal to $`i`$. The action of $`𝔤`$, considered as the zero-grade component of $`𝔤[t]`$, preserves this filtration and so has a well defined action on each of the components of the associated graded spaces $`\mathrm{gr}^i[F]F^i/F^{i1}`$. Each of the irreducible representations in the tensor product space will appear in one of these graded subspaces.
As an illustration, consider the simplest case $`𝔤=su(2)`$ where
$$[e,f]=h,[h,e]=+2e,[h,f]=2f.$$
(6.11)
In the usual physics spin-$`j`$ notation, we have the decomposition $`\frac{1}{2}\frac{1}{2}=10`$, or in our Dynkin index language
$$V_1V_1=V_2V_0.$$
(6.12)
The repeated action of the step down operator $`f`$ on the state
$$|v=||$$
(6.13)
(where $`|`$ and $`|`$ denote the states of weight $`1`$ and $`1`$ respectively) generates only the three states appearing in the spin-$`1`$ representation $`V_2`$, which therefore constitutes the space $`F^0`$. The action of $`ft`$, however, yields
$`(ft)|v`$ $`=`$ $`\zeta _1||+\zeta _2||,`$ (6.14)
$`=`$ $`\frac{1}{2}(\zeta _1\zeta _2)\left(||||\right)`$
$`+\frac{1}{2}(\zeta _1+\zeta _2)\left(||+||\right),`$
which, since $`||+||`$ lies in $`F_0`$, is equivalent in $`F^1/F^0`$ to
$$\frac{1}{2}(\zeta _1\zeta _2)\left(||||\right).$$
(6.15)
This last vector is the highest weight (indeed the only weight) in the spin-$`0`$ representation.
While the highest-weight vectors of each representation occurring in the tensor product must appear somewhere in this construction, their coefficients depend quite non-trivially on the $`\zeta _i`$, and some of these coefficients might vanish at non-generic points even for non-coincident $`\zeta _i`$. It is therefore not obvious that the grade at which a given representation first appears is independent of the choice of the $`\zeta _i`$. Feigin and Loktev conjecture that this is the case, and if this is true, the graded character
$$\mathrm{ch}_q(V_{\lambda _1}\mathrm{}V_{\lambda _N})=\underset{d}{}q^d\mathrm{ch}(\mathrm{gr}^d[F])$$
(6.16)
should be independent of the $`\zeta _i`$. In , we provide a proof of this conjecture by Feigin and Loktev in the spacial case where $`𝔤=su(r+1)`$ and when the $`\lambda _i`$ correspond to rectangular representations.
In the $`su(2)`$ example the $`\zeta _i`$ independence is manifest, and we have
$$\mathrm{ch}_q(V_1V_1)=\mathrm{ch}V_2+q\mathrm{ch}V_0.$$
(6.17)
Applying this construction to $`su(4)`$ we would find
$$\mathrm{ch}_q(V_{0,0,1}V_{0,1,0}V_{1,0,0})=\mathrm{ch}V_{1,1,1}+(q+q^2)\mathrm{ch}V_{0,1,0}+q\mathrm{ch}V_{2,0,0}+q\mathrm{ch}V_{0,0,2},$$
(6.18)
which looks like (6.3), but with $`q1/q`$ to reflect the fact that Feigin and Loktev define their grading in the opposite direction to the one which is natural for the affine algebra.
To obtain the characters for arbitrary representations of $`\widehat{su}(r+1)`$, we only need to consider the fusion product of $`r`$ rectangular representations $`V_1,\mathrm{},V_r`$, such that the highest weight $`\lambda _i`$ of $`V_i`$ is given by $`\lambda _i=n_i\omega _i`$, with $`n_i0`$. We will denote the $`q`$-polynomial coefficients in the decomposition of the character of this fusion product by $`𝒦_{𝐥,𝐧}(q)`$. Here, $`𝐥`$ is the vector whose entries are the Dynkin indices of the representations present in the fusion product. The entries of $`𝐧`$ are the $`n_i`$. In appendix B, we show how these $`q`$-polynomials can be calculated. The result is the fermionic formula (B.6). Details can be found in our paper , where we also showed that the $`q`$-polynomials are in fact generalized Kostka polynomials . Note that we will always assume that $`_in_ik`$, so we do not impose the level-restriction conditions in calculating the Kostka polynomials<sup>2</sup><sup>2</sup>2Level-restricted generalized Kostka polynomials are also important, but do not play a role in our calculations.. This restriction is always met in our case of obtaining characters of general $`\widehat{su}(r+1)`$ representations.
Thus, we have the following expression for the character of the fusion product of rectangular representations
$$\mathrm{ch}_q(V_1V_2\mathrm{}V_r)=\underset{𝐥}{}𝒦_{𝐥,𝐧}(q)\mathrm{ch}V_𝐥.$$
(6.19)
These $`q`$-polynomial $`𝒦_{𝐥,𝐧}(q)`$ coefficients that reduce for $`q=1`$ to the Littlewood-Richardson coefficients are generalized Kostka polynomials. As a result, we conclude that the sum in equation (6.19) is finite.
Classical Kostka polynomials are parametrized by two Young diagrams, $`\lambda `$ and $`\mu `$, such that $`|\lambda |=|\mu |`$. Generalized Kostka polynomials , which we consider in this paper, are parametrized by a Young diagram $`\lambda `$ and a set of rectangular diagrams $`\{\mu _i\}`$, such that the total number of boxes in the set $`\{\mu _i\}`$ is the same as the number of boxes in $`\lambda `$. Note that these Young diagrams are associated to $`𝔤𝔩_{r+1}`$. In our notation used above, the Young diagrams are associated to $`𝔰𝔩_{r+1}`$, which can be obtained from the $`𝔤𝔩_{r+1}`$ diagrams by ‘stripping off’ the columns of height $`r+1`$ (see for more details).
The classical Kostka polynomial is a special case of the generalized Kostka polynomial, where the set $`\{\mu _i\}`$ case is the set of single-row diagrams, each equal to a single row of the diagram $`\mu `$.
In both cases, the integer $`K_{\lambda ,\{\mu _i\}}(1)`$ is equal to the multiplicity of the $`𝔤𝔩_n`$-module $`V_\lambda `$ in the tensor product of representations $`V_{\mu _1}\mathrm{}V_{\mu _N}`$. The Kostka polynomial is therefore a refinement of tensor product multiplicities, or a grading.
## 7. Wess-Zumino-Witten conformal theory
Although the geometric, or fusion, co-product and the evaluation representation were initially defined for the finite-dimensional Lie algebra $`𝔤`$, they are motivated by the action of the affine Lie algebra $`\widehat{𝔤}`$ in Wess-Zumino-Witten (WZW) conformal field theory .
We will focus on the holomorphic half of a level-$`k`$ WZW model defined on the Riemann sphere. The symmetry algebra of this model is $`𝔤(𝜻)`$, consisting of $`𝔤`$-valued meromorphic functions with possible poles at the points $`\zeta _i`$ (where $`𝜻=(\zeta _1,\mathrm{},\zeta _N)`$).
Let $`\lambda `$ denote the highest weight of a finite-dimensional representation $`V_\lambda `$ of the finite-dimensional algebra $`𝔤`$, and suppose that the corresponding (in that its top graded component is $`V_\lambda `$) representation $`H_{\widehat{\lambda }}`$ of the level-$`k`$ affine Lie algebra $`\widehat{𝔤}`$ is integrable. Then, the Wess-Zumino primary field $`\phi _\lambda (\zeta )`$ acts on the vacuum at the origin $`\zeta =0`$ there to create the highest weight state $`|\widehat{\lambda }=\phi _\lambda (0)|0`$ of $`H_{\widehat{\lambda }}`$. If
$$a(z)=\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}a[n]z^{n1}$$
(7.1)
is the WZW current associated with the element $`a𝔤`$, and $`\gamma `$ is any contour surrounding the origin then
$$a[n]=\frac{1}{2\pi i}_\gamma z^na(z)𝑑z$$
(7.2)
has its usual action as an element of $`\widehat{𝔤}`$ on this representation.
More generally, for integrable $`\widehat{\lambda }_i`$ we create highest weight states $`|\widehat{\lambda }_i_{\zeta _i}`$ at the points $`\zeta _i`$ by acting on the vacuum at these points with $`\phi _{\lambda _i}(\zeta _i)`$ on the vacuum at $`\zeta _i`$. The current algebra acts on the tensor product $`H_{\widehat{\lambda }_1}\mathrm{}H_{\widehat{\lambda }_N}`$ of these spaces. We recall how this comes about. Let $`f(t)`$ denote a function meromorphic on the Riemann sphere, and having poles at no more than the $`\zeta _i`$. If $`\mathrm{\Gamma }`$ is a contour surrounding $`\zeta _1,\mathrm{},\zeta _N`$, then we insert $`\frac{1}{2\pi i}_\mathrm{\Gamma }f(t)a(t)𝑑t`$ into a suitable correlation function. The contour can be deformed to a sum of contours $`\gamma _i`$ each enclosing only one of the $`\zeta _i`$, and the $`a[n]`$ acting on the representation space at $`\zeta _i`$ are the local mode-expansion coefficients $`a(z)=_na[n](z\zeta _i)^{n1}`$, or
$$a[n]=\frac{1}{2\pi i}(t\zeta _i)^na(t)𝑑t.$$
(7.3)
Consequently, if
$$f(t)=\underset{n=p_i}{\overset{\mathrm{}}{}}f_n(\zeta _i)(t\zeta _i)^n$$
(7.4)
is the Laurent expansion of $`f`$ about $`\zeta _i`$ then $`af`$ acts on a state $`|v_{\zeta _i}H_{\lambda _i}`$ as
$$\underset{n=p_i}{\overset{\mathrm{}}{}}f_n(\zeta _i)a[n]|v_{\zeta _i},$$
(7.5)
and on $`H_{\widehat{\lambda }_1}\mathrm{}H_{\widehat{\lambda }_N}`$ as
$$\mathrm{\Delta }_{\zeta _1,\mathrm{},\zeta _N}(af)(|v_1_{\zeta _1}\mathrm{}|v_N_{\zeta _N})=\underset{i}{}|v_1_{\zeta _1}\mathrm{}\left(\underset{n=p_i}{\overset{\mathrm{}}{}}f_n(\zeta _i)a[n]|v_i_{\zeta _i}\right)\mathrm{}|v_N_{\zeta _N},$$
(7.6)
where $`\mathrm{\Delta }_{\zeta _1,\mathrm{},\zeta _N}`$ is the “geometric” co-product
$$\mathrm{\Delta }_{\zeta _1,\mathrm{},\zeta _N}(af)=\underset{i}{}𝕀\mathrm{}\left(\underset{n=p_i}{}f_n(\zeta _i)a[n]\right)\mathrm{}𝕀.$$
(7.7)
This co-product makes the tensor products of any number of level-$`k`$ representations into a level-$`k`$ representation. With the conventional co-product, the level would be $`Nk`$.
In the particular case that $`f(t)=t^n`$, the Laurent expansion about $`\zeta _i`$ is
$$t^n=\underset{m=0}{\overset{n}{}}\left(\genfrac{}{}{0pt}{}{n}{m}\right)\zeta _i^{nm}(t\zeta _i)^m.$$
(7.8)
If $`|v_{\zeta _i}`$ is a top-component state on which $`a[n]`$ with $`n>0`$ acts as zero, the action on $`|v_{\zeta _i}`$ is as $`\zeta _i^na[0]|v_{\zeta _i}`$. Since the zero-grade elements of $`\widehat{𝔤}`$ form an algebra isomorphic to $`𝔤`$, we recognize our evaluation co-product action described in equation (6.7).
When the contour $`\mathrm{\Gamma }`$ encloses all insertions of $`\phi _\lambda `$’s on the Riemann sphere, it can be contracted away, and the algebra acts trivially. This observation motivates the definition of the space $`(𝜻,\{\lambda \})`$ of conformal blocks. A conformal block $`\mathrm{\Psi }(𝜻;\{\lambda \})`$ is a mapping
$$\mathrm{\Psi }:H_{\widehat{\lambda }_1}\mathrm{}H_{\widehat{\lambda }_N},$$
(7.9)
invariant under the action of any $`af`$, $`a𝔤`$, with $`f`$ having poles at no more that the $`\zeta _i`$. It may be shown that such a mapping is uniquely specified by its evaluation on the top graded component $`V_{\lambda _1}\mathrm{}V_{\lambda _N}`$, and that the dimension $`\mathrm{dim}\left((𝜻,\{\lambda \})\right)`$ of the space of conformal blocks is the multiplicity of the identity operator in the fusion product of the $`\phi _{\lambda _i}`$.
This definition coincides with the physicist’s picture where
$$\mathrm{\Psi }(\zeta _1,\mathrm{},\zeta _N)\mathrm{\Psi }\left(|\mu _1_{\zeta _1}\mathrm{}|\mu _N_{\zeta _N}\right)=0|R\{\phi _{\mu _1}(\zeta _1)\mathrm{}\phi _{\mu _N}(\zeta _N)\}|0,$$
(7.10)
with $`\mu _i`$ weights in $`V_{\lambda _i}`$, is thought of as a radial-ordered correlation function of the primary fields $`\phi _{\mu _i}`$. This correlator is not uniquely defined because the chiral primary fields are not mutually local. It is, however, a solution of the Knizhnik-Zamolodchikov equations, whose solution space is precisely the space of conformal blocks as identified above. In the operator language this comes about because the $`0|`$ could be the dual of any of the copies of the highest-weight state of the vacuum representation that occur in the fusion-product.
We now recognize what it is that our function-space character is counting. It is the number of linear independent functions of the form
$`F(\{z_j^{(i)}\})`$ $`=`$ $`v|R\{f_{\alpha _1}(z_1^{(1)})\mathrm{}f_{\alpha _r}(z_{m^{(r)}}^{(r)})\phi _{\lambda _1}(\zeta _1)\mathrm{}\phi _{\lambda _r}(\zeta _r)\}|0,`$ (7.11)
$``$ $`{\displaystyle \mathrm{\Psi }\left(|v^{}_{\mathrm{}}|v_1_{\zeta _1}\mathrm{}|v_N_{\zeta _N}\right)}.`$
Here the $`\phi _{\lambda _i}(\zeta _i)`$ create rectangular representations with highest weight $`\lambda _i=l_i\omega _i`$ (no sum on $`i`$) by acting on the vacuum at the points $`\zeta _i`$. The state $`v|`$, analogous to the $`v|`$ of (3.1), is the dual of $`|v`$, which lies in a weight space in one of the integrable representations $`H_{\widehat{\lambda }}`$ occurring in the fusion product. The sum is over the $`|v_j`$ produced by the action of the $`f_\alpha `$ on the $`|\widehat{\lambda }_i`$. The state $`|v^{}_{\mathrm{}}`$ is in the dual representation $`H_{\widehat{\lambda }}^{}`$ that is located at the point $`\mathrm{}`$ on the Riemann sphere. The (left action) $`𝔤`$ weight of $`|v^{}`$ is minus that of $`|v`$. (The correspondence $`v||v^{}_{\mathrm{}}`$ is explained in .)
The fusion product of primary fields in the WZW model coincides with the finite-dimensional algebra Littlewood-Richardson rules only when $`k`$ is infinite. For finite $`k`$ we must use the level-restricted fusion rules of the Verlinde algebra. The level restrictions have no effect, however, when we build up an integrable representation of $`\widehat{𝔤}`$ by concatenating integrable rectangular representations, as we are doing in this paper. The representations appearing in our WZW fusion product therefore coincide with those appearing in the $`𝔤[t]`$ fusion product, which in turn are given by the Littlewood-Richardson rules of $`𝔤`$.
The fusion product itself is not a graded space. However, we can define a filtration on this space, in the same way as in the Feigin-Loktev fusion product. This filtration is inherited from the filtration of the algebra $`U(𝐧_{}(t^1))`$, where we assign a degree zero to the product of primary fields. The associated graded space has the character (4.4), which we can regard as the character of the fusion product. It is a sum over irreducible characters with a certain shift in the overall degree. We can regard this as the statement that the vector $`v|`$ in the matrix element belongs to a highest weight representation, with a highest weight on which $`d`$ acts by some negative integer instead of 0. (For an alternative point of view, where the space of highest weight vectors is interpreted as a quotient of the integrable representation – which is naturally graded – see .)
We would like to stress again that the characters of this space do not depend on the locations $`\zeta _i`$ of the rectangular representations!
## 8. Characters for arbitrary $`\widehat{su}(r+1)_k`$ representations
In this section, we will combine the results of the previous section and give explicit character formulæ for arbitrary (integrable) representations of $`\widehat{su}(r+1)_k`$. There is one more ingredient needed to to this, namely an explicit formula for the (generalized) Kostka polynomials, which form an essential ingredient in the character formulæ, as explained before. In this section, we will give the character formulæ in terms of the generalized Kostka polynomials. How to obtain an explicit formula will be described in appendix B (we refer to for the details). The resulting formula for the generalized Kostka polynomials is stated in equation (B.6).
We can write the decomposition of the character $`_{𝐧;k}^{\mathrm{}}(q,𝐱)`$ in the following way.
$$\mathrm{ch}_{𝐧;k}^{\mathrm{}}(q,𝐱)=\underset{𝐥}{}𝒦_{𝐥,𝐧}\left(\frac{1}{q}\right)\mathrm{ch}H_{𝐥;k}(q,𝐱).$$
(8.1)
We need to make a few remarks about the sum in the decomposition (8.1). First of all, we would like to note that this sum is finite. To show this, we introduce the notion of the threshold level, which is the lowest level for which an highest weight $`\lambda `$ corresponds to a highest weight representation. We will denote this threshold level by $`k(𝐥)`$. For $`\widehat{su}(r+1)`$, it is simply given by $`k(𝐥)=_{i=1}^rl_i`$. The only $`𝐥`$ for which the Kostka polynomial $`𝒦_{𝐥,𝐧}(q)`$ is non-zero, are the $`𝐥`$ such that $`k(𝐥)k(𝐧)`$. The only non-zero $`𝒦_{𝐥,𝐧}(q)`$ with $`k(𝐥)=k(𝐧)`$ is when $`𝐥=𝐧`$, in which case $`𝒦_{𝐧,𝐧}(q)=1`$, see equation (B.7). Using these results, we can view the Kostka polynomials $`𝒦_{𝐧,𝐥}(q)`$ as elements of a square matrix $`𝐊`$, which is upper triangular and $`1`$’s on the diagonal. Hence, this matrix is invertible, if we order the ‘highest weights’ $`𝐥`$ according to increasing threshold level. Note that we did not specify an ordering for weights with the same threshold level, but any ordering of those weights gives an upper triangular matrix.
Now that we established that the Kostka matrix $`𝐊`$ is invertible, we can invert the relation (8.1) to obtain explicit character formulæ for arbitrary (integrable) highest weight representations of $`\widehat{su}(r+1)_k`$
$$\mathrm{ch}H_{𝐥;k}(q,𝐱)=\underset{𝐧}{}(𝐊^1)_{𝐧,𝐥}\left(\frac{1}{q}\right)\mathrm{ch}_{𝐧;k}^{\mathrm{}}(q,𝐱).$$
(8.2)
Again, the (finite) sum is over highest weights $`𝐧`$ with threshold level $`k(𝐧)k(𝐥)`$.
## 9. Conclusion
We have shown how the introduction of the fusion product representation allows us to generalize the Feigin-Stoyanovsky strategy for computing characters of affine $`\widehat{su}(r+1)`$ to all integrable representations. We found that for non-rectangular representations, the character is not of fermionic type, but can be written as a linear combination of fermionic characters. The coefficients are polynomials in $`1/q`$ and are related to the generalized Kostka polynomials.
It is rather straightforward to generalize the Feigin-Stoyanovsky construction to arbitrary affine Lie algebras. In that case, even characters of rectangular representations are not always of fermionic type. It turns out that the character of those representations can be written in terms of characters of the so called ‘Kirillov-Reshetikhin’ modules . Details will be provided in a forthcoming publication.
## 10. Acknowledgments
The work of EA was sponsored in part by NSF grants DMR-04-42537 and DMR-01-32990, the work of MS by NSF grant DMR-01-32990 and the work of RK by NSF grant DMS-05-00759.
## Appendix A The structure of the character of the principal subspace
In this appendix, we will explain the structure of the character of the principal subspace of rectangular representations of $`\widehat{su}(r+1)`$, eq. (3.13). To this end, we will explain the zeros of the function $`f(z)`$ which appears in (3.6).
The exponent of $`q`$ in the character is the total degree of the polynomials $`f(\{z_j^{(i)}\})`$, which ensure that the vanishing conditions hold, combined with the explicit poles in $`F(\{z_j^{(i)}\})`$. The strategy presented in and , which we will outline here briefly, is to find the minimal amount of zeros necessary to enforce the vanishing conditions.
Let us first focus on the case we have $`2k`$ variables. We need to find a polynomial which does not vanish when $`k`$ variables are set to the same location, but does vanish when $`k+1`$ variables are set to the same location. Of course, the condition (3.8) allows for functions which vanish when less than $`k+1`$ variables are set to the same location, but we will deal with those later. A symmetric polynomial in the variables $`z_1,z_2,\mathrm{},z_{2k}`$ with the lowest possible degree which satisfies this property is ($`𝒮`$ denotes symmetrization over all variables)
$$𝒮\left[(z_1z_{k+1})(z_2z_{k+1})(z_2z_{k+2})\mathrm{}(z_kz_{2k1})(z_kz_{2k})(z_1z_{2k})\right].$$
(A.1)
We can display these zeros nicely in a graphical way in terms of a Young diagram. The boxes on the top row correspond to the variables $`z_1,\mathrm{},z_k`$ while the boxes on the second row correspond to $`z_{k+1},\mathrm{},z_{2k}`$. A line connecting the boxes of $`z_i`$ and $`z_j`$ corresponds to the factor $`(z_iz_j)`$:
,
where ‘periodic boundary conditions’ are assumed, to obtain the zero $`(z_1z_{2k})`$. It easily follows that if we pick $`k`$ variables, and set all of them to the same value, there will always be a term in the symmetrization of eq. (A.1) which is not zero. However, setting $`k+1`$ variables to the same location gives a zero for each term in the symmetrization.
We will now discuss the general case, i.e. we allow for functions which vanish when fewer than $`k+1`$ variables are set equal, and we also allow for a general number of variables.
To do this, we need to label the polynomials by $`r`$ partitions (or Young diagrams), one for each type of variable. These are partitions of the integers $`m^{(i)}`$, with the property that the width of the rows is maximally $`k`$, because the polynomials should vanish when $`k+1`$ variables are set to the same value.
The number of rows of width $`a`$ of partition $`(i)`$ is given by $`m_a^{(i)}`$, i.e. the partitions have the form $`(k^{m_k^{(i)}},(k1)^{m_{k1}^{(i)}},\mathrm{},2^{m_2^{(i)}},1^{m_1^{(i)}})`$. It follows that we have the relation $`m^{(i)}=_{a=1}^kam_a^{(i)}`$. See figure 1 for an example where $`m^{(1)}=13`$.
To each box of the partition $`(i)`$, we associate a variable $`z_j^{(i)}`$. Let $`z_1,\mathrm{},z_a`$ be the variables associated to a row of length $`a`$, of partition $`(i)`$ and $`\overline{z}_1,\mathrm{},\overline{z}_a^{}`$ to a row of length $`a^{}`$, also from partition $`(i)`$, such that $`a^{}a`$.
We need to assign the following zeros corresponding to these variables
$$(z_1\overline{z}_1)(z_2\overline{z}_1)(z_2\overline{z}_2)(z_3\overline{z}_2)\mathrm{}(z_a^{}\overline{z}_a^{})(z_{a^{}+1}\overline{z}_a^{}).$$
We will identify $`z_{a+1}=z_1`$. Pictorially, we can represent these zeros as
.
Again, each box corresponds to a variable, and a line connecting two boxes indicates a zero when the two corresponding variables have the same value. We obtain all the zeros needed to satisfy the integrability condition by multiplying all the zeros associated to each pair of rows belonging to the same partition
$$\underset{i=1}{\overset{r}{}}\underset{\begin{array}{c}\mathrm{pairs}\mathrm{of}\mathrm{rows}\\ \mathrm{of}\mathrm{partition}(i)\\ \mathrm{of}\mathrm{length}a,a^{}\\ aa^{}\end{array}}{}\underset{j=1}{\overset{a^{}}{}}(z_j^{(i)}\overline{z}_j^{(i)})(z_{j+1}^{(i)}\overline{z}_j^{(i)}).$$
(A.2)
To clarify this, we will look at one particular variable (which will correspond to the black box in the diagram below), and show with which other variables it has a zero. These other variables are given by the gray boxes:
$$\text{}.$$
(A.3)
To obtain all the zeros, we need to take the product over all ‘black boxes’, but without double counting the zeros. Finally, we need to take the product over all partitions $`i=1,\mathrm{},r`$. It is not to hard to convince oneself that these zeros do imply the integrability condition.
We will now focus on the zeros we have to include in the polynomials $`f(\{(z_j^{(i)}\})`$ to ensure that the Serre relations are taken into account properly. Let $`z_j^{(i)}`$ be a variable associated to a row of length $`a`$ of the partition of $`m^{(i)}`$. This variable has zeros with variables $`z_j^{}^{(i+1)}`$ which belong to a row of the partition of $`m^{(i+1)}`$ of length $`a^{}`$. In particular, these zeros are $`_{j^{}j}(z_j^{(i)}z_j^{}^{(i+1)})`$. We obtain all the zeros associated to the pair of partitions $`(i)`$ and $`(i+1)`$ by multiplying all the zeros associated to each variable of partition $`i`$.
$$\underset{i=1}{\overset{r1}{}}\underset{\begin{array}{c}\mathrm{rows}\mathrm{of}\\ \mathrm{partition}(i)\end{array}}{}\underset{\begin{array}{c}\mathrm{rows}\mathrm{of}\\ \mathrm{partition}(i+1)\end{array}}{}\underset{\begin{array}{c}j,j^{}\\ j^{}j\end{array}}{}(z_j^{(i)}z_j^{}^{(i+1)}).$$
(A.4)
Pictorially, these zeros are given by
$$\text{}.$$
(A.5)
Again, we need to take the product over all ‘black boxes’ and over $`i=1,\mathrm{},r1`$. Again, it is easy to check that these zeros, combined with the zeros which enforce integrability, give rise to the Serre conditions.
Finally, we need to include the zeros to make sure the highest weight condition is satisfied. For a row of length $`a>l`$ of the partition of $`m^{(p)}`$, we need to include the zeros
$$\underset{\begin{array}{c}\mathrm{rows}\mathrm{of}\\ \mathrm{partition}p\end{array}}{}\underset{j>l}{}z_j^{(p)},$$
(A.6)
or, pictorially, the zeros are given by the gray boxes
$$\text{}.$$
(A.7)
Counting the total degree of the zeros implied by (A.3), (A.5) and (A.7), combined with the degree corresponding to the poles in (3.6) and the extra term $`_{i=1}^rm^{(i)}`$ in the definition (3.11) precisely give the exponent of $`q`$ in the character formula (3.13).
Apart from the zeros discussed above, $`f(\{z_j^{(i)}\})`$ can contain additional symmetric polynomials. The degree of these symmetric polynomials is taken into account by the factors $`\frac{1}{(q)_m}`$ in (3.13), see for the details. The factor $`\frac{1}{(q)_m}`$ is the generating function for the symmetric polynomials in $`m`$ variables. The coefficient of $`q^d`$ in the expansion of $`\frac{1}{(q)_m}`$ is the number of symmetric polynomials of degree $`d`$ in $`m`$ variables.
## Appendix B Obtaining the Kostka polynomials
In the main text, we explained how the characters of the fusion product of rectangular representations can be decomposed into characters of arbitrary highest-weight representations. In this appendix we will outline a strategy to obtain an explicit expression for the ‘expansion coefficients’, which turn out to be generalized Kostka polynomials. We refer to our paper for the details of the proof.
We need to consider matrix elements of the form
$$G_{\lambda ,𝝁}(𝜻)=u_\lambda |f_{\alpha _1}(z_1^{(1)})\mathrm{}f_{\alpha _r}(z_{m^{(r)}}^{(r)})|v_1(\zeta _1)\mathrm{}v_r(\zeta _r),$$
(B.1)
where $`u_\lambda |`$ is dual to the state $`|u_\lambda ^{}_{\mathrm{}}`$, as explained in the discussion following equation (7.11). The $`v_i(\zeta _i)`$ are the highest weights of the rectangular representations $`H_{n_i\omega _i}(\zeta _i)`$. Thus, the rectangular representations inserted at the points $`\zeta _i`$ are given by $`n_i\omega _i`$. We define $`𝐧=(n_1,\mathrm{},n_r)^T`$. In addition, we define $`𝝁=_in_i\omega _i`$ and the coefficients $`l_i`$ are determined by $`\lambda =_il_i\omega _i`$.
To obtain the a formula for the generalized Kostka polynomials, we should only consider the bra’s $`u_\lambda |`$ at infinity, which are the highest weights of the ‘bottom component’ of the representation located at infinity. Note that the representation which is located at infinity is turned ‘upside down’, because we chose $`\frac{1}{z}`$ as the local variable at infinity.
We need to find the conditions such that the matrix elements (B.1) are non-zero. First of all, we need that $`\lambda =𝝁_im^{(i)}\alpha _i`$, which translates into
$$𝐦=C_r^1(𝐧𝐥).$$
(B.2)
As usual, we also need to impose the highest-weight conditions $`f_{\alpha _i}^{n_i+1}[0]|v_i=0`$ and the Serre relations. Finally, we need to incorporate a new ingredient, which is a degree restriction on the functions in the function space. These degree restrictions can be derived by letting the $`f`$’s in the matrix elements (B.1) act to the left on $`u_\lambda `$. Now, $`f_\alpha [n]`$ acts trivially at $`\mathrm{}`$ if $`n0`$, which translates a degree restriction on the functions.
The space of functions which we need to consider is spanned by rational functions in the variables $`z_j^{(i)}`$, which are symmetric under the exchange $`z_j^{(i)}z_j^{}^{(i)}`$. There can be poles of order one at the positions $`z_j^{(i)}=z_j^{}^{(i+1)}`$ and at $`z_j^{(i)}=\zeta _i`$. That is
$$G(𝐳)=\frac{g(𝐳)}{_{i=1}^r_j(\zeta _iz_j^{(i)})_{i=1}^{r1}_{j,j^{}}(z_j^{(i)}z_j^{}^{(i+1)})}.$$
(B.3)
The polynomial $`g(𝐳)`$ (which is symmetric under the exchange of variables of the same color) vanishes when any of the following conditions holds
$`z_1^{(i)}`$ $`=z_2^{(i)}=z_1^{(i+1)},z_1^{(i)}=z_1^{(i+1)}=z_2^{(i+1)},i=1,\mathrm{},r1`$ (B.4)
$`z_1^{(i)}`$ $`=z_2^{(i)}=\mathrm{}=z_{l_i+1}^{(i)}=\zeta _i,i.`$ (B.5)
Note that we do not impose the integrability conditions, because $`k`$ will always be large enough, i.e. $`k_in^{(i)}`$. Finally, we need to impose the degree restriction on the function $`G(𝐳)`$. For each variable, we have that $`\mathrm{deg}_{z_j^{(i)}}2`$, which follows form the fact that we use $`\frac{1}{z}`$ as a local variable at infinity. In , we showed that this leads to the following polynomials
$$𝒦_{𝐥;𝐧}(q)=\underset{\begin{array}{c}m_a^{(i)}0\\ i=1,\mathrm{},r\\ 𝐦=(𝐂_r^1)(𝐧𝐥)\end{array}}{}q^{\frac{1}{2}m_a^{(i)}(𝐂_r)_{i,j}𝐀_{a,b}m_b^{(j)}}\underset{i,a}{}\left[\begin{array}{c}𝐀_{a,n^{(i)}}(𝐂_r)_{i,j}𝐀_{a,b}m_b^{(j)}+m_a^{(i)}\\ m_a^{(i)}\end{array}\right]_q,$$
(B.6)
where $`\left[\right]_q`$ denotes the $`q`$-binomial, which is defined to be $`\left[\begin{array}{c}n+m\\ m\end{array}\right]_q=\frac{(q)_{n+m}}{(q)_n(q)_m}`$ for $`m,n_0`$ and zero otherwise. The coefficient of $`q^d`$ in the expansion of $`\left[\begin{array}{c}n+m\\ m\end{array}\right]_q`$ is the number of symmetric polynomials of total degree $`d`$, in $`m`$ variables, where the degree of each variable is maximally $`n`$. These $`q`$-binomials arise because of the degree restrictions on the polynomials $`G(𝐳)`$. Note that the Kostka polynomials do not depend on the positions $`\zeta _i`$, which follows from the fact that the degree of the rational functions (B.3) does not depend on these positions.
In , we showed that the functions $`𝒦_{𝐥;𝐧}(q)`$ of (B.6) are polynomials in $`q`$, and are related in a simple way to the generalized Kostka polynomials of . In particular, by setting $`q=1`$, we obtain the Littlewood-Richardson coefficients.
Let us make a few remarks about the structure of the polynomials $`𝒦_{𝐥;𝐧}(q)`$. Let $`k(𝐥)=_il_i`$, which is the lowest level at which $`\lambda =_il_i\omega _i`$ corresponds to an integrable representation (i.e. $`k(𝐥)`$ is the threshold level). We then have the following results
$$\begin{array}{cc}\hfill 𝒦_{𝐥;𝐥}(q)& =1\hfill \\ \hfill 𝒦_{𝐥;𝐧}(q)& =0\mathrm{if}\{\begin{array}{cc}& k(𝐥)>k(𝐧)\hfill \\ & k(𝐥)=k(𝐧)\mathrm{and}𝐥𝐧\hfill \\ & _iil_i_iin^{(i)}\mathrm{mod}r+1.\hfill \end{array}\hfill \end{array}$$
(B.7)
These results are obtained by making use of the constraint $`𝐦=(𝐂_r)^1(𝐧𝐥)`$ and the fact that all the summation variables $`𝐦=(m_1,\mathrm{},m_r)`$ have to be non-negative integers in the sum in equation (B.6).
We can view the polynomials $`𝒦_{𝐥;𝐧}`$ as the entries of a (square) matrix $`𝐊`$, with entries $`(𝐊)_{𝐥;𝐧}(q)=𝒦_{𝐥;𝐧}(q)`$. The relations (B.7) imply that there is an ordering such that the matrix $`𝐊`$ is upper triangular with $`1`$’s on the diagonal. Thus, the matrix $`𝐊`$ is invertible.
Before we move on and give the character formulæ for arbitrary highest-weight representations, we will first give an example of the ‘Kostka matrix’ related to $`\widehat{su}(4)`$. We use the following ordering of the entries $`𝐥`$
$$(0,0,0);(1,0,1),(0,2,0);(2,1,0),(0,1,2);(4,0,0),(2,0,2),(1,2,1),(0,4,0),(0,0,4)$$
With this ordering, we obtain the following Kostka matrix
$$𝐊(q)=\left(\begin{array}{cccccccccc}1& q& 0& 0& 0& 0& q^2& 0& 0& 0\\ & & & & & & & & & \\ 0& 1& 0& q& q& 0& q& q^2& 0& 0\\ 0& 0& 1& 0& 0& 0& 0& q+q^2& 0& 0\\ & & & & & & & & & \\ 0& 0& 0& 1& 0& 0& 0& q& 0& 0\\ 0& 0& 0& 0& 1& 0& 0& q& 0& 0\\ & & & & & & & & & \\ 0& 0& 0& 0& 0& 1& 0& 0& 0& 0\\ 0& 0& 0& 0& 0& 0& 1& 0& 0& 0\\ 0& 0& 0& 0& 0& 0& 0& 1& 0& 0\\ 0& 0& 0& 0& 0& 0& 0& 0& 1& 0\\ 0& 0& 0& 0& 0& 0& 0& 0& 0& 1\end{array}\right)$$
(B.8)
The inverse is
$$𝐊^1(q)=\left(\begin{array}{cccccccccc}1& q& 0& q^2& q^2& 0& 0& q^3& 0& 0\\ & & & & & & & & & \\ 0& 1& 0& q& q& 0& q& q^2& 0& 0\\ 0& 0& 1& 0& 0& 0& 0& qq^2& 0& 0\\ & & & & & & & & & \\ 0& 0& 0& 1& 0& 0& 0& q& 0& 0\\ 0& 0& 0& 0& 1& 0& 0& q& 0& 0\\ & & & & & & & & & \\ 0& 0& 0& 0& 0& 1& 0& 0& 0& 0\\ 0& 0& 0& 0& 0& 0& 1& 0& 0& 0\\ 0& 0& 0& 0& 0& 0& 0& 1& 0& 0\\ 0& 0& 0& 0& 0& 0& 0& 0& 1& 0\\ 0& 0& 0& 0& 0& 0& 0& 0& 0& 1\end{array}\right).$$
(B.9)
Using the results of the sections 6 and 7, we can now write down explicit character formulæ for arbitrary integrable highest-weight representations of $`\widehat{su}(r+1)`$, in terms of the characters $`\mathrm{ch}_{𝐧;k}^{\mathrm{}}(q,𝐱)`$ and the inverse matrix $`𝐊^1\left(\frac{1}{q}\right)`$.
First of all, we have the result that we can decompose the characters $`\mathrm{ch}_{𝐧;k}^{\mathrm{}}(q,𝐱)`$ in terms of the characters of arbitrary integrable highest-weight representations $`\mathrm{ch}H_{𝐥;k}(q,𝐱)`$ as follows
$$\mathrm{ch}_{𝐧;k}^{\mathrm{}}(q,𝐱)=\underset{\begin{array}{c}𝐥\\ k(𝐥)k(𝐧)\end{array}}{}𝒦_{𝐥;𝐧}\left(\frac{1}{q}\right)\mathrm{ch}H_{𝐥;k}(q,𝐱).$$
(B.10)
We can invert this relation, to obtain explicit character formulæ for arbitrary integrable highest-weight representations of $`\widehat{su}(r+1)`$
$$\mathrm{ch}H_{𝐥;k}(q,𝐱)=\underset{\begin{array}{c}𝐧\\ k(𝐧)k(𝐥)\end{array}}{}(𝐊^1)_{𝐧;𝐥}\left(\frac{1}{q}\right)\mathrm{ch}_{𝐧;k}^{\mathrm{}}(q,𝐱).$$
(B.11)
We would like to note that we have a similar formula for the character of the principal subspace of general highest-weight representations
$$\mathrm{ch}W_{𝐥;k}(q,𝐱)=\underset{\begin{array}{c}𝐧\\ k(𝐧)k(𝐥)\end{array}}{}(𝐊^1)_{𝐧;𝐥}\left(\frac{1}{q}\right)\mathrm{ch}_{𝐧;k}(q,𝐱),$$
(B.12)
where the character $`_{𝐥;k}(q,𝐱)`$ is given by
$$\mathrm{ch}_{𝐥;k}(q,𝐱)\stackrel{\mathrm{def}}{=}(\underset{i=1}{\overset{r}{}}x_i^{l_i})\underset{\begin{array}{c}m_a^{(i)}_0\\ a=1,\mathrm{},k\\ i=1,\mathrm{},r\end{array}}{}\left(\underset{i=1}{\overset{r}{}}x_i^{_j(𝐂_r)_{ji}m^{(j)}}\right)\frac{q^{\frac{1}{2}m_a^{(i)}(𝐂_r)_{i,i^{}}𝐀_{a,a^{}}m_a^{}^{(i^{})}𝐀_{a,l_i}m_a^{(i)}}}{_{i=1}^r_{a=1}^k(q)_{m_a^{(i)}}}.$$
(B.13)
This character can be viewed as the ‘untranslated’ version of the character (4.4), because applying the affine Weyl translation (as explained in section 4), results in (4.4).
EA: Department of Physics, University of Illinois, 1110 W. Green St., Urbana, IL 61801.
ardonne@uiuc.edu
RK: Department of Mathematics, University of Illinois, 1409 W. Green Street, Urbana, IL 61801. rinat@uiuc.edu
MS: Department of Physics, University of Illinois, 1110 W. Green St., Urbana, IL 61801.
m-stone5@uiuc.edu
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# Spin dynamics and violation of the fluctuation dissipation theorem in a non-equilibrium ohmic spin boson model.
## Abstract
We present results for the dynamics of an impurity spin coupled to a magnetic field and to two ohmic baths which are out-of equilibrium due to the application of a bias voltage. Both the non-equilibrium steady state and the rate constants describing the approach to steady state are found to depend sensitively on the relative strengths of a magnetic field and a voltage dependent decoherence rate. Computation of physical quantities including the frequency dependent ratio of response to correlation functions and the probabilities of the two spin states allows the extraction of voltage dependent effective temperatures. The temperatures extracted from different quantities differ from one another in magnitude and their dependence on parameters, and in general are non-monotonic.
A fundamental question in quantum condensed matter physics is understanding properties of non-equilibrium many body systems, some examples being the Kondo effect in quantum dots neqex1 , ultra-cold gases with rapidly tunable interactions neqex2 , and strongly driven optical lattices neqexp3 . While there are a variety of non-perturbative techniques in place to study equilibrium systems, these methods cannot be extended to non-equilibrium systems in a straightforward way neqrev . The experimental accessability neqex1 ; neqex2 ; neqex3 of the non-equilibrium regime of strongly correlated quantum many body systems gives rise to the need for developing the formalism further.
Many body systems driven out of equilibrium are known neqKth ; mamprl to acquire a steady state that may be quite different in character from their ground state properties, with the details of the steady state depending on the nature of the correlations, as well as on the way in which the system is driven out of equilibrium. One may characterize a system in steady state by the response function $`\chi (tt^{})`$ describing changes induced by weak external probes, and by the correlation function $`S(tt^{})`$ describing the probabilities of observing various configurations of the system. An important and still incompletely understood issue is the manner in which $`\chi `$ and $`S`$ characterizing a non-equilibrium system differ from those describing an equilibrium one. In particular there has been considerable interest fdt1 ; fdt2 in the possibility of establishing a generalized fluctuation-dissipation theorem relating $`\chi (\omega )`$ to $`S(\omega )`$ and thereby characterizing the departures from equilibrium in terms of an effective temperature. Several systems have been identified fdt2 where such a generalized fluctuation dissipation theorem is found to hold, with the extracted temperature often sensitive to the observables being studied.
In this paper we study the dynamics of the out of equilibrium ohmic spin-boson model. This model describes a two state system (which we represent in spin notation) with level splitting $`2B`$ and tunneling rate $`\mathrm{\Delta }`$, coupled via a coupling $`J_z`$ to a spin-less resonant level (creation operator $`d^{}`$), which is itself connected to two leads ($`L`$ and $`R`$) that may be at different chemical potentials. The Hamiltonian is
$`H`$ $`=`$ $`S_zB+\mathrm{\Delta }S_x+J_zS_zd^{}d+H_{bath}`$ (1)
$`H_{bath}`$ $`=`$ $`{\displaystyle \underset{k,\alpha =L,R}{}}ϵ_kc_{k\alpha }^{}c_{k\alpha }+{\displaystyle \underset{k,\alpha =L,R}{}}(t_{k\alpha }c_{k\alpha }^{}d+h.c.)`$
We assume the leads are infinite reservoirs characterized by the correlators $`c_{k\alpha }^{}c_{q\beta }=\delta _{kq}\delta _{\alpha \beta }\left(e^{\beta (ϵ_k\mu _\alpha )}+1\right)^1`$, and a non-equilibrium state occurs when $`\mu _L\mu _R=V0`$. Crucial parameters of the model are the left and right channel phase shifts $`\delta _{L,R}`$ defined by $`\mathrm{tan}\delta _L=\frac{a_LJ_z}{1isgn(V)a_RJ_z}`$, $`\mathrm{tan}\delta _R=\frac{a_RJ_z}{1+isgn(V)a_LJ_z}`$ with $`a_{L,R}=\frac{\mathrm{\Gamma }_{L,R}}{\left(\mathrm{\Gamma }_L+\mathrm{\Gamma }_R\right)^2}`$, with $`\mathrm{\Gamma }_{L,R}=\pi \rho t_{L,R}^2`$ and $`\rho =\frac{dk}{dϵ_k}`$. We will study properties of $`H`$ at $`T=0`$ but out of equilibrium ($`V0`$) working to leading nontrivial order in $`\mathrm{\Delta }`$ but to all order in $`J_z`$. The conditions under which perturbation theory is justified will be discussed below. We will find the time scales characterizing the approach to steady state, the response and correlation functions and the generalized fluctuation dissipation ratio.
In order to study the spin-dynamics, the appropriate starting point is the density matrix for the full Hamiltonian
$`{\displaystyle \frac{d\rho (t)}{dt}}=i[H,\rho (t)]`$ (3)
from which the density matrix $`\rho _S`$ for the local spin is obtained from taking a trace over the electronic degrees of freedom.
$`\widehat{\rho }_S=Tr_{el}\rho `$ (4)
We adopt a spin language, writing
$$\widehat{\rho }_S=\frac{1}{2}\left(1+S_z\widehat{\sigma }_z\right)$$
(5)
and study $`S_z`$. When $`\mathrm{\Delta }=0`$, the Hamiltonian is exactly solvable both in and out of equilibrium. For non-zero $`\mathrm{\Delta }`$ one may expand Eq. 3 perturbatively in $`\mathrm{\Delta }`$. The key object in the analysis is the time-evolution operator separating two spin-flip processes yuvand , $`K_\pm (t)=Tr_{el}\left[e^{iH(\mathrm{\Delta }=0,S_z=\pm 1/2)t}e^{iH(\mathrm{\Delta }=0,S_z=(1)/2)t}\right]=e^{\pm iBt}e^{C_\pm (t)}`$ where $`C_+(t)=C(t)=\left(C_{}(t)\right)^{}`$ computed from the linked cluster theorem has the following form at zero temperature mamprl ,
$`C(t)=C^{}(t)+iC^{\prime \prime }(t)`$
$`C^{}(t)=\left({\displaystyle \frac{\delta _{eq}}{\pi }}\right)^2\mathrm{ln}\left(\xi t\right)+\varphi ^{}(Vt)`$ (6)
$`C^{\prime \prime }={\displaystyle \frac{\pi }{2}}\left({\displaystyle \frac{\delta _{eq}}{\pi }}\right)^2sgn(t)+\varphi ^{\prime \prime }(Vt)`$ (7)
Here $`\xi `$ is a short time cut-off, $`\delta _{eq}=\delta _L+\delta _R=\mathrm{arctan}\left(\frac{J_z}{\mathrm{\Gamma }_L+\mathrm{\Gamma }_R}\right)`$ is the equilibrium phase shift, and $`\varphi (Vt)`$ is a function describing the crossover from the short-time ($`Vt1`$) equilibrium behaviour characterized by xrayeq $`C(t)=\left(\frac{\delta _{eq}}{\pi }\right)^2\left[\mathrm{ln}\left(\xi t\right)+i\frac{\pi }{2}sgn(t)\right]`$, to the long time ($`Vt1`$) non-equilibrium behaviour characterized by Ng $`C(t)=\left(\frac{\delta _L^2+\delta _R^2}{\pi ^2}\right)\left[\mathrm{ln}\left(\xi t\right)+i\frac{\pi }{2}sgn(t)\right]+\mathrm{\Gamma }_{neq}t`$ with $`\mathrm{\Gamma }_{neq}=V\frac{|\delta _L^{\prime \prime }\delta _R^{\prime \prime }|}{2\pi }`$. Correct treatment for the intermediate and short time behaviour of $`\varphi `$ is essential for obtaining correct results for $`\chi (\omega ),S(\omega )`$. A general analytic expression for $`\varphi `$ does not exist, here we use perturbation theory to third order in $`J_z`$ to obtain mamprl ,
$`\varphi (Vt)`$ $`={\displaystyle \frac{|\delta _L^{^{\prime \prime }}\delta _R^{^{\prime \prime }}|}{2\pi }}Vt\left[\left({\displaystyle \frac{2}{\pi }}\right)\left(Si(Vt){\displaystyle \frac{1\mathrm{cos}(Vt)}{Vt}}\right)\right]`$
$`{\displaystyle \frac{2Re\delta _L\delta _R}{\pi ^2}}\left[\gamma _eCi(Vt)+\mathrm{ln}(Vt)\right]`$
$`{\displaystyle \frac{2iIm\delta _L\delta _R}{\pi ^2}}\left[{\displaystyle \frac{2}{\pi }}{\displaystyle _0^1}𝑑u\mathrm{sin}(uVt){\displaystyle \frac{\left[\left(1u\right)\mathrm{ln}(1u)+u\mathrm{ln}u\right]}{u^2}}\right]`$
Now let us return to the evaluation of various spin observables. In terms of the symmetric and anti-symmetric time-evolution operators $`K_{s,a}(t)=Re\left[K_+(t)\pm K_{}(t)\right]`$, the equation of motion for the variable $`S_z`$ parameterized in Eq. 5 to leading non-trivial order in $`\mathrm{\Delta }`$ is given by
$$\frac{dS_z}{dt}=_0^t𝑑t^{}\left[K_s(t,t^{})S_z(t^{})+K_a(t,t^{})\right]$$
(9)
The Laplace transform of the two scattering rates $`K_{s,a}`$ defined by $`\stackrel{~}{K}_{s,a}(\lambda )=_0^{\mathrm{}}𝑑tK_{s/a}(t)e^{\lambda t}`$ have the form,
$`\stackrel{~}{K}_s(\lambda )=\mathrm{\Delta }^2{\displaystyle _0^{\mathrm{}}}𝑑te^{\lambda t}e^{C^{}(t)}\mathrm{cos}Bt\mathrm{cos}C^{\prime \prime }(t)`$ (10)
$`\stackrel{~}{K}_a(\lambda )=\mathrm{\Delta }^2{\displaystyle _0^{\mathrm{}}}𝑑te^{\lambda t}e^{C^{}(t)}\mathrm{sin}Bt\mathrm{sin}C^{\prime \prime }(t)`$ (11)
The above equations capture the effect of two sources of decoherence on spin dynamics, one is a Korringa type decoherence existing even in equilibrium, while the second arising mathematically from $`C^{}(t)`$ is due to a non-zero voltage and is intrinsically non-equilibrium.
The solution to Eq. 9 can be written as wprl85
$$S_z(t)=\frac{1}{2\pi i}_{ci\mathrm{}}^{c+i\mathrm{}}\frac{d\lambda }{\lambda }e^{\lambda t}\frac{\lambda S_z(0)\stackrel{~}{K}_a(\lambda )}{\lambda +\stackrel{~}{K}_s(\lambda )}$$
(12)
Eq. 12 allows straightforward analysis of the long-time behaviour. As $`t\mathrm{}`$, the integral is dominated by the pole at $`\lambda =0`$ and the residue gives $`S_z^{\mathrm{}}=S_z\left(t\mathrm{}\right)=\frac{\stackrel{~}{K}_a(0)}{\stackrel{~}{K}_s(0)}`$. Consideration of $`S_z(t)S_z^{\mathrm{}}`$ then yields the rate at which the system approaches steady state. In the small $`\mathrm{\Delta }`$ limit, and if at least one of $`B,\mathrm{\Gamma }_{neq}`$ is not too small, the result is exponential relaxation with rate $`\mathrm{\Gamma }_{rel}=K_s(0)`$. The value of $`\mathrm{\Gamma }_{rel}`$ depends crucially on whether the dominant time scales in Eq. 10 are large or small relative to $`V^1`$. If both $`B`$ and $`\mathrm{\Gamma }_{neq}`$ are less than $`V`$, one finds (for compactness we write for the symmetric case $`t_L=t_R`$)
$`\mathrm{\Gamma }_{rel}=\stackrel{~}{K}_s(0)`$
$`={\displaystyle \frac{\pi }{2}}{\displaystyle \frac{\mathrm{\Delta }^2}{\xi }}{\displaystyle \frac{1}{\mathrm{\Gamma }\left(\alpha \right)}}{\displaystyle \frac{\mathrm{sin}\left[\frac{\pi \alpha }{2}+(1\alpha )\mathrm{arctan}\frac{\mathrm{\Gamma }_{neq}}{B}\right]}{\mathrm{sin}\frac{\pi \alpha }{2}}}\left({\displaystyle \frac{\sqrt{B^2+\mathrm{\Gamma }_{neq}^2}}{\xi }}\right)^{\alpha 1}`$ (13)
where the non-equilibrium exponent $`\alpha =\left(\frac{\delta _L}{\pi }\right)^2+\left(\frac{\delta _R}{\pi }\right)^2`$. The most interesting situation is the relatively small phase shift limit, in which $`\mathrm{\Gamma }_{neq}V`$ and Eq. 13 holds even when $`B\mathrm{\Gamma }_{neq}`$, provided $`BV`$. For $`BV`$, one should set $`\mathrm{\Gamma }_{neq}=0`$ in Eq. 13 and replace the non-equilibrium exponent $`\alpha `$ by the equilibrium exponent $`\alpha _{eq}=\left(\frac{\delta _{eq}}{\pi }\right)^2`$, yielding the familiar $`T=0`$ Korringa relaxation.The above results of an exponential relaxation to steady state is obtained from neglecting the $`\lambda `$ dependence of $`K_s(\lambda )`$, which is justified when
$$\frac{\mathrm{\Gamma }_{rel}}{\sqrt{\mathrm{\Gamma }_{neq}^2+B^2}}\frac{\mathrm{\Delta }^2}{\xi ^2}\left(\frac{\sqrt{B^2+\mathrm{\Gamma }_{neq}^2}}{\xi }\right)^{\alpha 2}1$$
(14)
and therefore holds in the perturbative in $`\mathrm{\Delta }`$ regime provided the voltage or the magnetic field is not too small. Analysis similar to the equilibrium case shows that Eq. 14 is also the condition for validity of perturbation theory in $`\mathrm{\Delta }`$.
At steady state and in the small $`\mathrm{\Delta }`$ limit, the density matrix for the full system is an incoherent superposition of spin up times the electronic state appropriate to spin-up and spin-down times the electronic state appropriate to spin-down, which may be expressed follows,
$$\rho =\rho _S\rho _{el}=\left(\begin{array}{cc}\rho _{}\rho _{el}^{}& 0\\ 0& \rho _{}\rho _{el}^{}\end{array}\right)$$
(15)
where $`\rho _{el}^{}`$ is the density matrix corresponding to Hamiltonian $`H`$ with $`S_z=1/2`$ and $`\mathrm{\Delta }=0`$. Likewise $`\rho _{el}^{}`$ is the density matrix for $`S_z=1/2`$ and $`\mathrm{\Delta }=0`$, while $`\rho _,=\frac{1}{2}\left(1\pm S_z\right)`$. We now calculate the response and correlation functions appropriate to this state and also study the fluctuation-dissipation relation between them.
The correlation function we study is,
$`S_{xx}(t_1,t_2)`$ $`=`$ $`i\{S_x(t_1),S_x(t_2)\}_+`$ (16)
$`=`$ $`Tr\left[\rho \{S_x(t_1),S_x(t_2)\}_+\right]`$
and the corresponding spin response function is,
$`\chi _{xx}(t_1,t_2)`$ $`=`$ $`i\theta (t_1t_2)[S_x(t_1),S_x(t_2)]`$ (17)
$`=`$ $`i\theta (t_1t_2)Tr\left[\rho [S_x(t_1),S_x(t_2)]\right]`$
where the density matrix $`\rho `$ is evaluated to leading order in the spin flip term $`\mathrm{\Delta }`$ and hence given by Eq. 15. The Fourier transform of the imaginary part of the response and correlation functions are,
$`\chi _{xx}^{\prime \prime }(\omega )`$ $`=`$ $`\rho _{}\left[I(B+\omega )I(B\omega )\right]`$ (18)
$``$ $`\rho _{}\left[I(B\omega )I(B+\omega )\right]`$
$`iS_{xx}(\omega )`$ $`=`$ $`\rho _{}\left[I(B+\omega )+I(B\omega )\right]`$ (19)
$`+`$ $`\rho _{}\left[I(B\omega )+I(B+\omega )\right]`$
where
$$I(B)=Re\left[_0^{\mathrm{}}𝑑te^{iBt}e^{C(t)}\right]$$
(20)
and $`\frac{\rho _{}}{\rho _{}}=\frac{I(B)}{I(B)}`$. We also consider the fluctuation-dissipation ratio (also mentioned in hooleyspin )
$$h(\omega )=\frac{\chi ^{\prime \prime }(\omega )}{iS_{xx}(\omega )}$$
(21)
In equilibrium and at $`T>0`$ (when $`\varphi (0)=1`$ and $`C(t)=\left(\frac{\delta _{eq}}{\pi }\right)\mathrm{ln}\left(\frac{\xi }{\pi T}\mathrm{sinh}\pi tT\right))`$, $`h(\omega )=\mathrm{tanh}\frac{\omega }{2T}`$.
Out of equilibrium (and at $`T=0`$), $`h(\omega )`$ has the form shown in Fig. 1 which differs from the equilibrium solution $`sgn(\omega )`$. The calculated $`h(\omega )`$ is not a $`\mathrm{tanh}`$ function (compare dashed line), and therefore a generalized fluctuation dissipation theorem encompassing all frequencies does not exist. However we may define a frequency dependent effective temperature via $`T_{eff}(\omega )=\frac{\omega }{2tanh^1h}`$. This function is plotted in the inset of Fig. 1 and is seen to have a strong $`\omega `$ dependence and is indeed not monotonic. For $`\omega <V/2`$, $`T_{eff}`$ is seen to be of the order of $`V`$ and to depend weakly on $`\omega `$. For $`\omega >V/2`$, $`T_{eff}`$ drops sharply and at high $`\omega `$ approaches the equilibrium value (here,$`T=0`$).The results presented in Fig. 1 show that no unique definition of ”non-equilibrium effective temperature” exists, the value obtained depends on the quantity examined. Two obvious definitions are (i) from the $`\omega 0`$ limit of $`h(\omega )`$ glass , (ii) from the population ratio $`\rho _{}/\rho _{}`$. The effective temperature from definition (i) is obtained by expanding Eq. 18 and Eq. 19 for small $`\omega `$,
$$\frac{1}{T_{eff}^h}=\frac{h(\omega )}{\omega }|_{\omega =0}=\underset{\sigma =\pm }{}\frac{\mathrm{ln}I(x)}{x}|_{x=\sigma B}$$
(22)
while definition (ii) for the effective temperature leads to the expression
$$\frac{1}{T_{eff}^\rho }=\frac{1}{B}\mathrm{ln}\frac{\rho _{}}{\rho _{}}$$
(23)
Fig. 2 shows the dependence of these two measures of effective temperature on magnetic field. We see that the two curves differ in magnitude and in dependence on parameters; the variation in general is non-monotonic. The inset shows that the magnitude and field variation of the effective temperature (plots are for $`T_{eff}^h`$) also depends on coupling constant.
The non-monotonic behaviour as a function of $`B/V`$ may be understood as follows. For $`B\mathrm{\Gamma }_{neq}`$ and for $`\frac{\delta _L^{\prime \prime }\delta _R^{\prime \prime }}{2\pi }1`$ so that $`\mathrm{\Gamma }_{neq}V`$, the integrand in Eq. 20 is dominated by $`t1/\mathrm{\Gamma }_{neq}1/V`$. In this regime, Eq. 13 applies; from this one sees that the decoherence rate for the spin increases with magnetic field causing the initial downturn in Fig. 2. This behaviour may also be understood as originating from the opening up of an additional scattering channel on application of a magnetic field that corresponds to the relaxation of the higher energy spin state by creating particle-hole excitations in the leads. We make this more precise by studying Eq. 20 perturbatively in tunneling amplitude ($`t_{L,R}`$) to find
$$\frac{1}{T_{eff}^h}=\frac{a_L^2+a_R^2+a_La_R}{\left(a_L^2+a_R^2\right)|B|+a_La_R\left(|B|+V\right)}\frac{1}{|B|V}$$
(24)
For the special case of symmetric couplings $`a_L=a_R`$ (which corresponds to the case in Fig. 2), and for $`BV`$, one finds $`\frac{1}{2T_{eff}^h}\frac{2}{V}\left(1\frac{2B}{V}\right)`$ which captures the initial downturn in the plot for $`T_{eff}`$.
For $`B\mathrm{\Gamma }_{neq}`$ on the other hand the integrals in Eq. 20 is dominated by $`t1/B1/V`$. In this regime $`T_{eff}`$ approaches the equilibrium value of $`T_{eff}0`$ and therefore results in an upturn in Fig. 2. The $`BV`$ behaviour of the integral $`I(B)`$ was found mamprl to be $`I(BV)\left(\frac{V}{B}\right)^{\frac{B}{V}}`$, the physical significance of which may be understood as follows. $`I(B)`$ represents the population of the high energy spin state where the energy for populating it is supplied by the voltage source. The ratio $`n=\frac{B}{V}`$ represents the optimal number of bath electrons that can be transmitted across the voltage source in order to excite the higher energy spin state, while $`I(B)=\left(\frac{V}{B}\right)^n`$ is simply the probability for doing so. Plugging this expression for $`I(B)`$ into Eq. 22, the effective temperature is found to approach zero as $`T_{eff}^h\frac{V}{\mathrm{ln}\frac{B}{V}}`$ in the regime $`B/V1`$.
Let us now turn to the discussion of the spin response function itself. The imaginary part of the response function is plotted for different coupling strengths and ratio of $`B/V`$ in Fig. 3. The line-shape (main panel: Fig. 3) in addition to having the familiar asymmetric form of an x-ray response function Ng , now has a finite weight at $`|\omega |<|B|`$, which is forbidden at zero temperatures in equilibrium. The coupling constant (main panel) and voltage (inset panel) dependence of the broadening is illustrated in Fig 3. $`\chi ^{\prime \prime }(\omega )`$ is linear in $`\omega `$ for small $`\omega `$, with a slope that is inversely related to the long-time relaxation rate of the density matrix $`\mathrm{\Gamma }_{rel}=T_{eff}^h`$, while at large frequencies $`\omega B`$, $`\chi ^{\prime \prime }(\omega )\frac{1}{\omega ^{1(\frac{\delta _{eq}}{\pi })^2}}`$. These two different frequency regimes appear as a change in slope of the plots in the inset of Fig. 3.
In conclusion, we have studied the non-equilibrium ohmic spin-boson model including a non-vanishing level splitting and orthogonality effects exactly. Previous work hooleyspin studied the zero level splitting limit, treating the orthogonality effects perturbatively. Our results agree with previous results in the appropriate limit, but provide significant new information including the non-monotonic effective temperature and the line-shape at non-vanishing level splitting. The calculated spin dynamics reveal that the non-equilibrium regime can be quite complex because of the interplay between various voltage and magnetic field dependent relaxation mechanisms. While departures from equilibrium are qualitatively similar to a non-zero temperature (e.g. permitting sub-threshold absorption of Fig. 3), the analogy cannot be pushed too far. The ”fluctuation-dissipation” ratio is not a hyperbolic tangent and indeed is not characterized by a unique effective temperature (c.f. Fig. 1), and the low-frequency effective temperature is itself a non-trivial function of the control parameters (c.f. Fig. 2), and is different depending on the quantity used to evaluate it. In equilibrium, the spin boson and Kondo models are related by the simple mapping $`\frac{\delta _{eq}}{\pi }\sqrt{2}\left(1\frac{\delta _{eq}}{\pi }\right)`$. Our finding that non-equilibrium effects enter into different parameters in different ways suggests that the mapping will not be so simple in the non-equilibrium case. A direction for future research is to extend the analysis in this paper to arbitrary number of spin flip processes, and to perform an Anderson-Yuval-Hamann type renormalization group treatment for the out of equilibrium spin boson and Kondo models mmp . This work was supported by NSF DMR-0431350.
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# The fingerprint of binary intermediate mass black holes in globular clusters: supra-thermal stars and angular momentum alignment
## 1 Introduction
A number of different observations suggest that large black holes (BHs) may exist in nature, with masses between $`20M_{}10^4M_{}.`$ Heavier than the stellar-mass BHs born in core-collapse supernovae ($`3M_{}20M_{}`$; Orosz 2002), these intermediate mass black holes (IMBHs) are expected to form from the direct collapse of very massive stars, or in dense stellar systems through complex dynamical processes. Filling the gap between stellar-mass black holes (BHs) and super-massive black holes (SMBHs), they are of crucial importance in establishing the potential physical link between these two classes.
IMBHs may plausibly have formed in the early universe as remnants of the first generation of metal free stars (Abel, Bryan & Norman 2002; Heger et al. 2003; Schneider et al. 2002). If this is true, IMBHs can participate the cosmic assembly of galaxies and be incorporated in larger and larger units, becoming seeds for the formation of SMBHs (Madau & Rees 2001; Volonteri, Haardt & Madau 2003). IMBH may still be forming in young dense star clusters vulnerable to unstable mass segregation, via collisions of massive stars (Portegies Zwart 2004; Portegies Zwart et al. 2004; Portegies Zwart & McMillan 2002; Gürkan et al. 2004; Freitag et al. 2005a, 2005b; see van der Marel 2004 for a review). Unambiguous detections of individual IMBHs do not exist yet, but there are observational hints in favor of their existence, from studies of ultra-luminous X-ray sources in nearby star-forming galaxies (Fabbiano 2004; Mushotzky 2004; Miller & Hamilton 2002).
It has also been long suspected that globular clusters (GCs) may hide IMBHs in their cores (Shapiro 1977; Shapiro & Marchant 1978; Marchant & Shapiro 1979, 1980; Duncan & Shapiro 1982; see Shapiro 1985 for a review), despite the fact that their formation root is poorly known: runaway mergers among the most massive stars, at the time of cluster formation, can lead to an IMBH, similarly to what has been conjectured to occur in young dense star clusters (Portegies Zwart et al. 2004; Fregeau et al. 2004). An alternative pathway is based on the idea that IMBHs form later in the evolution of GCs through mergers of stellar-mass BHs. Segregating by dynamical friction in the core, these BHs are captured in binaries that form through dynamical encounters with stars and BHs. The picture is that hardening by subsequent interactions with BHs lead them to merge, emitting gravitational waves (Miller & Hamilton 2002). BHs more massive than stellar-mass BHs are thus created. However, the interactions that produce hardening also provide recoil, causing the ejection of BHs (Kulkarni, Hut & McMillan 1993; Sigurdsson & Hernquist 1993; Portegies Zwart & McMillan 2000), but the heaviest black holes may remain, if their mass exceeds a still uncertain value between $`50M_{}`$ and $`300M_{}`$ (Miller & Hamilton 2002; Colpi, Mapelli & Possenti 2003; Gultekin, Miller & Hamilton 2004).
At present, optical observations of GCs hint in favor of the existence of IMBHs. In particular Gebhardt, Rich & Ho (2002) suggest the presence of a $`2_{0.8}^{+1.4}\times 10^4M_{}`$ IMBH, to explain the kinematics and the surface brightness profile of the globular cluster G1 in M31. Gerssen et al. follow the same method to indicate the possible presence of a $`1.7_{1.7}^{+2.7}\times 10^3M_{}`$ IMBH in the galactic globular cluster M15 (Gerssen et al. 2002, 2003). These indications are still controversial, as Baumgardt et al (2003a, 2003b) showed that both the measurements of M15 and those of G1 can be explained with a central concentration of compact objects (neutron stars and white dwarfs) instead of the presence of a massive black hole. However, a more recent paper by Gebhardt, Rich & Ho (2005) further supports the hypothesis of an IMBH in G1. In addition, there has been some claim of the existence of rotation in the core of a few GCs (Gebhardt et al. 1995; Gebhardt et a. 1997; Gebhardt et al. 2000) suggesting the presence of a steady source of angular momentum in their cores.
Radio observations can also help in discovering concentrations of under-luminous matter in the core of GCs, thanks to the presence of millisecond pulsars that can probe the underlying gravitational field of the cluster. Currently, there are more than 100 known millisecond pulsars in GCs (Possenti 2003; Ransom et al. 2005; Camilo & Rasio 2005) and some show large and negative period derivatives. This is an important peculiarity, because millisecond pulsars spin down intrinsically due to magnetic braking. Negative period derivatives come from the variable Doppler shift caused by the acceleration of the pulsar in the gravitational field of the cluster itself (Phinney 1993). When these negative period derivatives can be ascribed solely to the gravitational field of the cluster, these pulsars place a lower limit on the mass enclosed inside their projected distance from the globular cluster center.
Two pulsars in M15, at 1” from the cluster center were found (Phinney 1993) with negative values of their period derivatives, consistent with the mass distribution implied by the stellar kinematics. The recent discovery by D’Amico et al. (2002) of two accelerated pulsars in NGC 6752 at 6” and 7” from the center has highlighted the presence of about $`2\times 10^3M_{}`$ of under-luminous matter in the core of the cluster (Ferraro et al. 2003a), making NGC 6752 a special target for the search of an IMBH. NGC 6752 is interesting also because it hosts two pulsars in its halo. PSR-A is a millisecond pulsar in a binary system with a white dwarf (Bassa et al. 2003; Ferraro et al. 2003b), located at 3.3 half mass radii away: it is the farthest pulsar ever observed in a cluster. Colpi, Possenti & Gualandris (2002), and Colpi, Mapelli & Possenti (2003, 2004) have suggested that this pulsar has been propelled in the halo due to a dynamical interaction with a binary IMBH. In modeling the gravitational encounter, Colpi, Mapelli & Possenti (2003) found that a $`(50200,10)M_{}`$ binary IMBH is the preferred target for imprinting the large kick to the pulsar. Colpi et al. (2005) shortly reviewed the dynamical effects that a binary IMBH would imprint on cluster stars.
Considering all the recent hints provided by optical and radio observations, we explore, in this paper, an new way to unveil a binary IMBH in a GC, previously overlooked, that exploits the dynamical fingerprint left by a binary IMBH on cluster stars. In particular we like to address a number of questions: (i) What signature does a binary IMBH imprint on cluster stars ? Are the stars heated during the scattering process still remaining bound to the cluster ? (ii) Is there direct transfer of angular momentum from the binary IMBH to the stars to produce some degree of alignment, given the large inertia of the BHs ? (iii) Are prograde or retrograde orbits equally scattered ?
In this paper, we simulate 3-body encounters of cluster stars with a binary IMBH, studying the energy and angular momentum exchange, and reconstructing the trajectories of those stars that are scattered away from equilibrium by the binary IMBH. In Section 2 we outline the method used to simulate 3-body encounters between a binary IMBH and the cluster stars. In Section 3 we present our main results: the formation, in a GC hosting a binary IMBH, of a population of high velocity (bound) stars (that we call ”supra-thermal stars”), which tend to align their orbital angular momentum with that of the binary IMBH. In Section 4 we investigate on the expected number of these supra-thermal stars. In particular, we use an upgraded version of the code presented in Sigurdsson & Phinney (1995) to follow the dynamical evolution of these supra-thermal stars inside a GC model which reproduces the characteristics of NGC 6752. This study allows us to put some constraint on the detectability of this family of high velocity stars. Section 5 contains our conclusions.
## 2 The simulations
We simulated 3-body encounters involving a binary IMBH (composed of two BHs of mass $`M_1`$ and $`M_2`$) and a cluster star of mass $`m=0.5M_{}.`$ The dynamics of the encounter is followed solving the equations of motion with a numerical code based on a Runge-Kutta fourth order integration scheme with adaptive stepsize and quality control (explained in Colpi, Mapelli & Possenti 2003, hereafter CMP). For the target binary IMBH we considered an interval of masses between $`60M_{}`$ up to $`210M_{}`$. These are the favored masses of the hypothetical binary IMBH in NGC 6752. The binary has semi-major axis $`a`$ of 1, 10, 100, and 1000 AU and eccentricity<sup>1</sup><sup>1</sup>1 We have repeated select simulations for different values of the eccentricity and find only minor quantitative differences ($`\stackrel{<}{}`$10%) in our results. $`e=0.7`$, which is nearly the average binary eccentricity in statistical equilibrium (Hills 1975). The set of models is described in Table 1.
The initial conditions, that are Monte Carlo generated as in CMP, are sampled using the prescriptions indicated in Hut & Bahcall (1983). In particular, the relative velocity $`u^{in}`$ is distributed homogeneously between $`8.511`$ km s<sup>-1</sup> (according to the value of stellar velocity dispersion measured in NGC 6752 which is our reference cluster; Dubath, Meylan & Mayor 1997), and the impact parameter $`b`$ is drawn at random from a probability distribution uniform in $`b^2`$ and in a range going from 0 to a truncation value $`b_{max}`$. The truncation value $`b_{max}`$ is chosen by requiring that the simulations include all the encounters with non-negligible energetic exchange, i.e. with outgoing velocities of the interacting star significantly higher than the initial velocity (see Appendix A for a discussion of our choices of the impact parameter ranges). The three orientation angles and the phase of the binary are generated as indicated in Table 1 and 2 of Hut & Bahcall (1983). We initiate (terminate) integration when the distance between the incoming (outcoming) star and the center of mass of the binary is comparable to the radius of gravitational influence of the IMBHs ($`r_a2G(M_1+M_2)/\sigma ^2`$, where $`\sigma `$ is the 1-D stellar dispersion velocity). At the start of each simulation, we place the binary BH and the single star on their respective hyperbolic trajectories.
At the end of every scattering experiment we store the final binding energy $`E_{\mathrm{BH}}^{fin}`$ and angular momentum $`𝐉_{\mathrm{BH}}^{fin}`$ of the binary IMBH, together with the velocity at infinity (or post-encounter asymptotic velocity) $`u^{fin}`$, defined as the velocity of the interacting star after the encounter, estrapolated at infinity, and the angular momentum $`𝐉_{}^{fin}`$ of the interacting star. In this paper we want, in particular, to quantify the mean change in the absolute value $`J_{}`$ of the vector $`𝐉_{}`$ and the extent of alignment of $`𝐉_{}`$ in the direction of $`𝐉_{\mathrm{BH}}`$.
## 3 Results
### 3.1 Supra-thermal stars
We find that a sizable number of high velocity stars is created, under certain conditions, that may highlight the presence of an IMBH in the cluster. These stars gain kinetic energy and an excess velocity relative to the mean, remaining bound to the cluster. In Table 3 we give the fraction of bound stars, defined as those having a final velocity lower than the escape velocity from the core $``$35-40 km s<sup>-1</sup>, and the fraction of high velocity or supra-thermal stars, defined as having a post-encounter asymptotic (or ”at infinity”) speed between 20 and 40 $`\mathrm{km}\mathrm{s}^1`$. Depending on the characteristics of the binary IMBH, these stars can account for $`\stackrel{<}{}50\%`$ of all stars that have experienced an encounter.
Figure 1 shows the distribution of the velocity at infinity of stars scattering off a binary IMBH, for cases D1, D2, D3 and D4 (in Table 1). The shaded area indicates the strip of supra-thermal stars. The mean values of the speed at infinity $`u^{fin}`$ and of the fractional binding energy exchange $`\mathrm{\Delta }E_{\mathrm{BH}}/E_{\mathrm{BH}}^{in}=(E_{\mathrm{BH}}^{fin}E_{\mathrm{BH}}^{in})/E_{\mathrm{BH}}^{in}`$ are given in Table 2 together with the dimensionless factor $`\xi _E`$, defined by:
$$\mathrm{\Delta }E_{\mathrm{BH}}/E_{\mathrm{BH}}^{in}=\xi {}_{E}{}^{}[m/(M_1+M_2)].$$
(1)
As illustrated in Figure 1, dynamical encounters widen the velocity distribution of the stars. This effect is particularly severe when the binary is hard, having a higher binding energy which becomes available in the interaction. $`\xi _E`$ clusters around values between 0.6 and 4, and shows a maximum declining as the binary becomes extremely hard. This explains why case D1 has a peak velocity at infinity smaller than case D2 corresponding to a less hard binary (see Table 2). We tried to fit the behavior of $`\xi _E`$, as a function of the orbital separation $`a`$, with a parabola. We obtained $`\xi _E=b\text{Log}^2(a)+c\text{Log}(a)+d`$ for $`b`$=-1.2495, $`c`$=3.8615, $`d`$=0.958; this fit is only an approximation, since we have few numerical points (Fig. 2). It is interesting to notice that for $`a=1`$ AU $`\xi _E`$ is lower than for wider binaries. This does not represent a violation of the first statement of Heggie’s law (i.e. that hard binaries tend to become more energetic, Heggie 1975; Hills 1990); in fact all our hard binaries become harder and harder. But a $`\xi _E`$ decreasing for $`a<10`$ AU seems in partial disagreement with the second statement of Heggie’s law, for which the hardening rate of the binaries is ”approximately independent of their binding energy” (Heggie 1975). This difference is the result of a different approach with respect to that followed by Heggie. In fact, we are treating very massive binaries with respect to the incoming star, and such binaries require very large maximum impact parameters, nearly independently of their semi-major axis, to include in our simulations all the interactions which have significant energy exchange (i.e. $`\mathrm{\Delta }E_{\mathrm{BH}}/E_{\mathrm{BH}}^{in}>10^3`$), corresponding to post-encounter asymptotic velocities of the star in the supra-thermal range (see Appendix A for a complete discussion).
In Figure 1, case D1 and D3 give the highest number of supra-thermal stars among the different runs. These supra-thermal stars are the ones that have also acquired a sizable fraction of the orbital angular momentum of the binary IMBH. Thus in scattering off the binary IMBH they are preferentially launched in a halo orbit, inside the cluster. Their detection (discussed in $`\mathrm{\S }`$ 4) would be the sign of an IMBH hidden in the cluster.
### 3.2 Angular Momentum Transfer and Alignment
An interesting question to address is whether and how the orbital angular momentum of the binary IMBH couples to the star after a close gravitational encounter. The BHs can be sufficiently massive and the orbit be sufficiently wide that the orbital angular momentum of the binary exceeds that of the incoming star. In this case a direct transfer of orbital angular momentum can occur. The star coming close to the binary IMBH can be dragged into corotation, i.e. the star can emerge after the encounter with an angular momentum nearly aligned with the binary IMBH.
If we denote with $`\mu =M_1M_2/(M_1+M_2)`$ the reduced mass of the binary hosting the two BHs, the total angular momentum of the system is
$$\begin{array}{c}𝐉=𝐉_{\mathrm{BH}}{}_{}{}^{in}+𝐉_{}^{in}\hfill \\ =\mu \sqrt{aG(M_1+M_2)}𝐳+bu^{in}\left[\frac{m(M_1+M_2)}{M_1+M_2+m}\right]𝐳^{}\hfill \end{array}$$
(2)
where $`𝐳`$ and $`𝐳^{}`$ are the unit vectors indicating respectively the directions of $`𝐉_{\mathrm{BH}}^{in}`$ and $`𝐉_{}^{in}`$. Angular momentum transfer from the binary to the interacting star and partial alignment become important when $`J_{BH}^{in}J_{}^{in}`$, i.e., when
$$\frac{\mu }{bu^{in}m}\sqrt{(M_1+M_2)aG}1.$$
(3)
Our experiment confirms this fact. Figure 3 (solid lines) shows the post-encounter distribution of $`J_{}`$ for cases D2, D3 and D4. It is compared with the distribution of $`J_{}`$ of the incoming stars (dashed lines) to show that the binary IMBH transfers angular momentum to the stars widening the distribution of $`J_{}`$. Table 2 contains the averaged value of the fractional angular momentum increase (in modulus) $`\mathrm{\Delta }J_{}/J_{}{}_{}{}^{in}`$ for the complete series of runs. $`\mathrm{\Delta }J_{}/J_{}^{in}`$ is large and exceeds unity in correspondence to the runs where the bulk of the supra-thermal stars are produced. It has a non monotonic trend with the hardness of the binary: a harder binary has little excess of angular momentum in its initial state relative to that of the incoming star; so the efficiency of angular momentum transfer reduces. A less hard binary (having a larger $`J_{\mathrm{BH}}^{in}`$) is also less efficient, since it does not alter much the post-encounter asymptotic velocity.
We can further quantify the importance of angular momentum transfer from the binary to the star, by computing the fraction of stars which increase $`J_{}`$ after an interaction; this is given in the last column of Table 4. The fraction of stars that do so, over the total (in the sample of the bound stars), is significantly high: it is above 70$`\%`$ in many cases.
Is alignment induced in the scattering process ? In Table 4 we compare the percentage of bound stars which were corotating (i.e. for which the scalar product between their angular momentum and the angular momentum of the binary is positive, that is $`𝐉_{}𝐉_{\mathrm{BH}}>0`$) before the 3-body interaction with the percentage of those stars which are corotating after the 3-body interaction. Whereas the fraction of corotating stars before the interaction is $``$50% (as one can expect given the initial sampling of the data), the fraction of corotating stars after the interaction is often over 60%, indicating a tendency of stars to align their angular momentum with that of the binary IMBH. This tendency is greater if the binary has a big reduced mass (cases with $`M_1=100,M_2=50M_{}`$ and with $`M_1=50,M_2=50M_{}`$) and if the binary is moderately hard (i.e. its orbital separation is neither too small, because in this case $`J_{\mathrm{BH}}^{in}`$ is comparable with $`J_{}^{in}`$, nor too large, since in this case the binary would be too soft, and so the kinetic energy exchange would be negligible). In case D3 the bound stars which after the interaction are corotating with the binary are $``$70%. This means that we can observe a family of supra-thermal stars which have an anisotropy in their orbital angular momentum relative to the center of the cluster, 70% of them rotating in one direction and the remaining 30% rotating in the other. Obviously, we have to take into account that orbits initially are isotropically distributed and that we observe them in projection; then the detection of this phenomenon is not so immediate.
Table 4 indirectly shows also another interesting effect. The first column indicates that generally less than 50% of the initially corotating stars remain bound to the cluster, while the post-encounter percentage of corotating stars is more than 50%. This means not only that a fraction of initially counter-rotating stars becomes corotating, but also that an initially corotating star is more easy ejected from the cluster than a counter-rotating star. This fact can be intuitively explained considering that, when a counter-rotating star interacts with the binary, its relative velocity with respect to the lighter BH is higher than in the case of a corotating star, and, then, the cross section is lower.
To better constrain orbit alignment in the direction of $`𝐉_{\mathrm{BH}}`$, we considered separately orbits which before the encounter (as initial condition) were corotating (Table 5) and orbits which before the encounter were counterrotating (Table 6). For both we investigated the final distributions, and in particular we compared the fraction of stars initially corotating which after the interaction remain corotating (Table 5, 2nd column) with the fraction of stars initially counterrotating which after the interaction become corotating (Table 6, 2nd column). We found that most (at least 70%) of the initially corotating stars remains corotating, even if during very energetic interactions -with the hardest binaries- the star tends to forget its initial angular momentum orientation (case D2). In contrast, a high percentage of initially counterrotating stars becomes corotating, flipping their angular momentum. The combination of these two tendencies (i.e. the tendency to remain corotating for initially corotating stars and the tendency to become corotating for initially counterrotating stars) provides the strongest evidence of an alignment between $`𝐉_{\mathrm{BH}}^{in}`$ and $`𝐉_{}^{in}`$ after the interaction (Fig. 4).In Fig. 5 we give a complete overview of the cases considered.
## 4 Detecting supra-thermal stars
We want now to estimate how many supra-thermal stars are produced by interactions with a binary IMBH and if they are observable with present instrumentation.
### 4.1 Estimating the number of supra-thermal stars
Since supra-thermal stars have excess kinetic energy and angular momentum, they are not in thermodynamic equilibrium with respect to the rest of the cluster. Thus, two-body relaxation will try to obliterate such anisotropies. Friction will return them to the core where they become unrecognizable, unless some memory of the angular momentum alignment (e.g. Section 3.2) is retained during the relaxation process.
The existence of supra-thermal stars requires:
* the survival of the IMBH as a binary in the cluster core;
* the persistence of the supra-thermal stars in the GC halo.
We can investigate these questions through a comparison between relevant timescales. The first timescale to consider is the hardening time of the binary IMBH (Quinlan 1996), given the stellar number density in the core $`n`$ and the binary separation $`a`$
$$t_{hard}=\frac{\sigma }{2\pi a\xi {}_{E}{}^{}Gmn}$$
(4)
The second relevant timescale is the relaxation time at the half mass radius
$$t_{rh}=\frac{0.14N}{\mathrm{ln}(0.4N)}\left(\frac{r_h^3}{GM}\right)^{1/2}$$
(5)
where $`N`$ and $`M`$ are respectively the number of stars and the total mass of the cluster. In Figure 6 we compare the two timescales for $`\sigma =8.5\mathrm{km}\mathrm{s}^1`$, $`\xi _E1`$ (see Table 2) and $`n=10^5`$ stars pc<sup>-3</sup>. This plot shows that, when the binary is not very hard, i.e., when $`a>\mathrm{\hspace{0.17em}4}`$ AU, the hardening time is shorter than the half-mass relaxation time. This implies that we can simultaneously observe all the supra-thermal stars in the halo for a time comparable to $`t_{rh}`$ with no dilution. On the other hand, when the binary has orbital separation $`a\mathrm{\hspace{0.17em}4}`$ AU, $`t_{rh}`$ is shorter than $`t_{hard}`$, and, then, we can not simultaneously observe all the produced supra-thermal stars in the halo, a part of them being already thermalized. However, this effect is balanced by the fact that the binary remains in this very hard configuration for most of its life, and, therefore, we can observe a considerable fraction ($`60`$%) of the total number of supra-thermal stars even when the binary is very hard ($`a\mathrm{\hspace{0.17em}4}`$ AU).
For completeness, we also considered a third timescale: the gravitational wave timescale, $`t_{gw}`$, i.e. the characteristic time for a binary IMBH to coalesce due to gravitational wave emission. When this timescale, defined as (Peters 1964; Quinlan 1996)
$$t_{gw}=\frac{5}{256}\frac{c^5a^4(1e^2)^{7/2}}{G^3M_1M_2(M_1+M_2)},$$
(6)
becomes shorter than $`t_{hard}`$, the binary would shrink, due to gravitational wave emission, at a rate shorter than what we expected from $`t_{hard}`$. Figure 6 shows the timescale $`t_{gw}`$ computed for a constant eccentricity $`e=0.7`$ (the initial eccentricity we chose for our runs) and for the range of masses considered in this paper. We expect that the assumption of constant eccentricity is correct for star-binary IMBH interactions, since, when $`M_1+M_2m`$, the binary eccentricity suffers very small changes. We also checked whether this assumption is correct on the basis of our simulations. We found that the mean eccentricity $`e`$ of the binary IMBH after a 3-body encounter is always 0.70, with a very narrow spread (the standard deviation $`\sigma `$ being always $`15\times 10^2`$), and the maximum post-encounter eccentricity is always $`\stackrel{<}{}0.8`$, in agreement with our statement.
From Figure 6 we note that $`t_{gw}`$ is longer than $`t_{hard},`$ for $`a\stackrel{>}{}0.50.3`$ AU. This implies that the lifetime of a binary IMBH is controlled by 3-body scattering on the hardening time $`t_{hard},`$ as long as the binary separation $`a`$0.1 AU. Therefore, we can neglect the effect of gravitational wave emission, since we consider binaries wider than 0.1 AU.
The expected number of supra-thermal stars is a fraction $`f`$ of the total number $`N`$ of stars strongly interacting with a binary IMBH. The latter is given by (cfr. CMP):
$$N=\frac{(M_1+M_2)}{m\xi _E}\mathrm{ln}\left(\frac{a_0}{a_{\mathrm{st}}}\right)$$
(7)
where $`a_0`$ is the initial semi-major axis of the binary and $`a_{\mathrm{st}}`$ is the minimum semi-major axis for the encounter to give a supra-thermal star. In eq. (7) we impose $`a_0=`$ 2000 AU, corresponding to the orbital separation below which the binary becomes reasonably hard to generate supra-thermal stars, and $`a_{\mathrm{st}}=`$ 0.1 AU, the orbital separation below which the cross section for 3-body encounters becomes negligible. We then calculate the number of stars that remain bound to the cluster and the number of supra-thermal stars using the statistics derived in our simulations (see the percentages given in Table 3, respectively in first column -for the bound stars- and in fourth column -for the supra-thermal). The results are shown in Table 7. The total number of bound stars is always greater than 400 and the number of stars stirred up to supra-thermal velocities is always greater than 100, even for the lightest system ($`M_1`$= 50 $`M_{}`$, $`M_2`$= 10 $`M_{}`$).
Supra-thermal stars may be recognized from their proper motion and/or Doppler line shift. Hence we have to take into account projection effects of velocity vectors. These effects depend in turn on the orientation of the angular momentum of the binary IMBH $`𝐉_{\mathrm{BH}}`$ with respect to the line of sight. We find that about 80% of the supra-thermal stars can be recognized as such in the best case, i.e. when the line-of sight happens to be nearly parallel to the orbital angular momentum of the binary IMBH $`𝐉_{\mathrm{BH}}`$ (if we are measuring proper motions) or when the line-of-sight is nearly perpendicular to $`𝐉_{\mathrm{BH}}`$ (if we are measuring Doppler-shifts). On the contrary, the percentage reduces to about 50% when we are in the most unlucky cases, i.e. when the line of sight happens to be nearly perpendicular to the orbital angular momentum of the binary IMBH $`𝐉_{\mathrm{BH}}`$ (if measuring proper motions), or when the line of sight is parallel to $`𝐉_{\mathrm{BH}}`$ (if measuring Doppler shifts). Assuming these percentages, we find that the number of recognizable supra-thermal stars over the entire cluster is $``$60-100 for a ”light” binary IMBH ($`M_1`$= 50 $`M_{}`$, $`M_2`$= 10 $`M_{}`$) and $``$130-210 for a ”moderately massive” binary IMBH ($`M_1`$= 100 $`M_{}`$, $`M_2`$= 50 $`M_{}`$).
Our calculation refers to all stars, including also compact remnants (neutron stars, white dwarfs) and stellar types which are too faint for allowing a measurement of proper motion or radial velocity. Using the numerical code which will be described in the next Section, we find that the stars with mass from 0.6 to 0.9 $`M_{}`$ (a mass range which includes red giant (RGB), horizontal branch (HB) and bright enough main sequence (MS) stars) represent about the 45% of the stars enclosed within<sup>2</sup><sup>2</sup>2We consider the region within 0.1 $`r_c`$, because we are interested in those stars which have the highest probability of interacting with the central binary IMBH. 0.1 $`r_c`$ of a GC like NGC 6752. This means that, in the case of a binary of (100$`M_{}`$ , 50$`M_{}`$), we are able to recognize only 50-100 supra-thermal stars.
There is a further problem concerning the number of surviving supra-thermal stars. Nearly all the proposed mechanisms of formation of binary IMBHs in GCs (Miller & Hamilton 2002; Sigurdsson & Hernquist 1993) predict that such binaries are born within the first Gyr since the formation of the GC itself. If this is true, considering the hardening time $`t_{hard}`$ plotted in Figure 6, the binary IMBHs which still survive today must have orbital separation $`\stackrel{<}{}1`$ AU (for a cluster with $`n=10^5`$ pc<sup>-3</sup>). This means that all the supra-thermal stars produced when the binary had orbital separation $`a\stackrel{>}{}3`$ AU have already been thermalized, because a time much longer than $`t_{rh}`$ has elapsed (Fig. 6). Then, we can observe only the supra-thermal stars produced in the last $`2`$ Gyr (about 3 $`t_{rh}`$, because a supra-thermal star can have apo-center distance larger than the half mass radius and, therefore, relaxation time longer than $`t_{rh}`$). Luckily, the orbital separation of a IMBH binary remains in the range 0.1-3 AU for a long time, because the hardening rate is lower for such a hard binary, and a large fraction of supra-thermal stars ($``$60%) is produced in this stage. The number of these ”surviving” supra-thermal stars is reported in Table 7, 6th column. Let us consider again the IMBH binary of (100,50) $`M_{}`$. The total number of supra-thermal stars produced by this binary in the last 2 Gyrs is $``$160. This means that, taking into account projection effects and considering only bright enough stars (i.e. in the mass range 0.6-0.9 $`M_{}`$), we are finally left with only 30-60 recognizable supra-thermal stars.
In this discussion we have not considered the effective instrumental errors so far. We now briefly report on them, without entering in details. Even if STIS is no more operative, its accuracy remains a good lower limit for future spectrographs. Observing with STIS stars in a globular cluster with distance from the Sun of the order of 4-10 kpc , one can expect an error of 1-2 km s<sup>-1</sup> in the determination of radial velocity. For example, van der Marel et al (2002), reported spectra of about 130 stars in the core of M15 (distance from the Sun about 10 kpc) with an observational error of the order of 1.3 km s<sup>-1</sup>. A somewhat higher error can be estimated for proper motion measurements with the HST/WFPC2. Drukier et al. (2003) combined two sets of observations with the WFPC2 (one taken in 1994, the second in 1999) for a sample of 1281 stars in NGC 6752, estimating a median error of 0.31 mas yr<sup>-1</sup> with a mode of 0.17 mas yr<sup>-1</sup>. For the distance of NGC 6752 this means a median error of $``$6 km s<sup>-1</sup> with a mode of $``$3 km s<sup>-1</sup>, which is still an acceptable accuracy to distinguish supra-thermal stars. Similar considerations hold for HST/ACS (see e.g. Anderson 2002; Anderson & King 2003). Since we defined supra-thermal the stars with projected velocity higher than $``$ 12 km s<sup>-1</sup> (16 km s<sup>-1</sup> for proper motion measurements), an error of 1-2 km s<sup>-1</sup> (3-6 km s<sup>-1</sup> for proper motion measurements) is sufficient to distinguish supra-thermal from other cluster stars, in clusters like NGC 6752. In summary, the main problem in detecting supra-thermal stars is not the error on the single measurement but the possibility of observing a sufficient large sample of stars, since supra-thermal stars are expected to be a very small fraction of cluster stars.
### 4.2 Spatial distribution of the supra-thermal stars
Since supra-thermal stars are only few tens, the possibility of recognizing them may be significantly enhanced if their radial distribution shows some characteristic feature. Thus, it is of interest to study how supra-thermal stars evolve in the cluster and what is their radial distribution.
To this purpose we have explored the dynamics of supra-thermal stars in a cluster under the action of dynamical friction and the influence of two-body relaxation effects using an updated version of the code described in Sigurdsson & Phinney 1995. The code first generates a cluster background model, i.e. a multi-mass King density profile which in our case reproduces that of NGC 6752, and we inject in this background supra-thermal stars whose initial positions and velocities are those obtained by our 3-body simulations. We followed the dynamical evolution of these stars for a random time $`t`$ uniformly distributed in the range $`0<t3t_{rh}`$. In this way we can see how supra-thermal distribute in the cluster before thermalization. After evolution for a time t$`\stackrel{<}{}3t_{rh}`$ we select the stars that still are supra-thermal. We applied this procedure to the cases D1, D2, D3 and D4.
The stars which were still supra-thermal when the simulation stopped and that can be observed as supra-thermal, taking into account two-dimensional (one-dimensional) projection effects and ejections, are about 70% (60%) of the initial sample. These stars present a radial distribution like that shown in Figure 7 (for the case D1). This distribution is peaked around a radius $`r_{peak}`$ which is of the order of few core radii for all the considered systems. For a binary IMBH with $`M_1=100M_{}`$, $`M_2=50M_{}`$ and $`a=1`$ AU (D1) $`r_{peak}=3r_c`$ and we have calculated that it is enough to measure radial velocities within a distance of 6 $`r_c`$ from the cluster center, in order to avoid poissonian fluctuations to smear out the peak.
Taking again into account both projection effects and luminosity criteria, only $``$30 supra-thermal stars are located within 6 $`r_c`$. We found that, if the velocity distribution is Maxwellian, $``$10 stars with velocity in the supra-thermal range are predicted to inhabit within a distance of 6 $`r_c`$ from the center of NGC 6752. Then, the number of supra-thermal stars produced by interactions with a binary IMBH prevails on the high velocity tail of the Maxwellian distribution. On the other hand the presence of this tail further dilutes the signature of supra-thermal stars.
In Fig. 7 we have finally compared the radial distribution of 1984 red giant stars (dotted line) observed in NGC 6752 (591 of them from HST/WFPC2 observations, 1393 from the ESO/MPG WFI. See Sabbi et al. 2004) with that calculated for supra-thermal stars (solid line). RGB stars appear to be more concentrated toward the center of the cluster (the maximum being between 0 and 2 $`r_c`$); but the two distributions are quite similar.
### 4.3 Detectability of angular momentum alignment
We now discuss the observability of the signature of the angular momentum alignment effect. Column 2 of Table 4 tells us that in the most favorite case about 70% of the stars which remain bound are corotating with the binary IMBH, whereas 30% are counterrotating. This means that, if there is a binary IMBH of $`M_1`$= 100 $`M_{}`$ and $`M_2`$= 50 $`M_{}`$, we should have about 90-150 corotating and only 40-60 counterrotating supra-thermal stars. On the other hand, angular momentum alignment effects are visible only for sufficiently wide binaries ($`a>10`$ AU). Thus, unless a binary IMBH is formed recently in a cluster, the alignment effect today is completely washed out by dynamical friction.
However, we notice that a very low probability mechanism of binary IMBH formation in GCs in recent epochs exists. In fact, there is some possibility that a ”last single BH” (Sigurdsson & Hernquist 1993), ejected from the core in the early stages of the GC life, remains in the halo for several Gyrs. The relaxation timescale out in the halo is long, and if the orbits can circularize there (maybe due to time varying galactic tidal field), then the return time to the core is long. When this BH comes back to the core, there is a high probability that it forms a binary with the central IMBH. Such a binary IMBH would be originated in late epochs. However, the probability for this process to occur is low, because it requires the BH to receive just the fine-tuned post-encounter velocity needed to remain in the outer halo: if the velocity is slightly too high, the BH will be ejected from the entire cluster; whereas, if the velocity is low, the return time to the core will be too short.
The process described here is an unmistakable feature of a binary IMBH, which can not be produced by a distribution of stellar binaries interacting via 3-body encounters with cluster stars. Instead, we need dedicated 3-body simulations to check whether even stellar binaries can produce effects like those described in $`\mathrm{\S }3.1`$ and $`\mathrm{\S }4.2`$.
## 5 Summary
In our work we explored the effects that the presence of a binary IMBH imprints on the velocity and angular momentum of cluster stars, due to 3-body encounters with it. The main results of our analysis can be summarized as follows.
1. A binary IMBH generates a family of supra-thermal stars, i.e. of stars which remain bound to the cluster, but have a velocity higher than the dispersion velocity. ($`\mathrm{\S }`$ 3.1).
2. A fraction of stars tend to align their angular momentum, $`𝐉_{}`$, with that of the binary IMBH, $`𝐉_{\mathrm{BH}}`$. This fraction depends on the reduced mass $`\mu `$ and on the semi-major axis $`a`$ of the binary IMBH, but it is always of the order of 55-70%. This means that the angular momentum distribution of stars that have suffered 3-body interaction with a binary IMBH presents a slight anisotropy ($`\mathrm{\S }`$ 3.2).
3. The present theoretical models of binary IMBH formation indicate that these systems are produced in the first Gyr of the GC life. This means that today, in dense clusters like NGC 6752, binary IMBHs have, due to hardening, orbital separation $`a\stackrel{<}{}`$ 1 AU. As a consequence, supra-thermal stars produced when the binary was wider (more than 2 Gyrs ago for the case of NGC 6752) have now already thermalized. That reduces the current population of residual supra-thermal stars about of a factor two with respect to the total population produced during the entire IMBH lifetime ($`\mathrm{\S }`$ 4.1). On the other hand these ”surviving” supra-thermal stars do not show any signature of angular momentum alignment ($`\mathrm{\S }`$ 4.3), because this effect is evident only when the stars interact with a relatively wide binary ($`a>10`$ AU).
4. How many supra-thermal stars are expected to be produced by a binary IMBH and still surviving in a GC like NGC 6752? We estimated few hundreds of supra-thermal stars. Once subtracted the fraction of compact remnants and faint stars and taking into account projection effects and thermalization of the oldest supra-thermal stars (see §4.1) we calculate that there may be at most few tens of supra-thermal stars suitable to be recognized as such in the whole cluster. It can be interesting to notice that there is already a claim for the detection of high velocity stars in 47 Tuc (Meylan, Dubath & Mayor 1991).
5. Given the aforementioned relatively small number of expected supra-thermal stars it is important to study their spatial distribution in order to improve the chances of recognizing their origin. To this purpose (see $`\mathrm{\S }`$ 4.2), we followed the dynamical evolution of supra-thermal stars in the cluster potential (assuming a cluster model which reproduces the observational characteristics of NGC 6752), taking into account dynamical friction and two body relaxation. We found that, before thermalization, supra-thermal stars tend to cluster within 6 $`r_c`$.
6. Although we have not performed a detailed simulation of the observability of the few tens of expected supra-thermal stars with present detectors, simple considerations (see §4.1 and §4.2) suggest that the various observational biases (mainly related with the limited field-of-view of the most sensitive instruments) may lead to the detection of only a sub-sample of the population.
In the light of these findings, the search of supra-thermal stars and of anisotropies in their angular momentum distribution can hardly provide useful indications on the existence of binary IMBHs in GCs with present instruments. However, it may become a feasible task for future instruments having much bigger collecting area and equipped with detectors having much larger field of view, like (Gilmozzi 2004) the Thirty Meter Telescope (TMT) or the OverWhelmingly Large Telescope (OWL).
## 6 Acknowledgments
We thank Francesco Ferraro for useful discussions on the issue of supra-thermal star detection. We also thank Marc Freitag, Cole Miller and the anonimous Referee for a critical reading of the manuscript. This work found first light during the workshop ”Making Waves with Intermediate-Mass Black Holes” (20-22 May 2004, Penn State University). M. C. and M. M. acknowledge the Center for Gravitational Wave Physics for kind hospitality. S. S. acknowledges support by the Center for Gravitational Wave Physics funded by the NSF under cooperative agreement PHY 01-14375 and NSF grant PHY 02-03046. M. C. and A. P. acknowledge financial support from the MURST, under PRIN03.
## Appendix A Choice of impact parameters
We chose impact parameters in order to select all the interactions which have not negligible energetic exchange (about $`\mathrm{\Delta }E_{\mathrm{BH}}/E_{\mathrm{BH}}^{in}\stackrel{>}{}10^3`$), to allow a complete coverage of all the interactions that enter in the supra-thermal domain (20-40 km s<sup>-1</sup>). We used the formula of the gravitational focusing (Sigurdsson & Phinney 1993), for which
$$b_{max}\left(\frac{2GM_Tp}{v_{\mathrm{}}^2}\right)^{1/2},$$
(8)
where $`b_{max}`$ is the maximum impact parameter, $`M_T`$ the total mass of the binary, $`v_{\mathrm{}}`$ the relative velocity at infinity and $`p`$ the distance of closest approach.
We firstly imposed $`p(10a)`$ for all the systems (where $`a`$ is the semi-major axis of the binary). Then we made some test-run to check what were the best maximum impact parameters to include all the strongly interacting stars (and then the supra-thermal stars). The results of these checks induced us to slightly modify $`b_{max}`$ as calculated with eq. (8). In fact, when a binary has a very large semi-major axis (case with $`a`$=100 or 1000 AU), the choice $`p(10a)`$ means that the incoming star will pass very far from the center of mass of the binary, will interact very weakly with the binary (unless it closely approaches the secondary component of the binary itself) and never obtain $`\mathrm{\Delta }E_{\mathrm{BH}}/E_{\mathrm{BH}}^{in}>10^3`$. The velocity at infinity peak is always $``$ 10 km/s.
On the other hand, when the binary has a very small semi-major axis ($`a`$1AU), the choice $`p(10a)`$ does not take into account all the strong interactions. In fact all the stars which have $`p(10a)`$ pass so closely to the very hard binary to become very fast ($`>`$ 100 km/s). But, if we consider only $`p(10a)`$, we miss most of the interacting stars. In fact, if we take $`p(50100)a`$, we still have stars that sufficiently approach to the center of mass of the binary to receive an outgoing velocity $``$ 20-40 km/s (exactly the range of supra-thermal stars). And, statistically, stars with $`p>10a`$ will be much more numerous than stars with $`p<10a`$. Then, if we take $`p\stackrel{<}{}(10a)`$, we miss the bulk of interacting stars belonging to the supra-thermal interval.
Then, both when we choose a small $`b_{max}`$ ($`p\stackrel{<}{}10a`$) for a very hard binary and a large $`b_{max}`$ ($`p\stackrel{>}{}10a`$) for a wide binary, we introduce a bias. In the former case we select only high velocity encounters ($`\mathrm{\Delta }E_{\mathrm{BH}}/E_{\mathrm{BH}}^{in}10^3`$), omitting the bulk of interacting stars; in the latter we risk to consider also unperturbing flybies ($`\mathrm{\Delta }E_{\mathrm{BH}}/E_{\mathrm{BH}}^{in}10^3`$).
Fig. A1 represents the post-encounter asymptotic velocity distribution of the star in the case of a binary with $`M_1`$=50 $`M_{}`$, $`M_2`$=10 $`M_{}`$, $`a`$=100 AU. The dotted line shows the case with $`b_{max}`$=2000 AU (corresponding to adopt $`p=10a`$ in eq. (8)): most of interactions are very weak flybies (velocity $`<`$ 8 km/s), and we cannot count them as strong interactions. Instead, the solid line shows the case with $`b_{max}`$=100 AU ($`p=a`$), i.e. the runs reported in the paper as case A3: this represents a valid statistic sample of all the strong interactions (the peak being at $``$ 13 km/s).
On the other hand, Fig. A2 shows the post-encounter asymptotic velocity distribution of the star in the case of a binary with $`M_1`$=50 $`M_{}`$, $`M_2`$=10$`M_{}`$, $`a`$=1 AU. The dotted line, showing the case with $`b_{max}`$=1 AU (corresponding $`p=1a`$), is peaked at very high velocities. Instead, the solid line, for $`b_{max}`$=60 AU (case A1), shows that, if we take into account also encounters with $`p>1a`$, but still with $`\mathrm{\Delta }E_{\mathrm{BH}}/E_{\mathrm{BH}}^{in}>10^3`$, we have a very different velocity distribution, peaked at 20-30 km/s. We think that the last one is the only statistically significant case, because it takes into account all the interacting stars with significant energetic exchange. The dotted line, instead, misses the most numerous class of interactions.
For these reasons we adopted new criteria for the choice of $`p`$ (to be substituted in eq. 8 to derive $`b_{max}`$) as:
* if $`a`$100 AU, then $`pa`$;
* if 1 AU$`<a<`$100AU, then $`p(10a)`$;
* if $`a`$=1 AU, then $`p(60100)a`$.
This means, that our selected maximum impact parameter is quite constant ($``$100 AU), independent of the semi-major axis a of the binary (at least if $`a<`$ 1000 AU).
We think that the big maximum impact parameter required by systems with $`a`$=1 AU is due to the uncommonly large mass ratio between the binary and the incoming star. In fact, when at least one of the two components of the binary is so massive, it plays an important role on the gravitational focusing, nearly independently of the semi-major axis. Then, the incoming star is attracted toward the center of mass of the binary, even if it starts with very large impact parameter. But a binary of $`a`$=1 AU has a very small geometrical cross-section; then only a very small fraction of interacting stars will have an effective pericenter $`pa`$. Thus, only a few stars will interact very strongly with the binary. Most of the stars will have a pericenter $`p>a`$; but the binary has a so large binding energy that also these stars receive non-negligible post-encounter velocity, and they can become supra-thermal.
Even the behavior of $`\xi _E`$ shown in Fig. 2, apparently different from what predicted by the Heggie’s law when $`a=`$1 AU, depends on our choice of $`b_{max}`$. In fact, if for a binary with semi-major axis $`a=`$ 1 AU we choose $`pa`$ (corresponding to $`b_{max}1`$ AU, instead of 60-100 AU), we obtain $`\xi _E4`$, higher than the value of $`\xi _E`$ derived for wider binaries and in agreement with Heggie’s law. On the other hand, we need to select $`pa`$, because we want to include all the interactions which have not negligible energetic exchange. But this choice lowers the average value of $`\xi _E`$ from 4 to about 1.
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# Dust Extinction in Compact Planetary Nebulae
## 1 Introduction
In general, PNs have axisymmetric structures. The classification of PNs into three basic types, round, elliptical and bipolar (Balick, 1987), is based on their axisymmetric appearances. However, many PNs also possess non-axisymmetric structures. Soker & Hadar (2002) proposed a scheme to classify PNs according to their departure from axisymmetry. They considered six categories: (1) PNs where the central star is not at the center of the nebula; (2) PNs having one side brighter than the other; (3) PNs having unequal size or shape of the two sides; (4) PNs where the symmetry axis is bent, e.g., the two lobes in a bipolar PN are bent toward the same side; (5) PNs where the main departure from axisymmetry is in the outer regions, e.g., an outer arc; (6) PNs showing no large-scale departure from axisymmetry. They found that about 50% of the analyzed sample has large-scale departures from axisymmetry.
The problem of departure from axisymmetry in PNs is usually investigated through the examination of optical images. This poses a problem because dust also plays a role in making such asymmetry. For example, in the sample analyzed by Soker & Hadar (2002), more than half of the PNs they found to have departure from axisymmetry are the types with unequal intensity or size. This kind of asymmetric appearance can also be caused by dust extinction, even if the intrinsic brightness of the nebula is axisymmetric. It is very difficult to draw conclusions about the real shapes of PNs based only on their apparent optical forms. In order to remove the bias caused by dust extinction while investigating the shapes of PNs, it is necessary to understand the dust distribution in PNs.
Asymptotic Giant Branch (AGB) stars eject large amount of dust and gas in the form of stellar winds. These stellar grains is a major source of dust grains in the Galaxy (Gehrz, 1989). Although the details of dust formation are still poorly understood, the properties of circumstellar dust have been extensively studied in the infrared, e.g., by the Infrared Astronomy Satellite (IRAS) and Infrared Space Observatory (ISO) missions. Since PNs descend from AGB stars, the remnants of the AGB dust envelopes should still be present in PNs (Kwok, 1982). The presence of dust in PNs was first detected in NGC 7027 via its infrared emission (Gillett, Low & Stein, 1967). This infrared excess peaks around 30 $`\mu `$m and has a color temperature of $``$100 K. Although the existence of dust in PNs first came as a surprise, its widespread presence is now confirmed by IRAS observations (Pottasch et al., 1984; Zhang & Kwok, 1991).
While the presence of dust in PNs is well established by infrared spectroscopy, very little is known about its spatial distribution. Recent development in mid-infrared cameras has just begun to explore the dust distribution of PNs (see, e.g., Volk & Kwok, 2003; Su et al., 2004). However, a full mapping of the dust distribution of PNs has to wait until the development of large-format mid-infrared cameras with high resolution and high sensitivity. An alternate and underutilized method to derive the dust distribution is by comparing the hydrogen recombination line map of a PN with its corresponding radio free-free continuum map. The recombination line (e.g. H$`\alpha `$) and free-free continuum emission are independent tracers of the distribution of the ionized gas, and both of their fluxes are proportional to $`n_e^2V`$ under optically thin conditions, where $`n_e`$ is the electron density and V is the volume of the nebula. If dust is present, they will have no effect on the long wavelength radio emission, but the H$`\alpha `$ emission will be attenuated by dust. Therefore the dust distribution can be indirectly obtained by comparing the H$`\alpha `$ and radio maps, assuming the dust is associated with the ionized gas. This method has been applied to different objects by various authors (e.g., Vázquez et al., 1999; Eyres et al., 2001). In this paper, we have derived optical depth maps of five compact PNs obtained with the same technique. Compact PNs are selected for this study because their small angular sizes allow the assumption that any foreground interstellar extinction is homogeneous over the entire object, so that the different degrees of extinction over a PN are only related to its internal dust distribution.
## 2 Derivation of the optical depth distribution
In a fully ionized region, both the free-free emission and the recombination fluxes are proportional to $`n_e^2V`$, where $`n_e`$ is the electron density. The ratio of the free-free flux to the H$`\beta `$ recombination line fluxes under Case B (where recombination to the ground state is excluded) is given by the following equation (Pottasch, 1984):
$$\frac{S_\nu (\mathrm{ff})}{F(\mathrm{H}\beta )}=2.51\times 10^7T_e^{0.53}\nu ^{0.1}Y\frac{\mathrm{Jansky}}{\mathrm{erg}\mathrm{cm}^2\mathrm{s}^1},$$
(1)
where $`T_e`$ is the electron temperature, and
$$Y=1+\frac{n(He^+)}{n_p}+3.7\frac{n(He^{++})}{n_p}$$
(2)
corrects for the contribution to the free-free flux from He ions. For a He to H number ratio of 0.11 and for He to be equally divided between singly and doubly ionized forms, $`Y`$ has a value of 1.258. Under Case B, the effective recombination coefficients for H$`\alpha `$ and H$`\beta `$ have a ratio of 2.85 (Hummer & Storey, 1987). For $`T_e=10^4`$ K, eq. 1 can be written as
$$F_{\mathrm{exp}}(\mathrm{H}\alpha )=6.85\times 10^{10}(\nu /\mathrm{GHz})^{0.1}(S_\nu /\mathrm{Jy})\mathrm{erg}\mathrm{cm}^2\mathrm{s}^1,$$
(3)
which can be used to calculate the expected flux of H$`\alpha `$ emission without dust extinction from the flux density of radio continuum emission.
In order to derive the extinction map, we begin with the radio continuum and H$`\alpha `$ maps of the same PN at similar angular resolutions. The expected H$`\alpha `$ emission at each pixel is derived from the radio continuum map, and then divided by the observed H$`\alpha `$ emission map for each object to obtain an extinction map, in which the dust optical depth at each pixel can be calculated. The intensity emitted by an object and received by the observer is related by the optical depth $`\tau `$ as
$$F_{\mathrm{obs}}=F_{\mathrm{exp}}e^\tau ,$$
(4)
where $`F_{\mathrm{exp}}`$ is the intensity emitted by the object, i.e., the intensity expected to be received by the observer if no dust is present. Therefore the dust optical depth for H$`\alpha `$ emission is obtained by taking the natural logarithm of the ratio of the expected and observed H$`\alpha `$ flux:
$$\tau _{\mathrm{H}\alpha }=\mathrm{ln}\frac{F_{exp}(\mathrm{H}\alpha )}{F_{obs}(\mathrm{H}\alpha )}.$$
(5)
## 3 Sample selection and data processing
In this paper, we have selected five compact PNs with known high infrared excesses. Figures 1 show the spectral energy distributions (SEDs) of these five PNs from optical to radio wavelengths, illustrating the IR excess caused by dust. The ground-based optical data are from Acker et al. (1992) and Acker, Marcout & Ochsenbein (1981). The near-infrared J, H, and K-band photometry are from the Two Micron All Sky Survey (2MASS), extracted from the 2MASS All-Sky Catalog of Point Sources (Cutri et al., 2003). The mid-infrared photometry includes IRAS photometry, IRAS low resolution spectra (LRS), and recent measurements from the Midcourse Space Experiment (MSX). The MSX measurements were extracted from the MSX Point Source Catalog (Egan et al., 2003). Radio photometric measurements (mostly at $`\lambda `$ 2 and 6 cm) are from Aaquist & Kwok (1990), Basart & Daub (1987), and this study. As clearly shown in Fig. 1, the infrared excess of the PNs is well represented by a single blackbody curve with a temperature $``$ 200 K.
High resolution H$`\alpha `$ and radio images of these PNs are available from the HST and the VLA respectively. The basic information for the five objects is listed in Table 1. With these high resolution images, the point-to-point spatial extinction distribution in PNs can be obtained with angular resolutions as high as $``$ 0.1 arcsec.
For IC 5117, M 1-61, M 2-43, and M 3-35, the HST WFPC2 H$`\alpha `$ images were obtained under the GO program 8307 (Kwok, Su & Sahai, 2003). For BD+303639, the HST WFPC2 H$`\alpha `$ image was first presented in Harrington et al. (1997) and the calibrated images were obtained from the HST archive. The only image processing needed was to remove cosmic rays with the task CRREJ in IRAF. The angular resolution FWHM of these images ranges from 0$`\stackrel{}{\mathrm{.}}`$085 to 0$`\stackrel{}{\mathrm{.}}`$11.
For BD+303639, the $`\lambda `$6 cm radio image was from Bryce et al. (1997). It was obtained by combining VLA and Multi-Element Radio Linked Interferometer Network (MERLIN) observations, resulting in a very high angular resolution of $``$ 0$`\stackrel{}{\mathrm{.}}`$08. For the remaining four objects, the radio continuum maps were observed by Aaquist & Kwok with the VLA at $`\lambda `$2 cm in A-configuration in March 1998. The unpublished data were retrieved from the VLA archive and reduced using the package AIPS following the standard procedures of continuum data reduction, as summarized in Appendix A of the AIPS Cookbook (NRAO, 2003). After the data were calibrated, the final CLEAN images were generated with robust weighting in order to obtain images with minimum-noise levels while maintaining the highest angular resolution. This ranges from $`0\stackrel{}{\mathrm{.}}12`$ to $`0\stackrel{}{\mathrm{.}}15`$ in the final CLEAN maps.
The operations taken to combine the radio and optical data and to derive the optical depth map are performed in the following sequence. The HST images were first oriented north-up with Karma software. Subsequent regridding and other steps were performed in AIPS. Because the coordinates of the HST images have a uncertainty of $`0\stackrel{}{\mathrm{.}}5`$, there is usually an offset in the positions of the same object in the radio and optical images. To correct this, the H$`\alpha `$ images were shifted and aligned by eye to match with the radio continuum maps using features common in both images. Table 3 gives the offset of the coordinates for each object; a positive sign indicates the optical image has been shifted eastward or northward. Each object’s H$`\alpha `$ image was convolved with an extended Gaussian to match the beam size of the corresponding radio image, except for BD+303639. Because the radio continuum image of BD+303639 has a higher resolution than its H$`\alpha `$ image, it was convolved with an extended Gaussian to match the FWHM of the H$`\alpha `$ image. The radio image was transformed so that its pixel coordinate grid was consistent with that of the HST image. Division and natural logarithm operations were then performed to obtain an optical depth map as in eq. 5.
## 4 Results and Analysis
Extinction maps showing the dust optical depth for the five compact PNs are presented in Fig. 2. For comparison, the H$`\alpha `$ brightness and extinction maps for each object are shown side by side, and the corresponding radio continuum contours are superposed on both maps. All images are oriented with north up and east to the left.
### 4.1 BD+303639
The radio continuum image of BD+303639 shows a ring-like structure with two peaks of roughly the same brightness in the north and south regions. The H$`\alpha `$ image also shows a ring-like structure with two peaks. However, the ring shows unequal brightness in the east and west regions, which indicates the ring suffers more dust extinction in the west as shown in the optical depth map. The two high-extinction regions in the northeast and southwest are in roughly the same positions as the pair of high-velocity molecular knots found with CO line imaging by Bachiller et al. (2000).
### 4.2 IC 5117
The radio continuum image of IC 5117 shows a shell structure with two emission peaks. Such a shape can be simply explained as a prolate ellipsoidal shell projected onto the plane of the sky (Aaquist & Kwok, 1991). The H$`\alpha `$ image shows a similar structure with two pairs of faint bipolar lobes oriented at different angles (see image displayed by Kwok, Su & Sahai, 2003). The lobes are not seen in the radio continuum because of relatively low sensitivity for the radio map. The extinction map shows a similar structure to that of the radio map, where the high-extinction region follows the region of high brightness in radio. However, the two extinction peaks fall slightly outside of the radio peaks. A similar situation has been found in the dust extinction map of NGC 7027 (Walton et al., 1988).
### 4.3 M 1-61
The radio continuum image shows a shell structure with two pronounced emission peaks of roughly the same brightness. In the H$`\alpha `$ image, the eastern peak is weaker than the western peak, and there is at least one pair of faint bipolar lobes in the outer region (see image displayed by Kwok, Su & Sahai, 2003). The extinction map shows the optical depth is higher in the region of brighter radio emission, with a strong extinction peak near the eastern radio peak.
### 4.4 M 2-43
The radio continuum image shows a shell structure, with the southern peak brighter than the northern one. The H$`\alpha `$ image shows similar structure with slightly extended faint emission outside the radio emission region. However, the northern peak of H$`\alpha `$ emission is brighter than the southern one. As a result, although the structure of the extinction map mainly follows that of the radio map, the optical depth in the southern part is significant higher than the northern part. The extinction peaks also falls slightly outside of the radio peaks.
### 4.5 M 3-35
The radio continuum image shows a deep central emission minimum and two surrounding peaks with concave-like contours at roughly the same brightness. The H$`\alpha `$ image shows two peaks at approximately the same location with two pairs of faint bipolar lobes outside (see image displayed by Kwok, Su & Sahai, 2003). The extinction map has similar structure to that of the radio map, with the high-extinction region coinciding with the concave-like region.
### 4.6 Uncertainties in the Optical Depth
Figure 3 show the uncertainties propagated from the radio and optical images. The uncertainty associated with the optical depth $`\tau `$ is given by
$$\mathrm{\Delta }\tau =\sqrt{\left(\frac{\sigma _{\mathrm{radio}}}{F_{\mathrm{radio}}}\right)^2+\left(\frac{\sigma _{\mathrm{H}\alpha }}{F_{\mathrm{H}\alpha }}\right)^2}$$
(6)
where $`\sigma _{\mathrm{radio}}`$ and $`\sigma _{\mathrm{H}\alpha }`$ are the rms noise of the radio continuum and optical H$`\alpha `$ images, respectively. Since the optical images have higher signal-to-noise ratios, the uncertainty of the optical depth is dominated by the radio image. Hence the uncertainty map of the optical depth and the radio continuum map have similar structure. The optical depth uncertainty ranges from $`\mathrm{\Delta }\tau `$ 0.01 to 0.4, where smaller values lie in the regions with stronger radio emission.
The calculations of optical depth above use several assumptions. To investigate whether the optical depth is sensitive to these assumptions, the uncertainties propagated from them should also be examined. Among these approximations, the electron temperature $`T_e`$ gives the most significant uncertainty. While most PNs have $`T_e`$ in the vicinity of 10,000 K, the electron temperatures of PNs cover a range of 5,000 - 15,000 K, with $`T_e`$ 20,000 K in some extreme cases. As a result, different choices of $`T_e`$ will give different values of optical depth.
To calculate the uncertainty propagated from the assumption that $`T_e=10,000`$ K, we estimate the uncertainty for $`T_e`$ to be roughly $`\mathrm{\Delta }T_e=\pm 5,000`$ K. Since $`\tau \mathrm{ln}F_{\mathrm{exp}}(\mathrm{H}\alpha )\mathrm{ln}T_e^{0.53}`$, the uncertainty propagated into the optical depth is $`\mathrm{\Delta }\tau 0.53\times \frac{\mathrm{\Delta }T_e}{T_e}0.53\times \frac{5,000}{10,000}0.27`$. This value is an order of magnitude smaller than the derived optical depth. This uncertainty from the electron temperature will only shift the values of the derived optical depths by a constant, but will not change the overall appearance of the extinction caused by its dust in a PN.
### 4.7 Dust distribution
All five PNs in this study show that the H$`\alpha `$ emission peaks and the radio continuum emission peaks are in the same regions. In four of them, the peaks of the dust extinction also reside in the same regions. These results suggest that the ionized gas and the dust have similar distributions, so at least part of the dust is mixed with the ionized gas.
On the quantative scale, the faintest parts of the optical image correlate with the peaks of the extinction distribution, as is expected. As an independent check, Kastner et al. (2002) compared the X-ray emission map of BD+303639 with its extinction map derived from optical and infrared images, and showed that the peak of X-ray emission also lies in the region with minimum extinction.
We note that the radio images, H$`\alpha `$ images, and the dust optical depth maps only represent the column density distribution along the line of sight. In order to remove projection effects and reconstruct the three dimensional ionized gas and dust distributions in PNs, we need to develop a 3-D model. The details of the model analysis are presented in the next section.
## 5 The Three Dimensional Dusty Ellipsoidal Shell (DES) Model
While the classification of PN morphology is based on appearance, this approach is vulnerable to projection effects. For example, elliptical and bipolar nebulae can appear round when the nebulae are viewed pole-on. In order to study their intrinsic properties, there have been many attempts to find the detailed three-dimensional structures of PNs (Khromov & Kohoutek, 1968). Since many PNs have the appearance of limb-brightened elliptical rings that are brightest on their minor axes, it is natural to consider models in which the nebulae have an elongated, hollow shell structure. Models with this simple geometry, such as the prolate ellipsoidal shell discussed by Masson (1989, 1990), are successful in reproducing many of the shapes observed in PNs. This was followed by the work of Aaquist & Kwok (1996) who developed a prolate ellipsoidal shell (PES) model to fit the radio images of PNs. Their model contains a spherical shell with both radial and latitude-dependent density gradients, and the shell is ionized by a central star to various depths at different latitudual directions. The ionized shell, when viewed at different angles, results in a variety of morphologies seen in PNs. The parameters derived from the PES model reveal the asymmetric distribution of the ionized gas in the nebula. The PES model was later called the ES model and explored extensively by Zhang & Kwok (1998) to investigate the statistical distributions of the asymmetry parameters. However, because of a typographical error in the source code, there are some uncertainties in the derived model parameters in Zhang & Kwok (1998).
In order to simulate dust extinction effects, we have modified the ES model to include a dust distribution with a constant gas-to-dust ratio, and we refer to this model as the dusty ellipsoidal shell (DES) model. The DES model has been used to analyze the observational images, including radio continuum, H$`\alpha `$, and dust extinction maps to constrain the ionized gas and dust distributions in compact PNs. These different images of the same object provide a way to study the distribution of dust in PNs, and also give better constraints on the model parameters than a single radio or optical image would. One advantage of using the model is that it helps to remove projection effects and to construct the three-dimensional structures of PNs. The DES model will give a quantitative classification of PNs morphology rather than the standard descriptive terms such as round, elliptical, or bipolar.
### 5.1 The Geometry of the DES model
The geometry of the DES model is shown in Figure 4. A single (hot) star is located inside a shell of ionized gas. The inner region of the nebula is assumed to be empty, supposedly swept clean by a wind from the central star. For simplicity, the empty region is assumed to be a prolate or oblate spheroid, depending on the choice of polar radius a and equatorial radius b. The shape of the shell is the result of the interacting winds process, where either the slow AGB wind is equatorially enhanced, or the fast central star wind is non-isotropic (e.g. collimated), or both. The distance from the star to the inner edge of the gas $`R_s`$ can be calculated from a, b and the polar angle $`\theta `$. The UV photons from the central star ionize the surrounding gas out to the distance $`R_i`$, which depends upon $`R_s`$ and the density structure of the surrounding nebula.
The model assumes the central star is emitting ionizing photons isotropically, and the ionized nebula is ionization-bounded along any radial directions from the central star. The depth of the ionization front along any radial vector is given approximately by
$$\frac{dL}{d\mathrm{\Omega }}=\alpha _B_{R_s}^{R_i}n^2(r,\theta ,\varphi )r^2𝑑r,$$
(7)
where $`dL/d\mathrm{\Omega }`$, the number of local ionizing photons ($`L`$) per unit solid angle ($`\mathrm{\Omega }`$), is set equal to the total number of recombinations within the local ionized region. In this equation, $`\alpha _B`$ is the recombination coefficient excluding captures to the ground level (Case B approximation), $`n(r,\theta ,\varphi )`$ is the density structure of the swept-up shell, and $`r,\theta ,\varphi `$ are the spherical coordinates. $`R_i`$ is a function of $`\theta `$ and $`\varphi `$ that depends on the density distribution. For simplicity, a separable density law is assumed, so that $`n(r,\theta ,\varphi )=n_0\eta _r\eta _\theta \eta _\varphi `$; here $`n_0`$ is the density at the inner edge, where $`(r,\theta ,\varphi )=(R_s,0,0)`$, and $`\eta _r,\eta _\theta ,\eta _\varphi `$ are dimensionless quantities describing the density relative to $`n_0`$:
$`\eta _r`$ $`=`$ $`(r/R_s)^\gamma `$ (8)
$`\eta _\theta `$ $`=`$ $`f_\theta \times \left\{\begin{array}{cc}(1\beta )\left({\displaystyle \frac{2\theta }{\pi }}\right)^\alpha +\beta ,\hfill & \hfill \text{if }0\theta {\displaystyle \frac{\pi }{2}}\\ & \\ (1\beta )\left({\displaystyle \frac{2\pi 2\theta }{\pi }}\right)^\alpha +\beta ,\hfill & \hfill \text{if }{\displaystyle \frac{\pi }{2}}\theta \pi \end{array}\right\}`$ (12)
$`\eta _\varphi `$ $`=`$ $`f_\varphi .`$ (13)
In this formulation, $`\gamma `$ gives the logarithmic radial density gradient, $`\beta `$ gives the pole-to equator density ratio, $`\alpha `$ governs the shape of the angular density function, and $`f_\theta ,f_\varphi `$ are additional asymmetry factors:
$`f_\theta `$ $`=`$ $`\delta _\theta +(1\delta _\theta )exp[\mathrm{cos}(\theta \theta _0)1];0\delta _\theta 1`$ (14)
$`f_\varphi `$ $`=`$ $`\delta _\varphi +(1\delta _\varphi )exp[\mathrm{cos}(\varphi \varphi _0)1];0\delta _\varphi 1,`$ (15)
where these control asymmetry in the latitude ($`\theta `$) and longitude ($`\varphi `$) directions, respectively. These two asymmetry factors are used to produce the departure from axisymmetry seen in many PNs. For simplicity, the phases $`\theta _0`$ and $`\varphi _0`$ are both assumed to be $`0^{}`$, so that $`f_\theta `$, $`f_\varphi `$ are maximized (=1) for $`\varphi =0^{}`$ and $`\varphi =0^{}`$, and minimized for $`\varphi =180^{}`$ and $`\varphi =180^{}`$.
With the given radial density law $`\eta _r`$, eq. 7 can be integrated to give a simple analytical solution:
$`{\displaystyle \frac{dL}{d\mathrm{\Omega }}}`$ $`=`$ $`\alpha _B{\displaystyle _{R_s}^{R_i}}(n_0\eta _r\eta _\theta \eta _\varphi )^2r^2𝑑r`$ (16)
$`=`$ $`(n_0\eta _\theta \eta _\varphi )^2\alpha _B{\displaystyle _{R_s}^{R_i}}r^{2\gamma +2}R_s^{2\gamma }𝑑r`$
$`=`$ $`(n_0\eta _\theta \eta _\varphi )^2\alpha _BR_s^3{\displaystyle \frac{(R_i/R_s)^{2\gamma +3}1}{2\gamma +3}}.`$
Once the parameters describing the inner boundary, the density law ($`\gamma ,\eta _\theta ,\eta _\varphi `$), and the equatorial axis shell thickness $`\mathrm{\Delta }b=[R_iR_s]_{\theta =90^{}}`$ are chosen, the outer ionized boundary $`R_i(\theta ,\varphi )`$ can be calculated from eq. 16.
The ionized shell is then projected onto the plane of the sky at different inclination angles i. For $`i=0^{}`$, the nebula is seen pole-on, i.e., the line of sight is perpendicular to the equatorial plane. The edge-on case has $`i=90^{}`$, with the line of sight in the equatorial plane. The projected emission from the ionized region is calculated by assuming that the intensity of the recombination H$`\alpha `$ line and the radio continuum emission is proportional to the integration of the square of the density function along the line of sight (eqs. 815). The calculations were done inside a cubic box containing the nebula.
To simulate optical extinction effect, we assume that the dust is well mixed with the gas in the ionized region. Thus, a constant gas-to-dust ratio is assigned throughout the ionized region, i.e., the amount of extinction is directly proportional to the density function at each integration point. The H$`\alpha `$ emission is then attenuated by the dust extinction using eq. 4.
### 5.2 The Model Images
The DES model allows us to explore how the morphology of a nebula changes with various parameters, especially $`\alpha `$ and $`\beta `$. A series of $`128\times 128`$ pixel model images has been produced with fixed parameters of $`a=20,b=18`$, and $`\mathrm{\Delta }b=6`$. The model images presented here have no radial density gradient, i.e., $`\gamma =0`$. This is a fairly good approximation, because most swept-up shells are thought to be relatively thin, allowing little density change with radius. We have computed the projected images for inclination angles $`i=0^{},30^{},60^{}`$, and $`90^{}`$ with different combinations of $`\beta `$ and $`\alpha `$ as shown in Figures 58.
Figure 5 shows the model images with emission calculated from the integration of the square of the density function along the line of sight. This simulates the radio continuum images of the long-wavelength free-free emission not affected by dust.
Figure 6 displays the input dust distribution. The images are computed from the integration of the amount of extinction along the line of sight all the way through the nebula. Because the gas-to-dust ratio is constant inside the ionized region, the dust distributions appear similar to the model radio images in Figure 5. If the dust temperature is uniform, and the cloud is optically thin for the dust thermal emission, the flux emitted by the dust is directly proportional to the absorption coefficient. Therefore, these dust distributions can mimic infrared images of the dust emission.
Figure 7 gives the model images when the dust extinction is taken into account. As in Figure 5, emission is calculated from the integration of the square of the electron density, but now it is attenuated by the dust in front of it along the line of sight. This simulates the optical H$`\alpha `$ images, which, unlike the radio continuum maps, are affected by dust. Some images clearly show an asymmetric appearance caused by dust extinction (e.g., the image of $`\alpha =1,\beta =0.1,`$ and $`i=30^{}`$).
Figure 8 displays output extinction optical depth maps, calculated as the natural logarithm of the ratio of the model radio and optical images (eq. 5). Because only dust in front of the emission can attenuate the light, the amount of dust in the extinction maps (Fig. 8) is less than the actual amount of dust in the nebulae (Fig. 6). Comparing these images with the input dust distributions in Figure 6, the amount of missing dust for different inclination angles and morphologies can be estimated. By comparing the radio and H$`\alpha `$ images, the nebular geometry can be estimated and the asymmetry effect caused by dust for this viewing geometry can be removed and reveal any intrinsic asymmetry of the gas distribution.
## 6 Comparison of Observed and Model Images
In order to constrain the parameters that describe the geometry and angular density distribution of the PNs presented, a set of model images were made to compare to each individual nebula’s observed radio continuum and H$`\alpha `$ images, extinction maps, and in one case (BD+303639), an infrared dust emission image (from Volk & Kwok, 2003). The comparison was made by visual inspection. The faint emission from the outer parts of each nebula was chosen to match the observed H$`\alpha `$ images because their S/N ratios are much higher than those for the radio continuum images. For each object, a set of parameters was found that generate model images matching most of the observed images. The images comparing the best models to the observations for these five nebulae are shown in Figures 913. The parameters used to generate the model images are listed in Table 4.
### 6.1 BD+303639
The model of BD+303639 has a relatively small equatorial thickness $`\mathrm{\Delta }b`$ compared to the size of the inner empty region, which has a prolate shape with an axial ratio of $`a/b=1.2`$. The small value of $`\beta (=0.1`$) suggests that it has a large equatorial to polar density ratio. A small asymmetry $`\delta _\varphi =0.9`$ in the longitude direction is added to explain the unequal intensity of the two sides observed in the radio image. Dust well-mixed with ionized gas is able to reproduce the anti-correlation of the optical image and extinction map. Due to its low inclination angle ($`i=15^{}`$), it has a round appearance. This low inclination is consistent with the value of $`i10^{}`$ found by Li, Harrington & Borkowski (2002), who also fit the radio image and the spectra by approximating the nebula as an ellipsoid. But a comparison with the model image in Fig. 7 (second column from the right) suggests that BD+303639 would appear as a prominent bipolar nebula if viewed edge-on.
### 6.2 IC 5117
The model of IC 5117 also has a relatively small equatorial thickness compared to the inner empty region, with a prolate axial ratio $`a/b=1.67`$. It has a small pole-to-equator ratio of $`\beta =0.1`$. No asymmetry factor in the longitude direction is needed. The brighter region toward the east in the optical image and the higher extinction toward the west can be reproduced by a dust distribution with a constant gas-to-dust ratio. The model images have an intermediate inclination angle of $`i=40^{}`$.
### 6.3 M 1-61
The model of M 1-61 has an equatorial thickness comparable to the dimension of the inner empty region, which has a prolate axial ratio of $`a/b=1.2`$. It has a small pole-to-equator ratio of $`\beta =0.2`$. A small asymmetry $`\delta _\varphi =0.9`$ in the longitude direction is needed to explain the mildly brighter western peak in the radio image. The anti-correlation of the optical bright peak in the west and the strong extinction peak in the east can be reproduced by the well-mixed dust. However, it is not as prominent as the observed images. The model images has a high inclination angle of $`i=70^{}`$, and they suggest that a higher dynamic range image of M 1-61 will reveal a pair of bipolar lobes.
### 6.4 M 2-43
The model of M 2-43 also has a relatively small equatorial thickness compared to the inner empty region with a prolate axial ratio $`a/b=1.67`$. It has a relatively small pole-to-equator ratio of $`\beta =0.3`$. A mild asymmetry $`\delta _\varphi =0.7`$ in the longitude direction is added to reproduce the brighter southern peak in the radio image. Although a simple constant gas-to-dust ratio is able to generate an anti-correlation between the H$`\alpha `$ emission and optical depth, it cannot reproduce optical and extinction maps that match the observed images well. The model images have an intermediate inclination angle of $`i=45^{}`$.
### 6.5 M 3-35
The model of M 3-35 has an equatorial thickness comparable to the dimension of the inner empty region, which has an oblate axial ratio of $`a/b=0.83`$. This is the only object that requires an oblate model. It has a intermediate pole-to-equator ratio of $`\beta =0.4`$. No asymmetry factor in the longitude direction is needed to reproduce the radio image. No obvious anti-correlation between optical and extinction maps can be found. The model images are highly inclined to the line of sight, with $`i=80^{}`$. As in M 1-61, a higher dynamic range image is expected to reveal the bipolar lobes in the NE-SW directions.
## 7 Discussion
As we can see from Figures 913, the model images can generally simulate the observed images. The remaining differences could be accounted for by several effects. For example, the assumed model geometry may be too simple to reproduce all of the complex small-scale structure of PNs. Other causes could include clumpy dust distributions that do not follow a simple constant gas-to-dust ratio. Three out of five PNs show an asymmetric density distribution in the longitudinal direction, a departure from axisymmetry. As summarized by Soker & Rappaport (2001), four main processes can result in the departure from axisymmetry in PNs: interaction with the interstellar medium, mass-loss enhancements due to cool starspots, or a wide or a close binary companion. They investigated the effects of companions and estimated that about 27% of all PNs are expected to acquire non-axisymmetric structure from binary interactions. Nevertheless, the effects of the ISM are also worth considering. Since the majority of PNs are distributed near the Galactic plane where the mean density of the ISM is higher, significant interaction between PNs and the ISM is expected. Such phenomena can be seen from the disturbed emission on the outskirts of some nebulae (e.g., Xilouris et al., 1996) and from some apparently off-centered nuclei (e.g., Sh 2-216, Tweedy et al. 1995; KFR 1 and MeWe 1-4, Rauch et al. 2000).
The electron densities of these five PNs are all found to be higher near their equators than near their poles. The small pole-to-equator density ratios ($`\beta `$) of some nebulae indicate they would show a bipolar morphology if seen edge-on. The fraction of intrinsic bipolar PNs therefore are likely to be higher than suggested by classification by apparent morphology. Since the morphology of PNs has been proposed to relate to their physical properties, (e.g., bipolar nebulae have been suggested to correlate with nebulae with more massive central stars; Stanghellini, Corradi, & Schwarz, 1993), it is therefore important that we have a good understanding of the intrinsic strucures of PNs. By using models such as DES, we can more reliably estimate the fraction of bipolar nebulae and therefore assess any correlation with physical properties of the PNs.
## 8 Conclusions
We have derived the dust distribution in five compact PNs from their H$`\alpha `$ and radio images. From this study, we are able to determine the extent of contribution by dust extinction to the apparent departure from axi-symmetry in these PNs. By comparing these results with simulations generated by the ellipsoidal shell model, we can infer the intrinsic morphologies of the PNs. This method gives a more quantitative approach to the analysis of PN morphology than visual inspection of apparent morphology, and is therefore more capable to deducing proper statistics on the distribution of morphological types of PNs.
We thank Kate Y.L. Su, Kevin Volk, Orla Aaquist, Steven Gibson for helpful discussions. Some figures were made with the software package WIP (Morgan 1995). We are grateful to R. Gooch for his expansion of the capabilities of the Karma visualization software package (Gooch 1996)<sup>1</sup><sup>1</sup>1See also http://www.atnf.csiro.au/karma. This work was supported in part by grants to SK from the Natural Sciences and Engineering Research Council of Canada and by the National Science Council of Taiwan. SK acknowledges the award of a Killam Fellowship from the Canada Council for the Arts.
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# Extinction law variations and dust excitation in the spiral galaxy NGC 300
## 1 Introduction
NGC 300 is a gas-rich galaxy of Sd type, at a distance of only 2.1 Mpc (Freedman et al., 1992). Its total extent in the blue light is over 20′, and its inclination on the line of sight, derived from H I kinematics, is about 50° (Puche et al., 1990). It contains numerous bright H II complexes (Sérsic, 1966; Deharveng et al., 1988) and forms stars at a current total rate of about 0.15 M yr<sup>-1</sup> (Helou et al., 2004). NGC 300 was targetted for surveys of supernova remnants ; between 20 and 40 candidates were found by various selection methods (Blair & Long, 1997; Pannuti et al., 2000; Payne et al., 2004), evidence that star formation has been going on for an extended period at a significant rate. Large numbers of Wolf-Rayet stars were also found (Breysacher et al., 1997; Schild et al., 2003). The oxygen abundance in the nebular gas phase varies between about 1.4 to 0.4 times the solar neighbourhood value, from the center up to a radial distance of about 10′ (Deharveng et al., 1988).
Many bubble-like structures are seen in the H$`\alpha `$ emission, especially in the outer disk, where they reach diameters of $`600`$ pc. About 50% of the total H$`\alpha `$ flux arises from the diffuse ionized gas component (Hoopes et al., 1996). Judging from the morphology of H II regions and the position of the most likely exciting sources, the hydrogen gas may be ionized very far away from young clusters (in projected distance). This argues for the existence of aged star-forming sites and a high degree of porosity of the interstellar medium around them. The large number of resolved stellar clusters and the variety of their associated H II regions, observed at a resolution of a few tens of parsecs, make it possible to study in great detail general patterns in the extinction and dust heating within a single galaxy.
Deriving spatially-resolved extinction in galaxies of different types, and determining how the extinction law is affected by properties of the interstellar medium such as metallicity and geometric configuration, is a first step toward assessing the validity of common recipes used to correct for extinction the global ultraviolet emission of galaxies. This knowledge is fundamental to obtaining good prescriptions to estimate star formation rates and ages from current and future ultraviolet and infrared surveys.
In this paper, we analyse jointly images of the stellar light from the far-ultraviolet to the far-red, of the hydrogen recombination emission in the H$`\alpha `$ and H$`\beta `$ lines, and of the infrared emission from 3 to 160 $`\mu `$m. The data are described in Section 2. After deriving large-scale infrared-to-ultraviolet ratio and extinction maps (Section 3), we focus on individual star-forming complexes selected in H$`\alpha `$ and in the ultraviolet, and constrain the fundamental properties of ionizing stellar clusters, to assess how critically the extinction law and excitation of mid-infrared emitting dust at these spatial scales depend on the local geometry of the interstellar medium (Section 4). The results are discussed topic by topic in Section 5, and summarized in Section 6.
## 2 Data
### 2.1 Ultraviolet images
NGC 300 was observed with the GALEX space telescope in the FUV and NUV bands at effective wavelengths of 1516 and 2267 Å (Martin et al., 2005; Morrissey et al., 2005). GALEX has a circular field of view of 1.2° in diameter, and observations are performed in a spiral dither pattern. The data from the photon-counting, multichannel-plate detectors were processed with the analysis pipeline operated by the Caltech Science Operations Center, which converts the photon lists into flux-intensity images and applies flat-fielding and first-order geometric distortion corrections. The absolute calibration is derived from ground test data which show good agreement, to within 10%, with observations of white dwarf standard stars (Morrissey et al., 2005). See also the GALEX Atlas of galaxies (Gil de Paz et al., 2005), where NGC 300 is included, for more information.
The angular resolution is of the order of 5″ (FWHM), but varies with source brightness and position across the field. In particular, the point spread function in the FUV band is slightly elliptical (elongated along the south-east to north-west direction) in some parts of the map. A first distortion correction between the NUV and FUV maps was estimated using Legendre polynomials of order 4 and cross-correlating the positions of about 900 stars and clusters. Distortion residuals and point spread function variations were handled separately for each individual sub-field of size $`2`$′, by fitting nearby point sources with gaussians in the FUV, NUV and U bands.
### 2.2 Infrared images
Broadband infrared images were acquired with the Spitzer space telescope at effective wavelengths of 3.6, 4.5, 5.7, 8.0, 24, 71 and 156 $`\mu `$m, and were first discussed by Helou et al. (2004). The data processing and mosaicing is also described in this paper. The angular resolution (full width at half maximum) of the mosaics is 2.5 to 3″ between 3.6 and 8 $`\mu `$m, 5.7″ at 24 $`\mu `$m, 16″ at 70 $`\mu `$m and 38″ at 160 $`\mu `$m.
The mid-infrared bands record gradually varying contributions, as a function of wavelength, from stars and dust, the global 3.6 $`\mu `$m emission being dominated by photospheres, the global 8 $`\mu `$m emission deriving mainly from aromatic band carriers, and longer-wavelength bands covering essentially the thermal continuum from very small to big dust grains.
The photospheric component was subtracted from the 8 $`\mu `$m map as explained in Helou et al. (2004). A model for intermediate to old stellar populations was derived from the Starburst99 population synthesis model (Leitherer et al., 1999), and was found to be remarkably invariant for solar and half-solar metallicities, and a broad range of star formation histories (excluding very young populations). Assuming that the 3.6 $`\mu `$m emission is entirely of stellar origin, this map was then used to normalize the photospheric spectrum at each point, to which the relevant filter transmission curves were applied. This procedure may be inaccurate locally, because of the contribution from young stars, but the 8 $`\mu `$m emission is heavily dominated by dust in star-forming regions, so that neglecting young stellar populations should not introduce any significant error.
### 2.3 Optical images
Mosaics in the UBVI bands of $`20.5\mathrm{}\times 20.5`$′ were obtained from S.C. Kim (Kim et al., 2004, private communication). An R-band mosaic of $`25\mathrm{}\times 25`$′ and H$`\alpha `$+\[N II\] and H$`\beta `$ maps in narrowband filters (in a field of view of $`8.85\mathrm{}\times 8.85`$′) were acquired at the Las Campanas du Pont 2.5 m telescope with the direct CCD camera. The angular resolution in all the broadband optical images, after the mosaicing and the geometric transformations applied to align them, is of the order of 1.7″ to 2″ (FWHM).
We also made use of an H$`\alpha `$ mosaic covering the same field as the R-band mosaic to measure the H$`\alpha `$ flux of the outermost H II regions ; the residual background was subtracted locally in each subfield separately. The map acquired by Hoopes et al. (1996), which was kindly provided to us by Charles Hoopes, supplied a means of checking the H$`\alpha `$+\[N II\] flux calibration, by comparing aperture photometry measurements on several bright H II complexes in both maps ; fluxes from the map of Hoopes et al. (1996) are systematically lower by about 10%.
The H$`\beta `$ photometric accuracy, however, is much more uncertain, because this line has a much smaller equivalent width than the H$`\alpha `$ line, and larger errors are introduced by continuum subtraction and correction for stellar absorption features. Because of these uncertainties, if no flux correction is applied, the H$`\alpha `$/H$`\beta `$ decrement values are unphysical in the case B recombination assumption, and in disagreement with published values for individual H II regions. In order to correct them and derive an estimate of the nebular extinction from H$`\alpha `$ and H$`\beta `$ photometry, we scaled the decrement in such a way that extinction values derived by Webster & Smith (1983) from spectro-photometry for 5 regions could be reproduced. We used an intrinsic H$`\alpha `$/H$`\beta `$ ratio of 2.86, appropriate for electronic densities and temperatures $`10^3`$ cm<sup>-3</sup> and $`10^4`$ K, respectively, and replaced the extinction law assumed by Webster & Smith (1983) with the Cardelli et al. (1989) law. This procedure requires that the raw H$`\alpha `$/H$`\beta `$ flux ratio be underestimated by 30% on average. Applying the above correction, the derived global extinction (from total fluxes in the H$`\beta `$ field of view) is A(H$`\alpha `$) = 0.5 mag.
## 3 Large-scale UV/IR ratio and extinction
The ultraviolet disk reaches much farther out than the infrared disk. To quantify this statement, we convolved all maps to the resolution of the 160 $`\mu `$m map, and estimated the total infrared emission (TIR) from a linear combination of the 24, 70 and 160 $`\mu `$m powers (Dale & Helou, 2002). To reach the 160 $`\mu `$m full width at half-maximum, the maps at wavelengths up to 24 $`\mu `$m were first convolved with a model of the point spread function at 70 $`\mu `$m<sup>1</sup><sup>1</sup>1which can be found on the website: http://ssc.spitzer.caltech.edu/mips/psf.html . ; these and the 70 $`\mu `$m map were then convolved with gaussians of the appropriate widths. We computed surface brightness profiles in the FUV band and in the TIR, averaged in elliptical annuli to account for the inclination of NGC 300 (Fig. 1). The FUV band was preferred over the NUV band for this comparison, because while the ratio of the extinction in these two bands is close to 1 (smaller than 1.3 for various extinction laws), the intrinsic FUV/NUV ratio of young stellar populations is $`1`$, the FUV band avoids the variable 2175 Å feature of the interstellar extinction curve (see Section 5.4), and the FUV band is minimally contaminated by foreground stars. On average, the FUV/TIR ratio increases exponentially with deprojected distance from the center, with irregularities mainly due to the FUV emission being less smooth than the infrared emission. The FUV/TIR ratio increases by at least a factor 5, maybe up to a factor 10, between the center and the outermost regions where it is still measurable. Also shown in Figure 1 is the opposite of the average slope of the metallicity gradient as derived by Deharveng et al. (1988). The slope of the FUV/TIR gradient is very similar, although somewhat smaller. At any given distance from the center, however, the dispersion in FUV/TIR ratios among individual star-forming regions is at least a factor 2, whereas the standard deviation of abundances is only of 0.11 dex, which seems to indicate that metallicity effects are not solely responsible for the FUV/TIR gradient. The UV-brightest regions tend to be deficient in infrared emission compared with their fainter counterparts. From studies of large samples of galaxies, a clear decrease of the global ultraviolet to infrared flux ratio was found as a function of a measure of the star formation rate (Wang & Heckman, 1996), or as a function of the blue luminosity (Buat & Burgarella, 1998) or the stellar mass (Pierini & Möller, 2003). Interpreting the ultraviolet to infrared ratio variations within resolved galaxies could help constrain the underlying causes of this behavior.
A stellar extinction map of NGC 300 can be derived from a simple energy balance calculation, using all the maps between 1516 Å and 3.5 $`\mu `$m to reconstruct the stellar emission. The emerging fluxes in the FUV, NUV, UBVRI and IRAC1 (3.5 $`\mu `$m) bands were interpolated by realistic stellar spectral energy distributions, using the actual filter transmission curves. Assuming several different extinction laws, the power absorbed in this range is set equal to the total power emitted in the infrared. The derived A(FUV) ranges between 0.1 and 0.6 (A(V) between 0 and 0.23). Uncertainties of 20% in the flux calibration at both 70 and 160 $`\mu `$m could increase the A(FUV) extinction values by 0.05 mag.
The large H II complexes including shell-like structures correspond to regions of lower extinction, consistent with the idea that they are aged star-forming sites having partly dispersed the molecular and dust material around them. The correspondance with gas column density cannot be tested because no map in a molecular transition is available. The H I map of Puche et al. (1990), with a beam size of 17″, allows us to see partial correspondance with the H I column density, though, especially in the outer disk (Fig. 5) ; the fact that there is no robust correlation with the extinction map is expected because H I does not trace adequately the total gas mass, which is usually dominated by the molecular phase in the inner parts of the disk where star formation is taking place actively.
The radial profile of the ratio of the total stellar power (integrated between 1000 Å and 3.5 $`\mu `$m) to the total infrared power is notably different from the FUV/TIR profile: it shows a sharp decline from the center up to the pseudo-ring of most active star formation, reaching a minimum around a radial distance of 3 to 4′, then rises again to reach its central value at a radial distance close to 10′.
Several effects can contribute to the large-scale FUV/TIR gradient.
1) The most obvious possibility is the variations in metallicity, which are expected to cause direct variations in the dust to gas ratio. However, various lines of evidence suggest that other parameters play a more fundamental role. Although the global dust to gas ratio of galaxies seems correlated with the far-infrared surface density (Andreani et al., 1995), which traces both the column density of dust and the star formation rate surface density, the dust to H I mass ratio of irregular galaxies is highly variable but not correlated with metallicity (Hunter et al., 1989). Sauvage et al. (1990) also showed that the far-infrared to blue flux ratio of dwarf galaxies is not correlated with their oxygen abundance. Within NGC 300, the nebular gas phase metallicity gradient, as determined by Deharveng et al. (1988), is slightly steeper than the FUV/TIR gradient. Interestingly, Bot et al. (2004) found that the dust to gas ratio of the Small Magellanic Cloud is much lower than expected from a linear scaling with metallicity, although the situation may differ between star-forming regions and the diffuse interstellar medium.
2) The large-scale variations of the FUV/TIR ratio may also be affected significantly by the cumulative effects of local variations in the geometry of the sources. Several factors can contribute to systematically increase the distance between stellar clusters and the absorbing dust, notably: a decrease in the volume density of all the phases of the interstellar medium ; the maturity of the star forming regions, i.e. the outer disk could contain slightly more evolved young stellar populations which have had time to disperse gas and dust away. Aged star formation complexes indeed seem to efficiently increase the porosity of the surrounding interstellar medium, because they are associated with diffuse, shell-like H II regions, and are affected by lower extinction.
3) The spectral energy distribution of the heating radiation could change in such a way that the total dust heating efficiency decreases. The fraction of the total stellar power emitted in the FUV band varies from less than 2% in the central region of the galaxy to about 5% in the outermost part of the disk, at and beyond a radial distance of 10′ (the whole SED being corrected for extinction). The non-ionizing part of the stellar radiation thus becomes on average more energetic in the outer disk (the old populations have a smaller scale length than the younger populations emitting in the ultraviolet), and the dust heating balance should shift progressively from big grains emitting in the far-infrared to smaller grains emitting in the mid-infrared, which could explain part of the far-infrared flux deficit. We indeed observe a slight increase of the $`F_{24}/\mathrm{TIR}`$ fraction in the outer disk (by about 30-40% compared with the annulus of intense star formation). However, this increase is too modest to account for the FUV/TIR gradient, which is mirrored closely by a similar $`\mathrm{FUV}/F_{24}`$ gradient.
Figure 2 shows the radial profile of the FUV/NUV energy ratio, which decreases on average beyond a radial distance of $`10`$′. Such a decrease of the FUV/NUV ratio in the outermost part of the disk cannot be explained by extinction or metallicity effects. It thus seems reasonable to assume that a gradient exists in the age of the youngest stellar populations which radiate in the ultraviolet. The range of FUV/NUV colors observed at large scales corresponds to average stellar population ages between about 16 Myr at radial distances near 10′, and $`300`$ Myr in the outermost disk.
In order to assess the role played by the geometry of the interstellar medium in the FUV/TIR variations, and to distinguish between the causes mentioned above, we now turn to individual H II complexes showing a wide range of morphologies, presumably corresponding to very different geometries. We take advantage of the high spatial resolution of the data from the UV to the near-infrared, and characterize individual stellar clusters or small groups of clusters. If each one consists of coeval populations, then it is possible to lift the star-formation history degeneracy that affects the large-scale emission.
## 4 Small-scale properties of H II complexes
### 4.1 Selection of H II regions/young stellar clusters
We selected individual stellar clusters or groups of clusters associated with prominent H II regions, that are visible in all the optical bands and sufficiently bright in the FUV and NUV bands to allow accurate photometry. Because of the selection method, the sample is biased against highly obscured regions, for which the analysis described in Section 4.3 would not be possible. Bright 24 $`\mu `$m sources which are not detected in the ultraviolet are discussed in Section 5.7. We also avoided objects with very complicated structure (making the identification of the clusters ionizing a particular HII region ambiguous), with contamination by foreground stars or blending with nearby regions. The sample covers a large range of morphologies, mid-infrared colors, and UV to infrared ratios. They were labelled according to the H II region which they excite, in the catalog of Deharveng et al. (1988) (but the letter sublabels follow a different convention). Detailed notes, pertaining both to the morphology and to the modelling results, can be found in Appendix A. Images of the cluster fields in the U band, in H$`\alpha `$ and at 8 $`\mu `$m are shown in Figure 6.
Table 1 lists, for each group of clusters, the initial diameter of the stellar photometric aperture (see Sect. 4.2) and some observables: the H$`\alpha `$ flux (corrected for an average \[N II\] contribution to the total H$`\alpha `$+\[N II\] flux of 10%), the nebular extinction derived from the H$`\alpha `$/H$`\beta `$ decrement, and the H$`\alpha `$/R ratio measured within the stellar photometric area, R referring to line-subtracted broadband flux (i.e. pure stellar emission). H$`\alpha `$ equivalent widths can be obtained by multiplying H$`\alpha `$/R by 1262 Å (the equivalent width of the R filter). We also measured the $`F_{24}/F_8`$ ratio at the peak of the 24 $`\mu `$m emission (which can be offset from both the exciting clusters and the peak of the H$`\alpha `$ emission), within a resolution element (2.5 times the 24 $`\mu `$m FWHM, except for confused regions for which the aperture was reduced to the FWHM in order to avoid neighboring sources), after convolving the 8 $`\mu `$m map to the 24 $`\mu `$m angular resolution.
Table 1 also indicates the association of the H II regions with supernova remnants in the catalogs of Pannuti et al. (2000), Payne et al. (2004) and Blair & Long (1997). We fail to see any direct connection between the presence of a supernova remnant and the morphology of the H II region, as some of these are classified here as compact (see in Section 5.4 the discussion about a quantification of the compactness of the H II regions, and the significance of this parameter in interpreting the results). Among these catalogued remnants, weak X-ray sources are found close to D40, D45, and D53 ; the radio sources are all very weak, and some have flat radio spectral indices, consistent with thermal emission. Such is the case for the regions D61, D77, D84. The influence of these candidate young supernova remnants on the surrounding interstellar medium is thus unclear, and the shell-like structures visible in the H$`\alpha `$ emission may correspond to older events.
### 4.2 Photometry
Special care was devoted to refine the photometric method so as to obtain accurate spectral energy distributions for the young clusters. In order to remove the local background from old populations, we devised a method appropriate to the combination of UV and optical data. We had to take into account the facts that the underlying emission can vary significantly on the scale of the considered regions, that its brightness relative to that of the clusters is variable from one band to the other, increasing with wavelength and being very high in the R and I bands, and that confusion makes any simple aperture always contaminated by old stellar populations and foreground stars. We define two diameters, constraining the source to be contained within the smaller ($`D_1`$) and estimating the background within the larger ($`D_2`$). Using the U band, where the background and contamination by old populations are minimal, we then determine refined photometric “apertures” in the following way. The average brightness $`b`$ and standard deviation $`\sigma `$ of the background in $`D_2`$ are computed iteratively in a narrowing brightness interval ($`\pm 2\sigma _\mathrm{i}`$ where $`\sigma _\mathrm{i}`$ is the standard deviation at each step) until convergence is achieved. Then, the source flux is summed over the pixels above $`b+3\sigma `$ within $`D_1`$ and the background averaged over the pixels below $`b+3\sigma `$ within $`D_2`$ (not excluding $`D_1`$). Then, we use the exact same pixels to compute the source flux and background in all the other optical bands, which have been previously aligned to within 0.2 arcsec, and control the results visually to ensure that no visible foreground star contaminates the aperture, and that $`D_1`$ and $`D_2`$ were appropriately chosen. We find that this method produces significantly better results than by measuring the background in a concentric annulus or another aperture centered on a nearby empty patch. For clusters that appear single in the optical bands, we obtain good fits with a single-population model, which is not the case when we use the alternative methods above.
Since the point spread function of the GALEX images varies across the field of view and is slightly elliptical in the FUV band, and because of residual distortion between the FUV and NUV images, we had to refine the photometric technique for the GALEX bands. In each small field around a selected cluster, we use point sources to fit positional offsets with respect to the optical images, and we determine the FWHM of the cluster independently in the FUV and NUV bands. In a few instances, where the selected cluster is partially confused by nearby clusters, we fitted and removed the latter before performing aperture photometry on the target. In a few cases when we have to truncate the aperture to avoid confusion, we apply an aperture correction, computed assuming that the source is perfectly gaussian. The measured broadband fluxes are listed in Table 2.
### 4.3 Stellar population fits
We used the Starburst99 synthesis model (Leitherer et al., 1999) to compute spectra of instantaneously-formed stellar clusters, with the metallicity $`Z=0.008`$, well suited to NGC 300, and a Salpeter initial mass function between 0.1 and 120 M.
The transmission curves of the filters used in the observations were applied to the model spectra before doing any comparison between the models and the measured fluxes. For the optical filters, we used the definitions of the Johnson-Cousins system given by Fukugita et al. (1995) and the references therein. For the ultraviolet and infrared filters, we used the actual transmission curves of GALEX<sup>2</sup><sup>2</sup>2available from the website: http://galexgi.gsfc.nasa.gov/tools/Resolution\_Response/index.html . and Spitzer<sup>3</sup><sup>3</sup>3available from: http://ssc.caltech.edu/irac/spectral\_response.html
and http://ssc.caltech.edu/mips/spectral\_response.html ..
We fit the observed FUV-NUV-UBVRI flux densities using three alternate extinction laws – Cardelli et al. (1989) for the Milky Way, with the ratio of total to selective extinction $`R_\mathrm{V}=A(\mathrm{V})/E(\mathrm{B}\mathrm{V})=3.1`$, appropriate for the diffuse interstellar medium ; Fitzpatrick (1986) for the LMC average and for the 30 Doradus complex – and the attenuation law of Calzetti et al. (1994). The results should be interpreted while keeping in mind that the first three extinction curves were derived toward the line of sight of individual stars, and the last one is in principle only valid for the integrated emission of starburst galaxies, while we are applying them to spatial extractions within regions ranging between 60 and 300 pc in diameter. While Misselt et al. (1999) have shown that the type of extinction law attributed by Fitzpatrick (1986) to 30 Doradus is in fact mostly observed in a region adjacent to 30 Doradus, the LMC2 giant shell, we retain the terminology of Fitzpatrick (1986) for reference to the particular extinction curve that he derived. Updated extinction curves for the LMC have been published by Misselt et al. (1999). Given that we are interested in extracting qualitative variations in the apparent local extinction law (see Section 5.4), and that the discrete extinction curves used here are in fact part of a continuum in dust properties (Gordon et al., 2003), these curves are sufficient for our purposes.
We fit the cluster spectral energy distributions in a simple dust screen configuration, which is implicit in the four extinction curves considered here, i.e. ignoring the effects of scattering and radiative transfer through layers of dust mixed with the stars. These effects are extensively described by e.g. Witt et al. (1992) and Calzetti (2001). Allowing the relative geometry of stars and dust to be unconstrained would introduce insurmountable degeneracy with the fitted parameters of the stellar populations. But there is another motivation for assuming that the stars are not mixed with the dust. Our photometric technique, in effect, is designed to extract the emission from discrete stellar clusters and remove efficiently background or foreground extended emission (Section 4.2). We argue that the stellar clusters can thus be treated as a collection of point sources, each with its own dust shell. This is especially true of the class of clusters that we call compact (see Section 5.4 for a quantitative definition). The smallest full width at half maximum of the images represents a linear distance of 20 pc, to be compared with the full widths at half maximum of young clusters in spiral galaxies as quantified by Larsen (2004), up to 10 pc. Even though the dust screen assumption does not reflect the likely geometry of the sources, we will show in Section 5.4 that it is a good approximation for shell-type sources, for which the inclusion of radiative transfer effects (in particular scattering effects) does not alter our conclusions. Furthermore, the results concerning the variations of the attenuation law (which includes radiative transfer effects, as opposed to the pure extinction law) are independent of any assumptions about the source geometry. Note that scattering should be a negligible component of the Calzetti et al. (1994) law, because light scattered in the line of sight compensates on average light scattered out of the line of sight for regions of large size (several kiloparsec).
The flux calibration uncertainty is assumed to be 10% for all the filters, to which the formal measurement uncertainty is added. The fits were not constrained to reproduce the 3-5 $`\mu `$m data, but only not to overproduce them (within the error bars), even when these bands are dominated by stellar emission. We find in some cases that no single coeval population can reproduce the photometric data, in which case we allow two populations of different ages but same extinction. This occurs for the double cluster D53c and for some complex fields associated with very large H II regions (D137b, D115). All the other regions were fitted with a single population.
In H II regions, the H$`\alpha `$ and \[O II\] lines at 6563 Å and 3727 Å, respectively, contribute a large fraction of the total flux in the R and U bands (up to 50% in R, and up to 40% in U, respectively, within the chosen apertures), except in very diffuse H II regions. The contribution of the H$`\beta `$ and \[O III\] lines can be neglected because they are observed through low-transmission parts of the B and V filters, and are intrinsically weaker. We accounted for H$`\alpha `$ and \[O II\] fluxes by subtracting scaled versions of the H$`\alpha `$ map from the R and U maps, using the \[O II\]/H$`\alpha `$ ratios measured by Webster & Smith (1983) and d’Odorico et al. (1983) and correcting them for differential extinction. Whenever a measurement of the \[O II\]/H$`\alpha `$ ratio is lacking, we adopt the average ratio of the H II regions discussed by Webster & Smith (1983) and d’Odorico et al. (1983) (\[O II\]/H$`\alpha =1.13`$). Because the \[O II\] contribution estimated in this way is uncertain and because of possible aperture mismatch, we affect a smaller weight to the U-band data than to the other bands, recovering an updated \[O II\]/H$`\alpha `$ ratio from the SED fits (Table 3).
The output parameters of the fits are the age(s) of the clusters, their mass(es), and the monochromatic extinction at 1516 Å, the effective wavelength of the FUV filter. The extinction and ages take discrete values, varying in steps of 0.1 mag and 1 Myr, respectively. For each extinction, each age of the young clusters, and optionally each age of older clusters, the masses have to be solved from a set of 7 equations (for the 7 bands from the FUV to I). In order not to bestow a larger weight on any particular combination of bands, we solve separately each equation with the new equation resulting from the sum of all the fluxes, and then we compute the average of the 7 solutions, with a lower weight for the U band. The ionizing photon flux predicted by the best fit is compared with the ionizing photon flux estimated from the H$`\alpha `$ flux, corrected for extinction using the H$`\alpha `$/H$`\beta `$ decrement, but is left unconstrained by the observations. We usually find agreement to within 20%, with a few exceptions (Table 3). To determine the confidence intervals of the cluster parameters, the various fits are only constrained to reproduce the ionizing photon flux from the best fit to within 30% (we set a constraint on the error, but not on the absolute value).
We expect the error on the observed ionizing photon flux to be large, because of uncertainties on the H$`\alpha /\mathrm{H}\beta `$ decrement, because of uncertainties in the extent of the region directly under the influence of the ionizing stars, and because an unknown fraction of the ionizing photons may leak out of the H II regions. The few cases where the agreement between the model value and the observed value is poorer than 20% may be due to local errors in the subtraction of the continuum and stellar absorption from the H$`\beta `$ line, or to increased porosity of the surrounding interstellar medium. One of the regions for which the best-fit model overpredicts the measured ionizing photon flux by a large percentage is D137b (by 40% for the 30 Doradus law, and 100% for the Calzetti et al. (1994) law). This region also stands out as the most peculiar in terms of morphology of the hydrogen gas, suggesting that previous generations of stars in the recent past have injected lots of mechanical energy in the surrounding medium, carving a tunnel-shaped void. It can thus be hypothesized that a sizeable fraction of the ionizing photons are leaking out of the visible H II region. According to studies of the diffuse ionized gas in spiral galaxies, up to 50% of the ionizing photons may escape H II regions (Ferguson et al., 1996). Alternatively, the inclusion of a second population increases the degeneracy of the fits, and the fitted stellar extinction could be too high.
We assess the validity of the different extinction laws by first comparing the $`\chi ^2`$ values of the best fits. But since the absolute uncertainty was fixed arbitrarily to 10%, not taking into account the smaller relative uncertainty between the two GALEX bands, we consider that a $`\chi ^2`$ value close to unity does not necessarily indicate that the fit is acceptable. We also check its ability to reproduce colors, in particular the FUV-NUV color, which is the most useful discriminant between the different extinction laws. The ratio of the extinction in these two bands varies from $`1.00`$ for the Milky Way and $`1.10`$ for the LMC to $`1.27`$ for 30 Doradus and $`1.20`$ for the Calzetti et al. (1994) law.
Whenever some significant interstellar component at 3-5 $`\mu `$m was indicated by the population synthesis model, we confirmed this finding by examining the morphology of the dust and stellar clusters (i.e. by checking that the structure of the 3-5 $`\mu `$m emission is much more similar to that at 8 $`\mu `$m than to that of the stars).
## 5 Results and discussion
### 5.1 Ages
One characteristic of the fits assuming instantaneous star formation is that they nearly always predict a cluster age of 3 Myr, with the notable exception of D53a, which is fitted with a stellar population of 1 to 2 Myr, and is one of the regions with the highest $`F_{24}/F_8`$ flux ratios. This uniformity of age is likely an artifact of adopting too simplified a representation of the clusters (instantaneous star formation model with coarse time steps of 1 Myr), whereas their formation may go on for several Myr instead of being coeval. That the clusters are all young and observed before the bulk of supernova explosions can be inferred from the fact that they produce a significant amount of ionizing radiation, since most were selected to be associated with bright H II regions. The case of D53a is interesting, as it confirms that complexes containing younger stars tend to have higher ratios of very small grain emission to aromatic band emission (Sauvage et al., 1990; Roussel et al., 2001b; Dale et al., 2004).
### 5.2 Masses
When considering such small stellar populations as those extracted here, formed of one or a few young clusters, statistical fluctuations of the mass spectrum of massive stars may become severe enough to affect the results of the fits. For the adopted initial mass function, the ratio $`\mathrm{N}(\mathrm{M}>\mathrm{M}_{\mathrm{hot}})/\mathrm{M}_{\mathrm{tot}}`$, where the numerator is the number of stars more massive than the threshold M<sub>hot</sub> and the denominator is the total cluster mass as given in Table 3, is equal to $`0.1264\mathrm{M}_{\mathrm{hot}}^{1.35}0.0002`$, with masses in units of solar mass. The whole range of cluster masses in Table 3 thus corresponds to between 9 and 126 O stars with $`\mathrm{M}_{\mathrm{hot}}=16\mathrm{M}_{\mathrm{}}`$ (Vacca, 1994). Adopting as a working definition of ionizing stars $`\mathrm{M}_{\mathrm{hot}}=10\mathrm{M}_{\mathrm{}}`$ (Kennicutt, 1984), the clusters masses also correspond to between 17 and 245 ionizing stars. The sampling of the initial mass function is thus a concern for the smallest clusters analysed in this paper, in which a single ionizing star may account for up to 5% of the far-ultraviolet emission. Note that a few regions, for which we could not obtain any satisfactory fits, were discarded from our sample ; one of the possible reasons for this is small-number statistics of ionizing stars in low-mass stellar clusters.
### 5.3 Stellar and nebular extinctions
The fitted stellar extinction at the effective wavelength of the FUV filter varies between nearly zero to $`1.5`$ mag. Its correspondence with the nebular extinction is shown in Figure 9. The lack of a tight correlation between the two can be understood in terms of spatial offsets between H II regions and clusters, and highly variable geometric configuration from one region to another. Within this sample, the A(H$`\alpha `$)/A<sub>R</sub> ratio ranges from 1 to more than 3.5 . The deviation from the value usually adopted for actively star-forming galaxies (Calzetti et al., 1994) can thus be greater than $`\pm 50`$% for individual regions.
There is a good qualitative correspondance between the extinction map at a resolution of 38″ and the extinctions of individual regions derived from the population synthesis fits, considering that significant deviations are expected from the fact that extinction is highly variable on small scales. It should also be noted that while the population synthesis fits give the extinction of the young stellar clusters only, the extinction in Figure 4, by construction, includes all the stellar populations at any given location. Since intermediate-age and old populations are, on average, much less affected by extinction at any wavelength, because decoupled over time from their parent gas clouds, the extinction values derived from the global energy budget are expected to be much lower than those of the ionizing clusters ; this is indeed observed.
### 5.4 Extinction law and geometry of the interstellar medium
We find clear differences in which extinction laws are applicable, as a function of the geometry of the H II regions with respect to the exciting clusters. The most compact regions follow either the Milky Way or the LMC law, whereas the clusters associated with H II regions forming shells and diffuse structures follow either the 30 Doradus or the Calzetti et al. (1994) law. Usually, the degeneracy between the latter two cannot be lifted, but models using the Calzetti et al. (1994) attenuation law systematically overpredict the ionizing photon flux by much larger amounts than the 30 Doradus law. We therefore consider the 30 Doradus law to yield a more reliable representation of the data, even though the stellar SED fits are formally of comparable quality using either of these two laws.
To qualitatively ascertain the consequences of neglecting radiative transfer effects, we examine the predictions of the DIRTY model (Witt & Gordon, 2000; Gordon et al., 2001; Misselt et al., 2001) with empirically-derived dust properties, in the configuration of a dust shell surrounding a stellar cluster, appropriate for our sample (Section 4.3). The dust properties of the LMC2 giant shell (called 30 Doradus here), not included in the original models, were derived by Gordon et al. (2003)<sup>4</sup><sup>4</sup>4The data tables were retrieved from the DIRTY website: http://dirty.as.arizona.edu/ .. The results obtained from the pure extinction laws for the Milky Way, the LMC2 region and the SMC are compared with those obtained from the corresponding attenuation laws for two types of local dust distribution: homogeneous and clumpy, the latter using the parameters defined by Witt & Gordon (2000). While other geometric configurations are available within DIRTY, better applicable to the extended emission of large regions within galaxies, it is beyond the scope of this observational study to perform comprehensive modelling. We stress that the basic result of large variations in the attenuation law, whether they correspond to genuine variations in the extinction law or not, is independent of any assumptions about the relative geometry of stars and dust.
Figure 10 shows the extinction and attenuation curves in the model shell configuration, for a particular optical depth, which corresponds to $`A_{1516}=0.7`$ using the Cardelli extinction law, within the range of derived extinction values for our clusters. The inclusion of scattering produces emerging spectra that are bluer, with significantly higher FUV/NUV ratios, as more light is scattered in the line of sight than out. It is thus conceivable that the clusters which were well fitted with the Cardelli extinction law could be equally well fitted with a LMC or SMC-based attenuation law. However, even at the low optical depth considered in Figure 10, the attenuation curves applicable to the clumpy case are extremely grey, and thus would not provide any good fit to the observed spectral energy distributions, since they would predict too much intrinsic emission in the optical bands.
We now examine the aptitude of each attenuation law to reproduce the observed FUV/NUV flux density ratios. Figure 11 shows the predictions from the DIRTY model. The intrinsic FUV/NUV ratio of a 3 Myr-old cluster is 1.18, and it varies between 1.13 and 1.31 for ages between 1 and 5 Myr. This uncertainty is reflected in the error bars in Figure 11.
Cluters well fitted with the Cardelli or LMC extinction law:
The data points overlap only with the Milky Way extinction law and the LMC2 extinction and attenuation laws. However, the attenuation law using the LMC2 dust properties, in the homogeneous case, is not much different from the pure extinction law, i.e. the effects of scattering are almost negligible. It is in principle possible to reproduce the observed FUV/NUV ratios with the LMC2 or 30 Doradus laws, but only if the clusters are younger (1 or 2 Myr). This seems unlikely, since the predicted ionizing photon fluxes would then be in strong excess with respect to the observational estimates (which are more robust for the compact HII regions associated with these clusters than for the diffuse HII regions). We note that the effects of scattering are strongest for the Milky Way dust properties, and that none of the clusters would be well represented in the dust shell configuration with the Milky Way attenuation laws (i.e. including scattering), either in the homogeneous or in the clumpy case.
Cluters well fitted with the 30 Doradus extinction law:
The data points lie in between the LMC2 (30 Doradus) and the SMC laws. Again, the SMC-dust attenuation law, in the homogeneous case, is very close to the pure extinction law, and the effects of scattering are also very small in this configuration. But the extinction laws in the line of sight of these clusters may span the full range between the 30 Doradus and the SMC laws.
We conclude that the correct treatment of radiative transfer effects would not drastically alter our qualitative conclusions in what follows, unless stars had extended spatial distributions and were significantly mixed with the dust, which we argue is very unlikely for the stellar clusters that we classify as compact.
To quantify the compactness of the regions in a convenient way, we use the H$`\alpha `$/R flux ratio, measured in the area defining the stellar aperture (i.e. if the H II region extends far away from the clusters, only a fraction of the H$`\alpha `$ emission is considered in this ratio). While this is proportional to the H$`\alpha `$ equivalent width, we argue below that it is also a valid measure of the compactness of the H II region. Although it is potentially distorted by differential extinction between the nebular phase and the stars, this effect should not have a large impact. Using the results from Figure 9, we derive that for the highest extinctions, observed values of H$`\alpha `$/R close to 1 at one extreme or close to 0.2 at the other extreme may be shifted by up to 0.9 upward or 0.2 upward, respectively, when corrected for extinction. Since the regions with the highest extinctions are also those with the highest H$`\alpha `$/R values, correcting rigorously for extinction would preserve the trend discussed below. Figure 12 confirms that the H$`\alpha `$/R ratio provides a reliable quantification of the degree of compactness of the entire H II region. We prefer to adopt the H$`\alpha `$/R indicator, because it is a relatively unbiased measurement: unlike the ordinate of Figure 12, it does not imply any assumption about the total extent of the H II region, which is difficult to define owing to the presence of neighboring H II regions and the diffuse and filamentary aspect of the H$`\alpha `$ emission.
Figure 13 shows the dependence of the extinction law on the H$`\alpha `$/R ratio. Whenever a region sits in between two extinction laws, it means that its SED is formally as well represented with either law as with the other. Unconstrained regions are those whose SED is well fitted by at least three of the four extinction laws (usually because of very small extinction). To indicate cases when the measured ionizing flux is overestimated by a large amount by the Calzetti et al. (1994) attenuation law, the points are plotted closer to the 30 Doradus law locus. The main discriminant between the various extinction laws is the A(FUV)/A(NUV) ratio (Sect. 4.3). The increase in A(FUV)/A(NUV), going from the Cardelli et al. (1989) law to the LMC law to the 30 Doradus law, is due to both the gradual disappearance of the so-called 2175 Å bump, which is covered by the NUV band, and the steepening of the FUV rise. Although the Calzetti et al. (1994) law is devoid of the 2175 Å bump, the corresponding A(FUV)/A(NUV) ratio is slightly smaller than that of the 30 Doradus law, because the Calzetti et al. (1994) law has a flatter wavelength dependence.
That the extinction law should vary within a galaxy is not a complete surprise given that some significant variations have been observed within the Milky Way and the Magellanic Clouds. In the SMC, Lequeux et al. (1982) discussed the line of sight to a star with a Milky Way-type extinction law. In the LMC, Misselt et al. (1999) discussed a line of sight associated with a giant shell, with a much weaker 2175 Å bump than its immediate surroundings. In the Milky Way, LMC-type extinction laws have be found toward low-density lines of sight (Clayton et al., 2000), and a SMC-type extinction law was derived toward a region suggested to have been profoundly affected by supernova shocks (Valencic et al., 2003). Although the disappearance of the 2175 Å bump, together with the steepening of the far-ultraviolet rise, are usually ascribed to a decrease in metallicity, which could explain the differences between the average extinction laws of the Milky Way and Magellanic Clouds, Mas-Hesse & Kunth (1999) found no such dependence of the extinction curve on metallicity, for a sample of dwarf star-forming galaxies. As pointed out by Misselt et al. (1999), an apparent dependence on metallicity may be only an indirect effect, because the dust grain properties and geometry of the interstellar medium are expected to be on average related to the metallicity.
It has been suggested that the intensity of star formation in the local environment could be the main driver of the extinction law shape, causing the destruction or processing of the dust grains responsible for the 2175 Å feature, via radiation hardness and supernova shocks (Fitzpatrick & Massa, 1986; Gordon et al., 1997; Gordon & Clayton, 1998; Mas-Hesse & Kunth, 1999). This scenario is well in line with the attribution of part of the 2175 Å feature to aromatic hydrocarbons (Joblin et al., 1992; Duley & Seahra, 1998), as this dust species is known to be less resilient than very small grains and to be destroyed by intense radiation fields. Supporting this identification, Vermeij et al. (2002) have reported variations of the mid-infrared band ratios (caused by changes in molecular structure, or degree of ionization and hydrogenation) that are correlated with the strength of the 2175 Å bump. It has been proposed that aromatic hydrocarbons are also responsible for part of the far-ultraviolet rise (Puget & Léger, 1989), and very small carbonaceous grains (of three-dimensional structure) are sometimes assumed to be the main contributor to the 2175 Å bump (Draine & Lee, 1984; Hecht, 1986; Sakata et al., 1994).
The present work, however, supports the idea that geometry effects are much more important than the distinction between quiescent and actively star-forming regions. According to the SED fits, all the exciting clusters have very similar ages. Therefore, rather than being a pure indicator of the evolution of the clusters, as the equivalent widths of the H recombination lines are usually assumed to be, the H$`\alpha `$/R ratio can be interpreted here as mainly governed by the geometric evolution of the whole region, bearing traces of previous generations of stars in the shaping of the interstellar gas. Indeed, the giant shells of hydrogen that are observed in many locations across the disk of NGC 300, where they are illuminated by ionizing clusters, may have been carved by supernova explosions, in which case they signal past star formation at the same location, more than 5 Myr ago. The difference in geometry between the regions D53a-b and D53c, which are part of the same complex and whose exciting clusters have approximately the same age, is particularly significant. High $`F_{24}/F_8`$ ratios, indicative of intense radiation fields, can also be found both in compact and diffuse star formation complexes. The most obvious difference between clusters associated with a Milky Way-LMC extinction law and those associated with a 30 Doradus-Calzetti law is in the distribution of the material around them. For the latter, both the atomic hydrogen and the mid-infrared-emitting dust (aromatic hydrocarbons and very small grains) have been removed far away, presumably by stellar winds and supernova shocks from previous generations of stars at the same location. This does not preclude important differences as well in the processing of dust grains by supernova shocks, but if the modified dust is dissipated together with the gas, it is only weakly participating in the extinction of the clusters.
We can speculate that the effects of a varying extinction law within and among galaxies may resemble the effects of a varying initial mass function (Rosa & Benvenuti, 1994), since for a given optical depth, the emerging ultraviolet continuum is much flatter in the case of a 30 Doradus-Calzetti law than for a Milky Way law. An upper mass cutoff of the initial mass function is not supported in NGC 300, since Wolf-Rayet features have been observed throughout its disk (Schild et al., 2003).
Finally, we can expect the behavior of more massive galaxies and of starbursts, where H II regions may generally be more centrally distributed and more compact than in NGC 300, to be significantly different.
### 5.5 Interstellar emission in the near-infrared
In some regions, we detect a clear excess at 3-5 $`\mu `$m above the stellar model. The improvement with respect to the previous treatment in Helou et al. (2004) (where it was assumed that the 3.5 $`\mu `$m emission covered by the IRAC1 band is entirely stellar) allows us to derive more accurately the amount of interstellar emission at these wavelengths. However, our method restricts the detectable interstellar component at 3-5 $`\mu `$m to the exact location of the stellar emission ; if the hot dust in this wavelength range is significantly offset from the exciting source, it will not be detected. The 3.5 $`\mu `$m/4.5 $`\mu `$m color of young clusters is very similar to that of dust (close to 1), making it very difficult to separate the two components if extensive information at other wavelengths does not exist. Complexes where the interstellar component at 3-5 $`\mu `$m is high with respect to the stellar fluxes (within the stellar aperture) are among the most compact (Fig. 14 and 13).
To explore the link between the extinction law and the dust emission, we adopt measurements of the 3-5 $`\mu `$m interstellar emission, as the favorable angular resolution in these bands allows us to select the dust emission within the stellar “apertures”, directly in the line of sight of the clusters. Figure 14 confirms that switching from the H$`\alpha `$ equivalent width to the column density of near-infrared-emitting dust in the line of sight of the clusters gives the same trend as seen in Figure 13: the shape of the extinction law seems to depend importantly on the amount of interstellar material which re-emits the absorbed energy in the near- to mid-infrared. We cannot explore related variations in the mid-infrared to far-infrared colors, because the angular resolution at far-infrared wavelengths is too coarse.
Even though this 3-5 $`\mu `$m dust component may not be physically related to the other mid-infrared dust species, the aromatic band carriers and very small grains, Figure 15 shows that it scales reasonably well with the 8 $`\mu `$m aromatic band emission. These measurements suffer from the degraded angular resolution at 8 $`\mu `$m combined with the narrow stellar “apertures” used to fit and extract the interstellar component from the 3.6 $`\mu `$m band, but are nevertheless in line with the $`F_{3.6}/F_8`$ flux ratios derived by Lu (2004) for interstellar emission in star-forming galaxies, as well as for the general interstellar medium of the Milky Way.
### 5.6 Porosity of the interstellar medium and excitation of aromatic band carriers
Examination of the morphology of the H II complexes in H$`\alpha `$ provides evidence that the interstellar medium is highly porous in some regions, thus that the 8 $`\mu `$m-emitting dust can be heated by photons travelling as far as several hundred parsecs (since the NUV photons travel farther away than the ionizing photons, which are able to excite H$`\alpha `$ emission a few hundred parsec away from the clusters). Diffuse emission also exists at 24 $`\mu `$m, but is harder to detect, because it is intrinsically fainter than at 8 $`\mu `$m and observed at much lower angular resolution, increasing the confusion by nearby bright sources, and the 8 $`\mu `$m observations are deeper. Although the 8 $`\mu `$m/H$`\alpha `$ flux ratio varies dramatically from bright H II regions to more diffuse regions (Helou et al., 2004), and the 8 $`\mu `$m emission has a more filamentary aspect than the H$`\alpha `$ and 24 $`\mu `$m emission, virtually every 8 $`\mu `$m peak is associated with a close-by H$`\alpha `$ peak (see Figure 6).
The claim of Haas et al. (2002) that aromatic band emission is more closely correlated with very cold dust than with hot dust within galaxies does not withstand examination in NGC 300. Figure 16 shows the relation between the aromatic band emission at 8 $`\mu `$m and the emission in the far-infrared, at 70 and 160 $`\mu `$m. As a comparison, the relation between 24 $`\mu `$m and the far-infrared emission is also included. The 160 $`\mu `$m emission can be considered a tracer of the cold dust, and the 70 $`\mu `$m emission arises from dust that is on average hotter, with a significant contribution from impulsively-heated very small grains (Désert et al., 1990). Figure 16 indicates clearly that aromatic bands are not associated with cold dust, and are not excited by the same radiation as the big grains. In the 8 $`\mu `$m-70 $`\mu `$m diagram, at high surface brightnesses, a branch with a slightly lower slope than in the low-brightness part is visible. This branch is slightly more prominent in the 8 $`\mu `$m-24 $`\mu `$m diagram ; it is linked with the disappearance of the aromatic bands from inside the H II regions. Nevertheless, the 8 $`\mu `$m-70 $`\mu `$m and 8 $`\mu `$m-24 $`\mu `$m relations are much closer to linearity than the 8 $`\mu `$m-160 $`\mu `$m relation. The linearity between 8 $`\mu `$m and 850 $`\mu `$m surface brightnesses reported by Haas et al. (2002), obtained at low angular resolution in galaxies of moderate angular size, and more importantly at mediocre sensitivity at 850 $`\mu `$m, may simply be a consequence of extracting only the brightest and most active regions, which generally have dust temperature distributions shifted toward higher values, and discarding more diffuse regions, with cooler dust temperatures. Their result thus does not imply that the aromatic bands and the cold dust component are physically related, nor heated by the same radiation.
### 5.7 Embedded and bare sources
Within the inner disk, a few non-stellar sources bright at 24 $`\mu `$m have weak optical and no ultraviolet counterparts (Table 4). Two of these are probably background distant galaxies, judging from their morphology in optical high-resolution images. The five other sources coincide with very faint H$`\alpha `$ emission, have no other optical counterpart, and may be ultracompact H II regions. We consider it likely at least for the object catalogued by Deharveng et al. (1988) as the faint and small H II region D125. If the H$`\alpha `$/H$`\beta `$ decrement of D125 can be trusted, it indicates a lower limit of the extinction in the H$`\alpha `$ line of 2.8 mag, much higher than for the H II regions associated with the clusters examined here (see Table 1). Its mid-infrared counterpart has the highest $`F_{24}/F_8`$ ratio within the galaxy ($`4.5`$) and is still visible at 70 $`\mu `$m despite the degraded angular resolution.
Conversely, some far-ultraviolet and H$`\alpha `$ sources are not associated with any detectable mid-infrared dust emission. These regions are of low brightness and diffuse in the H$`\alpha `$ emission, suggesting that they are aged H II regions, excited by slightly evolved clusters (see the previous paragraph). Both types of sources contribute negligibly to the emission from the optically-bright disk, but the latter (bare sources) become predominant in the outer disk.
## 6 Summary and conclusion
The large-scale UV/IR ratio gradient observed in NGC 300 is probably the consequence of a combination of factors, all causing an increase in the mean free path of ultraviolet photons. In this paper, we have focussed on geometric factors, and shown that they cannot be neglected.
Throughout the disk of NGC 300, mid-infrared peaks are commonly offset from both the ionized gas peaks and stellar clusters, except for compact regions (Helou et al., 2004). This is because we are able to partially resolve H II regions from photodissociation regions. However, in the outer disk, even the far-infrared emission sometimes appears significantly displaced from the location of the ultraviolet sources, at a resolution of 38″ (FWHM) or about 400 pc. In the outer disk, dust may be removed away from star-forming sites more efficiently than in the inner disk, either because the volume density of the interstellar medium is decreased or because the exciting clusters are older on average and have injected more non-thermal energy into the interstellar medium.
There are several examples of very porous star-forming sites at distances of 3-4′ ($`2`$ kpc) from the center, from which both the hydrogen gas and the dust have been removed (such as D118 and D137) ; as they belong to the optically-bright part of the disk, where presumably the molecular gas density remains high, i.e. above the threshold for gravitational instabilities discussed by Kennicutt (1989), these individual H II complexes are unable to affect significantly the global UV/IR radial profile. However, such regions may become dominant in the outer disk.
Most of the H$`\alpha `$ emission in the outer disk is in the form of diffuse, very large shells, which supports either one of the factors invoked above for dust removal: either the gas is very rapidly pushed away, during the lifetime of ionizing stars, because the interstellar medium is more tenuous, or the exciting stellar populations are older on average, so that they had more time to blow the gas away by their winds and supernova explosions. Good supporting examples of the hypothesis of aged stellar populations are the clusters D115 and D147, which are bright in the FUV, but seem to contain slightly evolved stellar populations (of the order of 5 Myr, as opposed to 3 Myr for most clusters). Since it is not possible to measure and fit the SEDs of a reasonably large number of clusters associated with giant shells in the outer disk, however, we cannot constrain the importance of this scenario relatively to the others (but see also Figure 2).
From population synthesis fitting of individual young stellar clusters or small groups of clusters, we find that the stellar attenuation law varies significantly on the explored spatial scales of $`100`$ pc, and we argue that this empirical result likely corresponds to real variations in the extinction law as well. Just as the UV/IR ratio does not depend only on age and metallicity, there is ample evidence that the apparent correlation between the shape of the extinction law and the metallicity or the mode of star formation (active versus quiescent) is only a secondary effect of a more direct cause. We argue that the strength of the 2175 Å bump and the steepness of the far-ultraviolet rise, whatever the responsible dust species, are directly related to the degree of compactness of the star-forming regions, and to the column density of interstellar material in front of the clusters. While most of the young stellar clusters considered here have similar ages and have not produced supernova explosions yet, the only obvious difference between them is that some appear to have formed in regions deeply affected by previous generations of stars. These are associated with diffuse, shell-like H II regions of large dimensions. If we divide star-forming complexes in two categories, based on the degree of compactness of the H$`\alpha `$ emission, we find that the most compact regions follow either the Milky Way or the average LMC extinction law, whereas the most diffuse regions follow an extinction law with a decreased 2175 Å bump and/or a steeper far-ultraviolet rise. While we cannot ascribe these features to a particular dust species, we show, using the near-infrared dust component extracted from our SED fitting procedure, thanks to the favorable angular resolution, that the distinction between the two types of morphologies and the two types of extinction laws is related to the column density of near-infrared emitting dust lying in the line of sight of the clusters.
We can hypothesize that the main difference in dust nature causing the change of extinction law between the compact and diffuse regions is related to a changing proportion of dust grains very close to the H II region, which have been processed by hard ultraviolet radiation from the young stars, with respect to grains lying further away and which have not been processed in the same way. This photo-processing is likely occurring very fast, so that the observed dust properties of compact H II regions should be systematically different from the dust properties of diffuse regions where dust lies further away from the exciting stars. According to Boulanger et al. (1990), who discussed the ubiquitous and large variations in mid-infrared to far-infrared colors within and among nearby molecular clouds, such variations cannot be explained only by changes in the radiation field, but require actual variations in the abundance of small grains with respect to big grains. They proposed that these abundance variations are caused by photo-processing, on time scales comparable to those of the mixing of gas by turbulent motions in clouds. Habart et al. (2001) and Abergel et al. (2002) also discussed abundance variations of small grains and aromatic hydrocarbons in resolved photodissociation regions and molecular clouds. This speculation would be consistent with the fact that the dust emitting in the 3-5 $`\mu `$m range is likely excited inside or very close to H II regions (Helou et al., 2004), and that we have seen here that its excitation has a close correspondence to the shape of the extinction law.
## Appendix A Notes on individual regions (including the clusters for which the fits are shown in Figure 7)
D53a: Compact, associated with very hot dust at 24 $`\mu `$m and 3-5 $`\mu `$m. Since the emission is confused by nearby clusters in the GALEX bands, we had to subtract a gaussian representation of these contaminating clusters, and apply a large aperture correction to the FUV and NUV fluxes, because of the small size of the aperture.
D53b: Compact, associated with hot dust at 24 $`\mu `$m and 3-5 $`\mu `$m. Two distinct clusters can be distinguished, but are well fitted with a single stellar population.
D53c: Two clusters, one with a very red SED. The morphology in the various bands clearly indicates that interstellar emission is negligible at 3-5 $`\mu `$m. The fits included a second population but were not constrained to reproduce the 3-5 $`\mu `$m fluxes. The modelling confirms, however, the stellar nature of the 3-5 $`\mu `$m emission. The ionized gas is diffuse and forms an extended fragmentary shell of size $`200`$ pc.
D118: Complex region with several clusters, which can be fitted by a single population. The ionized gas and the dust are distributed in both compact regions and a giant shell.
D119: Same remarks as for D118.
D137b: This complex region stands out as very peculiar: the ionized gas and dust seem to form two walls on each side of the brightest cluster, separated by $`80`$ pc in projected distance. Star formation has thus had dramatic effects on the shaping of the surrounding interstellar medium. The H$`\alpha `$ and dust peaks are coincident with weak UV sources, presumably because of extinction effects, while the cluster dominating the UV emission is located in a zone nearly void of interstellar emission. Since the complex contains redder stellar sources besides these clusters, we allow for a second population in the fits. This increases dramatically the degeneracy of the younger population parameters.
D84: Compact region, associated with hot dust at 24 $`\mu `$m and 3-5 $`\mu `$m.
We thank Sang Chul Kim and Charles Hoopes for graciously providing their UBVI and H$`\alpha `$ images, respectively, of NGC 300. We are most grateful to Karl Gordon for computing new DIRTY models incorporating the LMC2 extinction law, ahead of schedule and upon our request ; the corresponding data tables are publicly available on the DIRTY webpage (http://dirty.as.arizona.edu/). This research benefited from the NASA/IPAC Extragalactic Database (NED), which is operated by the Jet Propulsion Laboratory of Caltech.
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# 1 Introduction
## 1 Introduction
Coefficient functions, or partonic cross sections, form the backbone of perturbative QCD. These quantities are defined in terms of power expansions in the strong coupling constant $`\alpha _\mathrm{s}`$. In general, only a few terms in this expansion can be calculated. It is however possible, and necessary, to resum the dominant contributions to all orders in $`\alpha _\mathrm{s}`$ close to exceptional kinematic points. Close to threshold, for example, where real emissions are kinematically suppressed, the resummation takes the form of an exponentiation in Mellin-$`N`$ space , with the moments $`N`$ defined with respect to the appropriate scaling variable, like Bjorken-$`x`$ in deep-inelastic scattering (DIS) and $`x_T=2p_T/\sqrt{S}`$ in direct photon and inclusive hadron production.
The resummation exponents are given by integrals over functions in turn defined by a power series in $`\alpha _\mathrm{s}`$. Besides by dedicated calculations, the corresponding expansion coefficients can be obtained by expanding the exponentials and comparing to the results of fixed-order calculations. Hence progress in the latter sector also facilitates improved resummation predictions. At present the next-to-leading order (NLO) is the standard approximation for many important observables, facilitating a resummation with next-to-leading logarithmic (NLL) accuracy. For recent introductory overviews see, for instance, Refs. . The next-to-next-to-leading order (NNLO) corrections have been completed so far only for the coefficient functions for inclusive lepton-proton DIS , the Drell-Yan process and the related Higgs boson production in proton-proton collisions. Consequently, the threshold resummation has been carried out at the next-to-next-to-leading logarithmic (NNLL) accuracy only for these processes .
Recently we have computed the complete three-loop coefficient functions for inclusive photon-exchange DIS . Moreover, in the course of the calculation of the third-order splitting functions governing the NNLO evolution of the parton distributions , we have also computed DIS by exchange of a scalar directly coupling only to gluons. Together these results enable us to extend two more universal functions entering the resummation exponents, the quark and gluon jet functions $`B_\mathrm{q}`$ and $`B_\mathrm{g}`$ collecting final-state collinear emissions, to the third order in $`\alpha _\mathrm{s}`$. In fact, already the second-order coefficient for $`B_\mathrm{g}`$ represents a new result, relevant for future NNLL resummations of processes with final-state gluons at the Born level. Making use also of the results of Refs. we can furthermore effectively, i.e., up to the small contribution of the four-loop cusp anomalous dimension, extend the threshold resummation for inclusive DIS to an unprecedented next-to-next-to-next-to-leading logarithmic (N<sup>3</sup>LL) accuracy.
The remainder of this article is organized as follows: after recalling the general structure of the resummation exponents in Section 2, we extend the required integrations in Section 3 to the fourth logarithmic (N<sup>3</sup>LL) order. In Section 4 we determine the relevant expansion coefficients by comparison to our three-loop results for DIS and illustrate the numerical effect of the N<sup>3</sup>LL contributions to the resummation exponent. In Section 5 we present the resulting predictions for the leading seven large-$`x`$ terms of the four-loop coefficient function and discuss higher-order effects in terms of the expansion in towers of threshold logarithms. Our results are briefly summarized in Section 6. Some basic relations for the integrations of Section 2 can be found in the Appendix.
## 2 The general structure
For processes with only one colour structure at the Born level, the resummed Mellin-space coefficient functions $`C^N`$ (defined in the $`\overline{\text{MS}}`$ scheme) are given by a single exponential
$$C^N(Q^2)/C_{\mathrm{LO}}^N(Q^2)=g_0(Q^2)\mathrm{exp}[G^N(Q^2)]+𝒪(N^1\mathrm{ln}^nN).$$
(2.1)
Here $`C_{\mathrm{LO}}^N`$ denotes the lowest-order coefficient function for the process under consideration, e.g., $`C_{\mathrm{LO}}^N=1`$ for DIS. The prefactor $`g_0`$ collects, order by order in the strong coupling constant $`\alpha _\mathrm{s}`$, all $`N`$-independent contributions. The exponent $`G^N`$ contains terms of the form $`\mathrm{ln}^kN`$ to all orders in $`\alpha _\mathrm{s}`$. Besides the physical hard scale $`Q^2`$ ($`=q^2`$ in DIS, with $`q`$ the four-momentum of the exchanged gauge boson), both functions also depend on the renormalization scale $`\mu _r`$ and the mass-factorization scale $`\mu _f`$. The reference to these scales will be often suppressed for brevity.
The exponential in Eq. (2.1) is build up from universal radiative factors for each initial- and final-state parton $`p`$, $`\mathrm{\Delta }_\mathrm{p}`$ and $`J_\mathrm{p}`$, together with a process-dependent contribution $`\mathrm{\Delta }^{\mathrm{int}}`$. For example, the resummation exponents for inclusive deep-inelastic scattering, Drell-Yan (DY) lepton-pair production and direct photon production via $`q\overline{q}g\gamma `$ and $`qgq\gamma `$ take the form
$`G_{\mathrm{DIS}}^N`$ $`=`$ $`\mathrm{ln}\mathrm{\Delta }_\mathrm{q}+\mathrm{ln}J_\mathrm{q}+\mathrm{ln}\mathrm{\Delta }_{\mathrm{DIS}}^{\mathrm{int}},`$
$`G_{\mathrm{DY}}^N`$ $`=`$ $`2\mathrm{ln}\mathrm{\Delta }_\mathrm{q}+\mathrm{ln}\mathrm{\Delta }_{\mathrm{DY}}^{\mathrm{int}},`$
$`G_{\mathrm{ab}\mathrm{c}\gamma }^N`$ $`=`$ $`\mathrm{ln}\mathrm{\Delta }_\mathrm{a}+\mathrm{ln}\mathrm{\Delta }_\mathrm{b}+\mathrm{ln}J_\mathrm{c}+\mathrm{ln}\mathrm{\Delta }_{\mathrm{ab}\mathrm{c}\gamma }^{\mathrm{int}}.`$ (2.2)
The radiation factors are given by integrals over functions of the running coupling. Specifically, the effects of collinear soft-gluon radiation off an initial-state parton $`p=q,g`$ are collected by
$$\mathrm{ln}\mathrm{\Delta }_\mathrm{p}(Q^2,\mu _f^{\mathrm{\hspace{0.17em}2}})=_0^1𝑑z\frac{z^{N1}1}{1z}_{\mu _f^{\mathrm{\hspace{0.17em}2}}}^{(1z)^2Q^2}\frac{dq^2}{q^2}A_\mathrm{p}(\alpha _\mathrm{s}(q^2)).$$
(2.3)
Collinear emissions from an ‘unobserved’ final-state parton lead to the so-called jet function,
$$\mathrm{ln}J_\mathrm{p}(Q^2)=_0^1𝑑z\frac{z^{N1}1}{1z}\left[_{(1z)^2Q^2}^{(1z)Q^2}\frac{dq^2}{q^2}A_\mathrm{p}(\alpha _\mathrm{s}(q^2))+B_\mathrm{p}(\alpha _\mathrm{s}([1z]Q^2))\right].$$
(2.4)
Finally the process-dependent contributions from large-angle soft gluons are resummed by
$$\mathrm{ln}\mathrm{\Delta }^{\mathrm{int}}(Q^2)=_0^1𝑑z\frac{z^{N1}1}{1z}D(\alpha _\mathrm{s}([1z]^2Q^2)).$$
(2.5)
The functions $`g_0`$ in Eq. (2.1) and $`A_\mathrm{p}`$, $`B_\mathrm{p}`$ and $`D`$ in Eqs. (2.3) – (2.5) are given by the expansions
$$F(\alpha _\mathrm{s})=\underset{l=l_0}{\overset{\mathrm{}}{}}F_l\frac{\alpha _\mathrm{s}^l}{4\pi }\underset{l=l_0}{\overset{\mathrm{}}{}}F_la_\mathrm{s}^l,$$
(2.6)
where $`l_0=0`$ with $`g_{00}=1`$ for $`F=g_0`$, and $`l_0=1`$ else. Here we have also taken the opportunity to specify the reduced coupling $`a_\mathrm{s}`$ employed for the rest of this article. The extent to which these functions are known defines the accuracy to which the threshold logarithms can be resummed.
The situation for inclusive DIS is actually simpler than indicated in Eq. (2), as the function (2.5) is found to vanish to all orders in $`\alpha _\mathrm{s}`$ ,
$$D_k^{\mathrm{DIS}}=\mathrm{\hspace{0.33em}0},\mathrm{\Delta }_{\mathrm{DIS}}^{\mathrm{int}}=\mathrm{\hspace{0.33em}1}.$$
(2.7)
On the other hand, the resummation of processes with four or more partons at the Born level, like inclusive hadron production in $`pp`$ collisions , is more complicated than Eq. (2.1) due to colour interferences and correlations in the large-angle soft gluon emissions . However, the process-independent functions $`A_\mathrm{p}`$ and $`B_\mathrm{p}`$ retain their relevance also for such cases.
## 3 The resummation exponent to fourth logarithmic order
After performing the integrations in Eqs. (2.3) – (2.5), the function $`G^N`$ in Eq. (2) takes the form
$$G^N(Q^2)=\mathrm{ln}Ng_1(\lambda )+g_2(\lambda )+a_\mathrm{s}g_3(\lambda )+a_\mathrm{s}^{\mathrm{\hspace{0.17em}2}}g_4(\lambda )+\mathrm{},$$
(3.1)
where $`\lambda =\beta _0a_\mathrm{s}\mathrm{ln}N`$. For the actual computation of the functions $`g_i`$ it is convenient to employ the following representation for the scale dependence of $`a_\mathrm{s}`$ up to N<sup>3</sup>LO :
$`a_\mathrm{s}(q^2)`$ $`=`$ $`{\displaystyle \frac{a_\mathrm{s}}{X}}{\displaystyle \frac{a_\mathrm{s}^{\mathrm{\hspace{0.17em}2}}}{X^2}}{\displaystyle \frac{\beta _1}{\beta _0}}\mathrm{ln}X+{\displaystyle \frac{a_\mathrm{s}^{\mathrm{\hspace{0.17em}3}}}{X^3}}\left[{\displaystyle \frac{\beta _1^2}{\beta _0^2}}(\mathrm{ln}^2X\mathrm{ln}X1+X)+{\displaystyle \frac{\beta _2}{\beta _0}}(1X)\right]`$ (3.2)
$`+{\displaystyle \frac{a_\mathrm{s}^{\mathrm{\hspace{0.17em}4}}}{X^4}}[{\displaystyle \frac{\beta _1^3}{\beta _0^3}}(2(1X)\mathrm{ln}X+{\displaystyle \frac{5}{2}}\mathrm{ln}^2X\mathrm{ln}^3X{\displaystyle \frac{1}{2}}+X{\displaystyle \frac{1}{2}}X^2)`$
$`+{\displaystyle \frac{\beta _3}{2\beta _0}}(1X^2)+{\displaystyle \frac{\beta _1\beta _2}{\beta _0^2}}(2X\mathrm{ln}X3\mathrm{ln}XX(1X))]+𝒪(a_\mathrm{s}^{\mathrm{\hspace{0.17em}5}})`$
with $`a_\mathrm{s}a_\mathrm{s}(\mu _r^{\mathrm{\hspace{0.17em}2}})`$ and the abbreviation $`X=1+a_\mathrm{s}\beta _0\mathrm{ln}(q^2/\mu _r^{\mathrm{\hspace{0.17em}2}}).`$ The terms up to the $`n`$-th order in $`a_\mathrm{s}`$ in Eq. (3.2) contribute to $`g_n`$ in Eq. (3.1). Thus the calculation of $`g_4`$ requires the highest known coefficient of the beta function of QCD, $`\beta _3`$ .
Generalizing the approach of Ref. , the functions $`g_i(\lambda )`$ can be obtained using well-known methods for Mellin transforms based on properties of harmonic sums and harmonic polylogarithms in addition to algorithms for the evaluation of nested sums . The basic relations for this approach, suitable for the evaluation of Eq. (3.1) to any accuracy, are presented in the Appendix. As a check we have also carried out the integrations along the lines of Ref. .
For the convenience of the reader, we first recall the known results for $`g_1`$, $`g_2`$ and $`g_3`$ . For brevity suppressing factors of $`\beta _0`$ (see below) and using the short-hand notations $`L_{\mathrm{qr}}=\mathrm{ln}(Q^2/\mu _r^{\mathrm{\hspace{0.17em}2}})`$ and $`L_{\mathrm{fr}}=\mathrm{ln}(\mu _f^{\mathrm{\hspace{0.17em}2}}/\mu _r^{\mathrm{\hspace{0.17em}2}})`$, these functions can be written as
$`g_1^{\mathrm{DIS}}(\lambda )`$ $`=`$ $`A_1(1\mathrm{ln}(1\lambda )+\lambda ^1\mathrm{ln}(1\lambda )),`$ (3.3)
$`g_2^{\mathrm{DIS}}(\lambda )`$ $`=`$ $`\left(A_1\beta _1A_2\right)(\lambda +\mathrm{ln}(1\lambda ))+{\displaystyle \frac{1}{2}}A_1\beta _1\mathrm{ln}^2(1\lambda )`$ (3.4)
$`\left(A_1\gamma _eB_1\right)\mathrm{ln}(1\lambda )+L_{\mathrm{qr}}A_1\mathrm{ln}(1\lambda )+L_{\mathrm{fr}}A_1\lambda ,`$
$`g_3^{\mathrm{DIS}}(\lambda )`$ $`=`$ $`{\displaystyle \frac{1}{2}}(A_1\beta _2A_1\beta _1^2+A_2\beta _1A_3)\left(1+\lambda {\displaystyle \frac{1}{1\lambda }}\right)`$ (3.5)
$`+A_1\beta _1^2\left({\displaystyle \frac{\mathrm{ln}(1\lambda )}{1\lambda }}+{\displaystyle \frac{1}{2}}{\displaystyle \frac{\mathrm{ln}^2(1\lambda )}{1\lambda }}\right)+\left(A_1\beta _2A_1\beta _1^2\right)\mathrm{ln}(1\lambda )`$
$`+(A_1\beta _1\gamma _e+A_2\beta _1B_1\beta _1)\left(1{\displaystyle \frac{1}{1\lambda }}{\displaystyle \frac{\mathrm{ln}(1\lambda )}{1\lambda }}\right)`$
$`\left(A_1\beta _2+{\displaystyle \frac{1}{2}}A_1(\gamma _e^{\mathrm{\hspace{0.17em}2}}+\zeta _2)+A_2\gamma _eB_1\gamma _eB_2\right)\left(1{\displaystyle \frac{1}{1\lambda }}\right)`$
$`+L_{\mathrm{qr}}[(A_1\gamma _eA_1\beta _1+A_2B_1)(1{\displaystyle \frac{1}{1\lambda }})+A_1\beta _1\left({\displaystyle \frac{\mathrm{ln}(1\lambda )}{1\lambda }}\right)]`$
$`+L_{\mathrm{fr}}A_2\lambda L_{\mathrm{qr}}^{\mathrm{\hspace{0.17em}2}}{\displaystyle \frac{1}{2}}A_1\left(1{\displaystyle \frac{1}{1\lambda }}\right)L_{\mathrm{fr}}^{\mathrm{\hspace{0.17em}2}}{\displaystyle \frac{1}{2}}A_1\lambda .`$
The dependence on $`\beta _0`$ is recovered by $`A_kA_k/\beta _0^k`$, $`B_kB_k/\beta _0^k`$, $`\beta _k\beta _k/\beta _0^{k+1}`$ and multiplication of $`g_3`$ by $`\beta _0`$. In the same notation the new function $`g_4`$ (to be multiplied by $`\beta _0^{\mathrm{\hspace{0.17em}2}}`$) is given by
$`g_4^{\mathrm{DIS}}(\lambda )`$ $`=`$ $`{\displaystyle \frac{1}{6}}A_1\beta _1^3{\displaystyle \frac{\mathrm{ln}^3(1\lambda )}{(1\lambda )^2}}+{\displaystyle \frac{1}{2}}(A_1\beta _1^2\gamma _e+A_2\beta _1^2B_1\beta _1^2){\displaystyle \frac{\mathrm{ln}^2(1\lambda )}{(1\lambda )^2}}+{\displaystyle \frac{1}{2}}(A_1\beta _1^3A_1\beta _1\beta _2`$ (3.6)
$`A_1\beta _1(\gamma _e^{\mathrm{\hspace{0.17em}2}}+\zeta _2)+A_2\beta _1^22A_2\beta _1\gamma _eA_3\beta _1+2B_1\beta _1\gamma _e+2B_2\beta _1){\displaystyle \frac{\mathrm{ln}(1\lambda )}{(1\lambda )^2}}`$
$`(A_1\beta _1^3A_1\beta _1\beta _2){\displaystyle \frac{\mathrm{ln}(1\lambda )}{1\lambda }}+({\displaystyle \frac{1}{2}}A_1\beta _1^3A_1\beta _1\beta _2+{\displaystyle \frac{1}{2}}A_1\beta _3)\mathrm{ln}(1\lambda )+(A_1\beta _1^3`$
$`A_1\beta _1\beta _2A_1\beta _1^2\gamma _e+A_1\beta _2\gamma _eA_2\beta _1^2+A_2\beta _2+B_1\beta _1^2B_1\beta _2)({\displaystyle \frac{1}{2}}{\displaystyle \frac{1}{1\lambda }}`$
$`+{\displaystyle \frac{1}{2}}{\displaystyle \frac{1}{(1\lambda )^2}})+{\displaystyle \frac{1}{2}}({\displaystyle \frac{1}{3}}A_1\beta _1^3{\displaystyle \frac{1}{6}}A_1\beta _1\beta _2{\displaystyle \frac{1}{6}}A_1\beta _3{\displaystyle \frac{1}{3}}A_1(3\gamma _e\zeta _2+\gamma _e^{\mathrm{\hspace{0.17em}3}}+2\zeta _3)`$
$`+A_2\beta _1\gamma _eA_2(\gamma _e^{\mathrm{\hspace{0.17em}2}}+\zeta _2){\displaystyle \frac{5}{6}}A_2\beta _1^2+{\displaystyle \frac{1}{3}}A_2\beta _2+{\displaystyle \frac{5}{6}}A_3\beta _1A_3\gamma _e{\displaystyle \frac{1}{3}}A_4B_2\beta _1`$
$`+B_1(\gamma _e^{\mathrm{\hspace{0.17em}2}}+\zeta _2)+2B_2\gamma _e+B_3)(1{\displaystyle \frac{1}{(1\lambda )^2}})+{\displaystyle \frac{1}{3}}(A_1\beta _1^32A_1\beta _1\beta _2+A_1\beta _3`$
$`+A_2\beta _2A_2\beta _1^2+A_3\beta _1A_4)\lambda +L_{\mathrm{qr}}[(A_1\beta _1^2A_1\beta _2)({\displaystyle \frac{1}{2}}{\displaystyle \frac{1}{1\lambda }}+{\displaystyle \frac{1}{2}}{\displaystyle \frac{1}{(1\lambda )^2}})`$
$`+\left({\displaystyle \frac{1}{2}}A_1(\gamma _e^{\mathrm{\hspace{0.17em}2}}+\zeta _2){\displaystyle \frac{1}{2}}A_2\beta _1+A_2\gamma _e+{\displaystyle \frac{1}{2}}A_3B_1\gamma _eB_2\right)\left(1{\displaystyle \frac{1}{(1\lambda )^2}}\right)`$
$`+(A_1\beta _1\gamma _e+A_2\beta _1B_1\beta _1){\displaystyle \frac{\mathrm{ln}(1\lambda )}{(1\lambda )^2}}{\displaystyle \frac{1}{2}}A_1\beta _1^2{\displaystyle \frac{\mathrm{ln}^2(1\lambda )}{(1\lambda )^2}}]L_{\mathrm{qr}}^{\mathrm{\hspace{0.17em}2}}[{\displaystyle \frac{1}{2}}(A_1\gamma _e`$
$`+A_2B_1)(1{\displaystyle \frac{1}{(1\lambda )^2}})+{\displaystyle \frac{1}{2}}A_1\beta _1{\displaystyle \frac{\mathrm{ln}(1\lambda )}{(1\lambda )^2}}]+L_{\mathrm{qr}}^{\mathrm{\hspace{0.17em}3}}{\displaystyle \frac{1}{6}}A_1(1{\displaystyle \frac{1}{(1\lambda )^2}})`$
$`+L_{\mathrm{fr}}A_3\lambda L_{\mathrm{fr}}^{\mathrm{\hspace{0.17em}2}}(A_2+{\displaystyle \frac{1}{2}}A_1\beta _1)\lambda +L_{\mathrm{fr}}^{\mathrm{\hspace{0.17em}3}}{\displaystyle \frac{1}{3}}A_1\lambda .`$
The results $`g_i^{\mathrm{DY}}`$ for the Drell-Yan process and, with slightly different coefficients $`A_i`$ and $`D_i`$, Higgs production via gluon-gluon fusion, are related to Eqs. (3.3) – (3.6) as follows: the function corresponding to Eq. (3.3) reads $`g_1^{\mathrm{DY}}(\lambda )=2g_1^{\mathrm{DIS}}(2\lambda )`$, while the functions $`g_2^{\mathrm{DY}},g_3^{\mathrm{DY}}`$ and $`g_4^{\mathrm{DY}}`$ are obtained from Eqs. (3.4) – (3.6) by replacing $`\lambda 2\lambda `$ everywhere and substituting $`B_iD_i/2`$ in all terms. Finally, the constants have to be changed according to $`\gamma _e2\gamma _e`$ and $`\zeta _n2^n\zeta _n`$. The generalization of Eqs. (3.3) – (3.6) to other processes involving Eqs. (2.3) – (2.5) is obvious.
The functions $`g_1(\lambda )`$ collecting the leading logarithms $`L(a_\mathrm{s}L)^k`$ depend on $`A_1`$ only and are finite for all $`\lambda `$. The N$`^{n1}`$LL contributions $`g_{n>1}(\lambda )`$ to Eq. (3.1) including $`A_n`$, $`B_{n1}`$ and $`D_{n1}`$, on the other hand, exhibit Landau poles at the moments $`N=\mathrm{exp}[1/(\beta _0a_s)]`$ for DIS and $`N=\mathrm{exp}[1/(2\beta _0a_s)]`$ for the DY case (and at both values in general). As obvious from Eq. (3.3) – (3.6) the strength of these singularities increases with the logarithmic order, reaching $`[1(2)\lambda ]^2`$ at the N<sup>3</sup>LL level.
## 4 Resummation coefficients and numerical stability
Since observables are independent, order by order in $`\alpha _\mathrm{s}`$, of the factorization scale, the functions $`A_\mathrm{p}`$ in Eqs. (2.3) and (2.4) are given by the large-$`N`$ coefficients of the diagonal splitting functions for $`\mu _r=\mu _f`$,
$$P_{\mathrm{pp}}(\alpha _\mathrm{s})=A_\mathrm{p}(\alpha _\mathrm{s})\mathrm{ln}N+𝒪(1),$$
(4.1)
which in turn are identical to the anomalous dimension of a Wilson line with a cusp . The first and second order coefficients have been known for a long time, the third order has been recently completed by us. The expansion coefficients (2.6) for the quark case read
$`A_{\mathrm{q},1}`$ $`=`$ $`4C_F`$
$`A_{\mathrm{q},2}`$ $`=`$ $`8C_F\left[\left({\displaystyle \frac{67}{18}}\zeta _2\right)C_A{\displaystyle \frac{5}{9}}n_f\right]`$
$`A_{\mathrm{q},3}`$ $`=`$ $`16C_F[C_A^{\mathrm{\hspace{0.17em}2}}({\displaystyle \frac{245}{24}}{\displaystyle \frac{67}{9}}\zeta _2+{\displaystyle \frac{11}{6}}\zeta _3+{\displaystyle \frac{11}{5}}\zeta _2^{\mathrm{\hspace{0.17em}2}})+C_Fn_f({\displaystyle \frac{55}{24}}+2\zeta _3)`$ (4.2)
$`+C_An_f({\displaystyle \frac{209}{108}}+{\displaystyle \frac{10}{9}}\zeta _2{\displaystyle \frac{7}{3}}\zeta _3)+n_f^{\mathrm{\hspace{0.17em}2}}({\displaystyle \frac{1}{27}})].`$
Here $`n_f`$ denotes the number of effectively massless quark flavours, and $`C_F`$ and $`C_A`$ are the usual colour factors, with $`C_F=4/3`$ and $`C_A=3`$ in QCD. The gluonic quantities are given by
$$A_{\mathrm{g},\mathrm{i}}=C_A/C_FA_{\mathrm{q},\mathrm{i}}.$$
(4.3)
The perturbative expansion of $`A_\mathrm{p}`$ is very benign. For $`n_f=4`$, for example, Eqs. (4) lead to
$$A_\mathrm{q}(\alpha _\mathrm{s})\mathrm{\hspace{0.33em}0.4244}\alpha _\mathrm{s}(1+\mathrm{\hspace{0.25em}0.6381}\alpha _\mathrm{s}+\mathrm{\hspace{0.25em}0.5100}\alpha _\mathrm{s}^2+\mathrm{}).$$
(4.4)
Consequently, already the effect of $`A_3`$ on the resummed coefficient functions is very small , and a simple estimate suffices for the presently unknown fourth-order coefficients $`A_4`$ entering $`g_4`$. With the \[0/2\] results differing by less than 10%, we will employ the \[1/1\] Padé approximants
$`A_{\mathrm{q},4}\mathrm{\hspace{0.33em}7849},\mathrm{\hspace{0.25em}\hspace{0.25em}4313},\mathrm{\hspace{0.25em}\hspace{0.25em}1553}\text{for}n_f=\mathrm{\hspace{0.33em}3},\mathrm{\hspace{0.25em}\hspace{0.25em}4},\mathrm{\hspace{0.25em}\hspace{0.25em}5},`$ (4.5)
corresponding to an estimate of $`+\mathrm{\hspace{0.25em}0.4075}\alpha _\mathrm{s}^3`$ for the next term in Eq. (4.4).
For the determination of the coefficients $`B_i`$ we also need the constant-$`N`$ piece, $`g_0`$ in Eq. (2.1), for inclusive DIS. The corresponding expansion coefficients $`g_{0k}`$ can be obtained by Mellin inverting the +-distribution and $`\delta (1x)`$ parts of the $`k`$-th order coefficient function $`c_{\mathrm{\hspace{0.17em}2},\mathrm{q}}^{(k)}`$. Again using the expansion parameter $`a_\mathrm{s}=\alpha _\mathrm{s}/(4\pi )`$, the presently known terms are given by
$`g_{01}^{\mathrm{DIS}}`$ $`=`$ $`C_F\left(92\zeta _2+2\gamma _e^{\mathrm{\hspace{0.17em}2}}+3\gamma _e\right),`$ (4.6)
$`g_{02}^{\mathrm{DIS}}`$ $`=`$ $`C_F^{\mathrm{\hspace{0.17em}2}}({\displaystyle \frac{331}{8}}{\displaystyle \frac{51}{2}}\gamma _e{\displaystyle \frac{27}{2}}\gamma _e^{\mathrm{\hspace{0.17em}2}}+6\gamma _e^{\mathrm{\hspace{0.17em}3}}+2\gamma _e^{\mathrm{\hspace{0.17em}4}}+{\displaystyle \frac{111}{2}}\zeta _218\gamma _e\zeta _24\gamma _e^{\mathrm{\hspace{0.17em}2}}\zeta _2`$ (4.7)
$`66\zeta _3+24\gamma _e\zeta _3+{\displaystyle \frac{4}{5}}\zeta _2^{\mathrm{\hspace{0.17em}2}})+C_AC_F({\displaystyle \frac{5465}{72}}+{\displaystyle \frac{3155}{54}}\gamma _e+{\displaystyle \frac{367}{18}}\gamma _e^{\mathrm{\hspace{0.17em}2}}`$
$`+{\displaystyle \frac{22}{9}}\gamma _e^{\mathrm{\hspace{0.17em}3}}{\displaystyle \frac{1139}{18}}\zeta _2{\displaystyle \frac{22}{3}}\gamma _e\zeta _24\gamma _e^{\mathrm{\hspace{0.17em}2}}\zeta _2+{\displaystyle \frac{464}{9}}\zeta _340\gamma _e\zeta _3+{\displaystyle \frac{51}{5}}\zeta _2^{\mathrm{\hspace{0.17em}2}})`$
$`+C_Fn_f\left({\displaystyle \frac{457}{36}}{\displaystyle \frac{247}{27}}\gamma _e{\displaystyle \frac{29}{9}}\gamma _e^{\mathrm{\hspace{0.17em}2}}{\displaystyle \frac{4}{9}}\gamma _e^{\mathrm{\hspace{0.17em}3}}+{\displaystyle \frac{85}{9}}\zeta _2+{\displaystyle \frac{4}{3}}\gamma _e\zeta _2+{\displaystyle \frac{4}{9}}\zeta _3\right)`$
and
$`g_{03}^{\mathrm{DIS}}`$ $`=`$ $`C_F^{\mathrm{\hspace{0.17em}3}}({\displaystyle \frac{7255}{24}}+{\displaystyle \frac{1001}{8}}\gamma _e+{\displaystyle \frac{187}{4}}\gamma _e^{\mathrm{\hspace{0.17em}2}}{\displaystyle \frac{93}{2}}\gamma _e^{\mathrm{\hspace{0.17em}3}}9\gamma _e^{\mathrm{\hspace{0.17em}4}}+6\gamma _e^{\mathrm{\hspace{0.17em}5}}+{\displaystyle \frac{4}{3}}\gamma _e^{\mathrm{\hspace{0.17em}6}}{\displaystyle \frac{6197}{12}}\zeta _2`$ (4.8)
$`+{\displaystyle \frac{579}{2}}\gamma _e\zeta _2+66\gamma _e^{\mathrm{\hspace{0.17em}2}}\zeta _236\gamma _e^{\mathrm{\hspace{0.17em}3}}\zeta _24\gamma _e^{\mathrm{\hspace{0.17em}4}}\zeta _2411\zeta _3346\gamma _e\zeta _360\gamma _e^{\mathrm{\hspace{0.17em}2}}\zeta _3`$
$`+48\gamma _e^{\mathrm{\hspace{0.17em}3}}\zeta _3{\displaystyle \frac{1791}{5}}\zeta _2^{\mathrm{\hspace{0.17em}2}}+84\gamma _e\zeta _2^{\mathrm{\hspace{0.17em}2}}+{\displaystyle \frac{8}{5}}\gamma _e^{\mathrm{\hspace{0.17em}2}}\zeta _2^{\mathrm{\hspace{0.17em}2}}+556\zeta _2\zeta _380\gamma _e\zeta _2\zeta _3+1384\zeta _5`$
$`240\gamma _e\zeta _5+{\displaystyle \frac{8144}{315}}\zeta _2^{\mathrm{\hspace{0.17em}3}}{\displaystyle \frac{176}{3}}\zeta _3^{\mathrm{\hspace{0.17em}2}})+C_AC^{\mathrm{\hspace{0.17em}2}}_F({\displaystyle \frac{9161}{12}}{\displaystyle \frac{16981}{24}}\gamma _e{\displaystyle \frac{5563}{36}}\gamma _e^{\mathrm{\hspace{0.17em}2}}`$
$`+{\displaystyle \frac{8425}{54}}\gamma _e^{\mathrm{\hspace{0.17em}3}}+{\displaystyle \frac{433}{9}}\gamma _e^{\mathrm{\hspace{0.17em}4}}+{\displaystyle \frac{44}{9}}\gamma _e^{\mathrm{\hspace{0.17em}5}}+{\displaystyle \frac{191545}{108}}\zeta _2{\displaystyle \frac{28495}{54}}\gamma _e\zeta _2{\displaystyle \frac{592}{3}}\gamma _e^{\mathrm{\hspace{0.17em}2}}\zeta _28\gamma _e^{\mathrm{\hspace{0.17em}4}}\zeta _2`$
$`{\displaystyle \frac{284}{9}}\gamma _e^{\mathrm{\hspace{0.17em}3}}\zeta _2{\displaystyle \frac{49346}{27}}\zeta _3+752\gamma _e\zeta _3+{\displaystyle \frac{640}{9}}\gamma _e^{\mathrm{\hspace{0.17em}2}}\zeta _380\gamma _e^{\mathrm{\hspace{0.17em}3}}\zeta _3+{\displaystyle \frac{11419}{27}}\zeta _2^{\mathrm{\hspace{0.17em}2}}`$
$`+{\displaystyle \frac{299}{3}}\gamma _e\zeta _2^{\mathrm{\hspace{0.17em}2}}+{\displaystyle \frac{142}{5}}\gamma _e^{\mathrm{\hspace{0.17em}2}}\zeta _2^{\mathrm{\hspace{0.17em}2}}828\zeta _2\zeta _3+96\gamma _e\zeta _2\zeta _3{\displaystyle \frac{3896}{9}}\zeta _5+120\gamma _e\zeta _5`$
$`{\displaystyle \frac{23098}{315}}\zeta _2^{\mathrm{\hspace{0.17em}3}}+{\displaystyle \frac{536}{3}}\zeta _3^{\mathrm{\hspace{0.17em}2}})+C^{\mathrm{\hspace{0.17em}2}}_AC_F({\displaystyle \frac{1909753}{1944}}+{\displaystyle \frac{599375}{729}}\gamma _e+{\displaystyle \frac{50689}{162}}\gamma _e^{\mathrm{\hspace{0.17em}2}}`$
$`+{\displaystyle \frac{4649}{81}}\gamma _e^{\mathrm{\hspace{0.17em}3}}+{\displaystyle \frac{121}{27}}\gamma _e^{\mathrm{\hspace{0.17em}4}}{\displaystyle \frac{78607}{54}}\zeta _2{\displaystyle \frac{18179}{81}}\gamma _e\zeta _2{\displaystyle \frac{778}{9}}\gamma _e^{\mathrm{\hspace{0.17em}2}}\zeta _2{\displaystyle \frac{88}{9}}\gamma _e^{\mathrm{\hspace{0.17em}3}}\zeta _2`$
$`+{\displaystyle \frac{115010}{81}}\zeta _3+{\displaystyle \frac{121}{27}}\gamma _e^{\mathrm{\hspace{0.17em}4}}{\displaystyle \frac{78607}{54}}\zeta _2{\displaystyle \frac{18179}{81}}\gamma _e\zeta _2{\displaystyle \frac{778}{9}}\gamma _e^{\mathrm{\hspace{0.17em}2}}\zeta _2{\displaystyle \frac{88}{9}}\gamma _e^{\mathrm{\hspace{0.17em}3}}\zeta _2`$
$`+{\displaystyle \frac{3496}{9}}\zeta _2\zeta _3+{\displaystyle \frac{176}{3}}\gamma _e\zeta _2\zeta _3{\displaystyle \frac{416}{3}}\zeta _5+232\gamma _e\zeta _5{\displaystyle \frac{12016}{315}}\zeta _2^{\mathrm{\hspace{0.17em}3}}{\displaystyle \frac{248}{3}}\zeta _3^{\mathrm{\hspace{0.17em}2}})`$
$`+C_F^{\mathrm{\hspace{0.17em}2}}n_f({\displaystyle \frac{341}{36}}+{\displaystyle \frac{2003}{108}}\gamma _e+{\displaystyle \frac{83}{18}}\gamma _e^{\mathrm{\hspace{0.17em}2}}{\displaystyle \frac{683}{27}}\gamma _e^{\mathrm{\hspace{0.17em}3}}{\displaystyle \frac{70}{9}}\gamma _e^{\mathrm{\hspace{0.17em}4}}{\displaystyle \frac{8}{9}}\gamma _e^{\mathrm{\hspace{0.17em}5}}{\displaystyle \frac{10733}{54}}\zeta _2`$
$`+{\displaystyle \frac{2177}{27}}\gamma _e\zeta _2+{\displaystyle \frac{112}{3}}\gamma _e^{\mathrm{\hspace{0.17em}2}}\zeta _2+{\displaystyle \frac{32}{9}}\gamma _e^{\mathrm{\hspace{0.17em}3}}\zeta _2+{\displaystyle \frac{10766}{27}}\zeta _3{\displaystyle \frac{20}{9}}\gamma _e\zeta _3+{\displaystyle \frac{8}{9}}\gamma _e^{\mathrm{\hspace{0.17em}2}}\zeta _3{\displaystyle \frac{8}{3}}\gamma _e\zeta _2^{\mathrm{\hspace{0.17em}2}}`$
$`{\displaystyle \frac{10802}{135}}\zeta _2^{\mathrm{\hspace{0.17em}2}}{\displaystyle \frac{40}{3}}\zeta _2\zeta _3{\displaystyle \frac{784}{9}}\zeta _5)+C_AC_Fn_f({\displaystyle \frac{142883}{486}}{\displaystyle \frac{160906}{729}}\gamma _e`$
$`{\displaystyle \frac{7531}{81}}\gamma _e^{\mathrm{\hspace{0.17em}2}}{\displaystyle \frac{1552}{81}}\gamma _e^{\mathrm{\hspace{0.17em}3}}{\displaystyle \frac{44}{27}}\gamma _e^{\mathrm{\hspace{0.17em}4}}+{\displaystyle \frac{33331}{81}}\zeta _2+{\displaystyle \frac{5264}{81}}\gamma _e\zeta _2+{\displaystyle \frac{56}{3}}\gamma _e^{\mathrm{\hspace{0.17em}2}}\zeta _2+{\displaystyle \frac{16}{9}}\gamma _e^{\mathrm{\hspace{0.17em}3}}\zeta _2`$
$`{\displaystyle \frac{21418}{81}}\zeta _3+{\displaystyle \frac{1976}{27}}\gamma _e\zeta _3+8\gamma _e^{\mathrm{\hspace{0.17em}2}}\zeta _3+{\displaystyle \frac{164}{135}}\zeta _2^{\mathrm{\hspace{0.17em}2}}{\displaystyle \frac{128}{15}}\gamma _e\zeta _2^{\mathrm{\hspace{0.17em}2}}{\displaystyle \frac{64}{9}}\zeta _2\zeta _3+{\displaystyle \frac{8}{3}}\zeta _5)`$
$`+C_Fn_f^{\mathrm{\hspace{0.17em}2}}({\displaystyle \frac{9517}{486}}+{\displaystyle \frac{8714}{729}}\gamma _e+{\displaystyle \frac{470}{81}}\gamma _e^{\mathrm{\hspace{0.17em}2}}+{\displaystyle \frac{116}{81}}\gamma _e^{\mathrm{\hspace{0.17em}3}}+{\displaystyle \frac{4}{27}}\gamma _e^{\mathrm{\hspace{0.17em}4}}{\displaystyle \frac{2110}{81}}\zeta _2{\displaystyle \frac{8}{9}}\gamma _e^{\mathrm{\hspace{0.17em}2}}\zeta _2`$
$`{\displaystyle \frac{116}{27}}\gamma _e\zeta _2+{\displaystyle \frac{80}{81}}\zeta _3+{\displaystyle \frac{64}{27}}\gamma _e\zeta _3{\displaystyle \frac{292}{135}}\zeta _2^{\mathrm{\hspace{0.17em}2}})+{\displaystyle \frac{d^{abc}d_{abc}}{n_c}}fl_{11}(64+160\zeta _2`$
$`+{\displaystyle \frac{224}{3}}\zeta _3{\displaystyle \frac{32}{5}}\zeta _2^{\mathrm{\hspace{0.17em}2}}{\displaystyle \frac{1280}{3}}\zeta _5).`$
Note the new flavour structure $`fl_{11}`$ in $`g_{03}`$. This contribution, introducing the colour factor $`d^{abc}d_{abc}/n_c`$, for the first time leads to a difference between the flavour-singlet and non-singlet coefficient functions for the photon-exchange structure function $`F_2`$ in the soft-gluon limit, with $`fl_{11}^{\mathrm{ns}}=3e`$ and $`fl_{11}^\mathrm{s}=e^2/e^{\mathrm{\hspace{0.17em}2}}`$, where $`e^k`$ represents the average of the charge $`e^k`$ for the active quark flavours, $`e^k=n_f^1_{i=1}^{n_f}e_i^k`$. Correspondingly, the large-$`N`$ coefficient functions for $`Z`$\- and $`W`$-exchange DIS will differ from each other and from $`F_2^{\mathrm{e}.\mathrm{m}.}`$ at order $`\alpha _\mathrm{s}^3`$. Based on the size of the $`fl_{11}`$ term in Eq. (4.8), however, we expect these differences to be numerically insignificant.
Now the coefficients $`B_{\mathrm{q},k}`$ entering the jet function (2.4) can be derived successively from the $`\mathrm{ln}N`$ terms of the $`k`$-loop DIS coefficient functions $`c_{\mathrm{\hspace{0.17em}2},\mathrm{q}}^{(k)}`$. Expansion of Eqs. (3.1) and (3.3) – (3.6) in powers of $`a_\mathrm{s}`$ yields
$`c_{\mathrm{\hspace{0.17em}2},\mathrm{q}}^{(1)}|_{\mathrm{ln}N}`$ $`=`$ $`A_1\gamma _eB_1`$
$`c_{\mathrm{\hspace{0.17em}2},\mathrm{q}}^{(2)}|_{\mathrm{ln}N}`$ $`=`$ $`{\displaystyle \frac{1}{2}}A_1\beta _0(\gamma _e^{\mathrm{\hspace{0.17em}2}}+\zeta _2)+A_2\gamma _eB_1\beta _0\gamma _eB_2+g_{01}(A_1\gamma _eB_1)`$
$`c_{\mathrm{\hspace{0.17em}2},\mathrm{q}}^{(3)}|_{\mathrm{ln}N}`$ $`=`$ $`{\displaystyle \frac{1}{3}}A_1\beta _0^2(\gamma _e^{\mathrm{\hspace{0.17em}3}}+3\gamma _e\zeta _2+2\zeta _3)+{\displaystyle \frac{1}{2}}A_1\beta _1(\gamma _e^{\mathrm{\hspace{0.17em}2}}+\zeta _2)+A_2\beta _0(\gamma _e^{\mathrm{\hspace{0.17em}2}}+\zeta _2)`$ (4.9)
$`+A_3\gamma _eB_1(\beta _1\gamma _e+\beta _0^2\zeta _2+\beta _0^2\gamma _e^{\mathrm{\hspace{0.17em}2}})2B_2\beta _0\gamma _eB_3`$
$`+g_{02}(A_1\gamma _eB_1)+g_{01}\left({\displaystyle \frac{1}{2}}A_1\beta _0(\gamma _e^{\mathrm{\hspace{0.17em}2}}+\zeta _2)+A_2\gamma _eB_1\beta _0\gamma _eB_2\right).`$
The coefficients of $`\mathrm{ln}^lN`$, $`2l2k`$ in $`c_{\mathrm{\hspace{0.17em}2},\mathrm{q}}^{(k)}`$, on the other hand, are completely fixed by lower-order resummation coefficients, thus providing an explicit $`k`$-loop check of the exponentiation formula. Comparison of the relations (4.9) with the corresponding results from the fixed-order calculations of Refs. , using Eqs. (4), (4.6) and (4.7), leads to
$`B_{\mathrm{q},1}`$ $`=`$ $`3C_F,`$ (4.10)
$`B_{\mathrm{q},2}`$ $`=`$ $`C_F^{\mathrm{\hspace{0.17em}2}}\left[{\displaystyle \frac{3}{2}}+12\zeta _224\zeta _3\right]+C_FC_A\left[{\displaystyle \frac{3155}{54}}+{\displaystyle \frac{44}{3}}\zeta _2+40\zeta _3\right]`$ (4.11)
$`+C_Fn_f\left[{\displaystyle \frac{247}{27}}{\displaystyle \frac{8}{3}}\zeta _2\right],`$
$`B_{\mathrm{q},3}`$ $`=`$ $`C_F^{\mathrm{\hspace{0.17em}3}}\left[{\displaystyle \frac{29}{2}}18\zeta _268\zeta _3{\displaystyle \frac{288}{5}}\zeta _2^{\mathrm{\hspace{0.17em}2}}+32\zeta _2\zeta _3+240\zeta _5\right]`$ (4.12)
$`+C_AC_F^{\mathrm{\hspace{0.17em}2}}\left[46+287\zeta _2{\displaystyle \frac{712}{3}}\zeta _3{\displaystyle \frac{272}{5}}\zeta _2^{\mathrm{\hspace{0.17em}2}}16\zeta _2\zeta _3120\zeta _5\right]`$
$`+C_A^{\mathrm{\hspace{0.17em}2}}C_F\left[{\displaystyle \frac{599375}{729}}+{\displaystyle \frac{32126}{81}}\zeta _2+{\displaystyle \frac{21032}{27}}\zeta _3{\displaystyle \frac{652}{15}}\zeta _2^{\mathrm{\hspace{0.17em}2}}{\displaystyle \frac{176}{3}}\zeta _2\zeta _3232\zeta _5\right]`$
$`+C_F^{\mathrm{\hspace{0.17em}2}}n_f\left[{\displaystyle \frac{5501}{54}}50\zeta _2+{\displaystyle \frac{32}{9}}\zeta _3\right]+C_Fn_f^{\mathrm{\hspace{0.17em}2}}\left[{\displaystyle \frac{8714}{729}}+{\displaystyle \frac{232}{27}}\zeta _2{\displaystyle \frac{32}{27}}\zeta _3\right]`$
$`+C_AC_Fn_f\left[{\displaystyle \frac{160906}{729}}{\displaystyle \frac{9920}{81}}\zeta _2{\displaystyle \frac{776}{9}}\zeta _3+{\displaystyle \frac{208}{15}}\zeta _2^{\mathrm{\hspace{0.17em}2}}\right].`$
Eq. (4.10) is, of course, a well-known result . Eq. (4.11) has been derived by us before , establishing $`D_2^{\mathrm{DIS}}=0`$ from the $`n_f\mathrm{ln}^2N`$ term at three loops. For our new result (4.12), on the other hand, we have to rely on the subsequent all-order proofs of Eq. (2.7) in Refs. . The QCD expansion of $`B_\mathrm{q}`$ analogous to Eq. (4.4) appears far less stable than that for $`A_\mathrm{q}`$,
$$B_\mathrm{q}(\alpha _\mathrm{s},n_f=4)0.3183\alpha _\mathrm{s}(1\mathrm{\hspace{0.25em}1.227}\alpha _\mathrm{s}\mathrm{\hspace{0.25em}3.405}\alpha _\mathrm{s}^2+\mathrm{}).$$
(4.13)
The ingredients for the resummation of inclusive DIS are now complete, and in the left part of Fig. 1 we show the corresponding LL, NLL, N<sup>2</sup>LL and N<sup>3</sup>LL approximations to the exponent (3.1) resulting, for $`\alpha _\mathrm{s}=0.2`$ and three flavours, from Eqs. (3.3) – (3.6), (4), (4.5) and (4.10) – (4.12). For these parameters the expansion (3.1) is stable in the $`N`$-range shown in the figure. For example, the relative N<sup>3</sup>LL corrections amount to 2% at $`N=10`$ ($`\lambda =0.33`$) and 4% at $`N=40`$ ($`\lambda =0.53`$), whereas the corresponding N<sup>2</sup>LL figures read 9% and 12%. The large third-order contribution to $`B_\mathrm{q}`$ actually stabilizes $`g_4(\lambda )`$: for $`B_{\mathrm{q},3}=0`$ the N<sup>3</sup>LL term at $`N=40`$ would instead reach 12%, i.e., the size of the previous order. The effect of both $`A_{\mathrm{q},4}`$ and $`\beta _3`$, on the other hand, is very small, as their respective nullification would change the result even at $`N=40`$ by only 0.6% and 0.1%.
In the right part of Fig. 1 and in Fig. 2 the exponentiated results are convoluted with the typical input shape $`xf=x^{\mathrm{\hspace{0.17em}0.5}}(1x)^3`$ for a couple of values for $`\alpha _\mathrm{s}`$ and $`n_f`$. The Mellin inversion is in principle ambiguous due to the Landau poles briefly addressed at the end of Section 3. We employ the standard ‘minimal prescription’ (thus adopting the usual fixed-order contour) of Ref. , to which the reader is referred for a detailed discussion. For total soft-gluon enhancements up to almost an order of magnitude, as shown in the figures, the resulting N<sup>3</sup>LL corrections remain far smaller than their N<sup>2</sup>LL counterparts and amount to less than 10% even for $`\alpha _\mathrm{s}=0.3`$. Note that the dependence on $`n_f`$ is larger than the effect of $`g_4`$. Thus, at this level of accuracy, a reliable understanding of heavy-quark mass effects is called for also in the limit $`x1`$.
The gluonic coefficients corresponding to Eqs. (4.10) – (4.12) can be obtained in the same manner from DIS by exchange of a scalar $`\varphi `$ with a pointlike coupling to gluons, like the Higgs boson in limit of a heavy top quark. We have derived the corresponding coefficient functions $`c_{\varphi ,\mathrm{p}}^{(k)}`$ up to $`k=3`$ already during the calculations for Ref. , as a process of this type is required to access the lower row of the flavour-singlet splitting function matrix. Comparing those results to Eqs. (4.9) for $`c_{\varphi ,\mathrm{g}}^{(k)}`$ yields
$`B_{\mathrm{g},1}`$ $`=`$ $`{\displaystyle \frac{11}{3}}C_A+{\displaystyle \frac{2}{3}}n_f=\beta _0,`$ (4.14)
$`B_{\mathrm{g},2}`$ $`=`$ $`C_A^{\mathrm{\hspace{0.17em}2}}\left[{\displaystyle \frac{611}{9}}+{\displaystyle \frac{88}{3}}\zeta _2+16\zeta _3\right]+C_An_f\left[{\displaystyle \frac{428}{27}}{\displaystyle \frac{16}{3}}\zeta _2\right]+\mathrm{\hspace{0.25em}2}C_Fn_f{\displaystyle \frac{20}{27}}n_f^{\mathrm{\hspace{0.17em}2}},`$ (4.15)
$`B_{\mathrm{g},3}`$ $`=`$ $`C_A^{\mathrm{\hspace{0.17em}3}}\left[{\displaystyle \frac{1492081}{1458}}+{\displaystyle \frac{60875}{81}}\zeta _2+{\displaystyle \frac{13796}{27}}\zeta _3{\displaystyle \frac{2596}{15}}\zeta _2^{\mathrm{\hspace{0.17em}2}}{\displaystyle \frac{128}{3}}\zeta _2\zeta _3112\zeta _5\right]`$ (4.16)
$`+C_A^{\mathrm{\hspace{0.17em}2}}n_f\left[{\displaystyle \frac{498329}{1458}}{\displaystyle \frac{21014}{81}}\zeta _2{\displaystyle \frac{296}{9}}\zeta _3+{\displaystyle \frac{568}{15}}\zeta _2^{\mathrm{\hspace{0.17em}2}}\right]C_F^{\mathrm{\hspace{0.17em}2}}n_f`$
$`+C_AC_Fn_f\left[{\displaystyle \frac{8579}{54}}16\zeta _2{\displaystyle \frac{832}{9}}\zeta _3{\displaystyle \frac{32}{5}}\zeta _2^{\mathrm{\hspace{0.17em}2}}\right]+C_Fn_f^{\mathrm{\hspace{0.17em}2}}\left[{\displaystyle \frac{47}{3}}+{\displaystyle \frac{32}{3}}\zeta _3\right]`$
$`+C_An_f^{\mathrm{\hspace{0.17em}2}}\left[{\displaystyle \frac{48829}{1458}}+{\displaystyle \frac{716}{27}}\zeta _2{\displaystyle \frac{176}{27}}\zeta _3\right]+n_f^{\mathrm{\hspace{0.17em}3}}\left[{\displaystyle \frac{200}{243}}{\displaystyle \frac{8}{9}}\zeta _2\right],`$
where Eqs. (4.15) and (4.16) are new results. For $`n_f=4`$ the numerical expansion of $`B_\mathrm{g}`$ reads
$$B_\mathrm{g}(\alpha _\mathrm{s})0.6631\alpha _\mathrm{s}(1\mathrm{\hspace{0.25em}0.7651}\alpha _\mathrm{s}\mathrm{\hspace{0.25em}2.696}\alpha _\mathrm{s}^2+\mathrm{}),$$
(4.17)
exhibiting an enhanced third order correction similar to that of $`B_\mathrm{q}`$ in Eq. (4.13).
The gluonic threshold resummation resulting from Eqs. (4.3) and (4.14) – (4.16) is illustrated in Fig. 3 using the practically irrelevant scalar-exchange process, with same parameters as in Fig. 1 for direct comparison. The soft and collinear radiation effects are much larger here due to the larger colour charge of the gluons, but the qualitative pattern is rather similar to ‘normal’ inclusive DIS.
As mentioned above, the (closely related) Drell-Yan process and Higgs boson production via gluon-gluon fusion presently represent the only other processes for which the NNLL threshold resummation is known, with $`D_1=0`$ and
$$D_2^{\{\mathrm{DY},\mathrm{H}\}}=\{C_F,C_A\}\left[C_A\left(\frac{1616}{27}+\frac{176}{3}\zeta _2+56\zeta _3\right)+n_f\left(\frac{224}{27}\frac{32}{3}\zeta _2\right)\right].$$
(4.18)
Extending a result given in Ref. , we notice the following conspicuous relation between these coefficients and $`B_{\mathrm{p},2}`$:
$`{\displaystyle \frac{1}{2}}D_2^{\mathrm{DY}}B_{\mathrm{q},2}P_{\mathrm{q},\delta }^{(1)}`$ $`=`$ $`7\beta _0C_F`$
$`{\displaystyle \frac{1}{2}}D_2^\mathrm{H}B_{\mathrm{g},2}P_{\mathrm{g},\delta }^{(1)}`$ $`=`$ $`{\displaystyle \frac{1}{3}}\beta _0\left(4C_A+5\beta _0\right),`$ (4.19)
where $`P_{\mathrm{p},\delta }^{(1)}`$ denotes the coefficients of $`\delta (1x)`$ in the diagonal two-loop splitting functions, and the colour structures on the right-hand sides are those of $`A_{\mathrm{p},1}`$ and $`B_{\mathrm{p},1}=P_{\mathrm{p},\delta }^{(0)}`$, multiplied by $`\beta _0`$. Note especially the non-trivial cancellation of all $`\zeta `$-function terms between the three contributions on the left-hand sides of Eqs. (4.19) and the vanishing of the right-hand sides for $`\beta _00`$.
## 5 Fourth-order predictions and tower expansion
Another manner to organize the all-order information encoded in Eqs. (2.1) – (2.5) is to re-expand the exponential,
$$C^N(Q^2)/C_{\mathrm{LO}}^N(Q^2)=\mathrm{\hspace{0.33em}1}+\underset{k=1}{\overset{\mathrm{}}{}}a_\mathrm{s}^k\underset{l=1}{\overset{2k}{}}c_{kl}\mathrm{ln}^{\mathrm{\hspace{0.17em}2}kl+1}N,$$
(5.1)
and retain only those terms in the second sum which are completely fixed by the available information on the expansion coefficients in Eq. (2.6). Using the notation
$$g_i(\lambda )=\underset{k=1}{\overset{\mathrm{}}{}}g_{ik}\lambda ^k$$
(5.2)
for the expansion of Eqs. (3.3) – (3.6) together with Eq. (2.6) for $`g_0`$, the quantities $`c_{kl}`$ in Eq. (5.1) receive contributions from the following coefficients:
$`c_{k1}`$ $`:`$ $`g_{11}`$
$`c_{k2}`$ $`:`$ $`+g_{12},g_{21}`$
$`c_{k3}`$ $`:`$ $`+g_{13},g_{22},g_{01}`$
$`c_{k4}`$ $`:`$ $`+g_{14},g_{23},g_{31}`$
$`c_{k5}`$ $`:`$ $`+g_{15},g_{24},g_{32},g_{02}`$
$`c_{k6}`$ $`:`$ $`+g_{16},g_{25},g_{33},g_{41}`$
$`c_{k7}`$ $`:`$ $`+g_{17},g_{26},g_{34},g_{42},g_{03}`$
$`c_{k8}`$ $`:`$ $`+g_{18},g_{27},g_{35},g_{43},g_{51}\mathrm{}.`$ (5.3)
The complete relations for the first four terms $`c_{k1}\mathrm{}c_{k4}`$ can be found in Ref. in a slightly different notation, $`g_{31}g_{32}`$ (i.e., the second index denoting the total power of $`a_\mathrm{s}`$ in Eq. (3.1)). Note that the quantities $`c_{kl}`$ vanish factorially for $`k\mathrm{}`$ and fixed $`l`$.
Taking into account Eq. (2.7) and considering the coefficient $`A_i`$ as either known or irrelevant, the function $`g_i`$ for inclusive DIS is completely specified by its leading term, obtained by matching to the $`i`$-th order calculation of the coefficient functions as in Eqs. (4.9). The same holds for other processes, at least for cases like Eqs. (2), once the required coefficients $`B_i`$ are known from DIS.
The leading three towers of logarithms, $`c_{kl}`$ for any $`k`$ and $`l=1,\mathrm{\hspace{0.25em}2},\mathrm{\hspace{0.25em}3}`$, are fixed by a one-loop calculation (providing $`g_{i1}`$ for $`i=0,\mathrm{\hspace{0.17em}1},\mathrm{\hspace{0.17em}2}`$) together with the NLL resummation (adding $`g_{1k}`$ and $`g_{2k}`$ for $`k2`$). This is the status for many important observables . Correspondingly, a two-loop computation of the process under consideration specifies $`g_{31}`$ and $`g_{02}`$ and hence fixes, together with the NNLL resummation, also the next two towers. This is the accuracy reached for the Drell-Yan process and Higgs production . Finally a three-loop computation combined with the N<sup>3</sup>LL resummation fixes the first seven towers, $`c_{kl}`$ for $`l=1,\mathrm{},\mathrm{\hspace{0.25em}7}`$. With the results of Ref. and Sections 3 and 4, we have now reached this point for the structure function $`F_2`$ in DIS.
The resulting four-loop predictions, in $`x`$-space expressed in terms of the coefficients of the +-distributions $`𝒟_k=[(1x)^1\mathrm{ln}(1x)]_+`$, for the six highest terms read
$`c_{\mathrm{\hspace{0.17em}2},\mathrm{q}}^{(4)}|_{𝒟_7}`$ $`=`$ $`{\displaystyle \frac{16}{3}}C_F^{\mathrm{\hspace{0.17em}4}},`$ (5.4)
$`c_{\mathrm{\hspace{0.17em}2},\mathrm{q}}^{(4)}|_{𝒟_6}`$ $`=`$ $`28C_F^{\mathrm{\hspace{0.17em}4}}{\displaystyle \frac{308}{9}}C_AC_F^{\mathrm{\hspace{0.17em}3}}+{\displaystyle \frac{56}{9}}C_F^{\mathrm{\hspace{0.17em}3}}n_f,`$ (5.5)
$`c_{\mathrm{\hspace{0.17em}2},\mathrm{q}}^{(4)}|_{𝒟_5}`$ $`=`$ $`C_F^{\mathrm{\hspace{0.17em}4}}\left[18128\zeta _2\right]+C_AC_F^{\mathrm{\hspace{0.17em}3}}\left[{\displaystyle \frac{998}{3}}48\zeta _2\right]+{\displaystyle \frac{1936}{27}}C_A^{\mathrm{\hspace{0.17em}2}}C_F^{\mathrm{\hspace{0.17em}2}}`$ (5.6)
$`{\displaystyle \frac{164}{3}}C_F^{\mathrm{\hspace{0.17em}3}}n_f{\displaystyle \frac{704}{27}}C_AC_F^{\mathrm{\hspace{0.17em}2}}n_f+{\displaystyle \frac{64}{27}}C_F^{\mathrm{\hspace{0.17em}2}}n_f^{\mathrm{\hspace{0.17em}2}},`$
$`c_{\mathrm{\hspace{0.17em}2},\mathrm{q}}^{(4)}|_{𝒟_4}`$ $`=`$ $`C_F^{\mathrm{\hspace{0.17em}4}}\left[210+600\zeta _2+{\displaystyle \frac{400}{3}}\zeta _3\right]+C_AC_F^{\mathrm{\hspace{0.17em}3}}\left[{\displaystyle \frac{27835}{27}}+{\displaystyle \frac{6800}{9}}\zeta _2+400\zeta _3\right]`$ (5.7)
$`+C_A^{\mathrm{\hspace{0.17em}2}}C_F^{\mathrm{\hspace{0.17em}2}}\left[{\displaystyle \frac{24040}{27}}+{\displaystyle \frac{440}{3}}\zeta _2\right]{\displaystyle \frac{1331}{27}}C_A^{\mathrm{\hspace{0.17em}3}}C_F+C_F^{\mathrm{\hspace{0.17em}3}}n_f\left[{\displaystyle \frac{4630}{27}}{\displaystyle \frac{1040}{9}}\zeta _2\right]`$
$`+C_AC_F^{\mathrm{\hspace{0.17em}2}}n_f\left[{\displaystyle \frac{8120}{27}}{\displaystyle \frac{80}{3}}\zeta _2\right]+{\displaystyle \frac{242}{9}}C_A^{\mathrm{\hspace{0.17em}2}}C_Fn_f{\displaystyle \frac{640}{27}}C_F^{\mathrm{\hspace{0.17em}2}}n_f^{\mathrm{\hspace{0.17em}2}}{\displaystyle \frac{44}{9}}C_AC_Fn_f^{\mathrm{\hspace{0.17em}2}}`$
$`+{\displaystyle \frac{8}{27}}C_Fn_f^{\mathrm{\hspace{0.17em}3}},`$
$`c_{\mathrm{\hspace{0.17em}2},\mathrm{q}}^{(4)}|_{𝒟_3}`$ $`=`$ $`C_F^{\mathrm{\hspace{0.17em}4}}\left[{\displaystyle \frac{113}{2}}+264\zeta _21072\zeta _3+{\displaystyle \frac{1392}{5}}\zeta _2^{\mathrm{\hspace{0.17em}2}}\right]+C_A^{\mathrm{\hspace{0.17em}3}}C_F\left[{\displaystyle \frac{55627}{81}}{\displaystyle \frac{968}{9}}\zeta _2\right]`$ (5.8)
$`+C_AC_F^{\mathrm{\hspace{0.17em}3}}\left[{\displaystyle \frac{1534}{3}}{\displaystyle \frac{41824}{9}}\zeta _2{\displaystyle \frac{8800}{9}}\zeta _3+{\displaystyle \frac{3128}{5}}\zeta _2^{\mathrm{\hspace{0.17em}2}}\right]`$
$`+C_A^{\mathrm{\hspace{0.17em}2}}C_F^{\mathrm{\hspace{0.17em}2}}\left[{\displaystyle \frac{2154563}{486}}{\displaystyle \frac{52912}{27}}\zeta _2{\displaystyle \frac{13024}{9}}\zeta _3+{\displaystyle \frac{864}{5}}\zeta _2^{\mathrm{\hspace{0.17em}2}}\right]`$
$`+C_F^{\mathrm{\hspace{0.17em}3}}n_f\left[{\displaystyle \frac{280}{3}}+{\displaystyle \frac{7216}{9}}\zeta _2+{\displaystyle \frac{1888}{9}}\zeta _3\right]+C_A^{\mathrm{\hspace{0.17em}2}}C_Fn_f\left[{\displaystyle \frac{9502}{27}}+{\displaystyle \frac{352}{9}}\zeta _2\right]`$
$`C_AC_F^{\mathrm{\hspace{0.17em}2}}n_f\left[{\displaystyle \frac{339134}{243}}{\displaystyle \frac{14096}{27}}\zeta _2{\displaystyle \frac{1216}{9}}\zeta _3\right]+C_AC_Fn_f^{\mathrm{\hspace{0.17em}2}}\left[{\displaystyle \frac{1540}{27}}{\displaystyle \frac{32}{9}}\zeta _2\right]`$
$`+C_F^2n_f^{\mathrm{\hspace{0.17em}2}}\left[{\displaystyle \frac{24238}{243}}{\displaystyle \frac{928}{27}}\zeta _2\right]{\displaystyle \frac{232}{81}}C_Fn_f^{\mathrm{\hspace{0.17em}3}},`$
$`c_{\mathrm{\hspace{0.17em}2},\mathrm{q}}^{(4)}|_{𝒟_2}`$ $`=`$ $`C_F^{\mathrm{\hspace{0.17em}4}}\left[{\displaystyle \frac{1299}{2}}2808\zeta _2+1392\zeta _31836\zeta _2^{\mathrm{\hspace{0.17em}2}}640\zeta _2\zeta _3+4128\zeta _5\right]`$
$`+C_AC_F^{\mathrm{\hspace{0.17em}3}}\left[{\displaystyle \frac{13990}{3}}+{\displaystyle \frac{30704}{3}}\zeta _2+{\displaystyle \frac{2716}{3}}\zeta _3{\displaystyle \frac{12906}{5}}\zeta _2^{\mathrm{\hspace{0.17em}2}}3648\zeta _2\zeta _3720\zeta _5\right]`$
$`C_A^{\mathrm{\hspace{0.17em}2}}C_F^{\mathrm{\hspace{0.17em}2}}\left[{\displaystyle \frac{2254339}{243}}{\displaystyle \frac{86804}{9}}\zeta _2{\displaystyle \frac{24544}{3}}\zeta _3+{\displaystyle \frac{4034}{3}}\zeta _2^{\mathrm{\hspace{0.17em}2}}+832\zeta _2\zeta _3+1392\zeta _5\right]`$
$`+C_A^{\mathrm{\hspace{0.17em}3}}C_F\left[{\displaystyle \frac{649589}{162}}+{\displaystyle \frac{4012}{3}}\zeta _2+1452\zeta _3{\displaystyle \frac{968}{5}}\zeta _2^{\mathrm{\hspace{0.17em}2}}\right]`$
$`+C_AC_F^{\mathrm{\hspace{0.17em}2}}n_f\left[{\displaystyle \frac{713162}{243}}{\displaystyle \frac{82004}{27}}\zeta _2{\displaystyle \frac{10600}{9}}\zeta _3+{\displaystyle \frac{3772}{15}}\zeta _2^{\mathrm{\hspace{0.17em}2}}\right]`$
$`+C_A^{\mathrm{\hspace{0.17em}2}}C_Fn_f\left[{\displaystyle \frac{17189}{9}}{\displaystyle \frac{5096}{9}}\zeta _2352\zeta _3+{\displaystyle \frac{176}{5}}\zeta _2^{\mathrm{\hspace{0.17em}2}}\right]`$
$`+C_F^{\mathrm{\hspace{0.17em}3}}n_f\left[{\displaystyle \frac{145}{9}}{\displaystyle \frac{5132}{3}}\zeta _2936\zeta _3+{\displaystyle \frac{1032}{5}}\zeta _2^{\mathrm{\hspace{0.17em}2}}\right]+C_Fn_f^{\mathrm{\hspace{0.17em}3}}\left[{\displaystyle \frac{940}{81}}{\displaystyle \frac{32}{9}}\zeta _2\right]`$
$`C_AC_Fn_f^{\mathrm{\hspace{0.17em}2}}\left[{\displaystyle \frac{7403}{27}}{\displaystyle \frac{688}{9}}\zeta _216\zeta _3\right]C_F^{\mathrm{\hspace{0.17em}2}}n_f^{\mathrm{\hspace{0.17em}2}}\left[{\displaystyle \frac{52678}{243}}{\displaystyle \frac{6104}{27}}\zeta _2{\displaystyle \frac{304}{9}}\zeta _3\right].`$
The seventh term with $`𝒟_1`$ is not exactly known, since the fourth-order contribution to $`A_\mathrm{q}`$ has not been computed so far. Inserting the numerical values for the $`\zeta `$-functions and the QCD colour factors, including $`d^{abc}d_{abc}/n_c=5n_f/18`$, the resummation prediction is given by
$$c_{\mathrm{\hspace{0.17em}2},\mathrm{q}}^{(4)}|_{𝒟_1}=286702+64219.0n_f2019.24n_f^{\mathrm{\hspace{0.17em}2}}+2.0166n_f^{\mathrm{\hspace{0.17em}3}}63.402fl_{11}n_f+A_{\mathrm{q},4}.$$
(5.10)
As mentioned above, numerically insignificant are both the uncertainty due to $`A_{\mathrm{q},4}`$ (estimated in Eq. (4.5)) and the singlet$`/`$non-singlet difference introduced by the $`fl_{11}`$ contribution of Eq. (4.8). It is also interesting to note that the fourth coefficient $`\beta _3`$ of the beta function with its quartic group invariants $`d_F^{abcd}`$ and $`d_A^{abcd}`$ only enters the eighth tower, starting at the fifth order in $`\alpha _\mathrm{s}`$.
The numerical values of the $`N`$-space coefficients $`c_{kl}`$ in Eq. (5.1) are presented in Table 1 for $`l7`$ and $`k10`$. Recall that also these coefficients refer to an expansion in $`a_\mathrm{s}=\alpha _\mathrm{s}/(4\pi )`$. Whatever the normalization of the expansion parameter, however, the coefficients in each column (tower) finally vanish for $`k\mathrm{}`$, as mentioned below Eq. (5.3). Thus the series (5.1) converges at all $`N0`$ for any finite number of towers $`l_{\mathrm{max}}`$, i.e., with the upper limit in the second sum replaced by $`l_{\mathrm{max}}`$. The Mellin inversion of the product with the parton distributions $`f^N`$ is therefore well-defined, in contrast to the fully exponentiated result discussed above.
Before we turn to the higher-order predictions, it is instructive to compare the approximations by the leading large-$`x`$ and large-$`N`$ terms to the completely known two- and three-loop coefficient functions . This is done in Fig. 4 for the successive approximations in terms of the +-distributions $`𝒟_k`$ defined above Eq. (5.4). The corresponding results for the expansion in powers of $`\mathrm{ln}N`$ are presented in Fig. 5. Obviously both expansions reproduce the exact large-$`x`$ behaviour (up to terms not increasing as $`x1`$ for the ratios shown in the figures) at order $`\alpha _\mathrm{s}^n`$ once all enhanced terms, $`𝒟_k`$ with $`k=0,\mathrm{},\mathrm{\hspace{0.17em}2}n1`$ or $`\mathrm{ln}^lN`$ with $`l=1,\mathrm{},\mathrm{\hspace{0.17em}2}n`$, have been taken into account. The $`x`$-space expansion, however, would lead to a gross overestimate if only the terms up to $`kn+1`$ were known. The convergence in the $`N`$-space approach, on the other hand, is much smoother, with a good approximation already reached after $`n`$ terms.
A similar pattern is found for the fourth-order coefficient function $`c_{\mathrm{\hspace{0.17em}2},\mathrm{q}}^{(4)}`$ illustrated in the same manner in the left part of Fig. 6: the expansion in decreasing powers of $`\mathrm{ln}N`$ stabilizes after the fourth term. Based on these results and the higher-order coefficients shown in Table 1, we expect that the first $`l`$ logarithms should provide a good estimate up to about the $`l`$-th order in $`\alpha _\mathrm{s}`$, but severely underestimate the effect of the coefficient functions of much higher orders. Consequently, the tower expansion should underestimate the corrections towards $`x1`$, where more and more orders become relevant. This is exactly the pattern shown in the right part of Fig. 6, where the predictions of all effects beyond order $`\alpha _\mathrm{s}^{\mathrm{\hspace{0.17em}3}}`$ are compared between the tower expansion and the full exponentiation (for the latter again using a ‘minimal-prescription’ contour ). Both approaches agree very well, for the chosen input parameters, at $`x<0.93`$, but start to diverge at $`x\stackrel{>}{}0.95`$ where the exponentiation is also intrinsically more stable.
## 6 Summary
We have extended the threshold resummation exponents for few-parton processes to the fourth logarithmic (N<sup>3</sup>LL) order collecting the terms $`\alpha _\mathrm{s}^{\mathrm{\hspace{0.17em}2}}(\alpha _\mathrm{s}\mathrm{ln}N)^k`$ to all orders in $`\alpha _\mathrm{s}`$. For our reference process, inclusive deep-inelastic scattering (DIS), the N<sup>3</sup>LL contributions are specified by two universal expansion parameters: the four-loop cusp anomalous dimension $`A_{\mathrm{q},4}`$ and the third-order quantity $`B_{\mathrm{q},3}`$ which defines the jet function resumming collinear radiation off an unobserved final-state quark. The former coefficient has not been computed so far, but can be safely expected to have a very small effect of less than 1%. In fact, the perturbative expansion up to $`A_3`$ does not exhibit enhanced higher-order corrections, and $`A_4`$ can be estimated by Padé approximations. We have calculated the more important second coefficient $`B_{\mathrm{q},3}`$ by comparison of the expanded resummation result to our recent third-order calculation of electromagnetic DIS .
The perturbative expansion of $`B_\mathrm{q}`$ seems to indicate, as far as this can be judged from the first three terms, the onset of a factorial enhancement of the higher-order coefficients. However, the rather large size of the coefficient $`B_{\mathrm{q},3}`$ actually stabilizes the logarithmic expansion of the coefficient functions. In fact, the N<sup>3</sup>LL corrections are very small at large scales $`Q^2`$, and even facilitate a reliable prediction of the soft-gluon effects at scales as low as $`Q^24\text{ GeV}^2`$ (corresponding to $`\alpha _\mathrm{s}\mathrm{\hspace{0.25em}0.3})`$ down to very small invariant masses $`W`$ of the hadronic final state in $`epeX`$, $`W^2m_p^2\mathrm{\hspace{0.25em}0.5}\text{ GeV}^2`$. Thus we expect our results to be useful also for low-scale data analyses using parton-hadron duality concepts.
The threshold resummation can also be employed to predict, order by order in $`\alpha _\mathrm{s}`$, the leading $`\mathrm{ln}N`$ contributions to the higher-order coefficient functions. At the level of accuracy reached in the present article for inclusive DIS, the exponentiation fixes the seven highest terms, $`\mathrm{ln}^nN`$ with $`2l6n2k`$, at all orders $`k4`$ of $`\alpha _\mathrm{s}`$. Already the highest $`k`$ powers of $`\mathrm{ln}N`$ provide a good estimate of the soft-gluon enhancement of the $`k`$-loop coefficient functions at least for $`k7`$, in contrast to the (expected, see Ref. ) worse behaviour of the corresponding expansion in $`x`$-space +-distributions. Except very close to threshold, where too many orders in $`\alpha _\mathrm{s}`$ become important, the summation of the above seven $`N`$-space logarithms to all orders yields a good agreement with the exponentiated coefficient function. This agreement further confirms the ‘minimal prescription’ used for defining the in principle ambiguous Mellin inversion of the resummation exponential.
Besides the standard (gauge-boson exchange) process, we have also considered DIS by exchange of a scalar directly coupling to gluons. By comparison of the resummation to our unpublished three-loop coefficient function for this process we have derived, for the first time, the second and third order contributions to the coefficient $`B_\mathrm{g}`$ governing the jet function of a final-state gluon. The quantity $`B_{\mathrm{g},2}`$ will be employed to extend the NNLL resummation to more processes, once the corresponding NNLO results required to fix the process-dependent large-angle soft contribution become available. Finally we would like to draw attention to the curious relation (4.19) which connects, for both the quark and gluon channels, the second-order splitting functions, jet functions and the two-parton (Drell-Yan) large-angle soft emissions in a non-trivial manner.
### Acknowlegments
We thank E. Laenen and W. Vogelsang for stimulating discussions. The work of S.M. has been supported in part by the Helmholtz Gemeinschaft under contract VH-NG-105 and by the Deutsche Forschungsgemeinschaft in Sonderforschungsbereich/Transregio 9. The work of J.V. has been part of the research program of the Dutch Foundation for Fundamental Research of Matter (FOM).
## Appendix
Here we show some key elements of the calculation of the resummation exponents $`g_i`$ presented in Section 3. Useful auxiliary relations (for $`|x|<1`$) are
$`{\displaystyle \frac{1}{(1x)^{nϵ}}}`$ $`=`$ $`{\displaystyle \underset{i=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{\mathrm{\Gamma }(nϵ+i)}{\mathrm{\Gamma }(nϵ)}}{\displaystyle \frac{x^i}{i!}},`$ (A.1)
$`{\displaystyle \frac{\mathrm{ln}^k(1x)}{(1x)^{nϵ}}}`$ $`=`$ $`\left({\displaystyle \frac{}{ϵ}}\right)^k{\displaystyle \frac{1}{(1x)^{nϵ}}}.`$ (A.2)
The Mellin transforms of the +-distributions follow from the results for harmonic polylogarithms and are given by
$$\underset{0}{\overset{1}{}}𝑑z\frac{z^N1}{1z}\mathrm{ln}^k(1z)=(1)^{k+1}k!S_{\underset{k+1}{\underset{}{1,\mathrm{},1}}}(N),$$
(A.3)
where $`S_{m_1,\mathrm{},m_k}(N)`$ denotes the harmonic sums . Eqs. (A.1) – (A.3) lead to the master formula for the derivation of the functions $`g_i`$,
$`{\displaystyle \underset{0}{\overset{1}{}}}𝑑z{\displaystyle \frac{z^N1}{1z}}{\displaystyle \frac{1}{(1+a\mathrm{ln}(1z))^{nϵ}}}=`$ (A.6)
$`{\displaystyle \underset{i=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{\mathrm{\Gamma }(nϵ+i)}{\mathrm{\Gamma }(nϵ)}}(a_s\beta _0)^iS_{\underset{i+1}{\underset{}{1,\mathrm{},1}}}(N){\displaystyle \underset{j=0}{\overset{\mathrm{}}{}}}\left(\begin{array}{c}i+j1\\ j\end{array}\right)\left(a_s\beta _0\mathrm{ln}{\displaystyle \frac{Q^2}{\mu _r^2}}\right)^j`$
with $`a=(a_s\beta _0)/(1+a_s\beta _0\mathrm{ln}Q^2/\mu _r^2)`$. The double sum in Eq. (A.6) can be solved to the desired logarithmic accuracy with the algorithms for the summation of nested sums coded in Form . The expansion of the Gamma function in powers of $`ϵ`$ for positive integers $`n`$ reads
$`{\displaystyle \frac{\mathrm{\Gamma }(n+1+ϵ)}{n!\mathrm{\Gamma }(1+ϵ)}}`$ $`=`$ $`1+ϵS_1(n)+ϵ^2(S_{1,1}(n)S_2(n))+ϵ^3(S_{1,1,1}(n)S_{1,2}(n)`$ (A.7)
$`S_{2,1}(n)+S_3(n))+ϵ^4(S_{1,1,1,1}(n)S_{1,1,2}(n)S_{1,2,1}(n)`$
$`+S_{1,3}(n)S_{2,1,1}(n)+S_{2,2}(n)+S_{3,1}(n)S_4(n))+𝒪(ϵ^5).`$
Finally, with $`\theta _{ij}=1`$ for $`ij`$ and $`\theta _{ij}=0`$ else, the sums $`S_{1,\mathrm{},1}(N)`$ are factorized according to
$`i!S_{\underset{i}{\underset{}{1,\mathrm{},1}}}(N)`$ $`=`$ $`(S_1(N))^i+{\displaystyle \frac{1}{2}}i(i1)S_2(N)(S_1(N))^{i2}+{\displaystyle \frac{1}{3}}i(i1)(i2)S_3(N)(S_1(N))^{i3}`$ (A.9)
$`+{\displaystyle \frac{1}{4}}i(i1)(i2)(i3)\left(S_4(N)+{\displaystyle \frac{1}{2}}(S_2(N))^2\right)(S_1(N))^{i4}+\mathrm{}`$
$``$ $`\theta _{i1}\mathrm{ln}^i\stackrel{~}{N}+{\displaystyle \frac{1}{2}}\theta _{i3}i(i1)\zeta _2\mathrm{ln}^{i2}\stackrel{~}{N}+{\displaystyle \frac{1}{3}}\theta _{i4}i(i1)(i2)\zeta _3\mathrm{ln}^{i3}\stackrel{~}{N}`$
$`+{\displaystyle \frac{1}{4}}\theta _{i5}i(i1)(i2)(i3)\left(\zeta _4+{\displaystyle \frac{1}{2}}\zeta _{2}^{}{}_{}{}^{2}\right)\mathrm{ln}^{i4}\stackrel{~}{N}+\mathrm{},`$
where $`\stackrel{~}{N}=Ne^{\gamma _e}`$ and the algebraic properties of harmonic sums have been used .
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# Coexistence of Sharp Quasiparticle Dispersions and Disorder Features in Graphite
(February 28, 2005)
## Abstract
Angle resolved photoelectron spectroscopy (ARPES) on azimuthally disordered graphite demonstrates that sharp quasiparticle dispersions along the radial direction can coexist with a complete lack of dispersion along the azimuthal direction. This paradoxical coexistence can be explained in terms of van Hove singularities in the angular density of states. In addition, non-dispersive features at the energies of band maxima and saddle points are observed and possible explanations are discussed. This work opens a new possibility of studying the electronic structure of novel layered materials using ARPES even when large single crystals are difficult to obtain.
The ability to sharply resolve crystal momentum values of single particle excitations has made ARPES a very powerful tool in addressing the electronic structure of solid, as has been successfully demonstrated on single crystalline samples over the past decades Hufner . Due to the translational symmetry along the surface of a single crystal, the crystal momentum parallel to the surface (k) is conserved during the photoemission process, allowing a complete momentum space map of the initial state. This holds despite the short photoelectron lifetime Hedin ; Smith ; Lindroos which can severely broaden the resolution of the momentum perpendicular to the surface (k<sub>z</sub>). Indeed, even in the limit of an extreme k<sub>z</sub> broadening that results in no resolution of k<sub>z</sub>, strong ARPES dispersions are expected as a function of k, since the one dimensional density of states (1D-DOS) D<sub>z</sub>(E) $``$ dk<sub>z</sub>/dE obtained by integrating over k<sub>z</sub> is dominated by contributions from van Hove singularities in high symmetry planes Lindroos .
On the contrary, for those systems characterized by orientationally disordered domains, i.e. polycrystalline materials, the translational symmetry is preserved only within each domain. As a consequence, the dispersion measured by ARPES is the average dispersion over different domains, or equivalently azimuthal angle $`\varphi `$, which in general leads to no dispersion. However, extending the 1D-DOS D<sub>z</sub>(E) scenario for k<sub>z</sub> discussed above further to the plane, there is an interesting possibility that a layered polycrystalline sample, with a strong azimuthal disorder, can nevertheless give distinct dispersions in the radial direction. This would happen if the average dispersion is dominated by those along the high symmetry directions due to van Hove singularities in the angular density of states D<sub>ϕ</sub>(E) $``$ d$`\varphi /`$dE. To date, this possibility has never been demonstrated experimentally and photoemission studies on disordered samples have focused on angle-integrated features without any momentum information.
In this paper, we report a high resolution ARPES study on the electronic structure of azimuthally disordered graphite. For the first time, we report clear evidence that sharp quasiparticle dispersions, in agreement with band structure calculation along the high symmetry directions, can coexist with a circular Fermi energy intensity map, a definitive signature of azimuthal disorder Santoni . In addition we report non-dispersive features at the energies of band maxima or saddle points, which are attributed to the loss of momentum information by indirect photoemission process or elastic scattering. A practical implication of this study is that more ARPES opportunities can be made available for layered materials even when high quality single crystals of large size are difficult to obtain.
ARPES data were collected at beam line 10.0.1 of the Advanced Light Source (ALS) at the Lawrence Berkeley National Laboratory, using an SES-R4000 analyzer. The wide angular mode with acceptance angle of 30 and angular resolution of 0.9 was utilized for most scans, while high resolution angular mode with acceptance angle of 14 and angular resolution of 0.1 was utilized for one scan. The total instrumental energy resolution was 15 meV at 25 eV photon energy and 25 meV for other photon energies used (40, 55, 60 eV). The sample used was a grade ZYA highly oriented pyrolytic graphite (HOPG), obtained commercially from Structure Probe Inc. The sample was cleaved in situ in an ultra high vacuum better than 1.0$`\times `$10<sup>-10</sup> Torr and measured at temperature 50 K.
Figure 1(a) shows an ARPES intensity map measured at the Fermi energy taken at 40 eV photon energy. Throughout this paper, we use a color scale such that black represents high intensity in the raw data and blue represents low intensity. According to band structure calculation, a constant k<sub>z</sub> cross section of graphite Fermi surface can be a small hole pocket, a small electron pocket, or a point located at the six corners of the hexagonal Brillouin zone (dashed line in figure 1(a)), depending on the value of k<sub>z</sub> BandStructure ; el-h . Experimentally, the predicted small electron or hole pockets have yet to be resolved, and measurements on single crystalline samples have shown only small dots of high intensity at these corners Santoni , schematically drawn as shaded circles in Figure 1(a). For the graphite sample under study, the Fermi energy intensity map, symmetrized by three fold rotations to fill the entire Brillouin zone, shows a perfectly circular pattern <sup>1</sup><sup>1</sup>1The intensity variation along the circle is attributed to photoemission matrix element and is not a major concern here., in contrast to what is expected for single crystalline graphite. This is attributed to the angular spread of the dots to a circle due to the azimuthal disorder of the sample Santoni .
Figure 1(b) shows an ARPES intensity map as a function of binding energy and in-plane momentum k, corresponding to the momentum cut shown as a solid line in Figure 1(a). Despite the strong azimuthal disorder giving a circular Fermi energy intensity map in Figure 1(a), we observe, surprisingly, very clear dispersions over the entire energy range. Furthermore, at the Fermi energy crossing point k<sub>F</sub>, a sharp coherent quasi-particle peak is observed. This is shown in the inset, where an energy distribution curve (EDC), energy cut at a constant momentum, is plotted. Here the half width of the EDC peak is 20 meV (50 meV FWHM due to the asymmetry of the line shape), defining the sharpest peak observed in graphite so far Kihlgren .
In Figure 1(c) we report the second derivative of the raw data of Figure 1(b) with respect to energy. The second derivative method has been used in the literature to enhance the direct view of the ARPES dispersion. Local density approximation (LDA) band dispersions along two high symmetry directions $`\mathrm{\Gamma }`$-K-M$`^{^{}}`$ (solid lines) and $`\mathrm{\Gamma }`$-M-$`\mathrm{\Gamma }^{^{}}`$ (dashed lines) are plotted in the same figure for a direct comparison. Despite a polycrystal-like sample implied by the Fermi energy intensity map, an excellent agreement is observed between the experiment and the theory. We can identify the dispersions between 4 eV and 23 eV as originating from the $`sp^2`$ orbitals with strong intra-layer $`\sigma `$ bonding (black lines), and the dispersions between Fermi energy and 11 eV as originating from the $`p_z`$ orbitals with weaker $`\pi `$ bonding (white lines). We note that the calculated dispersions were stretched by 20$`\%`$ in energy throughout this paper, as suggested in the literature Bandwidening ; Kihlgren ; Louie . The stretching of the LDA band dispersions is attributed to missing self-energy corrections in LDA, since ab initio quasiparticle calculations based on the GW method show that for graphite the quasiparticle band dispersion near the Fermi level is 15% largerLouie .
The direct comparison between Figure 1(a) and Figure 1(b,c) shows an apparent paradox in our data, namely, the coexistence of azimuthal disorder feature (Figure 1(a)) with single crystalline features (Figure 1(b,c)). This can be readily understood if we consider an angular average of the calculated dispersions. Such an angular average would be necessary if the sample consisted of many small single crystallites with strong azimuthal disorder.
Figures 2(a-c) show ARPES cuts for three azimuthal angles, $`\varphi `$=0, 10, 20. A direct comparison between panels a, b and c shows no appreciable angular dependence of the dispersions, establishing that the azimuthally invariant electronic structure of Figure 1(a) at the Fermi energy extends to the entire band width. Thus, these data strongly support the azimuthal disorder model described in the previous paragraph, and indicate that the ARPES data measured are actually a 1D-DOS D<sub>ϕ</sub>(E) along the azimuthal direction $`\varphi `$, in analogy with the well-known 1D-DOS D<sub>z</sub>(E) along the k<sub>z</sub> direction Lindroos . As in the latter case, then, one would expect that van Hove singularities arising from states along the high symmetry directions to contribute dominantly, and this gives an explanation why the measured dispersions accurately reflect the dispersions along the two high symmetry directions. The 1D-like van Hove singularities arise from states along high symmetry lines because these states have zero group velocity along the arc with constant k magnitude.
In Figure 2(d, e) we show LDA calculations supporting this reasoning. In panel d, we show the dispersion for single crystalline graphite along an arc from A to B with radius equal to $`\mathrm{\Gamma }`$K distance. Within this arc, the $`\mathrm{\Gamma }`$K direction corresponds to point A and the $`\mathrm{\Gamma }`$M direction to point B respectively. As expected, extrema in the band dispersion occur at the two high symmetry directions, A (shaded circles) and B (open circles). The calculated 1D-DOS D<sub>ϕ</sub>(E) over this arc, and thus over the entire azimuthal angle range by symmetry, is shown in panel e. One can see diverging 1D van Hove singularities as sharp peaks occurring at energies where bands cross points A and B, which completely dominate over other contributions. This nicely explains why well-defined sharp peaks with large dispersions can be observed in this azimuthally disordered sample, despite the fact that the observed data come from averaging over all azimuthal directions.
The data presented so far can be summarized as showing well-defined dispersions along the radial direction with a complete lack of dispersion along the azimuthal direction. Therefore, our data suggest that the graphite sample under study consists of finite size single crystalline grains much smaller than the analysis area ($``$ 100 $`\mu m`$) with a complete azimuthal disorder. However, each grain is large enough to allow for highly dispersive quasiparticles to exist. In addition, we have measured the dispersion perpendicular to the surface, k<sub>z</sub> dispersion, using photon energy range from 34 to 155 eV at beam line 12.0.1 of the ALS, with a perfect agreement with previous results Law . This indicates that the crystalline order remains coherent along this direction, i.e. perpendicular to graphene layers, over a length scale larger than the probing depth of ARPES (order of 10 Å).
We now discuss additional features at 2.9 (e<sub>1</sub>), 4.3 (e<sub>2</sub>) and 7.8 eV (e<sub>3</sub>) characterized by no dispersion at all. These features can already be observed in Figure 1(b,c) taken at 60 eV photon energy, appearing as sharp extended horizontal lines at the same energies. Figure 3 shows a detailed view of these same features observed at another photon energy of 55 eV. In panel a, we show the first derivative of the ARPES intensity map with respect to energy. The first derivative enhances rising or falling edges in the data, and thus is particularly useful for detecting non-dispersive peaks and edges. The energies for these non-dispersive features are marked by arrows on the right of this panel. In panel b, we show EDCs at $`\mathrm{\Gamma }`$ point (thin line) and k$`=0.4`$ Å<sup>-1</sup> (thick line). In these panels, in addition to highly dispersive features that we already discussed, one can clearly see the non-dispersive features, i.e. peaks at energies e<sub>1</sub> and e<sub>3</sub> and edge at energy e<sub>2</sub>. The non-dispersive nature of these features can be checked more explicitly using EDCs. For example, in panel c, a stack of EDCs over an extended momentum range clearly show the non-dispersive nature of peak at e<sub>1</sub> and edge at e<sub>2</sub>. Such an analysis shows that the features at e<sub>1</sub>, e<sub>2</sub> and e<sub>3</sub> are non-dispersive within $``$ 200 meV.
Similar non-dispersive features have been noted and attributed to surface states or emission from isolated atoms McGovern . However, it is important to point out that in our data there is a strong connection between these non-dispersive features and the highly dispersive features discussed above. Indeed, one can easily see from Figure 1(c) and Figure 3(a) that the energies of these non-dispersive features correspond to the band energy extrema. Namely, e<sub>1</sub> corresponds to the saddle point of $`\pi `$ band, e<sub>2</sub> the degenerate maxima of $`\sigma `$<sub>1</sub> and $`\sigma `$<sub>2</sub> bands, and e<sub>3</sub> the saddle point of the $`\sigma `$<sub>1</sub> band. This strongly suggests that they are related to the bulk band structure rather than surface state. We suggest that these non-dispersive features can be explained by loss of electron momentum information in two possible ways, both of which are associated with the large density of states at the band energy extrema Bianconi . One possibility is indirect transition, i.e. non-k-conserving transition in photoemission process Law ; Marchand , for example, phonon assisted transition. Another possibility is elastic scattering of electrons in either the initial state or the final state by inhomogeneity or disorder. In both these scenarios, the observed non-dispersive features correspond to density of states of graphite. This explanation is consistent with the EDC line shape observed in panels (b,c). That is, the non-dispersive feature at e<sub>2</sub> is sharp at low binding energy side because e<sub>2</sub> is the energy maximum of $`\sigma `$<sub>1</sub> and $`\sigma `$<sub>2</sub> bands, whereas the non-dispersive features at e<sub>1</sub> and e<sub>3</sub> are broad on both sides, because e<sub>1</sub> and e<sub>3</sub> are saddle point energies. We note that similar non-dispersive features are also observed in single crystalline graphite Law ; Strocov , which suggests that point defects definitely play an important role, in addition to the possible role by the azimuthal disorder we noted above. Interestingly, similar non-dispersive features are also observed in NaMo<sub>6</sub>O<sub>17</sub> Gweon and high T<sub>c</sub> superconductor Kaminski . Future studies of these features may be interesting, in view of the general lack of understanding of the role of inhomogeneity or disorder in ARPES, and the importance of this topic in complex materials.
In conclusion, the data presented here provide a novel example of the coexistence of sharp quasiparticle dispersions with disorder features. Also, several non-dispersive features are identified at band energy extrema. Our study brings up an interesting possibility that layered crystals with azimuthal disorder can be studied with ARPES to obtain useful information.
###### Acknowledgements.
We thank C.M. Jozwiak, A. Bill, K. McElroy, and D.R. Garcia for useful discussions and A.V. Fedorov for experimental assistance at beam line 12.0.1. This work was supported by the Director, Office of Science, Office of Basic Energy Sciences, Division of Materials Sciences and Engineering of the U.S Department of Energy under Contract No. DEAC03-76SF00098; by the National Science Foundation through Grant No. DMR03-49361 and Grant No. DMR04-39768. Support by the Sloan Foundation, the Hellman Foundation and computer resource at the NSF NPACI Center is also acknowledged.
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# Thermalization of Sputtered Particles as the Source of Diffuse Radiation from High Altitude Meteors
## 1 Introduction
A meteor is a luminous phenomenon produced by a meteoroid colliding with the earth’s atmosphere. The collision starts with atmospheric particles impinging on the meteoroid’s body, by which the surface heats up and sublimates. The cloud of vapor particles surrounding the meteoroid collides further with the atmosphere and becomes ionized and excited. The meteor light is produced by deexcitations, with metals from the meteoroid dominating the spectrum. When the meteor reaches altitudes where the Knudsen number $`Kn`$ (ratio between the atmospheric mean free path and the meteoroid size) is less than $``$100, vapors completely engulf the body and shield it form direct collisions with the atmosphere. After that the heating and evaporation of the meteoroid is caused by hot vapor and radiation from the shock front. During the whole process the meteoroid body is ablating, that is losing mass either as gas or as fragments. The emission of light exists only as a consequence of the ablation process.
The limiting altitude at which most meteors become luminous is $``$120 km or less, with very few starting at $``$130 km (e.g. Betlem et al., 2000; Campbell et al., 2000; Brown et al., 2002). The atmospheric mean free path is on the order of meters at these altitudes, which is enough to induce intense evaporation of volatiles from large (centimeters in size) fast meteoroids. It came as surprise, therefore, when meteors of much higher beginning heights were detected in 1995 and 1996 by Fujiwara et al. (1998) using TV cameras. They reconstructed trajectories of two Leonid fireballs of -4<sup>m</sup> and -7<sup>m</sup> visual magnitude and found that they started at $``$160 km altitude. This is about 30 km higher than the highest meteors detected by photographic systems (Betlem et al., 2000). The detection was attributed to the much higher sensitivity of TV cameras, especially the infrared region of 900 nm.
This result was confirmed and extended further in 1998 by Spurný et al. (2000a), using high-sensitivity TV systems with limiting magnitude +6.5<sup>m</sup>. They detected thirteen Leonid fireballs at altitudes above 130km. Three of them were above 180km, with the highest at 200 km. The atmospheric mean free path is over 200 m at an altitude of 200 km, and more than 50 m at 160 km, while the meteoroid body has a size of no more than tens of centimeters for the brightest detected meteors. It is puzzling how meteors can produce light in such conditions. The meteoroid surface is still too cold for sublimation or evaporation and, thus, these high altitude meteors can not be explained by the standard theory of meteor ablation.
Spurný et al. (2000b) described the morphological features of these meteors. Initially they appear as a several kilometers big diffuse spot that rapidly transforms into a diffuse comet-like V-shape structure with a head and tail. Sometimes structures like arcs, jets and streamers are visible in the tail. At altitudes of $``$130 km the meteor changes its appearance into an ordinary droplet-like head shape with trail. This also marks the moment when the standard ablation theory becomes applicable. High altitude meteors up to 150 km were also detected in other meteor streams: $`\eta `$-Aquarids, Lyrids, and Perseids (Koten et al., 2001).
The decreasing size of the diffuse structure with altitude suggests that the atmospheric mean free path is important as the light is induced by collisions. Spurný et al. (2000a) noted that the beginning meteor height grows with the peak brightness of meteors. This implies that the number of collisions increases with the meteoroid size. Considering these properties of high altitude meteors, we suggest that they are produced by a cascade of collisions in the atmosphere during thermalization of particles sputtered from the meteoroid surface. Sputtering is a process in which an energetic particle projectile ejects atoms from the surface of a solid material. It has been used in meteor physics only to study heating of the surface or slowing down of micrometeoroids (Opik, 1958; Lebedinets et al., 1973; Bronshten, 1983; Coulson, 2002). The luminosity has been associated exclusively with the phase of intense evaporation<sup>1</sup><sup>1</sup>1During the preparation of this manuscript, we became aware of a study by Hill, Rogers & Hawkes (accepted for publication in Earth. Moon, and Planets) describing sputtering as a possible source of luminosity in high altitude meteors., during which the number of released atoms is increased dramatically (Bronshten, 1983). Unlike vaporization, sputtering can exist at high altitudes because it does not depend much on the meteoroid surface temperature.
It is the purpose of this work to reproduce basic properties of high altitude meteors. We present our 3D Monte Carlo model that traces multiple particle collisions within the atmosphere and show modeling results in the form of theoretical images from a ground-based CCD detector (TV camera).
## 2 Sputtering from the meteoroid surface
Meteoroids in meteor showers enter the earth’s atmosphere with geocentric velocity $``$15 km/s to 71 km/s. From the point of view of the meteoroid surface the atmospheric molecules behave like projectiles with kinetic energies of hundreds of eV. For example, an oxygen molecule at 71 km/s has 840 eV of kinetic energy. Collision with the surface dumps this energy into displacements of atoms in the solid until the projectile stops. If the energy is high enough, some atoms are ejected from the surface, and we call these sputtered particles. If the projectile exits the solid then we call it a scattered particle.
The sputtering yield $`Y`$ is defined as the number of sputtered particles relative to the number of projectiles. The yield depends on the projectile energy, physical and chemical properties of the solid target, and on the type of the sputtered atom. The experimental data on sputtering at the energies of interest here are scarce, thus we mostly rely on numerical models. Since the high altitude meteors detected so far belong to cometary material, we can assume that they are made of a fragile (low binding energy) amorphous material. Calculations show that sputtering yields on such targets approach $`Y`$=1 for projectile energies of hundreds of eV and that the sputtered atoms have typically 5-10% of the impact energy (Field et al., 1997). In the case of silicates typically present in interstellar grains, the yield can be one or two orders of magnitude smaller (May et al., 2000). Since in this work we do not investigate details of the meteor spectrum, we do not differentiate between different species of sputtered particles. We focus only on the number of particles originating from the surface, and therefore describe both sputtered and scattered particles with $`Y`$.
The velocity of sputtered and scattered particles has one component equal to the meteoroid velocity. An additional (smaller) component is present because a particle picks up some additional energy from the impact. The total velocity is a vector sum of those two components. This makes it about 100 times faster than the average thermal speed in the atmosphere. Such a fast particle will create a cascade of collisions before slowing down to the thermal level. During this process of thermalization, collisions will excite atoms and dissociate molecules, which results in production of light by deexcitations.
We make an estimate of whether this process creates enough light to be detected by the TV cameras used in the observations of high altitude Leonid meteors. During the time exposure $`\mathrm{\Delta }t`$ of one video frame, a meteor of velocity $`v_\mathrm{m}`$ travels a distance $`v_\mathrm{m}\mathrm{\Delta }t`$. The number of collisions with atmospheric particles is then $`R_\mathrm{m}^2\pi v_\mathrm{m}\mathrm{\Delta }tn_{\mathrm{atm}}`$, where $`R_\mathrm{m}`$ is the meteoroid radius and $`n_{\mathrm{atm}}`$ is the atmospheric number density. Each particle originating from the meteoroid surface will result in many photons from excited (mostly atmospheric) atoms in the cascade. The total number of emitted photons is
$$N_{\mathrm{photon}}^{tot}=\gamma YR_\mathrm{m}^2\pi v_\mathrm{m}\mathrm{\Delta }tn_{\mathrm{atm}},$$
(1)
where $`\gamma `$ is the total number of emitted photons in one cascade.
These photons are emitted within the overall volume of $`4\pi (L/2)^3/3`$, where high altitude meteors have size, $`L`$, of several kilometers. On the other hand, one pixel in the CCD of a TV camera covers only a fraction of the meteor angular size on the sky. Spurný et al. (2000a) used an imaging setup with angular pixel size $``$2.5”$`\times `$2.5”. This covers an area, $`A`$, of 150$`\times `$150 m<sup>2</sup> at 200 km distance. One pixel collects photons from the volume $`AL`$. If photons are emited uniformly within the volume of the meteor, and asumming a meteor size of $`L`$=10 km, the fraction of photons contributing to the flux measured in one pixel is $`\mathrm{\Psi }=24A/4\pi L^2`$0.0004. However, photons are not emitted uniformly, but with a radial gradient. The central region, which we recognize as the meteor head, can emit $``$100 times more photons than the average, as we will see in §4 from numerical models. This puts the estimate on $`\mathrm{\Psi }`$ closer to 4%.
If the average photon energy is $`E_0`$ then the measured flux in one pixel is
$$F_\mathrm{m}=\frac{\mathrm{\Psi }E_0N_{\mathrm{photon}}^{tot}}{4\pi D^2\mathrm{\Delta }t},$$
(2)
where $`D`$ is distance to the meteor. To work in units of stellar magnitudes, we compare this flux with the flux of Sirius, $`F_{\mathrm{sirius}}`$. Although the exact stellar magnitude depends on the spectral sensitivity of the detector, we assume here a visual magnitude of $``$1<sup>m</sup> for Sirius. The apparent meteor magnitude is then
$$𝔐=1^m2.5\mathrm{log}\left(\frac{\gamma YR_\mathrm{m}^2v_\mathrm{m}n_{\mathrm{atm}}\mathrm{\Psi }E_0}{4D^2F_{\mathrm{sirius}}}\right).$$
(3)
Spurný et al. (2000a) and Fujiwara et al. (1998) argue that the sensitivity of TV cameras to wavelengths of $``$800 nm is one of the reasons why they can detect high altitude meteors. Photons of that wavelength have $`E_0`$=2.5$`\times `$10<sup>-19</sup> J. One collisional cascade can produce $`\gamma `$100 photons (see §4), as a product of excitation by collisions or dissociation of molecules. The highest Leonid detected by Spurný et al. (2000a) had $`R_\mathrm{m}`$=0.07 m and was first detected at distance $`D`$200 km. By using $`v_\mathrm{m}`$=71 km/s for Leonids, $`n_{\mathrm{atm}}`$10<sup>16</sup> m<sup>-3</sup> and $`F_{\mathrm{sirius}}`$2$`\times `$10<sup>-8</sup> W/m<sup>2</sup> (integrated flux within the sensitivity range of the camera), our estimate of the brightness is $`𝔐`$6.5<sup>m</sup>. This is exactly the limiting magnitude of the video system used by Spurný et al. (2000a).
For a given observational setup, the limiting magnitude is, by definition, the same as the meteor brightness at the beginning height, $`H_{begin}`$. Equation 3 shows that $`H_{begin}`$ depends on the meteoroid size. Larger meteoroids reach higher maximum brightness $`𝔐^{max}`$ during their disintegration below $``$90 km, thus a correlation between these two parameters is to be expected. Indeed, observations show that $`H_{begin}`$ is higher for brighter meteors (figure 1). In addition to the meteoroid size $`R_\mathrm{m}`$ and distance to the beginning point $`D`$, this dependence comes from variables in equation 3 that differ in two meteors from the same meteor stream. The size of a meteor scales with the atmospheric mean free path $`l_{\mathrm{atm}}`$, thus $`\mathrm{\Psi }ł_{\mathrm{atm}}^2`$. For altitudes between 160 km and 200 km, the atmosphere has $`ł_{\mathrm{atm}}\mathrm{exp}(H_{begin}/27.3\mathrm{km})`$ and $`n_{\mathrm{atm}}\mathrm{exp}(H_{begin}/26.8\mathrm{km})`$ (the U.S. Standard Atmosphere 1976).
In the single body approximation of ablation, the maximum meteor brightness $`𝔐^{max}`$ in magnitudes is (see Appendix)
$$𝔐^{max}=7.5\mathrm{log}R_\mathrm{m}2.5\mathrm{log}(\mathrm{cos}Z)+\mathrm{constant},$$
(4)
where $`Z`$ is the zenith angle of the meteoroid’s trajectory. Combining this with equation 3 yields the dependence of $`H_{begin}`$ on $`𝔐^{max}`$:
$$H_{begin}=5.55\mathrm{km}\left[𝔐^{max}\right]+\mathrm{constant},$$
(5)
where $`\left[𝔐^{max}\right]`$ is the reduced maximum brightness:
$$\left[𝔐^{max}\right]=𝔐^{max}+7.5\mathrm{log}\left(\frac{D}{200\mathrm{k}\mathrm{m}}\right)+2.5\mathrm{log}(\mathrm{cos}Z)$$
(6)
scaled to the distance of 200 km.
Observed values of $`H_{begin}`$ as a function of $`\left[𝔐^{max}\right]`$ are plotted in figure 2. Linear regression gives a slope of -5.9$`\pm `$0.5 km, which agrees with the theoretical slope in equation 5.
## 3 Monte Carlo model
Unusual morphology is another characteristic of high altitude meteors. In order to explain changes in their appearance with altitude, we preformed Direct Simulation Monte Carlo (DSMC) modeling. The simulation follows a sputtered particle and all high velocity particles produced by collisions during the thermalization process. The three-dimensional distribution of collisions provides information on the spatial distribution of photons produced.
Although DSMC is widely used for modeling interactions between rarefied gases and solids, including reentry of spacecraft (Oran et al., 1998), it has very rarely been used in meteor physics. Boyd (2000) modeled the flow around a 1 cm Leonid meteor at 95 km altitude. The model shows that meteoroid vapors surround the body and interact with the atmosphere on a larger scale than when the vapors are not produced. This confirms theoretical expectations of vapor shielding. Zinn et al. (2004) also performed, but not with Monte Carlo, a numerical study of Leonids at these and lower altitudes with a hydrodynamics code coupled with radiative transfer. They found that the meteoroid vapors are stopped very quickly by the air and moved into the meteor wake, where the vapor-air mixture expands and cools. This stopping process is the dominant source of energy deposition in the atmosphere. These models demonstrate the importance of correct numerical treatment of the vapor cloud.
All these numerical models have been applied to conditions in which Kn$`<`$10, typically for altitudes below $``$100 km. This is where the intense vaporization and compressed air in front of the meteoroid create conditions appropriate for the approximation of vapor shielding and the continuous flow regime. This choice of flow conditions is also influenced by the numerical codes used in simulations. They are not designed specifically for meteors and their applicability requires the hydrodynamic fluid regime. In contrast, our DSMC code is written and designed specifically for meteors under the molecular flow regime at higher altitudes.
In this study we focus solely on the sputtering effect at altitudes above $``$130 km as relevant to high altitude meteors. For simplicity, we do not differentiate between different atomic and molecular species. Other physical processes (like vaporization, ionization, dissociation, excitation, chemistry, etc.) can be added in the future to expand the applicability of the code. The Monte Carlo methodology used in the code is described by Xie and Mumma (1996) for the case of cometary atmospheres. Our modifications are related to the physical conditions in the earth’s atmosphere and the sputtering process. Here we describe the basic numerical procedure, while for the detailed description we refer to the original paper.
The numerical simulation follows the flight of a meteoroid through the atmosphere and calculates how many particle collisions with the air happen during one time step $`\delta t`$ of calculation. Each modeled particle in DSMC is a statistical representative of a set of particles. For $`N_{\mathrm{MC}}`$ model particles within time step $`\delta t`$ at altitude $`H`$, the number of real particles in one set is
$$N_{\mathrm{set}}(H)=\frac{R_\mathrm{m}^2\pi v_\mathrm{m}n_{\mathrm{atm}}(H)\delta t(H)}{N_{\mathrm{MC}}}.$$
(7)
If $`\delta t`$ is held constant during the simulation then the distance covered by the meteor during one time step would eventually become much larger than the atmospheric mean free path $`ł_{\mathrm{atm}}`$. In order to avoid this statistical problem, we move the meteoroid one $`ł_{\mathrm{atm}}`$ per step, thus $`\delta t(H)=ł_{\mathrm{atm}}(H)/v_\mathrm{m}`$. The number of Monte Carlo particles used in our computations is $`N_{\mathrm{MC}}`$=1500.
A model particle is ejected in a random direction from a random location on the meteoroid surface. We use a spherical meteoroid of $`R_\mathrm{m}`$=0.05 $`m`$ and isotropic sputtering for simplicity. A more general shape and sputtering angle distribution could be used in the future. We also assume a constant sputtering velocity $`V_{\mathrm{sp}}`$ of 20 km/s relative to the meteoroid. This corresponds to $``$10% of the collision energy if the masses of projectiles and sputtered particles are equal. We plan to introduce a more realistic distribution of sputtered velocities, such as a Maxwellian (Coulson, 2002) or Sigmund-Thompson (Thompson, 1968; Sigmund, 1969).
The history of collisions of an ejected particle and all subsequent collisionally produced fast air particles is recursively traced by the code. The basic assumption is that fast particles collide only with the background atmosphere, but never with each other. This is a good approximation as long as the local number density of ejected particles around the meteoroid is not comparable to the atmospheric density. This condition is satisfied in our case of sputtered particles.
Since all modeled particles are treated as the same species, there is only one collisional cross section $`\sigma `$. We use the Variable Hard Sphere model where the cross section is velocity dependent:
$$\sigma =\sigma _{\mathrm{ref}}\left(\frac{v_{\mathrm{ref}}}{v}\right)^s.$$
(8)
We use the reference values typical for air: $`s`$=0.25, $`\sigma _{\mathrm{ref}}`$=1.26$`\times `$10$`{}_{}{}^{19}m_{}^{2}`$ and $`v_{\mathrm{ref}}`$=(4989/M)$`{}_{}{}^{0.5}m/s`$ (Boyd, 2000), where M is the molar mass ($`kg/mol`$) of air at the altitude of collision. The time between collisions for a particle of velocity $`v`$ is calculated from a random number $``$ between zero and one:
$$t_{\mathrm{coll}}=\underset{H}{}\mathrm{\Delta }t(H).$$
(9)
$$\mathrm{ln}()=\underset{H}{}\mathrm{\Delta }t(H)\sigma vn_{\mathrm{atm}}(H)$$
This was derived from the assumption that air particles, whose velocities are much smaller than the velocity of modeled particles, follow the Maxwellian distribution. The sum over $`n_{\mathrm{atm}}(H)`$ is introduced because the atmospheric density is changing with altitude along the particle trajectory. To describe this change we use altitude steps of 1 km. The time that the particle spends crossing an altitude layer is $`\mathrm{\Delta }t(H)`$. Since neither $`t_{\mathrm{coll}}`$ nor $`H`$ is known in advance, the equation is solved iteratively.
A collision between particles requires random selection of the velocity vector $`v_t`$ of an air particle in the Maxwellian distribution. We choose cylindrical coordinates ($`\phi _t`$,$`v_{\rho t}`$,$`v_{zt}`$) where
$$\stackrel{}{v_t}=v_{\rho t}\mathrm{cos}\phi _t\widehat{i}+v_{\rho t}\mathrm{sin}\phi _t\widehat{j}+v_{zt}\widehat{k}.$$
(10)
From three random numbers $`_1`$, $`_2`$ and $`_3`$, the velocity is given as:
$$\phi _t=2\pi _1$$
(11)
$$v_{\rho t}=\overline{v}\sqrt{\mathrm{ln}_2}$$
(12)
$$\mathrm{Erf}\left(\frac{|v_{zt}|}{\overline{v}}\right)=2|_30.5|$$
(13)
$$\mathrm{sign}(v_{zt})=\mathrm{sign}(_30.5),$$
where $`\mathrm{Erf}`$ is the Error integral, $`\overline{v}^2`$=$`2RT(H)/`$M$`(H)`$, and $`T(H)`$ and M$`(H)`$ are the temperature and molar mass of the air at altitude $`H`$. The collisional scattering of two particles is isotropic, with velocity directions chosen randomly in the center-of-mass frame. When the kinetic energy of a particle drops below some pre-defined limit (10 eV in our simulations), its trajectory and collisions are not traced any more. This limit cannot be below the energy of photons that we expect to be produced. Another reason for losing a particle is collision with the meteoroid. We do not consider sputtering from such secondary collisions.
## 4 Model results
In order to visualize a meteor, we collect the total number of collisions visible from the ground by a TV camera at location $`\stackrel{}{r}_{\mathrm{cam}}`$. The coordinate system is shown in figure 3. The meteor’s direction of flight is $`\widehat{v}_\mathrm{m}=\mathrm{sin}Z\widehat{j}\mathrm{cos}Z\widehat{k}`$, where $`Z`$ is the zenith angle (45 in our model). The camera’s frames point at subsequent locations along the meteor path in order to track the meteor flight. When a collision happens at $`\stackrel{}{r}_{\mathrm{coll}}`$ in space and $`t_{\mathrm{coll}}`$ in time (starting from the first frame), the camera’s frame number $`f_{\mathrm{cam}}`$ is calculated (by ignoring decimal places) from
$$f_{\mathrm{cam}}=\frac{t_{\mathrm{coll}}+|\stackrel{}{r}_{\mathrm{coll}}\stackrel{}{r}_{\mathrm{cam}}|/c}{\mathrm{\Delta }t}$$
(14)
where $`c`$ is the speed of light. The number of collisions in one set $`N_{\mathrm{set}}`$ (equation 7) is added to the image pixel that points at the collision.
The first appearance of the meteor depends on the camera’s sensitivity. The pixel values can be transformed into apparent magnitudes if we define the frame in which the meteor first appears. This is the frame showing the meteor at its beginning height $`H_{\mathrm{begin}}`$. The brightest pixel in that frame is set to be the limiting magnitude $`𝔐_{\mathrm{limit}}`$ and its value, $`N_{\mathrm{limit}}`$, is used for calculating the apparent magnitude. If the distance of the meteor from the camera in this first frame is $`D_{\mathrm{limit}}`$ then the brightness of a pixel is
$$𝔐=𝔐_{\mathrm{limit}}2.5\mathrm{log}\left[\left(\frac{D_{\mathrm{limit}}}{D}\right)^2\frac{N_{\mathrm{collisions}}}{N_{\mathrm{limit}}}\right],$$
(15)
where $`N_{\mathrm{collisions}}`$ is the pixel value and $`D`$ is the meteor’s distance.
Figure 4 shows images of our meteor model from four different camera locations. The basic morphological features described by Spurný et al. (2000b) are successfully reproduced. The meteor first appears as a diffuse structure that grows in size, but then turns into a comet-like V-shape, followed by a transition into a more compact, elongated structure. The exact shape depends on the relative position between the camera and the meteor path. Since we do not consider photon emission from long-lived atomic states, the feature missing at lower altitudes, below $``$140 km, is the train that develops behind the meteor. Figure 4 also shows the number of collisions “visible” by each pixel. Collisions are concentrated at the center of the diffuse structure, which demonstrates the plausibility of the value used for $`\mathrm{\Psi }`$ in equation 3.
The efficiency of light production in high altitude meteors (described as $`\gamma `$ in equation 3) depends on the number of collisions. Collisional cascades differ in their total number of collisions. This number is mostly between 100 and 200, with $``$180 collisions being the most frequent (figure 5). The total energy of collision between two particles follows the distribution $`dN=Nf(E)dE`$ that can be calculated from
$$f(E)=\frac{\mathrm{\Delta }N}{\mathrm{\Delta }EN},$$
(16)
where $`\mathrm{\Delta }N`$ and $`\mathrm{\Delta }E`$ are small increments in the number of collisions and collisional energy, and the total number of collisions is $`N`$. The distribution obtained is shown in figure 6. For most of its range it has the functional dependence $`f(E)E^p`$, with $`p`$=0.66.
From the meteoroid’s point of view, the diffuse coma does not change much with altitude when measured in the units of atmospheric mean free path $`l_{\mathrm{atm}}`$ . Figure 7 shows the meteor viewed by a camera that moves with the meteor. Since collisions between the fast coma particles are not considered, the size of the coma is controlled exclusively by $`l_{\mathrm{atm}}`$. Decreasing altitude only increases the number of collisions per unit volume. If the assumption of negligible coma-coma interactions holds all the way to the end of the molecular flow regime, then the meteor would shrink to a few meters in size at 120 km altitude, and below one meter at $``$100 km. The fact that the observed meteor size is larger at these altitudes indicates that the coma eventually becomes dense enough to produce internal collisions. In addition, vaporization starts to produce particles in larger numbers than sputtering, making the coma denser than during sputtering.
Meteor images based on apparent magnitudes (upper rows in figure 4) can be used to calculate the meteor light curve. Total meteor brightness is the sum of all pixels brighter than the limiting magnitude. Hence the absolute meteor magnitude, defined as the brightness at 100 km distance, is
$$M=𝔐_{\mathrm{limit}}2.5\mathrm{log}\left(\frac{D_{\mathrm{limit}}}{100\mathrm{k}\mathrm{m}}\right)^22.5\mathrm{log}\left(\underset{m_{\mathrm{pix}}<6.5^m}{}\frac{N_{\mathrm{collisions}}}{N_{\mathrm{limit}}}\right),$$
(17)
where $`m_{\mathrm{pix}}`$ is the pixels’ apparent magnitude. Since $`M`$ is a function of $`D_{\mathrm{limit}}`$, the light curve depends on the beginning height $`H_{\mathrm{begin}}`$ and on the camera’s position. The dependence on $`H_{\mathrm{begin}}`$ is stronger, because $`N_{\mathrm{limit}}`$ is a function of altitude. Notice that according to equation 3, $`N_{\mathrm{limit}}`$ is a constant for a given meteor and detector. By letting $`N_{\mathrm{limit}}`$ vary with altitude, we declare it an unknown that can be determined experimentally from the pixel value of the limiting magnitude in the first meteor frame.
Figure 8 shows the light curves of our meteor model for two beginning heights: 200 km and 171 km. The cameras have different distances to the starting point of the meteor, thus they differ in the absolute magnitude at $`H_{\mathrm{begin}}`$. In the first few frames, the meteor brightness is slightly erratic as the meteor is a few pixels in size and susceptible to stochastic variations in pixel values. After that the light curve smoothes out, but it is not linear; it shows a change in slope, which is also an observed characteristic of high altitude meteors (Spurný et al., 2000b).
## 5 Conclusion
We have shown that thermalization of particles sputtered from a meteoroid surface by the impinging atmosphere can explain the phenomenon of high altitude meteors. A sputtered particle has a much larger kinetic energy than the surrounding atmospheric particles. It takes between 100 and 200 collisions between particles in the resulting collisional cascade to redistribute this initial energy down to the thermal level of the atmosphere. The volume filled with collisions surrounding the meteoroid is controlled by the atmospheric mean free path. These collisions are the source of the atomic excitations required for the light production.
We showed that this process can create enough photons to be detected by high-sensitivity TV cameras. The number of sputtered particles, and consequently the number of photons produced, is proportional to the meteoroid size. This explains why the observed beginning height correlates with the maximum meteor brightness. Our analytical description of this correlation agrees with the data. We developed a Monte Carlo code specifically for the problem of high altitude meteors. Simulations successfully reproduced the observed size and morphology of these meteors. We found that their shape also depends on the relative angle between the meteor path and the camera’s line of sight.
This model provides the grounds for further study of meteoroid properties and meteor microphysics. Features like streams and jets observed in high altitude meteors are probably created by nonuniform sputtering and meteoroid rotation. This will be addressed in the future by Monte Carlo modeling. It is also possible to trace different atomic species during thermalization, which is a highly nonequilibrium process, and to calculate their spectral lines. Such numerical simulations will remove the need for “fudge” factors typically used to hide underlying microphysics in meteor equations.
## Acknowledgment
This work was supported by NSF grant PHY-0070928. The author would like to thank Peter Jenniskens, SETI Institute, and Bruce Draine, Princeton University, for valuable discussions.
## APPENDIX: Maximum brightness from the single body ablation theory
Meteor altitude $`H`$ is described by the equation of motion:
$$\frac{dH}{dt}=v_\mathrm{m}\mathrm{cos}Z.$$
(18)
Ablation is changing the meteoroid mass as:
$$\frac{dm}{dt}n_{\mathrm{atm}}(H)m^{2/3}v_\mathrm{m}^3.$$
(19)
Combining and integrating these two equations, while keeping the meteor velocity constant (a good approximation for fast meteors like Leonids) and assuming an exponential atmosphere, we get the solution:
$$m^{1/3}(H)=m_\mathrm{m}^{1/3}\mathrm{constant}\times \frac{n_{\mathrm{atm}}(H)}{\mathrm{cos}Z},$$
(20)
where $`m_\mathrm{m}`$ is the initial meteoroid mass. At the end point of the trajectory $`H_{\mathrm{end}}`$, the meteor mass becomes zero, thus
$$n_{\mathrm{atm}}(H_{\mathrm{end}})m_\mathrm{m}^{1/3}\mathrm{cos}Z.$$
(21)
The maximum meteor light intensity is (e.g. Bronshten, 1983; Opik, 1958)
$$I_{\mathrm{max}}m_\mathrm{m}^{2/3}n_{\mathrm{atm}}(H_{\mathrm{end}}).$$
(22)
Since $`m_\mathrm{m}^{1/3}R_\mathrm{m}`$, from equation 22 and 21 it follows that the maximum brightness in magnitudes is
$$𝔐^{max}=7.5\mathrm{log}R_\mathrm{m}2.5\mathrm{log}(\mathrm{cos}Z)+\mathrm{constant}.$$
(23)
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# Negation and Involutive Adjunctions
## 1 Introduction
The goal of this note is to present a phenomenon of adjunction present in assumptions about an involutive negation connective, like classical negation. Proof-theoretical assumptions concerning such a negation make an adjoint situation that we call an *involutive adjunction*. The notion of involutive adjunction amounts, in a sense to be made precise, to adjunction where an endofunctor is adjoint to itself, which in is called *self-adjunction*.
In a series of papers, which starts with (see , and ), Dov Gabbay has been working on characterizations of negation in terms of assumptions about a consequence relation. Sometimes, as in this note, Gabbay concentrates on negation in the absence of any other connective. The context of the present note replaces Gabbay’s logical framework of a consequence relation by a consequence *graph*, as this is done in categorial proof theory. We do not have any more only a relation between premises and conclusions, but we have arrows between them, and there may be more than one such arrow. We are interested in equalities between these arrows. Often these equalities, which are proof-theoretically motivated, exemplify important notions of category theory. This note shows that with an involutive negation we fall on a particular notion of adjunction. This is yet another corroboration of Lawvere’s thesis that all logical constants are tied to adjoint situations (see ), and of Mac Lane’s slogan that adjunction arises everywhere (see , Preface).
## 2 Self-adjunctions
To fix notation and terminology, we will rely on the following definition of the notion of adjunction (cf. , Section IV.1, and , Section 4.1.3).
An *adjunction* is a sextuple $`𝒜,,F,G,\phi ,\gamma `$ where
$`𝒜`$ and $``$ are categories,
$`F`$ from $``$ to $`𝒜`$ and $`G`$ from $`𝒜`$ to $``$ are functors,
$`\phi `$ is a natural transformation of $`𝒜`$ from the composite functor $`FG`$ to the identity functor of $`𝒜`$, which means that the following equation holds in $`𝒜`$ for every arrow $`f:A_1A_2`$ of $`𝒜`$:
| ($`\phi `$ nat) | $`f\phi _{A_1}=\phi _{A_2}FGf`$, |
| --- | --- |
$`\gamma `$ is a natural transformation of $``$ from the identity functor of $``$ to the composite functor $`GF`$, which means that the following equation holds in $``$ for every arrow $`g:B_1B_2`$ of $``$:
| ($`\gamma `$ nat) | $`GFg\gamma _{B_1}=\gamma _{B_2}g`$, |
| --- | --- |
the following *triangular* equations hold in $`𝒜`$ and $``$ respectively:
| ($`\phi \gamma F`$) | $`\phi _{FB}F\gamma _B`$ | $`=\text{1}_{FB}`$, |
| --- | --- | --- |
| ($`\phi \gamma G`$) | $`G\phi _A\gamma _{GA}`$ | $`=\text{1}_{GA}`$. |
A *self-adjunction* is a quadruple $`𝒮,L,\phi ,\gamma `$ where $`𝒮,𝒮,L,L,\phi ,\gamma `$ is an adjunction (this notion is taken over from , Section 10). So, in a self-adjunction, $`L`$ is an endofunctor, and the equations ($`\phi `$ nat) and ($`\gamma `$ nat) become
| $`f\phi _{A_1}=\phi _{A_2}LLf`$, |
| --- |
| $`LLf\gamma _{A_1}=\gamma _{A_2}f`$, |
while the triangular equations become
| ($`\phi \gamma L`$) | $`\phi _{LA}L\gamma _A=L\phi _A\gamma _{LA}=\text{1}_{LA}`$. |
| --- | --- |
A $`𝒦`$-*self-adjunction* is a self-adjunction that satisfies the additional equation
| ($`\phi \gamma 𝒦`$) | $`L(\phi _A\gamma _A)=\phi _{LA}\gamma _{LA}`$, |
| --- | --- |
and a $`𝒥`$-*self-adjunction* is a self-adjunction that satisfies the additional equation
| ($`\phi \gamma 𝒥`$) | $`\phi _A\gamma _A=\text{1}_A`$ |
| --- | --- |
(these notions are also from , Section 10). It is easy to see that every $`𝒥`$-self-adjunction is a $`𝒦`$-self-adjunction (the converse need not hold).
A $`𝒥`$-self-adjunction that satisfies
| ($`\gamma \phi `$) | $`\gamma _A\phi _A=\text{1}_{LLA}`$ |
| --- | --- |
is called a *trivial* self-adjunction. Note that for trivial self-adjunctions it is superfluous to assume the equations ($`\gamma `$ nat) and ($`\phi \gamma G`$), or alternatively ($`\phi `$ nat) and ($`\phi \gamma F`$); these equations can be derived from the remaining ones.
The *free self-adjunction* $`𝒮,L,\phi ,\gamma `$ generated by $`\{p\}`$, where we call $`p`$ a *letter*, is defined as follows. The category $`𝒮`$ has as objects the formulae of the propositional language generated by $`\{p\}`$ with a unary connective $`L`$. We may identify the formulae $`p`$, $`Lp`$, $`LLp`$,$`\mathrm{}`$ of this language with the natural numbers 0, 1, 2,$`\mathrm{}`$
The arrow terms of $`𝒮`$ are defined inductively out of the primitive arrow terms
$$\text{1}_A:AA,\phi _A:LLAA,\gamma _A:ALLA,$$
for every object $`A`$ of $`𝒮`$, with the help of the operations of composition $``$ and the unary operation that assigns to the arrow term $`f:AB`$ the arrow term $`Lf:LALB`$. On these arrow terms we impose the equations of self-adjunctions. In the set of these equations we have of course all the equations $`f=f`$, and this set is closed under symmetry and transitivity of equality, and under the rules
$$(\text{cong })\frac{f=f_1g=g_1}{fg=f_1g_1}$$
$$(\text{cong L})\frac{f=g}{Lf=Lg}$$
We assume for $`f`$ and $`g`$ in $`(\text{cong })`$ that they have composable types, such that $`fg`$ is defined; the same assumption is made for $`f_1`$ and $`g_1`$.
We define analogously the free $`𝒦`$-self-adjunction, the free $`𝒥`$-self-adjunction and the free trivial self-adjunction generated by $`\{p\}`$, just by imposing additional equations.
## 3 Involutive adjunctions
Consider a category $`𝒜`$ and a contravariant functor $`\neg `$ from $`𝒜`$ to $`𝒜`$, which means that for $`f:AB`$ in $`𝒜`$ we have $`\neg f:\neg B\neg A`$ in $`𝒜`$, and the following equations are satisfied:
| ($`\neg 1`$) | $`\neg \text{1}_A=\text{1}_{\neg A}`$ |
| --- | --- |
| ($`\neg 2`$) | $`\neg (fg)=\neg g\neg f`$, for $`f:AB`$ and $`g:CA`$. |
The contravariant functor $`\neg `$ may be conceived either as a functor from the category $`𝒜^{op}`$ to $`𝒜`$, which we denote by $`\neg `$ too, or as a functor from $`𝒜`$ to $`𝒜^{op}`$, which we denote by $`\neg ^{op}`$.
Suppose that for every object $`A`$ of $`𝒜`$ we have an arrow $`n_A^{}:\neg \neg AA`$ of $`𝒜`$. The arrow $`n_A^{}`$ becomes the arrow $`n_A^{op}:A\neg \neg A`$ in $`𝒜^{op}`$.
We say that $`𝒜,\neg ,n^{}`$ is an $`n^{}`$*-adjunction* when
$$𝒜,𝒜^{op},\neg ,\neg ^{op},n^{},n^{op}$$
is an adjunction. This means that in $`𝒜`$ we have for every $`f:A_1A_2`$ the equation
| ($`n^{}`$ nat) | $`fn_{A_1}^{}=n_{A_2}^{}\neg \neg f`$, |
| --- | --- |
alternatively written $`fn_{A_1}^{}=n_{A_2}^{}\neg \neg ^{op}f`$, which also delivers ($`n^{op}`$ nat) in $`𝒜^{op}`$, and the equation
| ($`n^{}`$ triang) | $`n_{\neg A}^{}\neg n_A^{}=\text{1}_{\neg A}`$, |
| --- | --- |
which delivers both the equation ($`\phi \gamma F`$), i.e. ($`n^{}n^{op}\neg `$), in $`𝒜`$, and the equation ($`\phi \gamma G`$), i.e. ($`n^{}n^{op}\neg ^{op}`$), in $`𝒜^{op}`$.
Suppose now that we have as before a category $`𝒜`$ and a contravariant functor $`\neg `$ from $`𝒜`$ to $`𝒜`$, and that for every object $`A`$ of $`𝒜`$ we have an arrow $`n_A^{}:A\neg \neg A`$ of $`𝒜`$. The arrow $`n_A^{}`$ becomes the arrow $`n_A^{op}:\neg \neg AA`$ in $`𝒜^{op}`$.
We say that $`𝒜,\neg ,n^{}`$ is an $`n^{}`$*-adjunction* when
$$𝒜^{op},𝒜,\neg ^{op},\neg ,n^{op},n^{}$$
is an adjunction. This means that in $`𝒜`$ we have for every $`f:A_1A_2`$ the equation
| ($`n^{}`$ nat) | $`\neg \neg fn_{A_1}^{}=n_{A_2}^{}f`$, |
| --- | --- |
which also delivers ($`n^{op}`$ nat) in $`𝒜^{op}`$, and the equation
| ($`n^{}`$ triang) | $`\neg n_A^{}n_{\neg A}^{}=\text{1}_{\neg A}`$, |
| --- | --- |
which delivers both the equation ($`\phi \gamma F`$), i.e. ($`n^{op}n^{}\neg ^{op}`$), in $`𝒜^{op}`$, and the equation ($`\phi \gamma G`$), i.e. ($`n^{op}n^{}\neg `$), in $`𝒜`$. Note that what we call $`n^{}`$-adjunction is called *self-adjunction* in (Section 3.1; cf. also , Section I.8), which should not be confused with our notion of self-adjunction in the preceding section.
We say that $`𝒜,\neg ,n^{},n^{}`$ is an *involutive* adjunction when $`𝒜,\neg ,n^{}`$ is an $`n^{}`$-adjunction and $`𝒜,\neg ,n^{}`$ is an $`n^{}`$-adjunction.
A $`𝒦`$*-involutive* adjunction is an involutive adjunction that satisfies the additional equation
| ($`n^{}n^{}𝒦`$) | $`\neg (n_A^{}n_A^{})=n_{\neg A}^{}n_{\neg A}^{}`$, |
| --- | --- |
and a $`𝒥`$*-involutive* adjunction is an involutive adjunction that satisfies the additional equation
| ($`n^{}n^{}𝒥`$) | $`n_A^{}n_A^{}=\text{1}_A`$. |
| --- | --- |
It is easy to see that every $`𝒥`$-involutive adjunction is a $`𝒦`$-involutive adjunction (the converse need not hold).
A $`𝒥`$-involutive adjunction that satisfies
| ($`n^{}n^{}`$) | $`n_A^{}n_A^{}=\text{1}_{\neg \neg A}`$ |
| --- | --- |
is called a *trivial* involutive adjunction.
Note that for trivial involutive adjunctions it is superfluous to assume the equations ($`n^{}`$ nat) and ($`n^{}`$ triang), or alternatively ($`n^{}`$ nat) and ($`n^{}`$ triang); these equations can be derived from the remaining ones. In trivial involutive adjunctions we have the equations
| $`n_{\neg A}^{}=\neg n_A^{}`$, |
| --- |
| $`n_{\neg A}^{}=\neg n_A^{}`$. |
The *free involutive adjunction* $`𝒜,\neg ,n^{},n^{}`$ generated by $`\{p\}`$ is defined as follows. The category $`𝒜`$ has as objects the formulae of the propositional language generated by $`\{p\}`$ with a unary connective $`\neg `$. We may identify these formulae with the natural numbers.
The arrow terms of $`𝒜`$ are defined inductively out of the primitive arrow terms
$$\text{1}_A:AA,n_A^{}:\neg \neg AA,n_A^{}:A\neg \neg A,$$
for every object $`A`$ of $`𝒜`$, with the help of the operations of composition $``$ and the unary operation that assigns to the arrow term $`f:AB`$ the arrow term $`\neg f:\neg B\neg A`$. On these arrow terms we impose the equations of involutive adjunctions. In the set of these equations we have of course all the equations $`f=f`$, and this set is closed under symmetry and transitivity of equality, under the rule (cong $``$ ), and also under the rule
$$(\text{cong }\neg )\frac{f=g}{\neg f=\neg g}$$
We define analogously the free $`𝒦`$-involutive adjunction, the free $`𝒥`$-involutive adjunction and the free trivial involutive adjunction generated by $`\{p\}`$, just by imposing additional equations.
Note that the category of the free involutive adjunction generated by an arbitrary set having more than one letter would be the disjoint union of isomorphic copies of the category $`𝒜`$ of the free involutive adjunction generated by $`\{p\}`$. An analogous remark applies to the category of the free self-adjunction generated by an arbitrary set having more than one member: it would be the disjoint union of isomorphic copies of the category $`𝒮`$ of the free self-adjunction generated by $`\{p\}`$.
## 4 Self-adjunctions and involutive adjunctions
We are now going to prove that in the free self-adjunction $`𝒮,L,\phi ,\gamma `$ and the free involutive adjunction $`𝒜,\neg ,n^{},n^{}`$, both generated by $`\{p\}`$, the categories $`𝒮`$ and $`𝒜`$ are isomorphic categories.
First, we define $`\neg `$, $`n^{}`$ and $`n^{}`$ in $`𝒮`$ in the following manner. On objects we have that $`\neg `$ is $`L`$, while for the arrow term $`f:AB`$ of $`𝒮`$ we define the arrow term $`\neg f:\neg B\neg A`$ of $`𝒮`$ inductively as follows:
| $`\neg \text{1}_A`$ | $`=L\text{1}_A=\text{1}_{LA}=\text{1}_{\neg A}`$, |
| --- | --- |
| $`\neg \phi _A`$ | $`=L\gamma _A`$, |
| $`\neg \gamma _A`$ | $`=L\phi _A`$, |
| $`\neg (fg)=\neg g\neg f`$, |
| $`\neg Lf`$ | $`=L\neg f`$. |
That this defines an operation $`\neg `$ on the arrows of $`𝒮`$ is shown by verifying that if $`f=g`$ in $`𝒮`$, then $`\neg f=\neg g`$ in $`𝒮`$; we verify, namely, that the equations of $`𝒮`$ are closed under the rule (cong $`\neg `$) of the preceding section. This is done by a straightforward induction on the length of the derivation of $`f=g`$ in $`𝒮`$. For that we use the fact that for every arrow term $`f`$ of $`𝒮`$ the arrow term $`\neg f`$ is equal in $`𝒮`$ to an arrow term of the form $`Lf^{}`$.
Finally, we have
$$n_A^{}=_{\text{df}}\phi _A,n_A^{}=_{\text{df}}\gamma _A.$$
Next, we define $`L`$, $`\phi `$ and $`\gamma `$ in $`𝒜`$ in the following manner. On objects we have that $`L`$ is $`\neg `$, while for the arrow term $`f:AB`$ of $`𝒜`$ we define the arrow term $`Lf:LALB`$ of $`𝒜`$ inductively as follows:
| $`L\text{1}_A`$ | $`=\neg \text{1}_A=\text{1}_{\neg A}=\text{1}_{LA}`$, |
| --- | --- |
| $`Ln_A^{}`$ | $`=\neg n_A^{}`$, |
| $`Ln_A^{}`$ | $`=\neg n_A^{}`$, |
| $`L(fg)=LfLg`$, |
| $`L\neg f`$ | $`=\neg Lf`$. |
That this defines an operation $`L`$ on the arrows of $`𝒜`$ is shown by verifying that if $`f=g`$ in $`𝒜`$, then $`Lf=Lg`$ in $`𝒜`$; we verify, namely, that the equations of $`𝒜`$ are closed under the rule (cong $`L`$) of §2 above. This is done by a straightforward induction on the length of the derivation of $`f=g`$ in $`𝒜`$. For that we use the fact that for every arrow term $`f`$ of $`𝒜`$ the arrow term $`Lf`$ is equal in $`𝒜`$ to an arrow term of the form $`\neg f^{}`$.
Finally, we have
$$\phi _A=_{\text{df}}n_A^{},\gamma _A=_{\text{df}}n_A^{}.$$
We verify easily by induction on the complexity of the arrow term $`f`$ that both in $`𝒮`$ and in $`𝒜`$ we have the equation
| ($`LL\neg \neg `$) | $`LLf=\neg \neg f`$. |
| --- | --- |
Next we verify that the equations of involutive adjunctions hold for the defined $`\neg `$, $`n^{}`$ and $`n^{}`$ in $`𝒮`$. This is done in a straightforward manner by induction on the length of derivation. In the basis of this induction, we use ($`LL\neg \neg `$), ($`\phi `$ nat) and ($`\gamma `$ nat) to verify ($`n^{}`$ nat) and ($`n^{}`$ nat), while the equations ($`n^{}`$ triang) and ($`n^{}`$ triang) reduce to ($`\phi \gamma L`$). In the induction step, we rely on the closure of $`𝒮`$ under (cong $`\neg `$), which we established above.
We verify also that the equations of self-adjunctions hold for the defined $`L`$, $`\phi `$ and $`\gamma `$ in $`𝒜`$. This is done again in a straightforward manner by induction on the length of derivation. In the basis of this induction, we use ($`LL\neg \neg `$), ($`n^{}`$ nat) and ($`n^{}`$ nat) to verify ($`\phi `$ nat) and ($`\gamma `$ nat), while the equations ($`\phi \gamma L`$) reduce to ($`n^{}`$ triang) and ($`n^{}`$ triang). In the induction step, we rely on the closure of $`𝒜`$ under (cong $`L`$), which we established above.
We have a functor $`F_𝒜`$ from $`𝒮`$ to $`𝒜`$ that maps the object of $`𝒮`$ corresponding to the natural number $`n`$ to the object of $`𝒜`$ corresponding to $`n`$, and that maps every arrow of $`𝒮`$ to the homonymous arrow in the defined $`𝒮`$ structure of $`𝒜`$. For example,
$$F_𝒜\phi _{LLp}=\phi _{\neg \neg p}=n_{\neg \neg p}^{}.$$
We define analogously a functor $`F_𝒮`$ from $`𝒜`$ to $`𝒮`$. That $`F_𝒜`$ and $`F_𝒮`$ are indeed functors follows from what we established above.
It is trivial that on objects we have that $`F_𝒮F_𝒜A`$ is $`A`$, and that $`F_𝒜F_𝒮B`$ is $`B`$. We show next by induction on the complexity of $`f`$ that in $`𝒮`$ we have
$$F_𝒮F_𝒜f=f.$$
When $`f`$ is of the form $`Lf^{}`$, we make an auxiliary induction on the complexity of $`f^{}`$, in which we use ($`LL\neg \neg `$). We show analogously that in $`𝒜`$ we have
$$F_𝒜F_𝒮g=g.$$
This concludes the proof that $`𝒮`$ and $`𝒜`$ are isomorphic categories.
We demonstrate analogously that the categories of, respectively,
the free $`𝒦`$-self-adjunction and the free $`𝒦`$-involutive adjunction,
the free $`𝒥`$-self-adjunction and the free $`𝒥`$-involutive adjunction,
the free trivial self-adjunction and the free trivial involutive adjunction,
all generated by $`\{p\}`$, are isomorphic categories.
The interest of considering $`𝒦`$ and $`𝒥`$ versions of self-adjunctions and involutive adjunctions comes from connections with Temperley-Lieb algebras and the associated geometrical interpretation (see and references therein). Roughly speaking, $`𝒦`$ is what we find in Temperley-Lieb algebras, where only the number of circles (which correspond to $`\phi _A\gamma _A`$ or $`n_A^{}n_A^{}`$) counts, while in $`𝒥`$ circles are disregarded.
The free trivial self-adjunction, and hence also the free trivial involutive adjunction, are preorders; namely, all arrows with the same source and target are equal. This follows from the results of (unabridged version) or .
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# Phases of QCD: lattice thermodynamics and a field theoretical model11footnote 1Work supported in part by BMBF and INFN
## 1 Introduction
Recent years have seen an expansion of activities devoted to the study of the QCD phase diagram. Heavy-ion experiments are looking for signals of the Quark-Gluon Plasma. Large-scale lattice simulations at finite temperature have become a principal tool for investigating the pattern of phases in QCD. Accurate computations of lattice QCD thermodynamics in the pure gauge sector have been performed. First results at finite quark chemical potential are available. The equation of state of strongly interacting matter is now at hand as a function of temperature $`T`$ and in a limited range of quark chemical potential $`\mu `$. Improved multi-parameter re-weighting techniques , Taylor series expansion methods and analytic continuation from imaginary chemical potential provide lattice data for the pressure, entropy density, quark density and selected susceptibilities.
A straightforward interpretation of these data in terms of QCD perturbation theory does not work because of poor convergence at any temperature of practical interest . In order to overcome this problem, resummation schemes have been proposed, based for example on the Hard Thermal Loop (HTL) approach or on dimensionally reduced screened perturbation theory (DRSPT) . However, these approaches still give reliable results only for temperatures $`T2.5T_c`$, far above the critical temperature $`T_c`$ 0.2 GeV. At these high temperatures, the HTL approach motivates and justifies a picture of weakly interacting quasiparticles, as determined by the HTL propagators.
In order to extend such descriptions to lower temperatures closer to $`T_c`$, various models have been proposed. Early attempts were based on the MIT bag model . More sophisticated approaches became necessary when more precise lattice data appeared. Various aspects of QCD thermodynamics have been investigated in terms of quasiparticle models based on perturbative calculations carried out in the HTL scheme , in terms of a condensate of $`Z_3`$ Wilson lines , by refined quasiparticle models based on the HTL-resummed entropy and extensions thereof , by an improved version with a temperature-dependent number of active degrees of freedom , by an evaporation model of the gluon condensate , by quasiparticle models formulated in dynamical terms , and by hadron resonance gas models below the critical temperature (for a recent review see ).
In this paper, we study the thermodynamics of two-flavour QCD at finite quark chemical potential. Our investigation is based on a synthesis of a Nambu Jona-Lasinio (NJL) model and the non-linear dynamics involving the Polyakov loop . In this Polyakov-loop-extended (PNJL) model, quarks develop quasiparticle masses by propagating in the chiral condensate, while they couple at the same time to a homogeneous background (temporal) gauge field representing Polyakov loop dynamics.
The “classic” NJL model incorporates the chiral symmetry of two-flavour QCD and its spontaneous breakdown at $`T<T_c`$. Gluonic degrees of freedom are “integrated out” and replaced by a local four-point interaction of quark colour currents. Subsequent Fierz transformations project this interaction into various quark-antiquark and diquark channels. The colour singlet $`q\overline{q}`$ modes of lowest mass are identified with the lightest mesons. Pions properly emerge as Goldstone bosons at $`T<T_c`$. However, the local $`SU(N_c)`$ gauge invariance of QCD is now replaced by a global $`SU(N_c)`$ symmetry in the NJL model, so that the confinement property is lost. Consequently, standard NJL-type models are bound to fail in attempts to describe $`N_c=3`$ thermodynamics around $`T_c`$ (and beyond) for non-zero quark chemical potential $`\mu `$.
On the other hand the NJL model, with just the simplest possible one-parameter colour current-current interaction between quarks, is remarkably successful in reproducing the thermodynamics of $`N_c=2`$ Lattice QCD at finite $`\mu `$ . Encouraged by this result, the NJL quasiparticle concept does suggest itself as a useful starting point. However, whereas aspects of deconfinement are less significant in the $`N_c=2`$ case, they figure prominently for $`N_c=3`$. This motivates our extension towards the PNJL Lagrangian as a minimal approach incorporating both chiral symmetry restoration and deconfinement.
The deconfinement phase transition is well defined in the heavy-quark limit, where the Polyakov loop serves as an order parameter. This phase transition is characterized by the spontaneous breaking of the $`Z(3)`$ center symmetry of QCD . In the presence of dynamical quarks the center symmetry is explicitly broken. No order parameter is established for the deconfinement transition in this case , but the Polyakov loop still serves as an indicator of a rapid crossover towards deconfinement. The chiral phase transition, on the other hand, has a well-defined order parameter in the chiral limit of massless quarks: the chiral (or quark) condensate $`\overline{q}q`$. This condensate, and its dynamical generation, is the basic element of the original NJL model.
The primary aim of this paper is to test the effectiveness of the PNJL approach when confronted with Lattice QCD thermodynamics. The PNJL Lagrangian is derived in Section 2. Parameters are fixed in Section 3 by reproducing known properties of the pion and of the QCD vacuum in the hadronic phase, while the Polyakov loop effective potential is adjusted to pure gauge lattice results. Sections 4 and 5 deal with the thermodynamics derived from the PNJL Lagrangian in the mean-field approximation. This discussion includes a detailed comparison with Lattice QCD results at zero and at finite chemical potential.
## 2 The PNJL model
Following we introduce a generalized $`N_f=2`$ Nambu Jona-Lasinio Lagrangian with quarks coupled to a (spatially constant) temporal background gauge field representing Polyakov loop dynamics (the PNJL model):
$`_{PNJL}=\overline{\psi }\left(i\gamma _\mu D^\mu \widehat{m}_0\right)\psi +{\displaystyle \frac{G}{2}}\left[\left(\overline{\psi }\psi \right)^2+\left(\overline{\psi }i\gamma _5\stackrel{}{\tau }\psi \right)^2\right]𝒰(\mathrm{\Phi }[A],\overline{\mathrm{\Phi }}[A],T),`$ (1)
where $`\psi =(\psi _u,\psi _d)^T`$ is the quark field,
$`D^\mu =^\mu iA^\mu \mathrm{and}A^\mu =\delta _{\mu 0}A^0.`$ (2)
The gauge coupling $`g`$ is conveniently absorbed in the definition of $`A^\mu (x)=g𝒜_a^\mu (x)\frac{\lambda _a}{2}`$ where $`𝒜_a^\mu `$ is the SU(3) gauge field and $`\lambda _a`$ are the Gell-Mann matrices. The two-flavour current quark mass matrix is $`\widehat{m}_0=diag(m_u,m_d)`$ and we shall work in the isospin symmetric limit with $`m_u=m_dm_0`$. A local, chirally symmetric scalar-pseudoscalar four-point interaction of the quark fields is introduced with an effective coupling strength $`G`$.
The quantity $`𝒰(\mathrm{\Phi },\overline{\mathrm{\Phi }},T)`$ is the effective potential expressed in terms of the traced Polyakov loop<sup>2</sup><sup>2</sup>2more precisely: the Polyakov line with periodic boundary conditions and its (charge) conjugate,
$$\mathrm{\Phi }=(\mathrm{Tr}_cL)/N_c,\overline{\mathrm{\Phi }}=(\mathrm{Tr}_cL^{})/N_c.$$
(3)
The Polyakov loop $`L`$ is a matrix in colour space explicitly given by
$$L\left(\stackrel{}{x}\right)=𝒫\mathrm{exp}\left[i_0^\beta 𝑑\tau A_4(\stackrel{}{x},\tau )\right],$$
(4)
with $`\beta =1/T`$ the inverse temperature and $`A_4=iA^0`$. In a convenient gauge (the so-called Polyakov gauge), the Polyakov loop matrix can be given a diagonal representation .
The coupling between Polyakov loop and quarks is uniquely determined by the covariant derivative $`D_\mu `$ in the PNJL Lagrangian (1). Note that in the chiral limit ($`\widehat{m}_00`$), this Lagrangian is invariant under the chiral flavour group, $`SU(2)_L\times SU(2)_R`$, just like the original QCD Lagrangian.
The trace of the Polyakov loop, $`\mathrm{\Phi }`$, and its conjugate, $`\overline{\mathrm{\Phi }}`$, will be treated as classical field variables throughout this work. In the absence of quarks, we have $`\mathrm{\Phi }=\overline{\mathrm{\Phi }}`$ and the Polyakov loop serves as an order parameter for deconfinement. The phase transition is characterized by the spontaneous breaking of the $`Z(3)`$ center symmetry of QCD. The temperature dependent effective potential $`𝒰`$ has the following general features. At low temperatures, $`𝒰`$ has a single minimum at $`\mathrm{\Phi }`$=0, while at high temperatures it develops a second one which turns into the absolute minimum above a critical temperature $`T_0`$. In the limit $`T\mathrm{}`$ we have $`\mathrm{\Phi }1`$. The function $`𝒰(\mathrm{\Phi },\overline{\mathrm{\Phi }},T)`$ will be fixed by comparison with pure-gauge Lattice QCD. We choose the following general form in accordance with the underlying $`Z(3)`$ symmetry:
$`{\displaystyle \frac{𝒰(\mathrm{\Phi },\overline{\mathrm{\Phi }},T)}{T^4}}={\displaystyle \frac{b_2\left(T\right)}{2}}\overline{\mathrm{\Phi }}\mathrm{\Phi }{\displaystyle \frac{b_3}{6}}\left(\mathrm{\Phi }^3+\overline{\mathrm{\Phi }}^3\right)+{\displaystyle \frac{b_4}{4}}\left(\overline{\mathrm{\Phi }}\mathrm{\Phi }\right)^2`$ (5)
with
$$b_2\left(T\right)=a_0+a_1\left(\frac{T_0}{T}\right)+a_2\left(\frac{T_0}{T}\right)^2+a_3\left(\frac{T_0}{T}\right)^3.$$
(6)
A precision fit of the coefficients $`a_i,b_i`$ is performed to reproduce the lattice data (see section 3).
Using standard bosonization techniques the Lagrangian (1) can be rewritten in terms of the auxiliary field variables $`\sigma `$ and $`\stackrel{}{\pi }`$:
$`_{eff}={\displaystyle \frac{\sigma ^2+\stackrel{}{\pi }^2}{2G}}𝒰(\mathrm{\Phi },\overline{\mathrm{\Phi }},T)i\mathrm{Tr}\mathrm{ln}S^1,`$ (7)
where an irrelevant constant has been dropped and
$$S^1=i\gamma _\mu ^\mu \gamma _0A^0\widehat{M}$$
(8)
is the inverse quark propagator with
$$\widehat{M}=\widehat{m}_0\sigma i\gamma _5\stackrel{}{\tau }\stackrel{}{\pi }.$$
(9)
The trace in (7) is taken over colour, flavour and Dirac indices. The field equations for $`\sigma `$, $`\stackrel{}{\pi }`$, $`\mathrm{\Phi }`$ and $`\overline{\mathrm{\Phi }}`$ are then solved in the mean field approximation<sup>3</sup><sup>3</sup>3In the mean field approximation the fields are replaced by their expectation values for which, in later sections, we will continue using the notation $`\sigma `$, $`\mathrm{\Phi }`$ and $`\overline{\mathrm{\Phi }}`$ for simplicity and convenience.. The expectation value $`\stackrel{}{\pi }`$ of the pseudoscalar isotriplet field is equal to zero for isospin-symmetric systems.
The $`\sigma `$ field has a non-vanishing vacuum expectation value as a consequence of spontaneous chiral symmetry breaking. Solving the field equations for $`\sigma `$, the effective quark mass $`m`$ is determined by the self-consistent gap equation
$$m=m_0\sigma =m_0G\overline{\psi }\psi .$$
(10)
Note that $`\sigma =G\overline{\psi }\psi `$ is negative in our representation, and the chiral (quark) condensate is $`\overline{\psi }\psi =\overline{\psi }_u\psi _u+\overline{\psi }_d\psi _d`$. For later purposes we note that $`\mathrm{\Phi }`$ and $`\overline{\mathrm{\Phi }}`$ are two independent field variables in the general case of finite quark chemical potential $`\mu `$. They become equal in the limiting case $`\mu =0`$. Their (thermal) expectation values $`\mathrm{\Phi }`$ and $`\overline{\mathrm{\Phi }}`$ are both real but differ at non-zero $`\mu `$.
Before passing to the actual calculations, we summarize basic assumptions behind eq. (1) and comment on limitations to be kept in mind. In fact the PNJL model (1) is quite schematic in several respects. It reduces gluon dynamics to a) chiral point couplings between quarks, and b) a simple static background field representing the Polyakov loop. This picture cannot be expected to work beyond a limited range of temperatures. At large $`T`$, transverse gluons are known to be thermodynamically active degrees of freedom, but they are ignored in the PNJL model. To what extent this model can reproduce lattice QCD thermodynamics is nonetheless a relevant question. We can assume that its range of applicability is, roughly, $`T(23)T_c`$, based on the conclusion drawn in ref. that transverse gluons start to contribute significantly for $`T>2.5T_c`$.
## 3 Parameter fixing
### 3.1 Polyakov loop effective potential
The parameters of the Polyakov loop potential $`𝒰`$ are fitted to reproduce the lattice data for QCD thermodynamics in the pure gauge sector. Minimizing $`𝒰(\mathrm{\Phi },\overline{\mathrm{\Phi }},T)`$ one has $`\mathrm{\Phi }=\overline{\mathrm{\Phi }}`$ and the pressure of the pure-gauge system is evaluated as $`p(T)=𝒰(T)`$ with $`\mathrm{\Phi }(T)`$ determined at the minimum. The entropy and energy density are then obtained by means of the standard thermodynamic relations. In Fig. 1 we show the (scaled) pressure, energy density and entropy density as functions of temperature. The lattice data are reproduced extremely well using the ansatz (5, 6), with parameters summarized in Table 1. The critical temperature $`T_0`$ for deconfinement appearing in Eq. (6) is fixed at $`T_0=270`$ MeV in the pure gauge sector. The resulting effective potential is displayed in Fig. 2 for two different temperatures: $`T=200`$ MeV (below $`T_0`$) and $`T=320`$ MeV (above $`T_0`$).
With the same parametrization, we are also able to reproduce the lattice data for the temperature dependence of the Polyakov loop itself. A comparison between these data and our results is shown in Fig. 3. The Polyakov loop vanishes below the critical temperature $`T_0`$, at which point it jumps discontinuously to a finite value, indicating a first order phase transition. It tends to one at large temperatures, as expected.
### 3.2 NJL sector
The pure NJL model part of the Lagrangian (1) has the following parameters: the “bare” quark mass $`m_0`$, a three-momentum cutoff $`\mathrm{\Lambda }`$ and the coupling strength $`G`$. We choose to reproduce the known chiral physics in the hadronic sector at $`T=0`$ and fix the three parameters by the following conditions:
* The pion decay constant is reproduced at its empirical value, $`f_\pi =92.4`$ MeV. In the NJL model, $`f_\pi `$ is evaluated using the following equation:
$$f_\pi ^2=4m^2I_\mathrm{\Lambda }^{(1)}\left(m\right)\mathrm{where}I_\mathrm{\Lambda }^{(1)}\left(m\right)=iN_c\frac{d^4p}{\left(2\pi \right)^4}\frac{\theta \left(\mathrm{\Lambda }^2\stackrel{}{p}^2\right)}{\left(p^2m^2+iϵ\right)^2},$$
(11)
with the effective (constituent) quark mass $`m`$ determined self-consistently by the gap equation (10).
* The quark condensate becomes
$$\overline{\psi }_u\psi _u=4mI_\mathrm{\Lambda }^{(0)}\left(m\right)$$
(12)
with
$$I_\mathrm{\Lambda }^{(0)}\left(m\right)=iN_c\frac{d^4p}{\left(2\pi \right)^4}\frac{\theta \left(\mathrm{\Lambda }^2\stackrel{}{p}^2\right)}{p^2m^2+iϵ}.$$
(13)
Its “empirical” value derived from QCD sum rules is
$$\overline{\psi }_u\psi _u^{1/3}\overline{\psi }_d\psi _d^{1/3}=\left(240\pm 20\right)\mathrm{MeV}.$$
(14)
* The current quark mass $`m_0`$ is fixed from the Gell-Mann, Oakes, Renner (GMOR) relation which is satisfied in the NJL model:
$$m_\pi ^2=\frac{m_0\overline{\psi }\psi }{f_\pi ^2}.$$
(15)
In the chiral limit, $`m_0=0`$ and $`m_\pi =0`$.
The values of the NJL model parameters, together with the resulting physical quantities, are summarized in Table 2.
## 4 Results at finite $`T`$ and $`\mu `$
We now extend the model to finite temperature and chemical potentials using the Matsubara formalism. We consider the isospin symmetric case, with an equal number of $`u`$ and $`d`$ quarks (and therefore a single quark chemical potential $`\mu `$). The quantity to be minimized at finite temperature is the thermodynamic potential per unit volume:
$`\mathrm{\Omega }(T,\mu )`$ $`=`$ $`𝒰(\mathrm{\Phi },\overline{\mathrm{\Phi }},T){\displaystyle \frac{T}{2}}{\displaystyle \underset{n}{}}{\displaystyle \frac{d^3p}{\left(2\pi \right)^3}\mathrm{Tr}\mathrm{ln}\frac{\stackrel{~}{S}^1(i\omega _n,\stackrel{}{p})}{T}}+{\displaystyle \frac{\sigma ^2}{2G}}.`$
Here $`\omega _n=(2n+1)\pi T`$ are the Matsubara frequencies for fermions. The inverse quark propagator (in Nambu-Gorkov representation) becomes
$$\stackrel{~}{S}^1(p^0,\stackrel{}{p})=\left(\begin{array}{ccc}\gamma _0p^0\stackrel{}{\gamma }\stackrel{}{p}m\gamma _0\left(\mu +iA_4\right)& 0& \\ 0& \gamma _0p^0\stackrel{}{\gamma }\stackrel{}{p}m+\gamma _0\left(\mu +iA_4\right)& \end{array}\right).$$
(17)
Using the identity $`\mathrm{Tr}\mathrm{ln}\left(X\right)=\mathrm{ln}det\left(X\right)`$ we reduce the trace in (LABEL:omega) and find:
$`\mathrm{\Omega }`$ $`=`$ $`𝒰(\mathrm{\Phi },\overline{\mathrm{\Phi }},T)+{\displaystyle \frac{\sigma ^2}{2G}}`$
$``$ $`2N_fT{\displaystyle \frac{\mathrm{d}^3p}{\left(2\pi \right)^3}\left\{\mathrm{Tr}_c\mathrm{ln}\left[1+L\mathrm{e}^{\left(E_p\mu \right)/T}\right]+\mathrm{Tr}_c\mathrm{ln}\left[1+L^{}\mathrm{e}^{\left(E_p+\mu \right)/T}\right]\right\}}`$
$``$ $`6N_f{\displaystyle \frac{\mathrm{d}^3p}{\left(2\pi \right)^3}E_p\theta \left(\mathrm{\Lambda }^2\stackrel{}{p}^2\right)}`$
where we have introduced the quark quasiparticle energy $`E_p=\sqrt{\stackrel{}{p}^2+m^2}`$. The last term involves the NJL three-momentum cutoff $`\mathrm{\Lambda }`$. The second (finite) term does not require any cutoff. A small violation of the underlying chiral symmetry at $`T>0.4`$ GeV, resulting from this procedure, is of no practical relevance since the model is supposed to be applied only at temperatures and chemical potential well below $`\mathrm{\Lambda }`$.
The remaining colour trace is then performed with the result
$`\mathrm{ln}\mathrm{det}\left[1+L\mathrm{e}^{\left(E_p\mu \right)/T}\right]+\mathrm{ln}\mathrm{det}\left[1+L^{}\mathrm{e}^{\left(E_p+\mu \right)/T}\right]`$ (19)
$`=`$ $`\mathrm{ln}\left[1+3\left(\mathrm{\Phi }+\overline{\mathrm{\Phi }}\mathrm{e}^{\left(E_p\mu \right)/T}\right)\mathrm{e}^{\left(E_p\mu \right)/T}+\mathrm{e}^{3\left(E_p\mu \right)/T}\right]`$
$`+`$ $`\mathrm{ln}\left[1+3\left(\overline{\mathrm{\Phi }}+\mathrm{\Phi }\mathrm{e}^{\left(E_p+\mu \right)/T}\right)\mathrm{e}^{\left(E_p+\mu \right)/T}+\mathrm{e}^{3\left(E_p+\mu \right)/T}\right].`$
From the thermodynamic potential (LABEL:omega2) the equations of motion for the mean fields $`\sigma ,\mathrm{\Phi }`$ and $`\overline{\mathrm{\Phi }}`$ are derived through
$$\frac{\mathrm{\Omega }}{\sigma }=0,\frac{\mathrm{\Omega }}{\mathrm{\Phi }}=0,\frac{\mathrm{\Omega }}{\overline{\mathrm{\Phi }}}=0.$$
(20)
This set of coupled equations is then solved for the fields as functions of temperature $`T`$ and quark chemical potential $`\mu `$.
Fig. 4(a) shows the chiral condensate together with the Polyakov loop $`\mathrm{\Phi }`$ as functions of temperature at $`\mu =0`$ where we find again $`\mathrm{\Phi }=\overline{\mathrm{\Phi }}`$. One observes that the introduction of quarks coupled to the $`\sigma `$ and $`\mathrm{\Phi }`$ fields turns the first-order transition seen in pure-gauge Lattice QCD into a continuous crossover. The original 1st order transition in the pure-gauge system appears at a critical temperature $`T_0=270`$ MeV. With the introduction of quarks, the crossover transitions for the chiral condensate $`\overline{\psi }\psi `$ and for the Polyakov loop perfectly coincide at a lower critical temperature $`T_c220`$ MeV (see Fig. 4(b)). We point out that this feature is obtained without changing a single parameter with respect to the pure gauge case. The value of the critical temperature that we obtain is a little high if compared to the available data for two-flavour Lattice QCD which gives $`T_c=(173\pm 8)`$ MeV. On the other hand, it is presently being discussed that detailed continuum extrapolation of these data can increase this temperature up to 210 MeV . For quantitative comparison with existing lattice results we choose to reduce $`T_c`$ by rescaling the parameter $`T_0`$ from 270 to 190 MeV. In this case we loose the perfect coincidence of the chiral and deconfinement transitions, but they are shifted relative to each other by less than 20 MeV. When defining $`T_c`$ in this case as the average of the two transition temperatures we find $`T_c=180`$ MeV. This is also consistent with the observations reported in .
As we turn to non-zero chemical potential, we find that $`\mathrm{\Phi }`$ and $`\overline{\mathrm{\Phi }}`$ are different from each other, even if they are both real. They will finally coincide again at high temperatures, as can be seen in Fig. 5. This feature was already observed in .
With increasing chemical potential, the crossover pattern evolves to lower transition temperatures (see Fig. 6) until it turns to a first order transition around $`\mu 0.3`$ GeV. At this point Cooper pairing of quarks presumably sets in. A more detailed discussion of the critical point and its neighbourhood therefore requires the additional incorporation of explicit diquark degrees of freedom in the PNJL model. Further developments along these lines will be reported elsewhere.
## 5 Detailed comparison with Lattice QCD
A primary aim of this work is to compare predictions of our PNJL model with the lattice data available for full QCD thermodynamics at zero and finite $`\mu `$. Consider first the pressure of the quark-gluon system at zero chemical potential:
$$p(T,\mu =0)=\mathrm{\Omega }(T,\mu =0;\sigma (T,0),\mathrm{\Phi }(T,0),\overline{\mathrm{\Phi }}(T,0)),$$
(21)
where $`\sigma (T,0),\mathrm{\Phi }(T,0)`$ and $`\overline{\mathrm{\Phi }}(T,0)`$ are the solutions of the field equations at finite temperature and zero quark chemical potential. Our results are presented in Fig. 7(a) in comparison with corresponding lattice data. We point out that the input parameters of the PNJL model have been fixed independently in the pure gauge and hadronic sectors, so that our calculated pressure is a prediction of the model, without any further tuning of parameters. With this in mind, the agreement with lattice results is quite satisfactory. One must note that the lattice data are grouped in different sets obtained on lattices with temporal extent $`N_t=4`$ and $`N_t=6`$, both of which are not continuum extrapolated. In contrast, our calculation should, strictly speaking, be compared to the continuum limit. In order to perform meaningful comparisons, the pressure is divided by its asymptotic high-temperature (Stefan-Boltzmann) limit for each given case. At high temperatures our predicted curve should be located closer to the $`N_t=6`$ set than to the one with $`N_t=4`$. This is indeed the case. Furthermore, Fig. 7(b) shows the predicted “interaction measure”, $`(\epsilon 3p)/T^4`$, in comparison with lattice data for $`N_t=6`$. One should of course note that the lattice results have been produced using relatively large quark masses, with pseudoscalar-to-vector mass ratios $`m_{PS}/m_V`$ around 0.7, whereas our calculation is performed with light quark masses corresponding to the physical pion mass. We have investigated the dependence of the pressure and of the energy density on the quark mass and found that the critical temperature scales approximately as $`T_cT_c(m_\pi =0)+0.04m_\pi `$, in agreement with the behaviour found in . Once the rescaling of $`T_c`$ is taken into account, the curves plotted in Figs. 7(a) and 7(b) as functions of $`T/T_c`$ have negligible remaining dependence on the quark mass.
At non-zero chemical potential, quantities of interest that have become accessible in Lattice QCD are the “pressure difference” and the quark number density. The (scaled) pressure difference is defined as:
$`{\displaystyle \frac{\mathrm{\Delta }p(T,\mu )}{T^4}}={\displaystyle \frac{p(T,\mu )p\left(T,\mu =0\right)}{T^4}}.`$ (22)
A comparison of $`\mathrm{\Delta }p`$, calculated in the PNJL model, with lattice results is presented in Fig. 8. This figure shows the scaled pressure difference as a function of the temperature for a series of chemical potentials, with values ranging between $`\mu =0.2T_c^{(0)}`$ and $`\mu T_c^{(0)}`$. The agreement between our results and the lattice data is quite satisfactory.
A related quantity for which lattice results at finite $`\mu `$ exist, is the scaled quark number density, defined as:
$$\frac{n_q(T,\mu )}{T^3}=\frac{1}{T^3}\frac{\mathrm{\Omega }(T,\mu )}{\mu }.$$
(23)
Our results for $`n_q`$ as a function of the temperature, for different values of the quark chemical potential, are shown in Fig. 9 in comparison with corresponding lattice data . Also in this case, the agreement between our PNJL model and the corresponding lattice data is surprisingly good.
It is instructive to study the effect of the Polyakov loop dynamics on the behaviour of the quark density $`n_q`$. The coupling of the quark quasiparticles to the field $`\mathrm{\Phi }`$ reduces their weight as thermodynamically active degrees of freedom when the critical temperature $`T_c`$ is approached from above. At $`T_c`$ the value of $`\mathrm{\Phi }`$ tends to zero and the quasiparticle exponentials exp$`[(E_p\pm \mu )/T]`$ are progressively suppressed in the thermodynamic potential as $`TT_c`$. This is what can be interpreted as the impact of confinement in the context of the PNJL model. In contrast, the standard NJL model without coupling to the Polyakov loop does not have this important feature, so that the quark density leaks strongly into the “forbidden” domain $`T<T_c170`$ MeV, as demonstrated in Fig. 10.
It is a remarkable feature that the quark densities and the pressure difference at finite $`\mu `$ are so well reproduced even though the lattice “data” have been obtained by a Taylor expansion up to fourth order in $`\mu `$, whereas our thermodynamic potential is used with its full functional dependence on $`\mu `$. We have examined the convergence in powers of $`\mu `$ by expanding Eq.(LABEL:omega2). It turns out that the Taylor expansion to order $`\mu ^2`$ deviates from the full result by less than 10 % even at a chemical potential as large as $`\mu T_c`$. When expanded to $`𝒪(\mu ^4)`$, no visible difference is left between the approximate and full calculations for all cases shown in Figs. 8 and 9.
## 6 Summary and conclusions
We have studied a Polyakov-loop-extended Nambu and Jona-Lasinio (PNJL) model with the aim of exploring whether such an approach can catch essential features of QCD thermodynamics when confronted with results of lattice computations at finite temperature and non-zero quark chemical potential. This PNJL model represents a minimal synthesis of the two basic principles that govern QCD at low temperatures: spontaneous chiral symmetry breaking and confinement. The respective order parameters (the chiral quark condensate and the Polyakov loop) are given the meaning of collective degrees of freedom. Quarks couple to these collective fields according to the symmetry rules dictated by QCD itself.
Once a limited set of input parameters is fitted to Lattice QCD in the pure gauge sector and to pion properties in the hadron sector, the quark-gluon thermodynamics above $`T_c`$ up to about twice the critical temperature is well reproduced, including quark densities up to chemical potentials of about 0.2 GeV. In particular, the PNJL model correctly describes the step from the first-order deconfinement transition observed in pure-gauge Lattice QCD (with $`T_c270`$ MeV) to the crossover transition (with $`T_c`$ around 200 MeV) when $`N_f=2`$ light quark flavours are added. The non-trivial result is that the crossovers for chiral symmetry restoration and deconfinement almost coincide, as found in lattice simulations. The model also reproduces the quark number densities at various chemical potentials remarkably well when confronted with corresponding lattice data. Considering that the lattice results have been found by a Taylor expansion in powers of the chemical potential, this excellent agreement came as a surprise and indicates rapid convergence of the power series in $`\mu `$.
Further developments will be directed towards improvements to overcome some obvious limitations. First, the NJL model operates with a constant four-point coupling strength which supposedly averages the relevant running coupling over a limited low-energy kinematic domain, corresponding to temperatures $`T2T_c`$ and chemical potentials $`\mu 0.3`$ GeV. Contacts with the high-temperature limit of QCD and the HTL approaches need to be established. Secondly, in order to proceed into the range of larger chemical potentials, diquark degrees of freedom need to be explicitly involved. Also, the effective potential for the Polyakov loop field, determined so far entirely as a function of temperature by investigating the pure gauge sector, must be examined with respect to its dependence on the chemical potential. And furthermore, the extension to 2+1 flavours with inclusion of strange quarks must be explored.
Nevertheless, considering the simplicity of the PNJL model, the conclusion that can be drawn at this point is promising: it appears that a relatively straightforward quasiparticle approach, with its dynamics rooted in spontaneous chiral symmetry breaking and confinement and with parameters controlled by a few known properties of the gluonic and hadronic sectors of the QCD phase diagram, can account for essential observations from two-flavour $`N_c=3`$ Lattice QCD thermodynamics.
## Acknowledgements
We gratefully acknowledge stimulating discussions with Jean-Paul Blaizot, Ulrich Heinz, Volker Koch, Krishna Rajagopal and Helmut Satz. One of us (W.W.) thanks the nuclear theory group at the Lawrence Berkeley National Lab for their kind hospitality. This work was supported in part by INFN and BMBF.
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# Generic initial ideals and exterior algebraic shifting of the join of simplicial complexes
## Introduction
Algebraic shifting, which was introduced by Kalai, is a map which associates with each simplicial complex $`\sigma `$ another simplicial complex $`\mathrm{\Delta }(\sigma )`$ with special conditions. Nevo studied some properties of algebraic shifting with respect to basic constructions of simplicial complexes, such as union, cone and join. With respect to union and cone, algebraic shifting behaves nicely. However, with respect to join, Nevo found that algebraic shifting does not behave nicely contrary to a conjecture by Kalai .
First, we will recall Kalai’s conjecture and Nevo’s counter example. Let $`\sigma `$ be a simplicial complex on $`\{1,2,\mathrm{},k\}`$, $`\tau `$ a simplicial complex on $`\{k+1,k+2,\mathrm{},n\}`$, and let $`\sigma \tau `$ denote their join, in other words,
$$\sigma \tau =\{SR:S\sigma \text{and}R\tau \}.$$
Kalai conjectured that $`\mathrm{\Delta }(\sigma \tau )=\mathrm{\Delta }(\mathrm{\Delta }(\sigma )\mathrm{\Delta }(\tau ))`$, where $`\mathrm{\Delta }(\sigma )`$ is the exterior algebraic shifted complex of $`\sigma `$. However, Nevo found a counter example. We quote his example. Let $`\mathrm{\Sigma }(\sigma )`$ denote the suspension of $`\sigma `$, i.e., the join of $`\sigma `$ with two points. Nevo showed that if $`\sigma `$ is the simplicial complex generated by $`\{1,2\}`$ and $`\{3,4\}`$ then the $`2`$-skeleton of $`\mathrm{\Delta }(\mathrm{\Sigma }(\sigma ))`$ is $`\{\{1,2,3\},\{1,2,4\},\{1,2,5\},\{1,2,6\}\}`$ and that of $`\mathrm{\Delta }(\mathrm{\Sigma }(\mathrm{\Delta }(\sigma )))`$ is $`\{\{1,2,3\},\{1,2,4\},\{1,2,5\},\{1,3,4\}\}`$.
Next, we will recall Nevo’s conjecture. Let $`\sigma `$ and $`\tau `$ be simplicial complexes on $`[n]=\{1,2,\mathrm{},n\}`$ and $`_{rev}`$ the reverse lexicographic order $`_{rev}`$ induced by $`1<2<\mathrm{}<n`$. In other words, for $`S[n]`$ and $`R[n]`$ with $`SR`$, define $`S_{rev}R`$ if (i) $`|S|<|R|`$ or (ii) $`|S|=|R|`$ and the minimal integer in the symmetric difference $`(SR)(RS)`$ belongs to $`S`$. For an integer $`d0`$, we write $`\sigma _d=\{S\sigma :|S|=d+1\}`$. Define $`\sigma _{Rev}\tau `$ if the smallest element w.r.t. $`_{rev}`$ in the symmetric difference between $`\sigma _d`$ and $`\tau _d`$ belongs to $`\sigma _d`$ for all $`d0`$, i.e., $`\mathrm{min}_{_{rev}}\{S:S(\sigma _d\tau _d)(\tau _d\sigma _d)\}\sigma `$ for all $`d0`$.
Nevo conjectured that (\[9, Conjecture 6.1\]), for any simplicial complex $`\sigma `$, one has
$$\mathrm{\Delta }(\mathrm{\Sigma }(\sigma ))_{Rev}\mathrm{\Delta }(\mathrm{\Sigma }(\mathrm{\Delta }(\sigma ))).$$
(In the previous example, the symmetric difference is $`\{\{1,2,6\},\{1,3,4\}\}`$ and $`\{1,2,6\}\mathrm{\Delta }(\mathrm{\Sigma }(\sigma ))`$.) In this paper, we will prove a stronger result. For any subset $`S[n]`$, let
$$m_{_{rev}S}(\sigma )=|\{R\sigma :|R|=|S|\text{and}R_{rev}S\}|.$$
We will prove the following. (The definition of $`\mathrm{\Delta }_\phi (\sigma )`$ will be given in §3.)
###### Theorem 3.1.
Let $`\sigma `$ be a simplicial complex on $`[n]`$ and $`\phi GL_n(K)`$. Then, for any $`S[n]`$, one has
$$m_{_{rev}S}(\mathrm{\Delta }(\sigma ))m_{_{rev}S}(\mathrm{\Delta }(\mathrm{\Delta }_\phi (\sigma ))).$$
By using Theorem 3.1, we can easily prove the next corollary which implies Nevo’s conjecture.
###### Corollary 3.2.
Let $`\sigma `$ be a simplicial complex on $`\{1,2,\mathrm{},k\}`$ and $`\tau `$ a simplicial complex on $`\{k+1,k+2,\mathrm{},n\}`$. Then, for any $`S[n]`$, one has
$$m_{_{rev}S}(\mathrm{\Delta }(\sigma \tau ))m_{_{rev}S}(\mathrm{\Delta }(\mathrm{\Delta }(\sigma )\mathrm{\Delta }(\tau ))).$$
Now, we will explain why Corollary 3.2 implies Nevo’s conjecture. Let $`d0`$ be a positive integer, and let $`\left(\genfrac{}{}{0pt}{}{[n]}{d}\right)`$ and $`\left(\genfrac{}{}{0pt}{}{[n]}{d}\right)`$ be families of $`d`$-subsets of $`[n]`$ which satisfy $`m_{_{rev}S}()m_{_{rev}S}()`$ for all $`S\left(\genfrac{}{}{0pt}{}{[n]}{d}\right)`$. Set $`T=\mathrm{min}_{_{rev}}\{S:S()()\}`$. Then $`m_{_{rev}T}()`$ must be strictly larger than $`m_{_{rev}T}()`$ since $`\{S:S_{rev}T\}=\{S:S_{rev}T\}`$. Thus we have $`T`$. This fact together with Corollary 3.2 implies that $`\mathrm{\Delta }(\sigma \tau )_{Rev}\mathrm{\Delta }(\mathrm{\Delta }(\sigma )\mathrm{\Delta }(\tau )))`$ for all simplicial complexes $`\sigma `$ and $`\tau `$, and Nevo’s conjecture is the special case that $`\tau `$ consists of two points.
To prove Theorem 3.1, we need some techniques of generic initial ideals which have a close connection with algebraic shifting. Let $`K`$ be an infinite field, $`V`$ a $`K`$-vector space with basis $`e_1,e_2,\mathrm{},e_n`$ and $`E=_{d=0}^n^dV`$ the exterior algebra of $`V`$. For a graded ideal $`JE`$, we write $`\mathrm{Gin}_{}(J)`$ for the generic initial ideal of $`J`$ with respect to a term order $``$. For every monomial $`e_S=e_{s_1}e_{s_2}\mathrm{}e_{s_d}E`$ and for every term order $``$, we write
$$m_{e_S}(J)=|\{e_RJ:|R|=|S|\text{and}e_Re_S\}|.$$
We will use the following proposition to prove Theorem 3.1.
###### Proposition 2.4.
Let $`JE`$ be a graded ideal and $``$ and $`^{}`$ term orders. Then, for any monomial $`e_SE`$, one has
$$m_{e_S}(\mathrm{Gin}_{}(J))m_{e_S}(\mathrm{Gin}_{}(\mathrm{in}_{^{}}(J))).$$
Note that the same property as Proposition 2.4 for generic initial ideals over the polynomial ring was proved by Conca .
This paper is organized as follows: In §1, we will give the definition of generic initial ideals and recall some basic properties. In §2, we will prove Proposition 2.4. In §3, we will prove Theorem 3.1 and Corollary 3.2.
## 1. Generic initial ideals in the exterior algebra
Let $`K`$ be an infinite field, $`V`$ a $`K`$-vector space with basis $`e_1,e_2,\mathrm{},e_n`$ and $`E=_{d=0}^n^dV`$ the exterior algebra of $`V`$. For an integer $`d0`$, let $`\left(\genfrac{}{}{0pt}{}{[n]}{d}\right)`$ denote the family of $`d`$-subsets of $`[n]`$. If $`S=\{s_1,s_2,\mathrm{},s_d\}\left(\genfrac{}{}{0pt}{}{[n]}{d}\right)`$ with $`s_1<s_2<\mathrm{}<s_d`$, then the element $`e_S=e_{s_1}e_{s_2}\mathrm{}e_{s_d}`$ will be called a monomial of $`E`$ of degree $`d`$. We refer the reader to for foundations of the Gröbner basis theory over the exterior algebra. Let $``$ be a term order. In this paper, for $`f=_{S[n]}\alpha _Se_SE`$ with each $`\alpha _SK`$, we write $`\mathrm{in}_{}(f)=\mathrm{max}_{}\{e_S:\alpha _S0\}`$.
Let $`GL_n(K)`$ denote the general linear group with coefficients in $`K`$. For $`\phi =(a_{ij})GL_n(K)`$ and for $`f(e_1,\mathrm{},e_n)E`$, we define
$$\phi (f(e_1,\mathrm{},e_n))=f(\underset{i=1}{\overset{n}{}}a_{i1}e_i,\mathrm{},\underset{i=1}{\overset{n}{}}a_{in}e_i).$$
Also, for a graded ideal $`JE`$ and for $`\phi GL_n(K)`$, define $`\phi (J)=\{\phi (f):fJ\}`$. A fundamental theorem of generic initial ideals is the following.
###### Theorem 1.1 (\[1, Theorem 1.6\]).
Fix a term order $``$. Then, for each graded ideal $`JE`$, there exists a nonempty Zariski open subset $`UGL_n(K)`$ such that $`\mathrm{in}_{}(\phi (J))`$ is constant for all $`\phi U`$.
This monomial ideal $`\mathrm{in}_{}(\phi (J))`$ with $`\phi U`$ is called the generic initial ideal of $`J`$ with respect to the term order $``$, and will be denoted $`\mathrm{Gin}_{}(J)`$.
###### Definition 1.2.
Fix a term order $``$. Given an arbitrary graded ideal $`J=_{d=0}^nJ_d`$ of $`E`$ with each $`J_d^dV`$, fix $`\phi GL_n(K)`$ for which $`\mathrm{in}_{}(\phi (J))`$ is the generic initial ideal $`\mathrm{Gin}_{}(J)`$ of $`J`$. Recall that the subspace $`^dV`$ is of dimension $`\left(\genfrac{}{}{0pt}{}{n}{d}\right)`$ with a canonical $`K`$-basis $`\{e_S:S\left(\genfrac{}{}{0pt}{}{[n]}{d}\right)\}`$. Choose an arbitrary $`K`$-basis $`f_1,\mathrm{},f_m`$ of $`J_d`$, where $`m=dim_KJ_d`$. Write each $`\phi (f_i)`$, $`1im`$, of the form
$$\phi (f_i)=\underset{S\left(\genfrac{}{}{0pt}{}{[n]}{d}\right)}{}\alpha _i^Se_S$$
with each $`\alpha _i^SK`$. Let $`M(J,d)`$ denote the $`m\times \left(\genfrac{}{}{0pt}{}{n}{d}\right)`$ matrix
$$M(J,d)=(\alpha _i^S)_{1im,S\left(\genfrac{}{}{0pt}{}{[n]}{d}\right)}$$
whose columns are indexed by $`S\left(\genfrac{}{}{0pt}{}{[n]}{d}\right)`$. For each $`S\left(\genfrac{}{}{0pt}{}{[n]}{d}\right)`$, write $`M_S(J,d)`$ for the submatrix of $`M(J,d)`$ which consists of the columns of $`M(J,d)`$ indexed by those $`R\left(\genfrac{}{}{0pt}{}{[n]}{d}\right)`$ with $`RS`$ and write $`M_S(J,d)`$ for the submatrix of $`M_S(J,d)`$ which is obtained by removing the column of $`M_S(J,d)`$ indexed by $`S`$.
It is not hard to see that we can know generic initial ideals by using the rank of these matrices. Indeed, following properties are known.
###### Lemma 1.3 (\[8, Lemma 2.1\]).
Let $`e_S^dV`$ with $`S\left(\genfrac{}{}{0pt}{}{[n]}{d}\right)`$. Then one has $`e_S(\mathrm{Gin}_{}(J))_d`$ if and only if $`\mathrm{rank}(M_S(J,d))<\mathrm{rank}(M_S(J,d))`$.
###### Lemma 1.4 (\[8, Corollary 2.2\]).
The rank of a matrix $`M_S(J,d)`$, $`S\left(\genfrac{}{}{0pt}{}{[n]}{d}\right)`$, is independent of the choice of $`\phi GL_n(K)`$ for which $`\mathrm{Gin}_{}(J)=\mathrm{in}_{}(\phi (J))`$ and independent of the choice of the $`K`$-basis $`f_1,\mathrm{},f_m`$ of $`J_d`$.
###### Lemma 1.5 (\[8, Corollary 2.3\]).
Let $`JE`$ be a graded ideal and $`\psi GL_n(K)`$. Then one has $`\mathrm{rank}(M_S(J,d))=\mathrm{rank}(M_S(\psi (J),d))`$ for all $`S\left(\genfrac{}{}{0pt}{}{[n]}{d}\right)`$.
Also, the next lemma immediately follows from Lemma 1.3.
###### Lemma 1.6.
Let $`JE`$ be a graded ideal. For every $`S\left(\genfrac{}{}{0pt}{}{[n]}{d}\right)`$, one has
$$m_{e_S}(\mathrm{Gin}_{}(J))=\mathrm{rank}(M_S(J,d)).$$
## 2. proof of proposition 2.4
We will follow the basic technique developed in . (See also \[4, Chapter 15\].)
###### Lemma 2.1 (\[5, Corollary 1.7\]).
For any term order $``$ and for any finite set of monomials $`ME`$, there exist positive integers $`d_1,d_2,\mathrm{},d_n`$ such that for any $`e_S,e_RM`$ with $`|S|=|R|`$, one has $`e_Se_R`$ if and only if $`_{kS}d_k>_{kR}d_k`$.
For every ideal $`JE`$, a subset $`G=\{g_1,\mathrm{},g_m\}J`$ is called a Gröbner basis of $`J`$ with respect to $``$ if $`\{\mathrm{in}_{}(g_1),\mathrm{},\mathrm{in}_{}(g_m)\}`$ generates $`\mathrm{in}_{}(J)`$. A Gröbner basis always exists and is actually a generating set of $`J`$ (\[1, Theorem 1.4\]).
###### Lemma 2.2.
Let $`K[t]`$ be the polynomial ring. Fix a term order $``$. For every graded ideal $`JE`$, there is a subset $`G(t)=\{g_1(t),\mathrm{},g_m(t)\}EK[t]`$ which satisfies the following conditions.
* One has $`g_i(0)=\mathrm{in}_{}(g_i(t_0))`$ for all $`t_0K`$;
* Let $`J(t_0)`$ with $`t_0K`$ be the ideal generated by $`G(t_0)`$. If $`t_00`$, then there exists $`\phi _{t_0}GL_n(K)`$ such that $`\phi _{t_0}(J)=J(t_0)`$;
* For all $`t_0K`$, $`G(t_0)`$ is a Gröbner basis of $`J(t_0)`$ with respect to $``$ and $`\mathrm{in}_{}(J(t_0))=\mathrm{in}_{}(J)`$.
###### Proof.
Let $`G=\{g_1,g_2,\mathrm{},g_m\}`$ be a Gröbner basis of $`J`$ with respect to $``$, where each $`g_j`$ is homogeneous. Let $`ME`$ be the set of monomials. Since $`M`$ is a finite set, Lemma 2.1 says that there exist positive integers $`d_1,d_2,\mathrm{},d_n`$ such that, for any $`e_S,e_RM`$ with $`|S|=|R|`$,
(1) $`e_Se_R\text{ if and only if }{\displaystyle \underset{kS}{}}d_k>{\displaystyle \underset{kR}{}}d_k.`$
Let $`e_{S_i}=\mathrm{in}_{}(g_i)`$. For each $`R[n]`$, write $`d(R)=_{kR}d_k`$. Set
$$g_i(e_1,\mathrm{},e_n)(t)=t^{d(S_i)}g_i(t^{d_1}e_1,\mathrm{},t^{d_n}e_n)$$
and $`G(t)=\{g_1(t),\mathrm{},g_m(t)\}`$. We will show that this set $`G(t)`$ satisfies conditions (i), (ii) and (iii).
First, we will show (i). Each $`g_i(t)`$ can be written in the form
$$g_i(t)=\alpha _{S_i}e_{S_i}+\underset{RS_i}{}\alpha _Rt^{\{d(R)d(S_i)\}}e_R,$$
where $`\alpha _{S_i}K\{0\}`$ and each $`\alpha _RK`$. Then (1) says that $`d(R)d(S_i)>0`$ for all $`R`$ with $`\alpha _R0`$. Thus we have $`g_i(0)=\mathrm{in}_{}(g_i)=\mathrm{in}_{}(g_i(t_0))`$ for all $`t_0K`$ as desired.
Second, for each $`t_0K\{0\}`$, define a matrix $`\phi _{t_0}GL_n(K)`$ by
$$\phi _{t_0}(e_i)=t_{0}^{}{}_{}{}^{d_i}e_i\text{ for }i=1,2,\mathrm{},n.$$
Then the construction of $`g_i(t)`$ says that $`\phi _{t_0}(g_i)=t_0^{d(S_i)}g_i(t_0)`$, and therefore we have $`\phi _{t_0}(J)=J(t_0)`$ for all $`t_0K\{0\}`$. Thus (ii) is satisfied.
Finally, we will show (iii). Since we already proved $`\mathrm{in}_{}(g_i(t_0))=\mathrm{in}_{}(g_i)`$ for all $`t_0K`$, what we must prove is $`\mathrm{in}_{}(J(t_0))=\mathrm{in}_{}(J)`$. The inclusion $`\mathrm{in}_{}(J(t_0))\mathrm{in}_{}(J)`$ follows from $`\mathrm{in}_{}(g_i(t_0))=\mathrm{in}_{}(g_i)`$. Recall that $`J`$ and $`\mathrm{in}_{}(J)`$ have the same Hilbert function, i.e., we have $`dim_K(J_d)=dim_K(\mathrm{in}_{}(J)_d)`$ for all $`d>0`$. Then we have
$$dim_K(\mathrm{in}_{}(J(t_0))_d)=dim_K(J(t_0)_d)=dim_K(J_d)=dim_K(\mathrm{in}_{}(J)_d).$$
Hence we have $`\mathrm{in}_{}(J(t_0))=\mathrm{in}_{}(J)`$. Thus $`G(t_0)`$ is a Gröbner basis of $`J(t_0)`$ for all $`t_0K`$. ∎
###### Lemma 2.3.
Let $`JE`$ be a graded ideal. For every $`t_0K`$, let $`J(t_0)E`$ be the ideal given in Lemma 2.2. Then, for all $`d>0`$, there exists a subset $`G_d(t)EK[t]`$ such that $`G_d(t_0)`$ is a $`K`$-basis of $`J(t_0)_d`$ for all $`t_0K`$.
###### Proof.
Let $`G(t)=\{g_1(t),\mathrm{},g_m(t)\}`$ be a subset of $`EK[t]`$ which is given in Lemma 2.2. Let
$$\stackrel{~}{G}_d(t)=\{e_Sg_i(t):\mathrm{deg}(e_Sg_i(0))=d\text{ and }e_Sg_i(0)0\}.$$
For every $`t_0K`$, since $`G(t_0)`$ is a Gröbner basis of $`J(t_0)`$ and $`g_i(0)=\mathrm{in}_{}(g_i(t_0))`$, the set $`\{\mathrm{in}_{}(h(t_0)):h(t_0)\stackrel{~}{G}_d(t_0)\}`$ spans $`\mathrm{in}_{}(J(t_0))_d=\mathrm{in}_{}(J)_d`$. Also, Lemma 2.2 (i) says that $`h(0)=\mathrm{in}_{}(h(t_0))`$ for all $`t_0K`$. Thus there is a subset $`G_d(t)\stackrel{~}{G}_d(t)`$ such that $`G_d(0)`$ is a $`K`$-basis of $`\mathrm{in}_{}(J)_d`$. On the other hand, for any $`t_0K`$, since each $`h(t_0)G_d(t_0)`$ has a different initial monomial, the set $`G_d(t_0)`$ is linearly independent. Thus we have
$$dim_K(\mathrm{span}\{G_d(t_0)\})=dim_K(\mathrm{span}\{G_d(0)\})=dim_K(\mathrm{in}_{}(J)_d)=dim_K(J(t_0)_d),$$
where $`\mathrm{span}(A)`$ denotes the $`K`$-vector space spanned by a finite set $`AE`$. Hence $`G_d(t_0)`$ is a $`K`$-basis of $`J(t_0)_d`$ for all $`t_0K`$. ∎
###### Proposition 2.4.
Let $`JE`$ be a graded ideal and $``$ and $`^{}`$ term orders. Then, for any monomial $`e_SE`$, one has
$$m_{e_S}(\mathrm{Gin}_{}(J))m_{e_S}(\mathrm{Gin}_{}(\mathrm{in}_{^{}}(J))).$$
###### Proof.
First, by Lemma 1.6, we have
$$m_{e_S}(\mathrm{Gin}_{}(J))=\mathrm{rank}(M_S(J,d))$$
and
$$m_{e_S}(\mathrm{Gin}_{}(\mathrm{in}_{^{}}(J)))=\mathrm{rank}(M_S(\mathrm{in}_{^{}}(J),d)).$$
Thus what we must prove is $`\mathrm{rank}(M_S(J,d))\mathrm{rank}(M_S(\mathrm{in}_{^{}}(J),d))`$.
Let $`m=dim_K(J_d)`$, $`G_d(t)=\{g_1(t),\mathrm{},g_m(t)\}EK[t]`$ the subset given in Lemma 2.3 w.r.t. the term order $`^{}`$ and $`J(t_0)`$, where $`t_0K`$, the ideal given in Lemma 2.2. Then, for each $`t_0K\{0\}`$, there exists $`\phi _{t_0}GL_n(K)`$ such that $`\phi _{t_0}(J)=J(t_0)`$. Thus Lemma 1.5 says that we have
(2) $`\mathrm{rank}(M_S(J,d))=\mathrm{rank}(M_S(J(t_0),d))\text{ for all }t_0K\{0\}.`$
Let $`AK`$ be a finite set with $`0A`$ and $`|A|0`$. Then Theorem 1.1 says that, for each $`aA`$, there exists a nonempty Zariski open subset $`U_aGL_n(K)`$ such that $`\mathrm{Gin}(J(a))=\mathrm{in}_{}(\phi (J(a)))`$ for all $`\phi U_a`$. Since $`U=_{aA}U_a`$ is also a nonempty Zariski open subset of $`GL_n(K)`$, we have $`\mathrm{Gin}_{}(J(a))=\mathrm{in}_{}(\phi (J(a)))`$, for all $`\phi U`$ and all $`aA`$.
Fix $`\phi U`$. Each $`\phi (g_i(t))`$, where $`1im`$, can be written in the form
$$\phi (g_i(t))=\underset{S\left(\genfrac{}{}{0pt}{}{[n]}{d}\right)}{}\alpha _i^S(t)e_S,$$
where $`\alpha _i^S(t)K[t]`$. Define the matrix $`\stackrel{~}{M}_S(J,d,t)=(\alpha _i^R(t))_{1im,RS}`$ by the same way as Definition 1.2. Recall that Lemma 2.3 says that $`G_d(a)`$ is a $`K`$-basis of $`J(a)_d`$ for all $`aA`$. Since Lemma 1.4 says that $`\mathrm{rank}(M_S(J(a),d))`$ is independent of the choice of a $`K`$-basis and independent of the choice of $`\phi U_a`$, it follows that
(3) $`\mathrm{rank}(\stackrel{~}{M}_S(J,d,a))=\mathrm{rank}(M_S(J(a),d))\text{ for all }aA.`$
Let $`l=\mathrm{max}\{\mathrm{deg}(\alpha _i^S(t)):1im,S\left(\genfrac{}{}{0pt}{}{[n]}{d}\right)\}`$. Recall that the rank of matrices is equal to the maximal size of nonzero minors. In this case, each minor of $`\stackrel{~}{M}_S(J,d,t)`$ is a polynomial of $`K[t]`$ and has at most degree $`lm`$. Furthermore, the number of nonzero minors of $`\stackrel{~}{M}_S(J,d,t)`$ is finite. Since integers $`l`$ and $`m`$ do not depend on $`A`$ and $`|A|`$ is sufficiently large, there exists $`a_00A`$ such that $`f(a_0)0`$ for all nonzero minors $`f(t)`$ of $`\stackrel{~}{M}_S(J,d,t)`$. In particular, we have $`\mathrm{rank}(\stackrel{~}{M}_S(J,d,a_0))\mathrm{rank}(\stackrel{~}{M}_S(J,d,t_0))`$ for all $`t_0K`$. Recall that Lemma 2.2 (i) and (iii) say that $`J(0)=\mathrm{in}_{^{}}(J)`$. Then, by (2) and (3), we have
$`\mathrm{rank}(M_S(J,d))`$ $`=`$ $`\mathrm{rank}(M_S(J(a_0),d))`$
$`=`$ $`\mathrm{rank}(\stackrel{~}{M}_S(J,d,a_0))`$
$``$ $`\mathrm{rank}(\stackrel{~}{M}_S(J,d,0))=\mathrm{rank}(M_S(\mathrm{in}_{^{}}(J),d)),`$
as desired. ∎
## 3. exterior shifting of the join of simplicial complex
Let $`\sigma `$ be a simplicial complex on $`[n]`$. The exterior face ideal $`J_\sigma `$ of $`\sigma `$ is the ideal of $`E`$ generated by all monomials $`e_S`$ with $`S\sigma `$. For every $`\phi GL_n(K)`$ and for every simplicial complex $`\sigma `$, the simplicial complex $`\mathrm{\Delta }_\phi (\sigma )`$ is defined by $`J_{\mathrm{\Delta }_\phi (\sigma )}=\mathrm{in}_{_{rev}}(\phi (J_\sigma ))`$. The exterior algebraic shifted complex $`\mathrm{\Delta }(\sigma )`$ of $`\sigma `$ is the simplicial complex defined by $`J_{\mathrm{\Delta }(\sigma )}=\mathrm{Gin}_{_{rev}}(J_\sigma )`$.
###### Theorem 3.1.
Let $`\sigma `$ be a simplicial complex on $`[n]`$ and $`\phi GL_n(K)`$. Then, for any $`S[n]`$, one has
$$m_{_{rev}S}(\mathrm{\Delta }(\sigma ))m_{_{rev}S}(\mathrm{\Delta }(\mathrm{\Delta }_\phi (\sigma ))).$$
###### Proof.
By the definition of exterior face ideals, for any $`d`$-subset $`S\left(\genfrac{}{}{0pt}{}{[n]}{d}\right)`$ and for any simplicial complex $`\tau `$, we have
(4) $`m_{_{rev}S}(\tau )`$ $`=`$ $`|\tau _{d1}||\{R\tau _{d1}:R_{rev}S\}|`$
$`=`$ $`|\tau _{d1}||\{R\left({\displaystyle \genfrac{}{}{0pt}{}{[n]}{d}}\right):R_{rev}S\}|+|\{e_RI_\tau :R_{rev}S\}|.`$
On the other hand, Proposition 2.4 says that, for any $`S[n]`$, one has
(5) $`m_{_{rev}e_S}(\mathrm{Gin}_{_{rev}}(\phi (J_\sigma )))m_{_{rev}e_S}(\mathrm{Gin}_{_{rev}}(\mathrm{in}_{_{rev}}(\phi (J_\sigma )))).`$
Also, Lemma 1.5 says that
(6) $`J_{\mathrm{\Delta }(\sigma )}=\mathrm{Gin}_{_{rev}}(J_\sigma )=\mathrm{Gin}_{_{rev}}(\phi (J_\sigma ))`$
and
(7) $`J_{\mathrm{\Delta }(\mathrm{\Delta }_\phi (\sigma ))}=\mathrm{Gin}_{_{rev}}(\mathrm{in}_{_{rev}}(\phi (J_\sigma ))).`$
Then, equalities (4), (5), (6) and (7) say that
$$m_{_{rev}S}(\mathrm{\Delta }(\sigma ))m_{_{rev}S}(\mathrm{\Delta }(\mathrm{\Delta }_\phi (\sigma ))),$$
as desired. ∎
###### Corollary 3.2.
Let $`\sigma `$ be a simplicial complex on $`\{1,2,\mathrm{},k\}`$ and $`\tau `$ a simplicial complex on $`\{k+1,k+2,\mathrm{},n\}`$. Then, for any $`S[n]`$, one has
$$m_{_{rev}S}(\mathrm{\Delta }(\sigma \tau ))m_{_{rev}S}(\mathrm{\Delta }(\mathrm{\Delta }(\sigma )\mathrm{\Delta }(\tau ))).$$
###### Proof.
Let $`l=|\{k+1,k+2,\mathrm{},n\}|.`$ Then there exist $`\phi \mathrm{GL}_k(K)`$ and $`\psi \mathrm{GL}_l(K)`$ such that $`\mathrm{\Delta }_\phi (\sigma )=\mathrm{\Delta }(\sigma )`$ and $`\mathrm{\Delta }_\psi (\tau )=\mathrm{\Delta }(\tau )`$. For any $`\phi \mathrm{GL}_k(K)`$, we define $`\overline{\phi }GL_n(K)`$ by
$`\overline{\phi }(e_i)=\{\begin{array}{c}\phi (e_i),\text{if}i\{1,2,\mathrm{},k\},\hfill \\ e_i,\text{otherwise}.\hfill \end{array}`$
Also, for any $`\psi \mathrm{GL}_l(K)`$, define $`\overline{\psi }GL_n(K)`$ in the same way. Then we have
$$\mathrm{\Delta }_{\overline{\phi }}(\mathrm{\Delta }_{\overline{\psi }}(\sigma \tau ))=\mathrm{\Delta }_{\overline{\phi }}(\sigma \mathrm{\Delta }(\tau ))=\mathrm{\Delta }(\sigma )\mathrm{\Delta }(\tau ).$$
Then Theorem 3.1 says that
$$m_{_{rev}S}(\mathrm{\Delta }(\sigma \tau ))m_{_{rev}S}(\mathrm{\Delta }(\mathrm{\Delta }_{\overline{\phi }}(\mathrm{\Delta }_{\overline{\psi }}(\sigma \tau ))))=m_{_{rev}S}(\mathrm{\Delta }(\mathrm{\Delta }(\sigma )\mathrm{\Delta }(\tau ))),$$
as desired. ∎
###### \[Remark\].
Proposition 2.4 holds for an arbitrary term order. However, Theorem 3.1 and Corollary 3.2 only hold for the reverse lexicographic order. Recall that Nevo’s example says that
$$\mathrm{\Delta }(\mathrm{\Sigma }(\sigma ))=\{\{1,2,3\},\{1,2,4\},\{1,2,5\},\{1,2,6\}\}$$
and
$$\mathrm{\Delta }(\mathrm{\Sigma }(\mathrm{\Delta }(\sigma )))=\{\{1,2,3\},\{1,2,4\},\{1,2,5\},\{1,3,4\}\},$$
where $`\sigma `$ is the simplicial complex generated by $`\{1,2\}`$ and $`\{3,4\}`$. In this case, $`\{1,2,5\}`$ and $`\{1,2,6\}`$ are larger than $`\{1,3,4\}`$ with respect to the lexicographic order $`_{lex}`$ induced by $`1<2<\mathrm{}<n`$. This implies that $`m_{_{lex}\{1,3,4\}}(\mathrm{\Delta }(\mathrm{\Sigma }(\sigma )))=2`$ but $`m_{_{lex}\{1,3,4\}}(\mathrm{\Delta }(\mathrm{\Sigma }(\mathrm{\Delta }(\sigma ))))=3`$.
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# Branch dependence in the “consistent histories” approach to quantum mechanics
## I Introduction
The consistent histories approach to quantum mechanics griffiths84 ; gellmann90 ; gellmann93 ; omnes94 ; dowker\_kent96 ; kent2000 ; griffiths2003 studies the dynamics of closed quantum systems as a stochastic process within a framework of alternative possible histories. Such a framework, or *family of histories*, must consist of pairwise exclusive and jointly exhaustive descriptions of the system’s dynamics.
This intuitive characterization does not yet state what a family of histories is mathematically. There are two formal definitions in the literature: So-called *product families* are straightforward generalizations of one-time descriptions of a system’s properties in terms of projectors to the case of multiple times griffiths84 . The so-called *history projector operator* formalism, introduced by Isham isham94 , is vastly more general. However, histories in that formalism do not necessarily have an intuitive interpretation in terms of temporal sequences of one-time descriptions.
In our paper we make formally precise the notion of a *branching* family of histories, which is meant to balance generality and intuitiveness: histories in such a family do correspond to temporal sequences of one-time descriptions, yet branching families are much more general than product families. In the context of quantum histories, the notion of branching, or branch dependence, was originally proposed by Gell-Mann and Hartle gellmann90 . It is invoked informally in many publications, but a formal definition is so far lacking. Our definition of *branching families of histories* is based on the theory of branching temporal logic. Apart from providing a precise reading of a useful concept, our approach allows us to comment on the relation between the consistency condition for families of histories and probability measures on such families. It turns out that the consistency condition is best viewed not as a precondition for introducing probabilities, as some authors suggest, but as the requirement that interference effects be absent from the description of a system’s dynamics. Our definition allows for consistent as well as inconsistent families of histories. It is therefore neutral with respect to the discussion about the pros and cons of consistency dowker\_kent96 ; griffiths98 ; kent2000 , and we refrain from taking a stance in that discussion.
Our paper is structured as follows: In section II, we introduce some basic facts about consistent histories and probabilities and review the definition of product families of histories and the history projection operator approach. We also sketch the intuitive motivation for a notion of branch-dependent families. In our central section III, we give our formal definition of branch-dependent families of histories and prove some relevant properties of such families. In the final section IV, we discuss the relation between our new definition and the two mentioned definitions of families of histories, and we comment on the consistency condition.
## II Consistent histories
### II.1 Histories, chain operators, and weights
In the consistent histories approach, a *history* is specified via properties of the system in question at a finite number of times.<sup>1</sup><sup>1</sup>1Continuous extensions of the theory have also been studied isham\_etal98 , but these will not be considered in this paper.—This section closely follows the notational conventions of griffiths2003 , which book provides a detailed and readable introduction to consistent histories. The system’s properties are expressed through orthogonal projectors<sup>2</sup><sup>2</sup>2More generally, one can specify POVMs or completely positive maps; cf. peres2000a . We will only consider projectors in this paper. on (closed) subspaces of the system’s Hilbert space $``$, i.e., operators $`P`$ for which
$$PP=P^{}=P.$$
(1)
Thus, a single history $`Y^\alpha `$ consists of a number of projectors $`P_\alpha ^i`$ at given times $`t_i`$, $`i=1,\mathrm{},n`$:
$$Y^\alpha =P_\alpha ^1P_\alpha ^2\mathrm{}P_\alpha ^n.$$
(2)
So far, the symbol “$``$” should be read as “and then”; in Isham’s *history projection operator* version of the history formalism, the symbol can be read as a tensor product (cf. section II.5).
A *family of histories* $`𝔉`$ (sometimes also called a *framework*) is an exhaustive set of alternative histories. In line with most of the literature on consistent histories, we will only consider finite families in this paper.
Associated with a history $`Y^\alpha `$ is a *chain operator* $`K(Y^\alpha )`$, which is formed by multiplying together the projectors $`P_\alpha ^i`$ associated with the times $`t_i`$ ($`i=1,\mathrm{},n`$), interleaved with the respective unitary time development operators $`T(t_i,t_{i+1})`$. Employing the convention of griffiths2003 , we define
$$K^{}(Y^\alpha )=P_\alpha ^1T(t_1,t_2)P_\alpha ^2\mathrm{}T(t_{n1},t_n)P_\alpha ^n.$$
(3)
These chain operators are often taken to be *representations* of the respective histories. This is appropriate in that $`K(Y^\alpha )`$ correctly describes the successive action of the projectors forming $`Y^\alpha `$ on the systen. However, the representation relation is in general many-one, which may be seen as a disadvantage in that one cannot recover a history from the associated chain operator uniquely.
The system’s dynamics explicitly enters the definition of the chain operators through the time development operators. Thus, a history $`Y^\alpha `$ can have a zero chain operator even though the history involves no zero projectors; such histories are *dynamically impossible*, i.e., ruled out by the system’s dynamics.
We assume that the initial state of the system is described by a density matrix $`\rho `$, which might be proportional to unity if no information is given.<sup>3</sup><sup>3</sup>3Many definitions implicitly take $`\rho =I`$ in the case of lacking information, which will lead to wrong scaling of inner products and thus, of weights, by a factor of dim($``$). The inner product of two operators, $`K_1,K_2_\rho `$, given $`\rho `$, is defined via
$$K_1,K_2_\rho =Tr[\rho K_1^{}K_2].$$
(4)
The chain operators allow us to associate with any history $`Y^\alpha `$ a *weight* $`W(Y^\alpha )`$, which is the inner product of the history’s chain operator with itself:
$$W(Y^\alpha )=K(Y^\alpha ),K(Y^\alpha )_\rho .$$
(5)
In general, one would hope that these weights correspond to probabilities for histories from a given family. We will comment on that issue below, but first we review a few notions from probability theory.
### II.2 Probabilities
A probability space is a triple $`𝔓=S,A,\mu `$, where $`S`$ is the sample space (the set of alternatives), $`A`$ is a Boolean $`\sigma `$-algebra on $`S`$, and $`\mu `$ is a normalized, countably additive measure on $`A`$, i.e., a function
$$\mu :A[0,1]\text{s.t.}\mu (S)=1,$$
(6)
and such that for any countable family $`(a_j)_{jJ}`$ of disjoint elements of $`A`$, $`\mu `$ is additive:
$$\mu \left(\underset{jJ}{}a_j\right)=\underset{jJ}{}\mu (a_j).$$
(7)
For a finite set $`S`$, the algebra $`A`$ is isomorphic to the so-called power set algebra, i.e., the Boolean algebra of subsets of $`S`$, with minimal element $`\mathrm{}`$ and maximal element $`S`$; the operations of join, meet, and complement are set-theoretic union, intersection, and set-theoretic complement, respectively. In the finite case, $`\mu `$ is uniquely specified by its value on the singletons (atoms), and normalization is expressed by the condition
$$\underset{sS}{}\mu (\{s\})=1.$$
(8)
For our treatment of quantum histories, which follows the literature in assuming finite families, these latter, simplified conditions are sufficient: A finite probability space is completely specified by giving a finite set $`S`$ of alternatives and an assignment $`\mu `$ of nonnegative numbers fulfilling (8). The algebra $`A`$ is then given as the power set algebra of $`S`$, and $`\mu `$ is extended to all of $`A`$ via (7).
### II.3 Families of histories: Product families
Intuitively, a family $`𝔉`$ of histories should consist of an exhaustive set of exclusive alternatives. Thus any one history $`h𝔉`$ should rule out all other histories from $`𝔉`$, and $`𝔉`$ should have available enough histories to describe any possible dynamic evolution of the system in question. Exclusiveness must in some way be linked to orthogonality of projectors. This rules out taking $`𝔉`$ to be the set of all time-ordered sets of projectors of the form (2)—such a family would be far too large. The simplest way to ensure the requirements of exclusiveness and exhaustiveness is to fix a sequence of time points
$$t_1<t_2<\mathrm{}<t_n$$
(9)
and to specify, for each time $`t_i`$, a decomposition of the identity operator $`I`$ on the system’s Hilbert space $``$ into $`n_i`$ orthogonal projectors $`\{P_1^i,\mathrm{},P_{n_i}^i\}`$, so that
$$P_j^iP_j^{}^i=\delta _{jj^{}}P_j^i,\underset{j=1}{\overset{n_i}{}}P_j^i=I.$$
(10)
In this case, the index $`\alpha `$ specifying a history $`Y^\alpha `$ can be taken to be the list of numbers $`(\alpha _1,\mathrm{},\alpha _n)`$, where $`1\alpha _in_i`$. The size $`|𝔉|`$ of the family $`𝔉`$ is given by
$$|𝔉|=n_1\times n_2\times \mathrm{}\times n_n.$$
(11)
Histories in $`𝔉`$ are pairwise exclusive, since any two different histories use different, orthogonal projectors at some time $`t_i`$. Such a family is also exhaustive, as at every time, the decomposition of the identity specifies an exhaustive set of alternatives. Formally, $`𝔉`$ corresponds to a cartesian product of decompositions of the identity at different times. Product families are thus the obvious generalization of one-time descriptions of a system’s properties in terms of projectors to the case of multiple times.
### II.4 Branch-dependent families of histories
While the construction of a product family of histories is the simplest way to ensure exclusiveness and exhaustiveness, that construction is by no means the only possibility. Many authors have noted that in applications, it will often be necessary to choose a time point $`t_{i+1}`$, or the decomposition of $`I`$ at $`t_{i+1}`$, dependent on the projector $`P_\alpha ^i`$ employed at time $`t_i`$. Thus, e.g., in order to describe a “delayed choice” quantum correlation experiment aspect82 , one chooses the direction of spin projection at time $`t_2`$ depending on the outcome of a previous selection event at time $`t_1`$.
It is intuitively quite clear what such “branch dependence” would mean; eqs. (22)–(25) below give an example of a branch-dependent family of histories. However, no formally rigorous description of such families of histories is available so far. Before we go on to give such a description in section III, we introduce Isham’s isham94 generalized definition of families of histories in terms of history projection operators (HPOs), with which our definition of branch dependence will be compared below.
### II.5 Isham’s history projection operators (HPO)
The guiding idea of the history projection operator framework is to single out, as for product families, $`n`$ times $`t_1,\mathrm{},t_n`$ at which the system’s properties will be described. One then forms the $`n`$-fold tensor product of the system’s Hilbert space:
$$\stackrel{~}{}=\mathrm{}\text{(}n\text{ times)}.$$
(12)
In that large history Hilbert space $`\stackrel{~}{}`$, a history $`Y^\alpha `$ is read as a tensor product of projectors:
$$Y^\alpha =P_\alpha ^1P_\alpha ^2\mathrm{}P_\alpha ^n=P_\alpha ^1P_\alpha ^2\mathrm{}P_\alpha ^n.$$
(13)
That tensor product operator $`Y^\alpha `$ is itself a projection operator on $`\stackrel{~}{}`$, a so-called *history projection operator*, fulfilling
$$Y^\alpha Y^\alpha =(Y^\alpha )^{}=Y^\alpha .$$
(14)
Along these lines one can give an abstract definition of a family of histories $`𝔉=\{Y^1,\mathrm{},Y^n\}`$ as a decomposition of the history Hilbert space identity $`\stackrel{~}{I}`$:
$$Y^\alpha Y^\beta =\delta _{\alpha \beta }Y^\alpha ,\underset{Y^\alpha 𝔉}{}Y^\alpha =\stackrel{~}{I}.$$
(15)
This generalization is formally rigorous, and it allows for further (e.g., continuous) extensions. However, the generality comes at a certain cost, since there is no condition that would ensure that a history projector $`Y^\alpha `$ should factor into a product of $`n`$ projectors on the system’s Hilbert space at the $`n`$ given times, as in (2). History projectors that do factor in this way are called *homogeneous histories*. As the main motivation for introducing histories is given in terms of homogeneous histories, some authors have expressed doubts as to whether the full generality of HPO is really appropriate (griffiths2003, , p. 118).
One possibility for constructing a narrower framework is to restrict the HPO formalism to the homogeneous case by requiring that all the $`Y^\alpha 𝔉`$ be products of projectors of the form (2). Such a restriction may be implicitly at work in griffiths2003 . However, such a restriction is to some extent alien to the HPO formalism. Nor does it single out a useful class of families of histories: as we will show below (eqs. (53)–(56)), not all homogeneous families are branch-dependent families, and some homogeneous families do not admit an interpretation of weights in terms of probabilities.
We now describe an alternative approach in which branch dependence is the natural result of an inductive definition. Furthermore, the framework allows one to distinguish two different types of coarse–graining for families of histories.
## III A formal framework for branch-dependent histories
The idea of a branching family of histories is based on the theory of branching temporal logic that originated in the work of Prior prior67 .<sup>4</sup><sup>4</sup>4Branching temporal logic has already found many applications in computer science, linguistics, and philosophy. The interested reader is referred to belnap2001 for an overview.
### III.1 Branching structures
For the purpose of constructing finite families of branching histories, branching temporal logic boils down to the following inductive definition of a *branching structure*, which is a set $`M`$ of *moments* $`m_i`$ partially ordered by $``$:<sup>5</sup><sup>5</sup>5One should think of $``$ in analogy to “less than or equal” ($``$), so there is a companion strict order (excluding equality), which is denoted by $``$. However, note that $``$ is only a partial order, i.e., some elements of $`M`$ may be incomparable.
###### Definition 1 (Branching structure)
(i) A singleton set $`M=\{m_0\}`$ together with the relation $`m_0m_0`$ is a branching structure. (ii) Let $`M,`$ be a branching structure and $`mM`$ a maximal element, and let $`m_1^{},\mathrm{},m_n^{}`$ be new elements. Let $`^{}`$ be the reflexive and transitive closure of the relation $``$ together with the new relations $`m^{}m_1^{}`$, …, $`m^{}m_n^{}`$. Then the set $`M\{m_1^{},\mathrm{},m_n^{}\}`$ together with the relation $`^{}`$ is again a branching structure.
By taking a finite number of steps along this definition, one constructs a finite branching tree in the form of a partially ordered set with the unique root element $`m_0`$.<sup>6</sup><sup>6</sup>6Note that the construction ensures the following formal features: The ordering $``$ is transitive (if $`xy`$ and $`yz`$, then $`xz`$), reflexive ($`xx`$) and antisymmetric (if $`xy`$ and $`yx`$, then $`x=y`$). Furthermore, the ordering fulfills the axioms of “no backward branching” (if $`xz`$ and $`yz`$, then either $`xy`$ or $`yx`$) and of “historical connection”, meaning the existence of a common lower bound for any two elements (for any $`x,yM`$, there is $`zM`$ s.t. $`zx`$ and $`zy`$).—In a more general approach to branching temporal logic, these formal features are taken as axioms of a logical framework. The maximal nodes in the tree are called “leaves”. Figure 1 illustrates the inductive process. Except for the root element $`m_0`$, each node has a unique direct predecessor, and except for the leaves, each node $`m`$ has one or more direct successors, which correspond to branching at $`m`$. Paths in the tree, i.e., maximal linearly ordered subsets, extend from the root to one of the leaves and are thus in one-to-one correspondence with the leaves. These paths are often called *histories* by logicians, and they can indeed be given an interpretation in terms of quantum histories, as we will show in the next section.
### III.2 Branching families of histories
A branching family of histories can be viewed as a quantum-mechanical interpretation of a branching structure $`M,`$. We assume that a system with Hilbert space $``$ (identity operator $`I`$) is given. The interpretation is given by two functions $`\tau `$ and $`P`$ that associate times and projectors with the elements of $`M`$, respectively. Formally, we define:
###### Definition 2 (Branching family of histories)
A *branching family of histories* is a quadruple
$$𝔉=M,,\tau ,P,$$
(16)
where $`M,`$ is a finite branching structure and $`\tau `$ is a function from $`M`$ to the real numbers respecting the partial ordering $``$:
$$\text{if }mm^{}\text{, then }\tau (m)<\tau (m^{}).$$
(17)
$`P`$ is a function from $`M`$ to the set of projectors on $``$ that assigns projection operators to the elements of $`M`$ in the following way: If $`mM`$ is not a maximal element, and $`m_1,\mathrm{},m_{n_m}`$ are the $`n_m`$ immediate successors of $`m`$, then a set of orthogonal projectors $`P_1^m,\mathrm{},P_{n_m}^m`$ forming a decomposition of the identity,
$$P_i^mP_j^m=\delta _{ij}P_i^m,\underset{i=1}{\overset{n_m}{}}P_i^m=I,$$
(18)
is assigned to the $`m_1,\mathrm{},m_{n_m}`$ via $`P(m_i)=P_i^m`$.
The number $`\tau (m)`$ is the time associated with $`mM`$. While the same time may be assigned to moments in different histories, we require that $`\tau `$ respect the partial ordering, as expressed via (17).<sup>7</sup><sup>7</sup>7The time parameter will only be needed in specifying the time development operators in the definition of the chain operators later on. Working in the Heisenberg picture, the function $`\tau `$ would not be needed; the condition on $`\tau `$ would be replaced by requiring an appropriate temporal ordering of the Heisenberg operators specified through $`P`$. This already points towards a relativistic generalization of the current approach to quantum histories, which is currently under preparation muller\_qmbst . As regards the assignment of projection operators, note that $`P(m_i)=P_i^m`$ means that at $`m`$ (the unique predecessor of $`m_i`$, not $`m_i`$ itself), the system had the property expressed by $`P_i^m`$.<sup>8</sup><sup>8</sup>8The slightly awkward reference to the previous node in (18) can be avoided if one assigns projectors more properly not to nodes, but to *elementary transitions* in the branching structure; cf. xu97 ; belnap2005 for details. We stick to our simplified exposition in order not to clutter this paper with technicalities—which will, however, be relevant for an extension to infinite structures, or to a relativistic version employing *branching space-times* belnap92 ; muller2005 ; muller\_qmbst . The function $`P`$ thus associates decompositions of the Hilbert space identity with instances of branching. In this way, each maximal path $`\alpha `$ of length $`n_\alpha +1`$ in $`M,`$,
$$m_0^\alpha m_1^\alpha \mathrm{}m_{n_\alpha }^\alpha ,$$
(19)
corresponds to the $`n_\alpha `$ elements long chain of projection operators
$$P(m_1^\alpha )P(m_2^\alpha )\mathrm{}P(m_{n_\alpha }^\alpha ).$$
(20)
Here, $`m_0^\alpha =m_0`$ is the root node, and $`m_{n_\alpha }^\alpha `$ is one of the maximal elements. The projectors give information for the $`n_\alpha `$ many times
$$\tau (m_0^\alpha )<\tau (m_1^\alpha )<\mathrm{}<\tau (m_{n_\alpha 1}^\alpha ).$$
(21)
Our Definition 2 captures the intuitive notion of branch dependence described in section II.4 in a formally exact manner, as illustrated by Figure 2. In that example of a branching family of histories, the root node $`m_0`$ has two successors, $`m_1`$ and $`m_2`$, corresponding to the system’s having the properties corresponding to projection operators $`P_0^1`$ and $`P_0^2`$ at $`\tau (m_0)`$, respectively. The vertical position of $`m_1`$ vs. $`m_2`$ indicates that $`\tau (m_1)\tau (m_2)`$, which is one aspect of branch dependence that is not available in product families: the time for which the system’s property is described after $`\tau (m_0)`$ depends on the system’s property at $`\tau (m_0)`$. Furthermore, the decompositions of the Hilbert space identity at $`m_1`$ ($`\{P_1^1,P_1^2,P_1^3\}`$) and at $`m_2`$ ($`\{P_2^1,P_2^2\}`$) are different, thus exhibiting the second form of branch dependence that is not available in product families.
### III.3 Properties of quantum branching histories
We first note that every product family of histories (cf. section II.3) is a branching family of a very symmetric kind.
###### Lemma 1 (Branching vs. product families)
Every product family of histories is a branching family of histories, but not conversely.
*Proof:* In order to see that product families are branching families, let a product family corresponding to $`n`$ decompositions of the identity $`\{P_1^i,\mathrm{},P_{n_i}^i\}`$ at times $`t_1,\mathrm{},t_n`$ be given. An equivalent branching family is constructed in $`n+1`$ steps as follows: We start with a single node, $`M_0=\{m_0\}`$ (stage $`0`$). Then at each stage $`i`$ ($`1in`$), we add new nodes and enlarge the structure. We assign the time $`t_i`$ to all of the maximal elements of $`M_{i1}`$, and we add the decomposition $`\{P_1^i,\mathrm{},P_{n_i}^i\}`$ above all these maxima by introducing $`n_i`$ new elements above *each* maximum, thus arriving at the new set of nodes $`M_i`$. When $`M_n`$ has been constructed, we finally assign some time $`t^{}>t_n`$ to all the maximal elements in $`M_n`$. This construction yields a symmetrically growing tree that in the end (at stage $`n`$) corresponds to the original product family of histories.
For a branching family that is not a product family, let $``$ have dimension 2, and let $`\{|\varphi _1,|\varphi _2\}`$ and $`\{|\psi _1,|\psi _2\}`$ be two different orthonormal bases of $``$. The family
$`h_1`$ $`=`$ $`|\varphi _1\varphi _1||\varphi _1\varphi _1|`$ (22)
$`h_2`$ $`=`$ $`|\varphi _1\varphi _1||\varphi _2\varphi _2|`$ (23)
$`h_3`$ $`=`$ $`|\varphi _2\varphi _2||\psi _1\psi _1|`$ (24)
$`h_4`$ $`=`$ $`|\varphi _2\varphi _2||\psi _2\psi _2|`$ (25)
describes the system at two times $`t_1`$ and $`t_2`$, yielding a branching family, but not a product family. $`\mathrm{}`$
Branching families of histories $`𝔉=M,,\tau ,P`$ thus yield a natural generalization of product histories while retaining a strong link between the formalism and the intended temporal interpretation of the histories. Furthermore, the inductive definition of the structure makes it easy to prove two key properties of branching families of histories. Firstly, the construction immediately shows that $`𝔉`$ is exclusive and exhaustive: The one-element family has that property, and it is retained by adding inductively further (exclusive and exhaustive) decompositions of the identity at a maximal node.<sup>9</sup><sup>9</sup>9“Exhaustive” does not mean “maximally detailed”. The latter notion is quite dubious anyway, as for every finite description of a quantum system’s dynamics one can give a more detailed one. Secondly, one can show that the weights of histories in a branching family $`𝔉`$ always add up to one. Thus, the weights immediately induce probabilities on the power set Boolean algebra of $`𝔉`$ (cf. section II.2).
###### Lemma 2 (Weights in branching families)
In a branching family of histories $`𝔉`$, the weights sum to one:
$$\underset{h𝔉}{}W(h)=1.$$
(26)
*Proof:* Assume that an initial state of the system is given by a density matrix $`\rho _0`$.<sup>10</sup><sup>10</sup>10In the case of complete ignorance, $`\rho _0=I/\text{dim}()`$. For the trivial family consisting of only one history $`\{m_0\}`$, the sum of weights reduces to
$$W(\{m_0\})=Tr[\rho _0]=1.$$
(27)
Now assume that the property in question holds for the family $`𝔉`$ corresponding to $`M,,\tau ,P`$, let $`m`$ be a maximal node, $`m^{}`$ its (unique) direct predecessor, and let
$$P_i^mP_j^m=\delta _{ij}P_i^m,\underset{i=1}{\overset{n}{}}P_i^m=I,$$
(28)
be the decomposition of the identity that is to be employed at the $`n`$ new maximal elements $`m_1^{},\mathrm{},m_n^{}`$ that are to be added after $`m`$. Suppose that $`𝔉`$ had $`N`$ elements $`h_1,\mathrm{},h_N`$, whose weights add to one. In order to facilitate book-keeping, suppose further that $`m`$ is the final node of $`h_N`$. To the new quantum branching structure $`M^{},^{},\tau ^{},P^{}`$ there corresponds a family $`𝔉^{}`$ of $`N+n1`$ histories, where for $`i=1,\mathrm{},N1`$, $`h_i^{}=h_i`$, whereas $`h_N`$ is replaced by the $`n`$ new histories $`h_N^{},\mathrm{},h_{N+n1}^{}`$ ending in the new elements $`m_1^{},\mathrm{},m_n^{}`$. In order to show that in $`𝔉^{}`$, the weights still add to one, we only need to show that
$$W(h_N)=\underset{i=1}{\overset{n}{}}W(h_{N+i1}^{}),$$
(29)
i.e., the histories replacing old $`h_N`$ must together have the same weight as $`h_N`$. Now in terms of the chain operator $`K(h_N)`$ for $`h_N`$, the chain operators for the new histories are
$$K(h_{N+i1}^{})=P_i^mT(\tau (m),\tau (m^{}))K(h_N).$$
(30)
The initial density matrix $`\rho _0`$, evolved along $`h_N`$, becomes
$$\rho _m=K(h_N)\rho _0K^{}(h_N),$$
(31)
and the weight of $`h_N`$ can be expressed as
$$W(h_N)=Tr[K(h_N)\rho _0K^{}(h_N)]=Tr[\rho _m].$$
(32)
The weights for the new histories can then be written as
$`W(h_{N+i1}^{})=`$
$`Tr[P_i^mT(\tau (m),\tau (m^{}))\rho _mT(\tau (m^{}),\tau (m))(P_i^m)^{}]=`$
$`Tr[T(\tau (m^{}),\tau (m))P_i^mT(\tau (m),\tau (m^{}))\rho _m],`$
where we used the cyclic property of the trace and $`P^{}=P^2=P`$. Now as $`T`$ is unitary, the $`n`$ operators
$$\stackrel{~}{P}_i^m=T(\tau (m^{}),\tau (m))P_i^mT(\tau (m),\tau (m^{}))$$
(34)
are again projectors forming a decomposition of the identity, so that by the linearity of the trace,
$`{\displaystyle \underset{i=1}{\overset{n}{}}}W(h_{N+i1}^{})`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{n}{}}}Tr[\stackrel{~}{P}_i^m\rho _m]`$ (35)
$`=`$ $`Tr\left[{\displaystyle \underset{i=1}{\overset{n}{}}}\stackrel{~}{P}_i^m\rho _m\right]`$ (36)
$`=`$ $`Tr[I\rho _m]=W(h_N),`$ (37)
which was to be proved. $`\mathrm{}`$
## IV Coarse graining, probabilities and the consistency condition
The general idea of coarse graining is that it should be possible to move from a more to a less detailed description of a given system in a coherent way. If probabilities are attached to a fine-grained description, then an obvious requirement is that the probability of a coarse-grained alternative should be the sum of the probabilities of the corresponding fine-grained alternatives. Considerations of coarse graining are important for quantum histories because of the interplay between weights of histories and probability measures in a family of histories. Our discussion will show that one needs to distinguish two notions of coarse graining.
One notion of coarse graining comes for free in any probability space: Due to the additivity of the measure, if $`b^{}`$ is the disjoint union of $`b_1,\mathrm{},b_n`$ in the event algebra, then $`\mu (b^{})=_{i=1}^n\mu (b_i)`$.<sup>11</sup><sup>11</sup>11In the infinite case, that equation holds for countable unions. By Lemma 2, for any branching family (and thus, by Lemma 1, for any product family) of quantum histories, the weights $`W(h)`$ of the histories $`h𝔉`$ induce a probability measure on $`𝔉`$ via $`\mu (\{h\})=W(h)`$, which is extended to the power set Boolean algebra of $`𝔉`$ via (7):
$$\mu (\{h_1,\mathrm{},h_n\})=\underset{i=1}{\overset{n}{}}\mu (\{h_i\})=\underset{i=1}{\overset{n}{}}W(h_i).$$
(38)
If $`h_1,\mathrm{},h_n`$ are fine-grained descriptions of a system’s dynamics, eq. (38) shows that the coarse-grained description $`\{h_1,\mathrm{},h_n\}`$ automatically is assigned the correct probability. Thus, branching families of histories unconditionally and naturally support this notion of coarse graining.
If all branching families support probabilities and coarse graining, then what is behind the consistency condition? In the literature it is often suggested that the possibility of defining probabilities or the possibility of coarse graining for a family of histories is conditional upon the so-called *consistency condition*,
$$K(Y^\alpha ),K(Y^\beta )_\rho =0\text{if }\alpha \beta ,$$
(39)
which demands that the chain operators of different histories must be orthogonal. Condition (39) is also called “medium decoherence” gellmann93 .
The above considerations show that for branching families, both probabilities and one notion of coarse graining are unproblematic, independent of any condition like (39). However, that condition does play an important role with respect to a second, different notion of coarse graining.
That second notion of coarse graining is based on the idea of constructing from the histories $`h𝔉`$ not *sets of histories*, as in eq. (38), but *new histories*, by something like addition. Whether such additive combination of histories is possible at all, generally depends on what histories are mathematically. We will see below that additive combination is not always possible for histories in a branching family. Accordingly, in order not to suggest that addition of histories is always unproblematic, we will use the formal notation “$`sum(h_1,h_2)`$” when we wish to leave open the question whether that sum is in fact defined.
If we consider the additive combination of two histories, $`h_1`$ and $`h_2`$,
$$h=sum(h_1,h_2),$$
(40)
the idea behind coarse graining suggests that for $`h`$, which is a less detailed description than the two fine-grained histories, the probabilities should just add:
$$\mu (h)=\mu (sum(h_1,h_2))=\mu (h_1)+\mu (h_2).$$
(41)
Even apart from the question of whether $`sum(h_1,h_2)`$ can be defined, eq. (41) is problematic as it stands: No family of histories can contain both two histories $`h_1`$ and $`h_2`$ and their sum $`h`$, as that would violate the requirement of exclusiveness. Thus, $`\mu `$ in (41) cannot be a probability measure in a single family of histories. At this point the idea of weights $`W(h)`$ as probabilities enters. Assuming that $`h=sum(h_1,h_2)`$ is indeed a history, $`W`$ is defined for all three of $`h_1`$, $`h_2`$, and $`h`$, and the main idea of (41) can be reformulated as
$$W(h)=W(sum(h_1,h_2))=W(h_1)+W(h_2).$$
(42)
The validity of (42) is indeed linked to the consistency condition (39). However, depending on which type of family of histories one considers, there are some subtle issues, as the following sections point out.
### IV.1 Coarse graining in product families
If $`𝔉`$ is a product family of histories, the sketched idea of coarse graining makes immediate sense, as the sum of any two histories in a given product family can be defined. To consider the basic case, let two histories $`h_1,h_2𝔉`$ be given,
$$h_\alpha =P_\alpha ^1\mathrm{}P_\alpha ^n,\alpha =1,2,$$
(43)
such that they coincide everywhere except for the $`j`$-th position: $`P_1^i=P_2^i`$ for $`ij`$, $`P_1^jP_2^j`$. In this case, one can define their sum
$$h=sum(h_1,h_2):=P_1^1\mathrm{}P_1^{j1}(P_1^j+P_2^j)P_1^{j+1}\mathrm{}P_1^n,$$
(44)
i.e., at each time $`t_i`$ for which the histories $`h_1`$ and $`h_2`$ are defined, their sum, $`h`$, specifies either the same projector as each of $`h_1`$ and $`h_2`$, or gives a less detailed description in terms of the projector $`P_1^j+P_2^j`$. The assumption of a product family is crucial in this definition, as it guarantees that the $`P_1^i`$ and $`P_2^i`$ are both defined at the same times, and that $`P_1^j`$ and $`P_2^j`$ commute—for a branch-dependent family, $`P_1^j+P_2^j`$, even if defined, wouldn’t normally be a projector.
For a history like $`h`$ in (44), the weight function $`W`$ is naturally defined even though $`h𝔉`$, and it appears natural to demand that
$$W(h)=W(sum(h_1,h_2))=W(h_1)+W(h_2).$$
(45)
This equation does not hold in general in product families of histories. A family of histories $`𝔉`$ must satisfy the above-mentioned condition of *consistency* if it is to satisfy (45) for all $`h_1,h_2𝔉`$, and it was along these lines that Griffiths griffiths84 originally motivated the consistency condition for product families of histories.<sup>12</sup><sup>12</sup>12An extended discussion of questions of uniqueness conditions for probability assignments in product families is given in nistico99 .—The notion of consistency in griffiths84 , which corresponds to (45), is weaker than the consistency condition (39) formulated above: For (45) to hold, it is sufficient that the *real part* of $`K(Y^\alpha ),K(Y^\beta )_\rho `$ vanish for $`\alpha \beta `$. The latter condition is known as *weak consistency*. In what follows, we will not differentiate between medium and weak consistency. However, as the next section shows, the symmetric nature of product families hides an important asymmetry in adding histories.
### IV.2 Coarse graining in branching families
We have already seen that in product families $`𝔉`$, the formal addition $`sum(h_1,h_2)`$ can be defined for any $`h_1,h_2𝔉`$. For branching families, this is not always possible. In fact, we will see that with respect to formula (45) one should distinguish two types of coarse graining, which we call *intra-branch* and *trans-branch* coarse graining. *Intra-branch* coarse graining means that maximal nodes from an otherwise shared branch are added, whereas *trans-branch* coarse graining means adding “across branches”. The notion of intra-branch coarse graning and the respective summation of histories for the basic case can be defined as follows:
###### Definition 3 (Intra-branch coarse graining)
In a branching family $`𝔉`$, the formal summation $`h=sum(h_1,h_2)`$ of two histories $`h_1,h_2𝔉`$,
$$h_\alpha =P_\alpha ^1\mathrm{}P_\alpha ^{n_\alpha },\alpha =1,2,$$
(46)
is called *intra-branch coarse graining* iff $`n_1=n_2`$, the histories are defined at the same times, and for $`1i<n_1`$, $`P_1^i=P_2^i`$. In that case, the sum is defined to be
$$h=sum(h_1,h_2):=P_1^1\mathrm{}P_1^{n_11}(P_1^{n_1}+P_2^{n_1}).$$
(47)
Intra-branch coarse graining is both well-defined and probabilistically unproblematic for all branching families:
###### Lemma 3
For intra-branch coarse graining $`h=sum(h_1,h_2)`$ as in (47) in a branching family of histories, the weights add according to (45), i.e.,
$$W(h)=W(sum(h_1,h_2))=W(h_1)+W(h_2).$$
*Proof:* We can follow the lines of the inductive proof of Lemma 2: Let two histories $`h_1`$ and $`h_2`$ fulfilling Definition 3 be given, and assume that $`m`$ is the immediate predecessor of maximal nodes $`m_1^{}`$ and $`m_2^{}`$ of histories $`h_1`$ and $`h_2`$, respectively. Let $`m^{}`$ be the unique direct predecessor of $`m`$. Let $`K(h_m)`$ be the chain operator for the path from the root node to $`m`$, and set $`T=T(\tau (m),\tau (m^{}))`$. Then the chain operators for the histories $`h_\alpha `$ are:
$$K(h_\alpha )=P(m_\alpha ^{})TK(h_m).$$
(48)
The sum of the two final projectors, $`P(m_1^{})+P(m_2^{})`$, is again a projector in virtue of (18). Accordingly, the weight of the coarse-grained history $`h=h_1+h_2`$ is
$`W(h)`$ $`=`$ $`K(h),K(h)_\rho =K(h_1+h_2),K(h_1+h_2)_\rho `$
$`=`$ $`Tr[(P(m_1^{})+P(m_2^{}))TK(h_m)\rho `$
$`K^{}(h_m)T^{}(P(m_1^{})+P(m_2^{}))^{}]`$
$`=`$ $`Tr[T^{}P(m_1^{})TK(h_m)\rho K^{}(h_m)]+`$
$`Tr[T^{}P(m_2^{})TK(h_m)\rho K^{}(h_m)]`$
$`=`$ $`W(h_1)+W(h_2),`$
where we employed $`P^{}(m_\alpha ^{})=P(m_\alpha ^{})=(P(m_\alpha ^{}))^2`$, the fact that $`P(m_1^{})P(m_2^{})=0`$ (18), and linearity and the cyclic property of the trace. $`\mathrm{}`$
So all is well probabilistically if histories are formed as sums of histories that differ only at the last node, i.e., via intra-branch coarse graining. Note that this result carries over to product families of histories: for intra-branch coarse graining in branching families and in product families, eq. (45) holds automatically, without having to presuppose a consistency condition like (39).
What about trans-branch coarse graining? For a product family, eq. (44) shows how to build histories from other histories quite generally, and we have mentioned the fact that the validity of eq. (45) for trans-branch coarse graining in a product family generally depends on a consistency condition like (39) griffiths84 ; griffiths2003 . For the more general case of branching families, so far the summation of histories for trans-branch coarse graining has not been defined. It is possible to define that type of summation in the extended framework of Isham’s HPO formalism, but this amounts to discarding the intuitive idea that histories are temporal sequences of one-time descriptions (cf. the next subsection).
The problem of defining trans-branch coarse graining while holding on to the intuitive interpretation of a history may be illustrated by considering a two-dimensional Hilbert space and the two histories $`h_1`$ and $`h_3`$ from eq. (22) and eq. (24), respectively, taken to be defined at the two times $`t_1`$ and $`t_2`$. How should one define $`sum(h_1,h_3)`$? Surely one can coarse-grain by considering the *set* of histories $`\{h_1,h_3\}`$, and the sum of $`h_1`$ and $`h_3`$ in Isham’s formalism amounts to this exactly. But no temporal interpretation in terms of a single history is forthcoming, as the projectors involved at $`t_2`$ do not commute.
Thus, for trans-branch coarse graining, the formal summation $`sum(h_1,h_2)`$ generally must remain undefined. With respect to the coarse-graining criterion (45), this means that in all cases in which it generally makes sense to ask whether it is satisfied, it is satisfied automatically: Eq. (45) is generally defined only for intra-branch coarse graining, and for that case, Lemma 3 has shown the equation to hold unconditionally.
This result does of course not mean that all branching families of histories are consistent. The test of eq. (39) can still be applied for any branching family, and many branching families will be classified as inconsistent. However, in these cases, the link with eq. (45), which holds for product families, can no longer be made in our branching histories framework. This points to a somewhat different interpretation of the consistency condition: In a branching family, that condition should not be read as a precondition for the assignment of probabilties via weights (which is unproblematic in view of Lemma 2), but rather as the condition that the different descriptions of the system’s dynamics given by the histories in the family amount to wholly separate, interference-free alternatives. This is just what it means for he chain operators to be orthogonal, i.e., to satisfy eq. (39). We hold it to be an advantage of the branching family formalism proposed here that by leaving trans-branch coarse graining undefined, it forces us to rethink the interpretation of the consistency condition. The formalism of product histories is deceptively smooth in treating all times in the same way, thus blurring the distinction between intra-branch and trans-branch coarse graining.
### IV.3 Coarse graining in Isham’s HPO
The comparison between branching families and Isham’s HPO scheme is illuminating in another respect. HPO is more general and more abstract than branching families. However, that abstractness comes at a price: we will argue that much of the intuitiveness of branching families is lost by moving to the HPO scheme.
In Isham’s HPO scheme, histories are themselves represented as projectors on the (large) history Hilbert space, not as chain operators on the (much smaller) system Hilbert space. As an HPO-family $`𝔉`$ must correspond to a decomposition of the history identity operator, sums of histories in $`𝔉`$ will again correspond to history projectors (even though not from the given family); the formal addition $`h=sum(h_1,h_2)`$ has a direct interpretation as the literal addition of history projectors. Thus one can form a Boolean algebra with elements
$$Y=\underset{Y^\alpha 𝔉}{}\pi _\alpha Y^\alpha ,\pi _\alpha \{0,1\},$$
(50)
that is isomorphic to the power set algebra of $`𝔉`$ (the $`\pi _\alpha `$ playing the role of characteristic functions)—just like in the first type of coarse graining considered at the beginning of this section. From one point of view, addition here always forms like objects from like objects: sums of history projectors are again history projectors. From another point of view, however, addition remains problematic: Even if the $`Y^\alpha `$ are homogeneous histories, i.e., have an intuitive interpretation as temporally ordered sequences of projectors, that will generally not be so for their sums. Thus in general, the elements (50) of the Boolean algebra are inhomogeneous histories without an intuitive interpretation.
At the level of abstraction of eq. (50), one need not distinguish between different kinds of coarse graining. Accordingly, it seems more natural to demand additivity of weights,
$$W\left(\underset{Y^\alpha 𝔉}{}\pi _\alpha Y^\alpha \right)=\underset{Y^\alpha 𝔉}{}\pi _\alpha W(Y^\alpha ),$$
(51)
which corresponds to satisfaction of the consistency condition (39). However, at the same level of abstraction, one can note that $`W`$ is a quadratic function, whereas (51) demands linearity—not a natural demand at all. The motivation for additivity was, after all, given in terms of the time-ordered sequences of projectors that formed the basis of the history framework, not in terms of some abstract algebra. Furthermore, even when HPO is restricted to families of homogeneous histories, a probability interpretation of the weights is not forthcoming generally; families of HPO histories can violate a number of seemingly straightforward assumptions isham\_linden94 . As an example, consider the following family of histories: Let $``$ have dimension 2, and let $`\{|\varphi ,|\psi \}`$ be an orthonormal basis, so that $`\varphi |\psi =0`$. Then, define
$$|\chi =\frac{1}{\sqrt{2}}(|\varphi +|\psi );|\chi ^{}=\frac{1}{\sqrt{2}}(|\varphi |\psi ).$$
(52)
Note that $`\chi |\chi ^{}=0`$. We now construct a family of two-time histories; the corresponding history Hilbert space $`\stackrel{~}{}`$ has dimension 4:
$`h_1`$ $`=`$ $`|\chi \chi ||\psi \psi |`$ (53)
$`h_2`$ $`=`$ $`|\chi ^{}\chi ^{}||\psi \psi |`$ (54)
$`h_3`$ $`=`$ $`|\varphi \varphi ||\varphi \varphi |`$ (55)
$`h_4`$ $`=`$ $`|\psi \psi ||\varphi \varphi |`$ (56)
These four histories are pairwise orthogonal, and their sum is the identity operator in $`\stackrel{~}{}`$. Thus, $`\{h_1,\mathrm{},h_4\}`$ is a homogeneous family of histories. Now, taking the initial density matrix $`\rho `$ to be the pure state $`\rho =|\varphi \varphi |`$, one can compute the following:
$`K(h_1)\rho K^{}(h_1)={\displaystyle \frac{1}{4}}|\psi \psi |;W(h_1)`$ $`=`$ $`{\displaystyle \frac{1}{4}}`$ (57)
$`K(h_2)\rho K^{}(h_2)={\displaystyle \frac{1}{4}}|\psi \psi |;W(h_2)`$ $`=`$ $`{\displaystyle \frac{1}{4}}`$ (58)
$`K(h_3)\rho K^{}(h_3)=|\varphi \varphi |;W(h_3)`$ $`=`$ $`1`$ (59)
$`K(h_4)\rho K^{}(h_4)=0;W(h_4)`$ $`=`$ $`0`$ (60)
Thus, the sum of weights in this HPO family of histories is $`3/2`$, barring any straightforward probability interpretation. This family is *not* a branching family of histories in the sense of section III, proving that branching families of histories are a subclass even of homogeneous HPO families.<sup>13</sup><sup>13</sup>13Note that the temporally reversed family *is* a branching family, and the weights do sum to unity (as only the mirror image of $`h_3`$, which is $`h_3`$ itself, contributes a non-zero weight).
### IV.4 Discussion
The consistent history approach to quantum mechanics offers a view of quantum mechanics that honours many classical intuitions while remaining, of course, faithful to the empirical predictions of orthodox quantum mechanics. While the initial motivation of the approach in terms of product families makes good pedagogical sense, it is too narrow for applications. Furthermore, as we have shown in section IV.2, product families may be misleading because they fail to distinguish between two importantly different notions of coarse-graining.
A number of applications demand branching families of histories. However, so far there has not been available a formally rigorous definition of that class of families. Isham’s HPO formalism, while offering the necessary generality, is in danger of losing touch with the intuitive motivation of the history approach. To be sure, this does not amount to any fundamental criticism, but it gives additional support for providing a less general definition that stays closely tied to the intuitive motivation of branch-dependent histories.
Through our definition of branching families of histories we have here provided the sought-for formal framework. Branching families are more general than product families, and they are general enough for applications while retaining a natural interpretation.
Figure 3 gives a graphical overview of the various kinds of families of histories that we considered in this paper. Branching families always admit an interpretation of weights in terms of probabilities, as the weights of the histories in such a family add to one. Consistent branching families in addition are free from interference effects. The formal framework presented here is neutral with respect to the question of whether families of histories that are not consistent in this sense can be put to good physical use or not.
###### Acknowledgements.
I would like to thank audiences at Oxford and Pittsburgh for stimulating discussions. Special thanks to Nuel Belnap, Jeremy Butterfield, Jens Eisert, Robert Griffiths, and Tomasz Placek. Thanks to Przemysław Gralewicz and to four anonymous reviewers for additional helpful comments. Support by the Polish KBN and by the Alexander von Humboldt-Stiftung is gratefully acknowledged.
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# A self-adjusted Monte Carlo simulation as a model for financial markets with central regulation
## 1 Introduction
During the past decades, the financial markets around the world have become more and more interconnected. The financial globalization has changed the organization structure of the stock markets by creating new risks and challenges for market participants. Dramatic expansion of the cross-border financial flows as well as domestic flows within the countries due to very rapid increase of telecommunication and computer-based products have emerged. As a consequence, the world financial markets are increasingly efficient today than ever before.
The integration and globalization process opens discussion about the activity of the central economic institutions that focus attention to financial market monitoring. The subsequent regulatory legislative yields asymptotically to dynamic balance between coalitions based on agreements of the cooperating companies and on the other hand competitors that enhance a diversity of the price trends. The feedback signals of different reliability are monitored and daily evaluated by central economic institutions. Consequently, this affects correlated decisions of the market participants. An expected effect of the global self-regulation is a portfolio efficiency.
Recently, there have been many attempts to gain a control over the dynamics of complex systems. The approaches based on the principles of feedback and adaptivity emerges as a perspective branch Melby2000 . The importance of the positive feedback mechanism in the economic context has been studied in Ref. YoussHolyst2000 . An approach based on the feedback control of spatially extended optical system has been employed in the Ref. Bleich1996 . The feedback mechanism is also exploited in Monte Carlo (MC) algorithm introduced by us recently in Ref. HorvathGmitra2004 . The algorithm is formulated as a random walk of temperature variable biased by feedback that mixes current and past stochastic signals, originating at the platform of extended statistical system. The time integrated signals are used to build up the actions shifting an instant temperature towards the critical temperature value. Due to limited memory depth, the convergence process is necessarily accompanied by uncertainty. Since the feedback treats a partial information only, many parallels with risky decisions within the business world can be found. Therefore, the MC feedback model may be used for simulation of decision-making in the stochastic environment.
The plan of the paper is as follows. In the next section the self-adjusted MC algorithm is reintroduced briefly. Its relationship to non-equilibrium spin updates is discussed in Sec. 3. The Sec. 4 deals with the statistics of spin clusters near the critical conditions and the self-adjusted algorithm is related to cellular automaton models of self-organized criticality (SOC). In Sec. 5, the applications of mentioned MC dynamics to models of financial markets is presented.
## 2 The self-adjusted random walk near criticality
Before proceeding to the introduction of financial model, we recall the salient points of the self-adjusted MC model HorvathGmitra2004 that is based on the random walk of the temperature variable $`T_t`$ defined by the one-step recurrence
$$T_{t+1}=T_t+\xi _t\mathrm{\Delta }\text{sign}(F_t),$$
(1)
where $`\mathrm{\Delta }`$ parameterizes a maximum step length $`|T_{t+1}T_t|`$, randomized by $`\xi _t`$ that is uniformly distributed within the $`1,1`$ range. The subscript $`t`$ accounts for MC steps per $`N`$ degrees of freedom. The feedback regulatory response
$$F_t=\frac{E^3_t3E_tE^2_t+2E_t^3}{T_t^4N}2\frac{E^2_tE_t^2}{T_t^3N},$$
(2)
mixes an influence of bias and randomness involved in the energy series $`\{E_t\}`$ that represents outcome of some extended statistical system. Eq.(2) approximates a temperature derivative of the specific heat $`\frac{}{T}[(E^2_{\mathrm{eq}}E_{\mathrm{eq}}^2)/(T^2N)]`$ at the canonical equilibrium (eq) ensemble defined by the constant temperature $`TT_t`$ Binder1988 . In Eq.(2) the canonical averages are approximated by instantaneous running averages of the energy cumulants $`E^p_t`$, $`p=1,2`$ calculated with the assistance of linear filter Principe2000 as it follows
$$E^p_t=(1\eta )E^p_{t1}+\eta E_t^p.$$
(3)
Here the plasticity parameter $`\eta (0,1)`$ reweights an energy samples $`E_t`$ according to delay. For $`\eta 0`$, running averages vary slowly and long-memory processes dominates, whereas in the limit $`\eta 1^{}`$ the memorizing becomes faster but less efficient.
The feedback is applied to maintain the critical regime of the 2d Ising spin micromodel of ferromagnet. This is defined for $`i=1,2,\mathrm{},N=L^2`$ sites of $`L\times L`$ square lattice characterized by the exchange energy $`E=_{\text{nn}}S_iS_j`$, $`S_i\{\frac{1}{2},\frac{1}{2}\}`$, where exchange coupling constant rescales an instant temperature $`T_t`$. The nearest neighbor (nn) interaction is considered. This interaction picture is supplemented by the periodic boundary conditions in all directions. The standard Metropolis single-spin-flip algorithm Binder1988 ; Metropolis where two states are linked by the flip acceptance probability $`p_{\mathrm{acc}}=\mathrm{min}\{1,\mathrm{exp}\left(\delta E_i/T_t\right)\}`$ depends on the energy change $`\delta E_i`$ at $`i`$th site. The technical issue of the proposed implementation is that $`T_t`$ is hold constant during $`N=L^2`$ spin updates (i.e. 1 MC step per sample). The example which illustrates $`T_t`$ dynamics is shown in Fig.1. As follows from Eq.(3), a filtered information serves to provide $`T_{t+1}`$ near criticality. In compliance with a previous facts, the feedback role in dynamics can be summarized as follows: a) it receives series of energy $`\{E_t\}`$ values from extended statistical system (or only sufficient part of the series is served for this purpose); b) it treats the series via the linear filtering that generates running averages $`E^p_t`$; c) it transmits the regulatory signals $`\text{sign}(F_t)`$ via the temperature given by Eq.(1).
After the transient time the thermal noise overcomes a regular bias involved in stochastic $`F_t`$ and system with feedback maintain a stationary regime. Due to fluctuations of $`T_t`$ a non-equilibrium distribution is generated. Therefore, the expectation mean values, worth our interest are given by
$$T^p=\frac{1}{n_1n_0}\underset{t^{}=n_0}{\overset{n_1}{}}T_t^{}^p,p=1,2.$$
(4)
where $`n_1n_0`$. If the number of excluded steps $`n_0`$ is larger than the transient time, the average $`T`$ estimates critical point of the extended system. Thus, the observation time is defined by $`\tau _{\mathrm{obs}}=n_1n_0`$. In our simulations we used $`n_010^5`$, $`n_110^7`$.
## 3 The finite-size effects of $`T`$
There are many examples of the simple non-equilibrium spin-flip stochastic models Buonsante2002 ; Garrido87 giving rise to the various dynamical phases. To simulate nearly equilibrium systems, the specific non-equilibrium approaches have been suggested. For instance, the Ref. Tomita2001 exploits the controlled growth of percolating spin cluster.
The postulation of artificial dynamics clearly places novel time scales into original equilibrium problem. In our case the combination of two parameters $`\eta `$ and $`\mathrm{\Delta }`$ affect the distance between a stationary non-equilibrium statistics and canonical equilibrium. The basic properties of stationary non-equilibrium regime can be summarized as follows:
1. The stationary regime with a ferro-paramagnetic order and positive $`T`$ is stabilized by a sufficiently slow dynamics limited to $`\mathrm{\Delta }<\mathrm{\Delta }_{\mathrm{tr}}(\eta )`$, where $`\mathrm{\Delta }_{\mathrm{tr}}(\eta )`$ is a threshold value. The irreversible drop to negative $`T`$ (antiferromagnetic order) occurs at $`\mathrm{\Delta }=\mathrm{\Delta }_{\mathrm{tr}}`$.
Similar transition can be found for the coupled map lattice where antiferromagnetic order is generated spontaneously by ferromagnetic couplings accompanied by sufficiently fast synchronous dynamics Angelini2003 . For $`L=10`$, $`\eta =10^3`$ we estimated $`\mathrm{\Delta }_{\mathrm{tr}}2.2\times 10^3`$.
2. The canonical equilibrium limit $`TT_\mathrm{c}(L)`$, where $`T_\mathrm{c}(L)`$ is the pseudocritical temperature of specific heat, is reached for the sufficiently slow dynamics $`1\mathrm{\Delta }>0`$, when the feedback memory is deep ($`1\eta `$) and parameters admit the scale separation $`\eta \mathrm{\Delta }`$. In addition, the stationary averages have to be accumulated for $`\tau _{\mathrm{obs}}1/\eta `$.
3. For stationary regime (see point 1) the nonzero variance $`(TT)^2`$ gives rise to the broad energy distribution in comparison to Boltzmann one.
The question of practical relevance is whether true $`T_\mathrm{c}`$ can be attained by non-equilibrium dynamics. The result of simulation realized for $`\eta =10^5`$, $`\mathrm{\Delta }=10^9`$ is depicted in Fig.2. The standard finite-size dependence of $`T_\mathrm{c}(L)T=T_\mathrm{c}+b/L`$ Binder1988 (where $`b`$ denotes the thermal coefficient) has been obtained from the fit. Without additional tuning the simulation provides estimate $`T_\mathrm{c}0.5679`$ of the equilibrium exact value $`T_\mathrm{c}^{\mathrm{ex}}[2\mathrm{ln}(1+\sqrt{2})]^10.56729`$ ExactSolOnsager1944 .
As one can see in Fig.3a follows that substantial finite-size corrections occur for a short-time memory. In that case, the temperature mean value can be interpolated by
$$T=T_\mathrm{c}+\frac{b}{L}A\eta ^q,\mathrm{\Delta }\eta <\eta _\mathrm{c}(L),$$
(5)
where cut-off $`\eta _\mathrm{c}(L)`$ is roughly interpolatated by $`\eta _\mathrm{c}(L)=\eta _{\mathrm{}}0.343/L`$ with $`\eta _{\mathrm{}}=0.468`$. From the fit carried out within the range $`\eta \mathrm{\hspace{0.17em}0},\frac{\eta _\mathrm{c}}{2}`$ we have determined the exponent $`q1.5`$. Assuming that the lowest mean value $`TT_{\mathrm{min}}T_2(L)`$, where $`\eta _{\mathrm{}}`$ is achieved, for $`L\mathrm{}`$ we may write $`A=(T_\mathrm{c}T_{\mathrm{min}})/\eta _{\mathrm{}}^q`$. The presence of the temperature lower bound is associated with the condition that peak of specific heat is formed. This implies the requirement of a residual acceptance of single-spin-flips at very low temperatures. The bound can be explained in the frame of two-level specific heat model denoted as $`C_2`$. The model takes into account only single branch of ground state $`L^2/2`$ (since the ergodicity is broken, the ground state is counted once) and $`L^2`$-fold single-spin excitation of energy $`2L^2/2`$. The condition for extreme $`\mathrm{d}C_2/\mathrm{d}T|_{T=T_2}=0`$ yields equation $`(T_21)\mathrm{exp}(2/T_2)+L^2(1+T_2)=0`$. The last relation predicts a monotonous decrease of $`T_2`$ with $`L`$ and solution can be simply approximated by $`T_2\frac{1}{\mathrm{ln}L}`$ for infinite $`L`$ asymptotics. To estimate $`T_2`$ we have carried out simulation for $`L=90`$ and $`\eta (0.45,0.55)`$ that gives $`T`$ from $`(0.25,0.3)`$ range. The mean value $`T`$ is approximated by $`T_2=0.212\frac{1}{\mathrm{ln}(90)}=0.222`$.
The assumption about the energy gap corresponding to two lowest levels is consistent with the presence of the crossing over memory depth $`1/\eta _\mathrm{c}2`$ (two MC steps). To maintain a stationary regime the substantial memory reduction for $`\eta >\eta _\mathrm{c}`$ have to be compensated by the enlarged $`T`$ value. We observed that for $`\eta \eta _\mathrm{c}`$ the ergodic time $`\tau _\mathrm{e}`$ between the plus and minus magnetization stages diverges as $`\tau _\mathrm{e}L^z`$, where $`z2`$ is the dynamical exponent LandauBinder2000 . If $`\tau _{\mathrm{obs}}`$ is smaller than $`\tau _\mathrm{e}`$, the broken symmetry can be detected. Clearly, for $`\eta >\eta _\mathrm{c}`$ the ”jaggy” $`T(\eta )`$ dependence indicates the occupancy of only few lowest energy levels.
## 4 The relationship to SOC and percolation models
The concept of SOC proposed by Bak Bak88 represents the unifying theoretical framework of systems that drive themselves spontaneously to critical regime. The sand pile, forest-fire Drossel92 and game-of-life Alstrom94 are well-known examples of cellular automaton SOC models. Their stationary regime can be characterized by invariant probability distribution functions (pdf’s) of spatial or temporal measures of the dissipative events called avalanches. The corresponding pdf’s include a power-law intervals with the specific non-equilibrium critical exponents.
According to Ref. Kadanoff , the stationary SOC regime implies the operation of inherent feedback mechanism. On the contrary, our model has explicitely defined global feedback coupled to the MC dynamics. Such coupling is completely different from the standard models of SOC, but the arguments given in this section suggest that additional coupling anchors statistical properties indistinguishable from SOC models.
In the next, we focus our attention to the problem of partial information regulation to answer the question, how important can be a concrete feedback form. We consider spatially truncated form of feedback defined by energy cumulants calculated for the rectangular $`[\varphi L]\times L`$ lattice segment, where $`\varphi 0,1`$ and $`[x]`$ denotes an integer value of $`x`$. The results are depicted in Fig.3b for $`L=60`$. The simulation has confirmed that stationary regime is robust with respect to the partial information regulation represented by $`\eta `$ and $`1\varphi `$ terms. As one can be see in Fig.3, both sufficiently small $`\eta `$, ($`\eta <\eta _\mathrm{c}(L))`$, and $`1\varphi `$ have similar influence on diminishing the stationary $`T`$ value.
As pointed out in Ref. Sornette , the pdf’s of spin clusters at a second-order critical point resembles pdf’s of the spatial extent of avalanches. Admittedly, the correspondence between spins and clusters is not unique and thus numerous mapping schemes could be examined. We have utilized the well-known mapping to bond percolation model Fortuin72 via Wolf’s algorithm Wolff89 to identify cluster without overturning. To interpret the cluster statistics, we suggest the universal scaling in the form
$$P(s,L)=L^\beta g\left(\frac{s}{L^D}\right),$$
(6)
where $`s`$ is the number of spins belonging to the same cluster. In Fig.4 are depicted spin cluster pdf’s. As one can see, the distribution functions for clusters $`ss_{\mathrm{max}}(L)`$ clearly obey the scaling given by Eq.6, where $`s_{\mathrm{max}}(L)`$ is proportional to the size of a spanning cluster. It is associated with the peak $`P(s_{\mathrm{max}},L)`$ in $`P(s,L)`$ distribution. Using the scaling relation $`s_{\mathrm{max}}(L)L^D`$, the fractal dimension $`D1.98`$ has been determined from the fit, see inset in Fig. 4a. Assuming an asymptotic form $`g(x)x^\tau `$ in the limit $`x0`$ we obtained $`P(s,L)L^{\beta +\tau D}s^\tau `$ that implies independence on $`L`$ if $`\beta =\tau D`$. Surprisingly, the dependencies in Fig.4b indicate that small-scale power-law scenario survives for full $`\eta (0,1)`$ range. It doesn’t matter how distant an assembly is from the thermal equilibrium except a spanning cluster that grows with increasing $`\eta `$.
Since the static bond percolation model pays no attention to temporal variations, the alternative maps should be examined, for instance those taking into account a short-time memory of cluster evolution Mouritsen92 . The question whether differences in a spin - cluster map definition might affect a non-equilibrium exponent $`\beta `$, introduced in Eq.6, is left for future research.
## 5 The relationship to financial market models
### 5.1 Statistical properties relevant for financial market modeling
Performed simulations for pdf of $`F_t`$ defined as $`P_\mathrm{F}(F)\delta _{F,F_t}`$ confirm that flat tails exhibit within the full range of $`\eta `$, see Fig.5. Therefore, it is worth to check functional stability of aggregated regulatory responses
$$r_{n,t}=\underset{t^{}=t}{\overset{t+n}{}}F_t^{}.$$
(7)
The results of analysis depicted in Fig.5b has validated the expected scaling
$$P_\mathrm{F}(r_n)=r_n^{1/\alpha }f\left(r_nn^{1/\alpha }\right)$$
(8)
for $`n<10`$ with estimated exponent $`\alpha 1.3`$, where $`f()`$ is some universal function. Since $`F_t`$ is bounded, the crossing-over from Lévy to Gaussian scaling has been verified, see Fig.5b. More profound connection to the price dynamics is discussed in the following section.
The dynamics of self-adjusted system can be analyzed by means of autocorrelation functions of the time lag $`t`$:
$`\mathrm{C}_t^{(p)}`$ $`=`$ $`{\displaystyle \frac{F_t^{}^pF_{t^{}+t}^pF_t^{}^pF_{t^{}+t}^p}{(F_t^{})^{2p}(F_t^{}^p)^2}},p=1,2.`$ (9)
The MC model yields nearly exponential decrease of $`\mathrm{C}_t^{(1)}`$, see Fig.6. The power-law dependence of $`\mathrm{C}^{(2)}`$ which becomes more pronounced for larger $`L`$ is identified. Assuming the form $`\mathrm{C}_t^{(2)}t^q^{}`$ we have determined $`q^{}1.5`$.
### 5.2 The financial market model formulation
On the basis of previous facts we have proposed model based on the two postulates: (i) under the conditions of new economy the world market entities (participants, companies) are forced to evolve towards the longer auto-correlation times. It means that hardly predicable - volatile price statistics gets stuck in the regime which imply highest degree of uncertainty. It might be classified as an edge of the chaos. Such interpretation is fundamentally equivalent with the dynamics of the continuous phase transitions and dynamics of information processing Langton1990 ; (ii) the centralized institutions are established to satisfy the global information and legislative needs of market entities looking for the arbitrage opportunities. The action of participants leads to enhancement of market effectiveness in the sense of the postulate (i).
Our simulation have shown that presented spin model that involves feedback has much in common with Ising spin market models of the price dynamics Ponzi2000 . The feedback introduced in Bornholdt’s model Bornholdt2001 , which was later elucidated by Kaizoji, Bornholdt and Fujiwara (BKF), is linked to the local in time and global in lattice space magnetization $`m_t=(2/L^2)_{i=1}^{L^2}S_{i,t}`$. Both models are inspired by the paradigm of minority game Challet1997 . Speaking in game theory term, the feedback promotes the continuing tournaments between competing ordered (ferro) and disordered (para) magnetic phases. To maintain a balance between ordered ferromagnetic and disordered - paramagnetic spin phase, the feedback mechanism including a global cooling or heating is applied. The predicted temperature is simultaneously transmitted to all lattice sites. Clearly, the effect of such step has probabilistic nature, because a freedom exists between the market entities which can uphold or abolish older or establish new coalitions, respectively. Subsequently, any decision is reflected by correlations between entities.
Two principal situations may be observed from the view point of the proposed model. At the subcritical temperatures the clusters grow, which means that correlations among market entities become increase. A subcritical temperature can increase a mean cluster size until a spanning cluster come into existence. The ferromagnetic domain state can be understood as a market situation where competition is distorted due to monopoly, duopoly or trust formation effects. To prevent from highly correlated stocks, the feedback triggers the antitrust operations Wrigley1992 . The portfolio diversification is represented by the fragmentation of the spin clusters.
### 5.3 The comparison with other spin models of financial markets
The spin model presented in previous sections displays several interesting statistical properties that are particularly relevant for stochastic financial market dynamics. One of our challenges is to examine cross-links between our feedback model and BKF model. In BKF model the local spin behavior is given by the competition between the local and non-local interactions involved in local field $`h_{t,i}=2(_{\mathrm{nn}}S_{j,t}\kappa S_{i,t}|m_t|)`$, where $`\kappa `$ is proportionality factor of feedback instant magnetization. The probability of $`\pm \frac{1}{2}`$ lattice spin states is given by $`1/[1+\mathrm{exp}(2\stackrel{~}{\beta }h_i)]`$, where $`\stackrel{~}{\beta }`$ is the inverse constant temperature. The difference $`m_{t+1}m_t`$ is interpreted as a logarithm of the price return. Within the BKF model the instant price is proportional to $`\mathrm{exp}(\lambda m_t)`$ term, where $`\lambda `$ is related to the ratio of the interacting and fundamental traders.
In further we assume the equivalence of $`\mathrm{exp}(\lambda m_t)`$ and $`\mathrm{exp}(\lambda F_t)`$ terms considered for different models. The interpretation of $`\mathrm{exp}(\lambda F_t)`$ as a price term is in accord with requirement that price increases when the competition is weakened. The crucial items emphasizing the common aspects of models are listed in table 1 and illustrated in Figs. 7 and 8.
We see that both models yield Lévy distributions with the nearly common index $`\alpha _13.3`$ related to cumulative distributions with power law slope $`2.3=\alpha _1+1`$ that is consistent with the empirical data Gopikrishnan2000 . The bowl-shaped anomaly at small $`|F_{t+1}F_t|`$ is probably due to $`\text{sign}(F_t)`$ sharpness. The remarkable feature of this comparison is a similarity of the exponents obtained for pdf’s of the bull (bear) market durations. In addition, the similarity of temporal aspects can be seen in coincidence of stationary conditional probabilities of bull-bear switching listed in the Table 1. We see that $`F_t`$ feedback model acquires the exponential form typical for the positively correlated quantities, whereas the magnetization variations predicted by BKF are anticorrelated. Clearly, this difference can be explained providing energy cumulants temporal filtering. The simulation of BKF model reveals some gap between Hurst exponents Kim2004 . Contrary to the standard diffusive $`\delta t^{0.5}`$ behaviour which is typical for $`|m_{t+\delta t}m_t|`$, the superdiffusive short-time regime $`|F_{t+\delta t}F_t|\delta t^{0.7}`$ has been identified for $`F_t`$ feedback model.
### 5.4 The application to inter-stock correlations
In the subsection certain diverse application of our model is presented. As one sees, the MC model may be used to understand the basic features of the mechanism of regulation of bilateral stock cross-correlations. The role of central institutions, modeled via the feedback, is to perform large-scale eligible decisions that could enhance portfolio efficiency. These contribute to $`F_t`$ comprising a summary of the stock prices. We assume that different stocks are distinguished by lattice position and sign of dominant trend at given time. And for now the correspondence between spins and stocks follows a simple majority rule: $`S_i=\frac{1}{2}`$ ($`S_i=\frac{1}{2})`$ when the sell (buy) market orders prevail during elementary time lag (1 MC step/site), respectively. The cross-correlations among directly regulated - four ”nearest” stocks are related to the mean exchange energy term $`S_iS_{j\mathrm{nn}(i)}`$ optimized by $`F_t`$ to be extremely susceptible to regulation. The long-range cross-correlations stem from indirect regulation of the portfolio composition near to criticality.
The presence of correlations (and anticorrelations) between pairs of stocks play a key role in the theory of selecting of the most efficient portfolio composition of the financial goods Mantegna99 . For discussed spin model the properties of the pairs of stocks can be characterized by the cross-correlation coefficient
$$c_{ij,t}=\frac{4}{n_{\mathrm{av}}}\underset{t^{}=t}{\overset{n_{\mathrm{av}}+t}{}}S_{i,t^{}}S_{j,t^{}},i,j\{1,2,\mathrm{},L^2\}$$
(10)
ranging from $`1`$ to $`1`$. Here $`i,j`$ denote the pair of randomly chosen sites. In the simulations the temporal averages have been calculated for $`n_{\mathrm{av}}=10^3`$ steps and accumulated in pdf’s depicted in Fig.9. We see that the pair correlations are dominantly positive except a negative tail that agrees qualitatively with the empirical studies Mantegna99 . For $`\eta `$ larger than $`0.1`$ the negative tail vanishes and pdf of $`c_{ij,t}`$ is shifted towards $`1`$. Clearly, such ”coherent” regime is typical for inefficient portfolio structure.
## 6 Conclusion
In conclusion, we have presented application of adaptivity concept to critical phenomena. Despite of the principal difficulties, insufficient understanding of the systematic errors which arise proceeding from non-equilibrium to equilibrium, we expect that work might inspire numerous applications to other inverse statistical problems where constrains are written in terms of running averages.
We have demonstrated that MC model can be reinterpreted in macro-economic terms. The simulations indicated spontaneous emergence of the power-law distribution for feedback regulatory signals which is universal feature of the stationary non-equilibrium regime. Surprisingly, we found that $`F_{t+1}F_t`$ statistics contradicts to our previous expectations based on the finding $`\text{C}_t^{(2)}t^q^{}`$, which does not imply power-law fall-off of the autocorrelations of $`|F_{t+1}F_t|`$. The influence of the memory depth related to parameter $`\eta `$ has also been investigated extensively. The critical point related to the most efficient portfolio structure, is an abstract of ideal situation when $`\eta 0^+`$, which is unattainable in a real finite-memory conditions. Otherwise, we have found that the portfolio efficiency is weakened via short-time memory for $`\eta >0.5`$. The comparative study of two spin models have confirmed that presented MC dynamics based on the specific heat feedback offers an alternative description of financial markets. Despite of outstanding similarities, the comparison of exponents indicate non-trivial relationship between the universality classes of studied models.
Acknowledgment
The authors would like to express their thanks to Slovak Grant agency VEGA (grant no.1/2009/05) and agency APVT-51-052702 for financial support.
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# AB-INITIO STUDY OF STRUCTURE AND DYNAMICAL PROPERTIES OF CRYSTALLINE ICE
## 1. INTRODUCTION
Ice, the frozen form of liquid water, is one of the most common materials on earth and in outer space, and has important relevance to a large number of diverse fields such as astronomy, geophysics, chemical physics, life sciences, etc. (Petrenko and Whitworth, 1999). Ice-Ih, the hexagonal form of ice is the most commonly known phase of ice in which each oxygen atom has four neighbours at the corners of a tetrahedron. When water freezes the forces of interaction between H<sub>2</sub>O molecules win over their thermal motion and they form a most stable arrangement precisely with hexagonal symmetry. This is the reason why snow flake crystals are always hexagonal. The crystal structure of ice at the atomic level marks the way these crystals will look. The hydrogen atoms are covalently bonded to the nearest oxygen atoms to form water molecules which are linked to each other through hydrogen bonds. The high translational symmetry is not retained at the level of the crystallographic unit cell. The phase diagram of ice is very complex with over 12 known phases existing, (see Fig. 1). Besides the environmental importance, ice is also special because of the interesting phenomena contained within its structure. The crystal structure of ice is very unusual because, while the molecules lie on a regular crystal lattice, there is disorder in their orientations. This property leads to many interesting characteristics in electrical polarization and conductivity.
Numerous attempts have been made for the past decades to understand the nature of the lattice vibrations of ice, in particular, ice-Ih (Marchi, Tse and Klein, 1986). Experimental information has been obtained from infrared absorption (Whalley and Bertie, 1967), Raman scattering (Klug and Whalley, 1978) and both coherent and incoherent inelastic neutron scattering (Renker and Blanckenhagen, 1969; Prask and Trevino, 1972). On the theoretical side, three basic approaches have been adopted to study the lattice modes. The earliest studies involved the application of lattice dynamics to hypothetical proton ordered structures (Whalley and Bertie, 1967; Prask and Trevino, 1972). Later lattice dynamics was used to study more realistic orientationally disordered structures (Nielsen, Townsend and Rice, 1984). The most recent work has utilized the molecular dynamics simulation techniques (Tse, Klein and McDonald, 1984). Many of the these theoretical studies involved the use of classical modeled potential through empirical method in order to describe the interaction of the system (Marchi, Tse and Klein, 1986; Nielson, Townsend and Rice, 1984). As result of all these works, the overall features of the lattice mode spectrum in the translational region (0-300 cm<sup>-1</sup>) and in the librational region (450-950 cm<sup>-1</sup>) are reasonably well understood. Potential-based empirical modelling has had some success towards the end, but to date, there are no empirical potentials capable of describing the ice dynamics and related properties across its whole spectra range and describing certain key spectra features.
The ab-initio method has recently gained ground not only because of its reliability in the study of static and dynamics properties of ice (Cote, Morrison, Cui, Jenkins et al., 2003; Morrison and Jenkins, 1999) but also because the method allow to model some important features such as the ordered periodic ice structure (Lee, Vanderbilt, Laarsonen, Car et al., 1993), and the nature of hydrogen bond in different geometries (Xantheas and Dunning, 1993). Our present study to understand the microscopic nature and lattice vibrations of ice-Ih makes use of the Vienna Ab Initio Simulation Package (VASP) (Kresse and Futhmüller, 1996) which is designed to perform ab-initio quantum-mechanical molecular dynamics using pseudopotentials and a plane wave basis set and the recent implementation of the projected augmented wave (PAW) method.
The computational study presented in this work begins with the choice of unit cell of ice as described in Section 2, which is followed by molecular dynamics simulation of ice in a supercell (Section 3) with the unit cell replicated in all directions in order to obtain the structural behaviour in comparison with the available experimental data. In Section 4 and 4.1 we present the results of calculated phonon spectrum of ice in the translational region, which is done in a $`3\times 1\times 1`$ supercell in $`y`$-direction, as well as the corresponding integrated vibrational density of states in Section 4.2. The anomalous behaviour of this ice structure observed in the low energy region is presented in Section 4.3. The final Summary of this work is presented in Section 5.
## 2. COMPUTATIONAL DETAILS
As mentioned above, the calculations in this work were carried out by using the Vienna Ab Initio Simulation Package (VASP) (Kresse and Futhmüller, 1996) which has been designed for ab-initio quantum-mechanical molecular dynamics simulations using pseudopotentials and a plane wave basis set and the recent implementation of projected augmented wave (PAW) method. In order to investigate the structural properties of ice, we started with the construction of ice crystal using the Bernal and Fowler rule (Bernal and Fowler, 1933). The rule is based on a statistical model of the position of hydrogen atoms produced by Pauling (Pauling, 1935) using the six possible configurations of hydrogen atoms within Ice Ih. It is defined as ideal crystal based on the assumptions that:
* Each oxygen atom is bonded to two hydrogen atoms at a distance of 0.95 to form a water molecule;
* Each molecule is oriented so that its two hydrogen atoms face two, of the four, neighbouring oxygen atoms in the tetrahedral coordination;
* The orientation of adjacent molecules is such that only one hydrogen atom lies between each pair of oxygen atoms;
* Ice Ih can exist in any of a large number of configurations, each corresponding to a certain distribution of hydrogen atoms with respect to oxygen atoms.
The schematic drawing shown in Fig. 2 satisfies one out of the six possible orientations of protons of the central water molecule according to these rules. Each of the oxygen atoms can be linked to another oxygen by the combination of a covalent bond plus a hydrogen bond to form the tetrahedral arrangement of oxygen atoms. A unit cell of this ice was prepared in a cubic box according to Fig. 3 with 8 molecules of water. All the atomic degrees of freedom were relaxed using VASP with the projected-augmented wave (PAW) formalism at high precision. The optimum Monkhorst Pack $`4\times 4\times 4`$ $`k`$-point was used in addition to the generalized gradient approximation (GGA) of Perdew-Wang in order to describe the exchange-correlation and to give good description of the hydrogen bonding of water. We used a energy cut-off of 500 eV because the 2p valence electrons in oxygen require a large plane wave basis set to span the high energy states described by the wavefunction close to the oxygen nucleus and also the hydrogen atoms require a larger number of planes waves in order to describe localization of their charges in real space. The lattice constants of the unit cell were calculated from the plot of the energy against the volume. The estimated values of the lattice constants are $`a`$ = 6.1568 Å, $`b`$ = 6.1565 Å, $`c`$ = 6.0816Å. The values of $`abc`$ imply that the relaxed structure is tetragonal with $`c/a`$ ratio $``$ 0.988. The experimental lattice constant reported by Blackman et. al. (Blackman and Lisgarten, 1958) is 6.35012652 Å for the cubic geometry. The final geometry obtained was used in our molecular and lattice dynamical study of ice.
## 3. MOLECULAR DYNAMICS
In our molecular dynamics simulation, the final relaxed geometry in Fig. 2 (or Fig. 3 was replicated in all directions using the calculated values of the lattice parameters to produce 64 molecules of water which form the hexagonal structure shown in Fig. 4. Here the molecular dynamics simulation was done for the $`\mathrm{\Gamma }`$-point only in a box of dimension $`12.411356\times 12.411356\times 12.259675`$ Å<sup>3</sup> corresponding to the density 1.01 $`\mathrm{gcm}^3`$ to be compared to the real density of ice Ih and ice Ic which is 0.92 $`\mathrm{gcm}^3`$. Molecular dynamics was carried out for 3 ps using a time step of 0.5 fs. The angle H-O-H of the ice structure was compared with liquid water at the different temperatures as can be seen in Fig. 5. For ice structure, the H-O-H angle was found to be 107.8 compared to liquid water case at different temperature or isolated water molecule also calculated with VASP which is found to be 104.5 and 105.4 degree, respectively, as can be seen in Fig. 5. The angle shown by solid ice is an indication that the oxygen centre preserves its tetrahedral structural units. Also, the range of the angular distribution for the high temperature water (or supercritical water) is wider and broader than the corresponding liquid water at ambient temperature and the solid ice cases.
The radial distribution functions for the solid ice accumulated over a 3 ps run are plotted at two different temperatures, 100 K and 220 K. The $`g_{\mathrm{OO}}`$ and $`g_{\mathrm{OH}}`$ in Fig. 6, and $`g_{\mathrm{HH}}`$ in Fig. 7 at 100 K, all exhibit the long-range like order when compared to the short range order of liquid water. There is a well pronounced first peak of $`g_{\mathrm{OO}}`$ at 2.75 Å at 100 K and also the position of the first minimum is deeper when compared to the liquid water radial distribution functions. This result is comparable to the result obtained using the classical TIP4P modelled potential (Svishchev and Kusalik, 1994) The result of the radial distribution function of oxygen atoms, $`g_{\mathrm{OO}}`$, oxygen-hydrogen atoms $`g_{\mathrm{OH}}`$ (Fig. 8), and hydrogen atoms (Fig. 9) at 220 K were fairly comparable with available neutron diffraction scattering data of Soper (Soper, 2000). The position of the first peak (VASP calculation) shows very little difference from 100 K while the height of the peak is lower for 220 K due to the effect of entropy increase of the hydrogen atoms which results from re-orientation of protons that tends to force apart more oxygen atoms to a larger distance. The positions of the second minimum for 100 K and 220 K solid ice are, respectively, 4.8 Å and 4.9 Å. The experimental result is slightly shifted to the right to the value 5.1 Å for 220 K. The deviation from experimental value might be due to the phase of crystalline ice under consideration.
## 4. LATTICE DYNAMICS
In this study phonon dispersion curves are calculated by using the PHONON package developed by K. Palinski (Parlinskin, 2002) which has been designed to take input data of Hellmann-Feymann forces calculated with the help of an ab-initio electronic structure simulation program such as VASP. We carried out the lattice dynamics study of ice by using the geometry of eight-molecule primitive cell discussed in the Section 1. The calculations of force constants was carried out by considering a $`3\times 1\times 1`$ supercell containing 24 molecules of water which is obtained by matching 3 tetragonal unit cells. At the first step of the calculation, the PHONON package is used to define the appropriate crystal supercell for use of the direct method. As done for the primitive unit cell, all the internal coordinates were relaxed until the atomic forces were less than $`10^4`$ eV/Å. The relaxed geometry for a $`1\times 1\times 1`$ supercell from the initial configurations containing 8 molecules is shown in Fig. 3. The starting geometry of the molecules in the simulation box shown is such that no hydrogen bonds were present but the positions of oxygen atoms follow the tetrahedral orientation. After the relaxation, all the protons perfectly point to the right direction of oxygen atoms and make the required hydrogen bonds necessary as indicated by the dotted lines in Fig. 3 to preserve the tetrahedral orientation of the ice structure.
Figure 10(I) shows the Brillouin zone belonging to the relaxed structure of our (model) ice shown in Fig. 3. Figure 10(II) shows the Brillouin zone used in the analysis of the measured phonon spectra for the model structure of D<sub>2</sub>O ice. Let us say once again that the relaxed structure shown in Fig. 3 has the long-range orientational order while the actual structure of ice Ih has no long-range orientational order. Therefore, in the analysis of the measured phonon dispersion curves of D<sub>2</sub>O, one uses another model for ice shown in Fig. 4. We have to keep this in mind when comparing our calculated dispersion curves with the measured frequencies. For the evaluation of the phonon dispersion curves we have used the direct ab-initio force constant method (Parlinskin, Li and Kawazoe, 1997), whereby the forces are calculated with VASP via the Hellmann-Feymann theorem in the total energy calculations. Usually, the calculations are done for a supercell with periodic boundary conditions. In such a supercell, a displacement $`𝐮(0,k)`$ of a single atom induces forces $`𝐅(lk)`$ acting on all other atoms,
$`F_\alpha (lk)={\displaystyle \underset{l^{}k^{}\beta }{}}\mathrm{\Phi }_{\alpha \beta }(lk;l^{}k^{})u_\beta (l^{}k^{}).`$ (1)
This expression allows to determine the force constant matrix directly from the calculated forces (see Parlinski et. al.) (Parlinskin, Li and Kawazoe, 1997; Parlinskin, 2002). The phonon dispersion branches calculated by the direct method are exact for discrete wave vectors defined by the equation
$$\mathrm{exp}(2\pi ı𝐤_L𝐋)=1,$$
(2)
where $`𝐋=(L_a,L_b,L_c)`$ are lattice parameters of the supercell. A related technique has recently been used to obtain accurate full phonon dispersions in highly symmetric structures of $`\mathrm{Ni}_2`$GaMn (Zayak, Entel, Enkovaara, Ayuela et al., 2003).
In order to obtain the complete information of the values of all force constants, every atom of the primitive unit cells was displaced by 0.02 Å in both positive and negative non-coplanar, $`x`$, $`y`$ and $`z`$ directions to obtain pure harmonicity of the system. We use a $`3\times 1\times 1`$ supercell which implies that 3 points in the direction are treated exactly according to the direct method. The points are \[$`\zeta `$00\], with $`\zeta `$ = 1, 1/3, 2/3. We calculate forces induced on all atoms of the supercell when a single atom is displaced from its equilibrium position, to obtain the force constant matrix, and hence the dynamical matrix. This is then followed by diagonalization of the dynamical matrix which leads to a set of eigenvalues for the phonon frequencies and the corresponding normal-mode eigenvectors. The vibrational density of states (VDOS) is obtained by integrating over $`k`$-dependent phonon frequencies from the force-constant matrix in supercells derived from the primitive molecule unit cells.
## 4.1 Phonon dispersions of crystalline ice
The phonon dispersion curves calculated for our ice crystal in direction are shown in Fig. 11 for low lying energy vibrations. According to the geometry of the supercell, the low-frequency (or low-energy 0-50 meV) acoustic modes can be compared to Renker’s inelastic neutron scattering measurement (Renker, 1973) along the direction ($`\mathrm{\Gamma }`$A) of hexagonal symmetry shown in Fig. 10, taken from reference (Li, 1996), though our calculation was done only along the Cartesian direction of the cubic cell. Our transverse and longitudinal acoustic (LA/TA) dispersions are well behaved when compared to some other modelled calculations or experimental results (Renker, 1973) in the high symmetry directions $`\mathrm{\Gamma }A`$ of hexagonal ice as shown in Fig. 11(c). We can also compare our result to the Cote et. al. result in Fig. 11(b) (Cote, Morrison, Cui, Jenkins et al., 2003), where they have recently used the ab-initio method to obtain the phonon dispersions in the translational frequency range for the ice structure in the Brillouin zone of the orthorhombic eight-molecule unit cell. Altogether our LA and TA dispersions are better than Cote’s LA/TA in comparison to the experimental curves in Fig. 11(c). Our dispersion curves in direction can as well be compared to the dispersion curves obtained using a dynamical model with two force constants to describe the low frequencies of vibrations of hexagonal ice as proposed by Faure (Faure, 1969).
Although everywhere along the $`\mathrm{\Gamma }`$-point, our dispersions are completely degenerate in the optic region whereas Cote’s dispersions in Fig. 11(b) show some splittings, the so-called longitudinal/transversal optic (LO/TO) splittings whose origins is explained below, while at the zone boundary, some of the dispersions are non-degenerate unlike our results. Our inability to reproduce these splittings at the $`\mathrm{\Gamma }`$-point is due to the direct method approach in which absolute periodicity of the crystal according to Born-von Kármán conditions was considered. The splitting of LO and TO branches for long wavelengths occurs in almost all crystals which are heteropolar (partially ionic such as GaAs) or ionic (such as NaCl) at the $`\mathrm{\Gamma }`$-point, and only for infrared active modes (Srivastava, 1990). The long-range part of the Coulomb interaction causes the splitting of the $`k`$ = 0 optic modes raising the frequency of LO modes above those of TO modes. The long-range part of the Coulomb interaction corresponds to the macroscopic electric field arising from ionic displacements. Ice is a tetrahedrally covalently bonded polar system whose dipole-dipole interactions give rise to the electric field when they are disturbed. The origin of the splitting is therefore the electrostatic field created by long wavelength modes of vibrations in such crystals. Usually a microscopic electric field influences only the LO modes while TO modes remain unaltered. The field therefore breaks the Born-von Kármán conditions, as a consequence with a direct method only finite wave vector $`𝐤0`$ calculations are possible. Elongated sub-supercells are needed to recover the $`𝐤0`$ limit of the LO phonon branch (Parlinskin, 2002).
In our result there are two transverse acoustic branches which are highly degenerate and a longitudinal acoustic branch. The first optical branch of the dispersion curves is degenerate with the transverse acoustic branches at energy $``$9.0 meV. The transverse and longitudinal velocities of sound are calculated from the initial slopes of the corresponding transverse and longitudinal acoustic branches of the dispersion curves in the long wavelength limit. The experimental values of velocities reported in Table 1 are those of longitudinal and transverse sound waves propagating along the $`c`$-direction of single crystals of ice at 257 K. It is well known that velocities of sound depend much on the direction of propagation and also on the temperature. Inelastic X-ray scattering data from water at 5 C shows a variation of the sound velocity from 2000 to 3200 m/s in the momentum range of 1-4 $`\mathrm{nm}^1`$. The so-called transition from normal to fast sound in liquid water at $``$ 4 meV, the energy of sound excitations which is equal to the observed second weakly dispersed mode, was reported to be due to the reminiscent of a phonon branch of ice Ih of known optical character (Sette, Ruocco, Krisch, Masciovecchio et al., 1996). We can conclude that our calculated values of longitudinal velocity, $`v_L`$ is in a reasonable range of velocity of sound in ice along the $`[100]`$-direction chosen for our calculation. We must also stress the fact that our phonon dispersions were calculated at 0 K.
## 4.2 Vibrational density of states of crystalline ice
In order to understand the mode of collective vibration of molecules of water in ice from a spectroscopic point of view, we need to consider the three normal modes of an isolated water molecule shown in Fig. 12 as phases of water vapour; liquid water and ice consist of distinct H<sub>2</sub>O molecules recognized by Bernal and Fowler in 1933, which explained one of the factors leading to the Pauli model of the crystal structure of ice. The fact that the forces between the molecules are weak in comparison with the internal bonding results in a simple division of the lattice modes into three groups involving the internal vibrations, rotations, and translations of the molecules. The frequency of the first two groups depend primarily on the mass of the hydrogen or deuterium nuclei, and the frequencies of the translations depend on the mass of the whole molecule (Petrenko and Whitworth, 1999). A free H<sub>2</sub>O molecule has just three normal modes of vibration illustrated in Fig. 12. The comparatively small motions of the oxygen atoms are required to keep the centre of mass stationary, and these motions result in the frequency $`\nu _3`$ being slightly higher than $`\nu _1`$; these depend on the force constant for stretching the covalent O-H bond, while the bending mode $`\nu _2`$ depends on the force constant for changing the bond angle. In the vapour the free molecules have a rich rotation-vibration infrared spectrum (Benedict, Gailar and Plyler, 1956), from which the frequencies of the molecular modes are deduced to be:
$`\nu _1`$ = 3656.65 cm<sup>-1</sup> $``$ 453.4 meV,
$`\nu _2`$ = 1594.59 cm<sup>-1</sup> $``$ 197.7 meV,
$`\nu _3`$ = 3755.79 cm<sup>-1</sup> $``$ 465.7 meV.
For ice the band around 400 meV is thus ratified with the O-H bond stretching modes $`\nu _1`$ and $`\nu _3`$. The frequencies are thus lowered from those of the free molecules by the hydrogen bonding to the neighbouring molecules, but as a single molecule cannot vibrate independently, this coupling also leads to complex mode structures involving many molecules.
We can now discuss the vibrational density of states, VDOS, for H<sub>2</sub>O ice based on the lattice dynamics obtained from the results of our calculation and compare them to some of the well known spectra of ice such as infrared and Rahman spectra and inelastic neutron scattering data. The total VDOS calculated from the phonon dispersions for our ice structure is shown in Fig. 13. Also shown in Fig. 15 is the corresponding partial VDOS for both hydrogen and the oxygen atoms in the ice system. We note that these phonon DOS are not complete since the summation is not done over the whole Brillouin zone, but only in the direction of the cubic symmetry. The distribution of the partial DOS is given by
$`g_{\alpha ,k}(\omega )={\displaystyle \frac{1}{nd\mathrm{\Delta }\omega }}{\displaystyle \underset{𝐤,j}{}}|e_\alpha (k;𝐤,j)|^2\delta _{\mathrm{\Delta }\omega }(\omega \omega (𝐤,j)),`$ (3)
where $`e_\alpha (k;𝐤,j)`$ is the $`\alpha `$-th Cartesian component of the polarization vector for the $`k`$-th atom; $`n`$ is the number of sampling points and $`d`$ is the dimension of the dynamical matrix (Parlinskin, 2002). The total VDOS is calculated by summing all the partial contributions. Figure 13 shows the total VDOS together with the full phonon dispersion curves along the direction. Also shown in Fig. 14 is the enlargement of the intermolecular frequency range on which we superimpose the inelastic neutron-scattering spectra data extracted from Ref. (Cote, Morrison, Cui, Jenkins et al., 2003). The comparison is made with the ice Ih data though our ice structure is not perfectly hexagonal but still there is very little difference between the neutron data for ice Ih and ice geometry used in our calculation in the translational region as explained below. There are well defined separated peaks in the whole range of the vibrations. The illustrative discussion in Fig. 12 can be well understood if we consider the partial DOS in Fig. 15. The covalently O-H stretching mode of both phase and anti-phase, analogous to the frequencies $`\nu _1`$ and $`\nu _3`$ for isolated free water molecules, can be seen clearly in Fig. 15(b) in the energy range (350-410 meV) or frequency range (3010-3400 cm<sup>-1</sup>). We can notice that the collective motion of oxygen is almost static when compared to the collective contributions from the hydrogen atoms. According to the Rahman spectra, a strong peak is observed at 382.3 meV (3083 cm<sup>-1</sup> at 95 K) (Bertie and Whalley, 1967; Whalley and Bertie, 1967) for $`\mathrm{D}_2\mathrm{O}`$. If we take into account the mass difference between deuterium and hydrogen atoms (i.e., isotope effect), the peaks which are observed at 2950, 3000 (very short), 3250, and 3270 cm<sup>-1</sup> are in good range when compared to the experimentally observed values for $`\nu _1`$ and $`\nu _3`$. In the intra-molecular bending region, analogous to the frequency $`\nu _2`$ for an isolated water molecule (1580-1680 cm<sup>-1</sup>), there is an interesting feature. Our results show that, of all the contributions resulting from the collective motion of hydrogen atoms as contribution from collective motion of oxygen atoms is recessive, only one of the components of the collective motions of hydrogen atoms contributes to the intra-molecular bending modes and it is one that is dominant. Figure 15(b) gives an example of such contribution being dominated mainly by the $`y`$-component of the intra-molecular vibration of the O-H. This means that intra-molecular bending of the angular motion takes place mostly in one direction.
Tanaka has identified hydrogen-bond bending modes with negative expansion coefficients associated with this region (Tanaka, 1998). If we go further down to the low frequency region such as 600-1200 cm<sup>-1</sup> called the molecular librational region and then to (0-400 cm<sup>-1</sup>) called the molecular translational region, where we estimated the sound velocities from the corresponding phonon dispersion curve, the VDOS peaks of these modes of vibration agree very well with the experimental observation from inelastic neutron scattering data. The general agreement of the features in the translational optic region is good with all the three distinct peaks present at 400, 270 and 105 cm<sup>-1</sup> (Li, 1996).
## 4.3 The boson peak in ice
The high frequency (0.1-10 THz) or energy (0.4136-41.36 meV) excitations have been experimentally shown to have linear dispersion relations in the mesoscopic momentum region ($``$ 1-10 nm<sup>-1</sup>). So many amorphous materials display low temperature anomalies in their specific heats that they are generally regarded as being universal properties of the glassy state. These anomalies are usually of two kinds. The first concerns observation that, while many crystals obey the Debye law $`CT^3`$ for temperatures less than say, 1K, glasses with the same chemistry frequently display the law $`CT`$ at correspondingly low temperature (Phillips, 1996). This second observation is tied in with the appearance of the ubiquitous Boson peak (BP) in inelastic neutron and Raman spectra (Buchenau, Zhou, Nucker, Gilroy et al., 1987). The name Boson peak therefore refers to the fact that the temperature dependence of its intensity scales roughly with the Bose-Einstein distribution. Dove et. al. demonstrated that the Boson peak arises largely from a flattening of the dispersion of the transverse acoustic modes, in a manner similar to that which occurs in a crystalline material at its Brillouin zone boundary (Dove, Harris, Hannon, Parker et al., 1997). This idea is not new, being originally suggested by (Leadbetter, 1969), and recently developed further by (Elliott, 1996; Taraskin, 1997; Taraskin and Elliot, 1997). A standard way of extracting the boson peak (BP) from the vibrational spectrum is to plot $`g(E)/E^2`$ (as done in Fig. 16(b)), since in the Debye approximation $`g(E)E^2`$ at low energy. In Fig. 16 we show the plot of $`g(E)`$ and the corresponding $`g(E)/E^2`$ in the translational low energy range for the ice geometry in our calculation. According to the Debye law, it is expected that $`g(E)/E^2`$ should be constant for the whole range of energy. This constant relation is only obtained at energies larger than 40 meV in agreement with the experimental observation range for the BP as mentioned above. There is an anomalous sharp peak at 3.5 meV which can be ascribed to the region of low-energy excess vibrational excitation of the so-called BP. The reason for this peak is unknown since we have a crystalline structure for our ice geometry. The peak reveals the anomalous behaviour of hydrogen bonding in the crystal ice which shows similarity in the behaviour as in the results of an inelastic neutron scattering study of a crystalline polymorph of SiO<sub>2</sub> ($`\alpha `$-quartz), and a number of silicate glasses (pure silica, SiO<sub>2</sub>) with tetrahedral coordination (Harris et al., 2000). Also amorphous solids, most supercool liquids and the complex systems show this anomalous character (Grigera, Martin, Parisi and Verrocchio, 2003).
## 5. SUMMARY
Our MD simulation, carried out through the analysis of radial distribution functions of the ice crystal, shows fair agreement in the positions of peaks in comparison with neutron diffraction data.
The phonon dispersion calculations in direction shows a good result especially in the acoustic region in comparison with experiment. Our optic modes are degenerate at the $`\mathrm{\Gamma }`$-point due to application of the direct method used in this calculation which does not take into account the effect of polarization arising from dipole-dipole interactions of water molecules which is expected to yield a significant effect in the splitting of longitudinal and transverse optic modes at the $`\mathrm{\Gamma }`$-point. We intend to include this polarization effect in our future work through the calculation of effective charge tensors using the Berry phase approach to see actually if there will be splitting. Our calculated longitudinal acoustic velocity agrees well with the longitudinal acoustic velocity from inelastic neutron scattering data. The vibrational density of states reproduces all the features in covalently O-H stretching region, intra-molecular bending region, molecular librational region as well as in the molecular translational region, when comparing our results to some infrared spectra, Rahman spectra as well as inelastic neutron scattering results.
The analysis of the vibrational density of states shows a boson peak, a characteristic feature common to amorphous systems, at low energy of the translational region. The anomalous sharp peak of $`g(E)/E^2`$ at 3.5 meV, which can be ascribed to the region of low-energy excess vibrational excitation of the boson peak, might be due to the anomalous behaviour of hydrogen bonding since we have a perfect crystal. The origin of this anomalous behaviour of this abnormal peak in non-crystalline solids and supercool liquids is still subject of scientific debate.
Acknowledgments
We acknowledge the support by the Deutsche Forschungsgemeinschaft (Graduate College “Structure and Dynamics of Heterogeneous Systems”).
We thank Prof. Keith Ross for allowing us to use one of their phonon dispersions for comparison with our results. Also, we acknowledge Prof. Soper for providing us the neutron diffraction data of radial distribution functions of liquid water and ice for fitting our ab-initio simulation results.
TABLES
* Velocities of sound calculated from the initial slope of the phonon dispersion curves of ice in direction compared to the experimental result.
FIGURE CAPTIONS
* The phase diagram of the stable phases of ice.
* The tetrahedral unit from which the hexagonal ice is created.
* Initial and the relaxed geometry of the unit cell of ice. The ice structure was initially packed in a cubic unit cell with initial lattice constant taken from the literature (Lee, Vanderbilt, Laarsonen, Car et al., 1993) to be 6.35 Å. There are no hydrogen bonds in the initial prepared structure shown on the left but were perfectly formed after the relaxation. The relaxed geometry has the values of $`abc`$ which implies that the relaxed structure is tetragonal with $`c/a`$ ratio $``$ 0.988.
* Relaxed crystalline structure of ice produced from the geometry in Fig. 2 by replicating in all direction by the calculated lattice parameters.
* Distribution of H-O-H angle in ice and water (in arbitrary units). The comparison is done for the ice at 220 K and liquid water simulated at room temperature (298 K) and at high temperature in the supercritical regime. The H-O-H angles of ice structure are larger compared to those of liquid water.
* Radial distribution functions $`g_{\mathrm{OO}}`$ and $`g_{\mathrm{OH}}`$ at 100 K and 220 K obtained with VASP. There is a gradual loss of long-range order at 220 K as can be noticed in its 2nd and 3rd peaks of $`\mathrm{g}_{\mathrm{OO}}`$ when compared to the results for 100 K.
* Radial distribution functions $`g_{\mathrm{HH}}`$ at 100 K and 200 K obtained with VASP. There is a gradual “loss of peaks” observed for temperature 100 K at 2.5, 4.4, 4.8 and 5.4 Å due to the loss of long-range order as the temperature increases.
* Radial distribution functions $`g_{\mathrm{OO}}`$ and $`g_{\mathrm{OH}}`$ at 220 K obtained with VASP compared to neutron diffraction scattering data (NDS) (Soper, 2000).
* Radial distribution functions $`g_{\mathrm{HH}}`$ obtained with VASP at 220 K compared to neutron diffraction scattering data (NDS) (Soper, 2000).
* (I) Brillouin zone for the cubic symmetry used in our VASP calculation of the ice system. Our phonon dispersions are calculated in $`k_x`$-direction, (II) The first Brillouin zone for the structure of ice Ih with origin at the point $`\mathrm{\Gamma }`$. $`\mathrm{\Gamma }`$A =$`\frac{1}{2}c^{}`$ and $`\mathrm{\Gamma }`$M =$`\frac{1}{2}a^{}`$, where $`a^{}`$ and $`c^{}`$ are the vectors of the reciprocal lattice (Petrenko and Whitworth, 1999). The dispersion curves are commonly drawn along the lines of symmetry $`\mathrm{\Gamma }`$A, $`\mathrm{\Gamma }`$M, and $`\mathrm{\Gamma }`$K.
* (a) With VASP calculated phonon dispersion curves of ice in direction of the tetragonal unit cell compared to (b) dispersion relations in the translational frequency range for the ice structure plotted along the Cartesian directions from zone center to zone edge in the Brillouin zone of the orthorhombic eight-molecule unit cell (Cote et. al. (Cote, Morrison, Cui, Jenkins et al., 2003)) and (c) the experimental dispersion of D<sub>2</sub>O ice according to Renker’s model (Renker, 1973). The difference in scale of (c) from (a) and (b) is due to the isotopic effect because of difference in the masses of hydrogen and Deuterium atoms.
* The three normal modes of an isolated water molecule. Motion with frequency $`\nu _1`$ can be regarded as symmetric stretch, $`\nu _2`$ as bending and $`\nu _3`$ as anti-symmetric (Petrenko and Whitworth, 1999).
* Total vibrational density of states (VDOS) of ice based on the lattice dynamics together with full phonon dispersion curves. The VDOS show all the important regions such as the intermolecular translational, librational, bending and the stretching frequency range.
* Enlargement of calculated total vibrational density of states in Fig. 13 showing the intermolecular range of frequencies. The broken line is taken from the inelastic neutron scattering data available through Ref. (Cote, Morrison, Cui, Jenkins et al., 2003).
* Partial density of states of ice based on the lattice dynamics. (a) shows the $`x`$, $`y`$ and $`z`$ components of VDOS for the oxygen atoms in the H<sub>2</sub>O ice. (b) shows the $`x`$, $`y`$ and $`z`$ VDOS for corresponding hydrogen atoms for the whole range of frequencies. It is interesting to notice that in the intra molecular bending region (1580-1680 cm<sup>-1</sup>), only one of the components of the VDOS of hydrogen (say $`y`$) dominates while $`x`$ and $`z`$ contribute less. The contribution from the oxygen atoms can be regarded as being completely recessive in this region.
* Plot of the VDOS $`g(E)`$ and the corresponding $`g(E)/E^2`$ vs. $`E`$ for the region of translational mode. The boson peak is found in the low-energy region at 3.5 meV.
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# A matched expansion approach to practical self-force calculations
## 1 Introduction
There is currently considerable interest in calculating the self-force on a charged particle moving in a curved background spacetime. This interest is twofold. On one hand, self-force is an intrinsically interesting phenomenon, apparently both deep and subtle. On the other, there is increasing practical interest in understanding the motion of particles with small but non-negligible masses in curved geometries. In particular, there is wide-spread belief that accurate computation of gravitational self-force is needed to calculate gravitational waveform templates for key data analysis efforts associated with interferometric gravitational wave detectors. With such templates, the Laser Interferometer Space Antenna (LISA) is expected to produce a wealth of information about strongly-field gravity.
We define self-force to be any force on a particle of quadratic (or higher) order in the charge carried by that particle. For electrically charged particles in flat spacetime, the self-force is given by the Abraham-Lorentz-Dirac formula . DeWitt and Brehme derived the general expression for the self-force on an electrically charged particle in a curved background (however, they missed a term which was later supplied by Hobbs). The gravitational self-force was first calculated almost simultaneously by Mino, Sasaki and Tanaka and by Quinn and Wald . Later, Quinn derived the equivalent formula for a charge coupled to a minimally-coupled massless scalar field. These results have resolved many of the issues of principle in computing self-forces in curved spacetime. Poisson has recently written a comprehensive review article in which these expressions and their implications are discussed at length.
The self-force expressions for all fields of physical interest are of similar form. We henceforth restrict our attention, therefore, to the simplest case; a particle with mass $`\mu `$ and scalar charge $`q`$ coupled to a minimally-coupled massless scalar field in a curved background geometry. Lessons learned here should be extendible to other physical fields without major modification.
The field equation for a minimally-coupled massless scalar field $`\varphi `$ is
$$\mathrm{}\varphi =4\pi \rho .$$
(1)
Here, $`\mathrm{}`$ is the D’Alembertian of the curved background and $`\rho `$ is the charge density. We consider a point particle, in which case
$$\rho (x)=q\delta ^4(x,z(\tau ))𝑑\tau ,$$
(2)
where $`\delta ^4()`$ is a generalized Dirac distribution in four dimensions and $`z(\tau )`$ denotes the worldline of the particle.
For such a particle, Quinn has shown that the self-force is given by
$`f^\alpha =q^2[{\displaystyle \frac{1}{3}}(\dot{a}^\alpha a^2u^\alpha )+{\displaystyle \frac{1}{12}}(2R^{\alpha \beta }u_\beta +2R_{\beta \gamma }u^\beta u^\gamma u^\alpha Ru^\alpha )`$
$`+\underset{ϵ0^+}{lim}{\displaystyle _{\mathrm{}}^{\tau ϵ}}^\alpha G_{\mathrm{ret}}(z(\tau ),z(\tau ^{}))d\tau ^{}].`$ (3)
The quantities in equation (3) are defined as follows: $`u^\alpha `$ is the four-velocity of the particle and $`a^\alpha =u^\beta _\beta u^\alpha `$ its four-acceleration. $`\dot{a}^\alpha =u^\beta _\beta a^\alpha `$ denotes derivative of the acceleration with respect to its proper time, and $`a^2=a^\alpha a_\alpha `$ is the magnitude of the acceleration squared. The quantity $`R_{\alpha \beta }`$ is the Ricci tensor of the background spacetime, and $`R`$ is its scalar curvature. $`\tau `$ is the proper time of the particle at its current position, while $`\tau ^{}`$ denotes the proper time at any other point along the particles worldline. Finally, $`G_{\mathrm{ret}}(x,x^{})`$ is the retarded Green’s function for the scalar field equation (1), which satisfies the Green’s function equation
$$\mathrm{}G_{\mathrm{ret}}(x,x^{})=4\pi \delta ^4(x,x^{}),$$
(4)
and has support only when $`x^{}`$ is in the causal past of $`x`$.
Notice that the terms on the right-hand-side of (3) can be gathered into three groups. The first group are local terms involving the acceleration of the particle. These give the Abraham-Lorentz-Dirac force on an accelerating particle in flat spacetime. They arise because the accelerating particle radiates, and that radiation produces a “radiation reaction” or recoil force on the particle. It is interesting to note that for a particle in Minkowski spacetime with constant $`a`$ (whose worldline is a hyperbola), that $`\dot{a}^\alpha 0`$ (because the direction of the four-vector changes). In fact, $`a^\alpha =a^2u^\alpha `$ in this case, and the first group vanishes. This fact, which might seem surprising at first, can be understood as a consequence of the equivalence principle. We further note that for a particle freely falling in curved space (i.e. following a geodesic), $`a^\alpha =0`$, and thus these terms will also vanish.
The second group are also local terms, this time involving the background curvature. These terms involve only the Ricci curvature, and thus represent a self-force mediated by the matter content of the background. Interestingly, it is not necessary for the matter to interact directly with the particle for this to be true. We will largely ignore these first two groups of terms for the remainder of our discussion because for a particle following a geodesic in a vacuum background spacetime, which is the case of most practical interest, they vanish. Moreover, even when they do not vanish, they are easily calculated.
In contrast, there is considerable practical difficulty in calculating the single non-local term which constitutes the third group. Interestingly, this non-local self-interaction arises in curved backgrounds, but not in a flat background. This is because of two ways in which the propagation of massless fields differs in curved and flat backgrounds. In Minkowski space, massless fields propagate along null geodesics. Thus, a particle would have to be null-separated from some point on its past worldline to affect itself. However, the simple causal structure of Minkowski space does not allow this for a massive particle. The particle is restricted to a time-like geodesic, and no two points in Minkowski space can be connected both by a time-like geodesic and a null geodesic. Thus, any point with which the charged particle can interact, it cannot travel to, and vice versa. This is not true in a curved background, however, where the causal structure can be considerably more complicated. In this case, it is possible for the field, which “leaves” the particle along a null geodesic, to re-intersect that particle, which follows a timelike geodesic, at a later time.
If this were the only mechanism by which non-local self-interactions arose, then there would be no self-interactions arising from within the normal neighbourhood of the point at which the particle sits. Recall that the normal neighbourhood of a point $`x`$ is the set of all other points which are connected to $`x`$ via a unique geodesic. Recall also that such a neighbourhood is guaranteed to exist. Thus, if a particle at $`x`$ is massive and non-accelerating (i.e. following a timelike geodesic), then any point on the particle’s worldline connected to $`x`$ by a null geodesic is connected to $`x`$ by (at least) two geodesics. Therefore, by definition, it is not in the normal neighbourhood of $`x`$.
However, in general, fields do not propagate only along null geodesics in curved backgrounds. This is apparent from the Hadamard form of the retarded Green’s function,
$$G_{\mathrm{ret}}(x,x^{})=\mathrm{\Theta }[tt^{}]\left\{U(x,x^{})\delta [\sigma (x,x^{})]V(x,x^{})\mathrm{\Theta }[\sigma (x,x^{})]\right\},$$
(5)
where $`\mathrm{\Theta }[]`$ is the Heaviside step function, $`\delta []`$ is the Dirac delta distribution, $`\sigma (x,x^{})`$ is one half of the square of the geodesic distance between points $`x`$ and $`x^{}`$, and $`U(x,x^{})`$ and $`V(x,x^{})`$ are smooth functions which depend on the details of the background and field equation. Note that this expression is only well defined if one can unambiguously define geodesic distance between $`x`$ and $`x^{}`$, which implies that $`x^{}`$ is within the normal neighbourhood of $`x`$.
Recall that the Green’s function gives the field at position $`x`$ due to a charge at position $`x^{}`$. Now, the first term on the right-hand-side of equation (5), which is known as the direct part of the Green’s function, has support only when the geodesic distance between $`x`$ and $`x^{}`$ vanishes, or, in other words, when $`x`$ and $`x^{}`$ are separated by a null geodesic. This term is always present, but, as described above, cannot contribute to the self-force. However, the second term, which is known as the tail part of the Green’s function, does contribute whenever the square of the geodesic distance between $`x`$ and $`x^{}`$ is negative, or, in other words, when $`x`$ and $`x^{}`$ are separated by a timelike geodesic<sup>1</sup><sup>1</sup>1This fascinating fact was first elaborated upon by Hadamard, who described it as a failure of Huygens’ principle to hold in general for hyperbolic partial differential equations - c.f. the Paul Günther memorial edition of Zeitschrift für Analysen und ihre Anwendungen for a collection of articles addressing the proof of the modified Hadamard’s conjecture, which states that the only spacetimes for which standard wave equations obey Huygen’s principle are those conformal to Minkowski spacetime and one plane-wave family of spacetimes.. Since every point on the particles past worldline is timelike separated from the particle by a timelike geodesic (to wit, the worldline), through this term a self-force can be generated at every time within the normal neighbourhood.
Given that all points on a particles past worldline can interact with the particle, it may seem somewhat arbitrary to have singled out in the above discussion points on the past worldline which are null-separated from the particle. Indeed, in the literature, it is often stated that the self-force arises from the tail part of the field (or Green’s function) and left at that. This statement can be somewhat confusing, however. Outside of the normal neighbourhood, there is no clear distinction between tail and direct parts - is the field from a point connected by both timelike and null geodesics a direct field, a tail field, or both?
Furthermore, there is reason to believe that these points play a qualitatively different role in the self-interaction. In particular, when doing the integral over the Green’s function, distributions (i.e. $`\delta `$-functions and step functions) will be encountered when the source point and field point are null separated. In effect, the particle can “feel” its own direct field “sent” from points in the past. This is depicted in figure 1. In the specific case of the $`O[M]`$ Green’s function, such distributions do appear, and their contribution make up the entire self-force (see equation (27)).
Let us turn now to the matter of calculating the non-local part of the self-force,
$$f^\alpha =q^2\underset{ϵ0^+}{lim}_{\mathrm{}}^{\tau ϵ}^\alpha G_{\mathrm{ret}}(z(\tau ),z(\tau ^{}))𝑑\tau ^{}.$$
(6)
The key feature of this formula is that the integral is well behaved over the entire region of integration, i.e. it is integrable throughout its domain. The termination of the integral at $`\tau ϵ`$ (as opposed to integrating all the way to $`\tau `$) is a required for this to be true, however, because the Green’s function has an unintegrable singularity at the coincidence limit ($`ϵ=0`$). Nonetheless, this contribution is distributional, and the integrand is perfectly well behaved as long as we approach $`\tau ^{}=\tau `$ from the past.
At first glance, evaluating equation (6) seems like a relatively straightforward task - an approximate retarded Green’s function can be calculated, for instance, using mode-sum techniques. However, in practice, we can only sum a finite number of modes and thereby obtain an approximate Green’s function. Furthermore, as we will demonstrate in section 3, the number of modes needed to obtain a given accuracy grows without bound as epsilon approaches zero. This does not preclude the use of modes to calculate the non-local part of the self-force, and, as can be seen in the pages of this special issue, modes are indeed widely used. The modes must, however, first be regularized by some method.
Nonetheless, it is our purpose in this paper to explore an alternative method which may have some advantages over a regularized mode sum. It does not need to make use of regularized modes, although one might choose to do so. Rather, it is a method of matched expansions, in which one calculates the self-force using the tail within some portion of the normal neighbourhood of the particle and using an unregularized mode expansion for the remainder of the particle’s worldline. More precisely, we propose to express equation (6) as
$$f^\alpha =q^2_{\tau \mathrm{\Delta }\tau }^\tau ^\alpha V(z(\tau ),z(\tau ^{}))𝑑\tau ^{}+q^2_{\mathrm{}}^{\tau \mathrm{\Delta }\tau }^\alpha G_{\mathrm{ret}}(z(\tau ),z(\tau ^{}))𝑑\tau ,$$
(7)
where $`\mathrm{\Delta }\tau `$ is an interval of proper time. We require that $`\mathrm{\Delta }\tau `$ be chosen so that $`z(\tau ^{})`$ is within the normal neighbourhood of $`z(\tau )`$ for all $`\tau \mathrm{\Delta }\tau <\tau ^{}<\tau `$. In other words, $`\mathrm{\Delta }\tau `$ distinguishes the contribution to the self-force coming from the recent history of the particle from the contribution that comes from the more distant past. It is the choice of $`\mathrm{\Delta }\tau `$ and the feasibility of evaluating the two integrals that will occupy us for the remainder of this paper.
There are several notable features of equation (7). First, this expression is exactly equivalent to equation (6), since, as is evident from equation (5), $`G_{\mathrm{ret}}(x,x^{})V(x,x^{})`$ provided we restrict ourselves to the interior of the past light cone, as is always the case for this integral. Second, there is no longer a limit needed because $`V(x,x^{})`$ is regular everywhere. Finally, notice that we now have a new parameter, $`\mathrm{\Delta }\tau `$, which we are free to choose so long as it is not too large.
One might suspect that the contribution to the Green’s function “falls off” fast enough, in general, so that perhaps only the first integral needs to be evaluated; and, since it is restricted to the normal neighbourhood, only the Hadamard expansion is needed to compute the self-force. Surprisingly, there are simple situations in which this line of reasoning turns out to be false; the second integral gives a significant (perhaps, the dominant) contribution to the total force. We demonstrate this in section 3 using a Green’s function for spacetime with a large central mass $`M`$ where we only keep the terms of leading order in the central mass (the “O\[M\] Green’s function”). Using the O\[M\] Green’s function, none of the force originates from the tail contribution in the normal neighbourhood; all of the force comes times prior to the light reflection time.
Unfortunately, neither integrand in equation (7) can be calculated exactly. As mentioned above, $`G_{\mathrm{ret}}(x,x^{})`$ can be approximated in the second integrand by a mode sum expansion. The advantage here is that the mode sum does not need to be extended to the limit of the particle’s position. Thus, for any fixed finite precision required, a finite number of modes will be needed for the second integral in equation (7). As noted above, however, that number will grow as $`\mathrm{\Delta }\tau `$ decreases.
On the other hand, for $`V(x,x^{})`$, we can take advantage of the fact that we are within the normal neighbourhood, and can therefore define a Riemann normal coordinate system. In such a coordinate system, it is relatively straightforward to expand $`V(x,x^{})`$ in a covariant Taylor series. Such a series will have coefficients constructed of geometric quantities (notably, the particle’s four-velocity and the curvature of the background) evaluated at the particles position, and will be an expansion in geodesic distance from that position. Again, only a finite number of terms will be available, so only an approximation of $`V(x,x^{})`$ can be calculated. In contrast to the mode sum expansion of $`G_{\mathrm{ret}}(x,x^{})`$, the covariant Taylor series will, in general, require more terms to achieve a given accuracy as $`\mathrm{\Delta }\tau `$ increases.
The goal, then, is to calculate both series to sufficient accuracy within their applicable domains, and to then integrate and sum. We may take advantage or our parameter, $`\mathrm{\Delta }\tau `$, to adjust the number of terms required in each series to achieve the required accuracy. However, both series are non-trivial to calculate, and the difficulty of calculating each successive term is greater than for the last. There is, therefore, no guarantee that there is any value of $`\mathrm{\Delta }\tau `$ such that the number of terms needed in both series can be calculated in practice.
It is the main goal of the remainder of this paper to investigate if it is plausible that there would exist a choice of $`\mathrm{\Delta }\tau `$ which would make this scheme, first proposed by Poisson and Wiseman, viable for practical calculations. Because of the difficulty of performing both expansions, we will exploit existing results when possible. Furthermore, we will use the simplest and most transparent results that still bear upon the problem. We will not explicitly try to calculate the self-force for any specific geometry or particle motion, but will, rather, simply investigate the convergence of both series expansions.
The plan of the remainder of this paper is as follows: we will review and examine the existing results for the expansion of the $`V(x,x^{})`$ in the next section. Following that, in section 3, we will investigate the mode expansion. Finally, we will discuss our conclusions in section 4.
## 2 The quasi-local part of the self-force
Covariant normal neighbourhood expansions have a venerable history for calculations of both self-forces and quantum field effects in curved backgrounds. There is therefore a plethora of material from which to begin a calculation of the first term in equation (7), which we will call the quasi-local part of the self-force
$$f_{\mathrm{QL}}^\alpha =q^2_{\tau \mathrm{\Delta }\tau }^\tau ^\alpha V(z(\tau ),z(\tau ^{}))𝑑\tau ^{}.$$
(8)
There are, however, to our knowledge, only three papers in the literature which address the calculation of this part of the self-force. The first was by Roberts, in which he stopped just short of calculating the quasi-local part of the self-force for an electric charge in an arbitrary background. A trivial extension of his paper gives the leading order to the quasi-local part of the self-force in this case,
$$f_{\mathrm{QL}}^\alpha =\frac{q^2}{4}(\delta _\alpha ^\beta +u_\alpha u^\beta )C^\alpha {}_{\beta \gamma }{}^{\delta }{}_{;\delta }{}^{}u_{}^{\beta }u^\gamma \mathrm{\Delta }\tau ^2+O(\mathrm{\Delta }\tau ^3).$$
(9)
Here, $`C_{\alpha \beta \gamma \delta }`$ is the Weyl curvature tensor of the background at the particle’s position, $`u^\alpha `$ is the particle’s four-velocity, and $`\mathrm{\Delta }\tau `$, as previously, is a proper time interval along the particle’s past worldline.
The second result is due Anderson and Hu. They calculated the tail of the retarded Green’s function in the normal neighbourhood for a particle with a minimally-coupled massless scalar field in a Schwarzschild background geometry. By using a Hadamard-WKB expansion for the Euclidean Green’s function, they are able to expand $`V(x,x^{})`$ to sixth order in the geodesic separation. They do not go on to calculate the quasi-local self-force explicitly (although that is clearly the motivation for their paper), but it is a relatively straightforward matter to do so. Because their results are somewhat unwieldy, we refer the reader to their paper for the actual result, although we will be reproducing the self-force derived from it for two simple particle motions below.
Most recently, the quasi-local part of the gravitational self-force for a particle in an arbitrary curved background has been calculated by Anderson, Flanagan and Ottewill. They calculate the first two non-vanishing terms in the Taylor series, and find
$`f_{QL\alpha }(\tau ,\mathrm{\Delta }\tau )=\mu ^2(\delta _\alpha ^\beta +u_\alpha u^\beta )C_{\beta \gamma \delta \epsilon }C_{\sigma \rho }^{\gamma \epsilon }u^\delta u^\sigma u^\rho \mathrm{\Delta }\tau ^2`$
$`+\mu ^2(\delta _\alpha ^\beta +u_\alpha u^\beta )u^\gamma u^\delta [{\displaystyle \frac{1}{6}}C_{\gamma \mu \delta \nu }C_{\epsilon \sigma ;\beta }^{\mu \nu }u^\epsilon u^\sigma {\displaystyle \frac{3}{20}}C_{\beta \gamma \mu \delta ;\nu }C_{\epsilon \sigma }^{\mu \nu }u^\epsilon u^\sigma `$
$`+{\displaystyle \frac{1}{3}}\left({\displaystyle \frac{1}{2}}C_{\mu \nu \gamma \lambda }C_{\delta ;\beta }^{\mu \nu \lambda }+C_{\mu \epsilon \gamma \lambda }C_{\sigma \delta ;\beta }^{\mu \lambda }u^\epsilon u^\sigma \right)`$
$`{\displaystyle \frac{19}{60}}({\displaystyle \frac{1}{2}}C_{\mu \nu \gamma \lambda }C_{\delta \beta }^{\mu \nu ;\lambda }+C_{\mu \epsilon \gamma \lambda }C_{\sigma \delta \beta }^{\mu ;\lambda }u^\epsilon u^\sigma )]\mathrm{\Delta }\tau ^3+O(\mathrm{\Delta }\tau ^4).`$ (10)
They also provide explicit expressions for some particle motions in Schwarzschild and Kerr geometries.
For the purposes of evaluating the convergence of these expansions, we will focus on the results of Anderson and Hu, both because the scalar field case is inherently the simplest and because the provide the highest order expansion with which to work. Unfortunately, Anderson and Hu do not take the derivative of $`V(x,x^{})`$ necessary to calculate the self-force contribution. Furthermore, their result is expressed in terms of coordinate expansion rather than an expansion in geodesic distance. It is only for relatively simple particle motions that one can easily recast coordinate distance into geodesic distance. Fortunately, it is in the spirit of our explorations here to take such simple cases, and we shall.
The fundamental question we ask in this section, then, is “what is the rate of convergence of the expansion for the quasi-local part the self-force as a function of $`\mathrm{\Delta }\tau `$?” Clearly, as $`\mathrm{\Delta }\tau 0`$, the expansion becomes exact at any order. The more interesting limit is the one where $`\mathrm{\Delta }\tau `$ approaches the boundary of the normal neighbourhood. It might be reasonable to expect that the expansion ceases to converge at all in that case, since that is the boundary of its domain of validity. Indeed, there is no fundamental reason that the domain of convergence of the series expansion could not be much smaller than the normal neighbourhood. Addressing this question, however, is slightly complicated by the fact that it is not immediately clear what the value of $`\mathrm{\Delta }\tau `$ at the boundary of the normal neighbourhood is.
In the case of a Schwarzschild black hole spacetime, at least, one can get some insight into the normal neighbourhood boundary by studying null geodesics intersecting a particle in a circular geodesic at radius $`R`$. For such a particle, the angular displacement is related to the particle’s proper time by
$$\mathrm{\Delta }\varphi _{\mathrm{particle}}=\sqrt{\frac{M}{R3M}}\frac{\mathrm{\Delta }\tau }{R},$$
(11)
where $`M`$ is the mass of the black hole. On the other hand, the angular deflection (as seen at $`r=R`$) of a null geodesic passing within a nearest distance $`R_0>3M`$ of the black hole is
$$\mathrm{\Delta }\varphi _{\mathrm{photon}}=2_{1/R}^{1/R_0}𝑑u\sqrt{\frac{R_0^3}{R_02MR_0^3u^2(12Mu)}},$$
(12)
(note that if $`R_03M`$, there is no turning point for the geodesic and it does not return to larger radii). Finally, the proper time, as measured by the particle, for a photon to descend from a distance $`R`$ to within a distance $`R_0`$ of the black hole and return to a distance $`R`$ is given by
$$\mathrm{\Delta }\tau =2\left(\frac{R_02M}{R0}\right)_{R_0}^R𝑑r\frac{r}{r2M}\sqrt{\frac{r^3(R_02M)}{rR_0(r^2R_0^2)2M(r^3R_0^3)}}.$$
(13)
What we would like to find is the minimum value of $`R_0`$ (and hence $`\mathrm{\Delta }\tau `$) for which a null geodesic intersecting the particle’s circular geodesic at a given $`R`$ can re-intersect it at a the same $`R`$ but a different time. For a particle orbiting at $`R3M`$ from the black hole, the condition for this to occur is
$$\mathrm{\Delta }\varphi _{\mathrm{particle}}+\mathrm{\Delta }\varphi _{\mathrm{photon}}=2\pi .$$
(14)
Assuming this and substituting equations (11 \- 13) into equation (14), we can choose a value of $`R`$ and solve equation (14) numerically for $`R_0`$. We give some values obtained in this way in the table 1.
There are a number of noteworthy features of table 1. First, at at approximately $`R_0=3.2M`$, the angular deflection of the null geodesic is $`2\pi `$, and this therefore represents the value of $`R_0`$ corresponding to a particle orbit at $`R\mathrm{}`$. Next, we note that the values of $`\mathrm{\Delta }\tau `$ take values between $`3R`$ for the closest orbits and $`2R`$ for the most distance orbits, asymptoting to $`\mathrm{\Delta }\tau =2R`$ in the limit $`R\mathrm{}`$. In other words, for orbits at large distances the time is dominated by the time for the photon to reach the black hole and return.
The $`\mathrm{\Delta }\tau `$’s quoted in table 1 represent approximate upper bounds on the extent of the normal neighbourhood along the past worldline of the particle since they demarcate two intersections of the particle’s geodesic with a null geodesic. Thus, for a particle at a fixed radius $`R`$, we need not worry about the convergence of the expansion of the quasi-local part of the self-force beyond $`2R`$ to the past.
Let us now consider two such expansions. First, we consider the expansion for a static particle at radius $`R`$ coupled to a minimally-coupled massless scalar field by scalar charge $`q`$. We take the expression for $`V(x,x^{})`$ given by Anderson and Hu and take partial derivatives with respect to the Schwarzschild coordinates. Next, we convert these coordinates into proper time using the coordinate parameterization
$$\mathrm{\Delta }\varphi =0,\mathrm{\Delta }\theta =0,\mathrm{\Delta }r=0,\mathrm{\Delta }t=\sqrt{1\frac{2M}{r}}(\tau \tau ^{}),$$
(15)
which is appropriate for a static particle. We also, without loss of generality, set $`\theta =\pi /2`$. We can then integrate the expansion of $`_\alpha V(x,x^{})`$ with respect to proper time $`\tau `$ to obtain
$`f_{\mathrm{QL}}^r={\displaystyle \frac{9}{2240}}{\displaystyle \frac{q^2M^2}{R^{15}}}\left(4R11M\right)\left(R2M\right)^5\mathrm{\Delta }\tau ^5`$
$`+{\displaystyle \frac{1}{3360}}{\displaystyle \frac{q^2M^2}{R^{20}}}\left(20R^3195MR^2+598M^2R585M^3\right)\left(R2M\right)^6\mathrm{\Delta }\tau ^7`$ (16)
$`+O(\mathrm{\Delta }\tau ^8),`$
which is the only non-vanishing component of the quasi-local part of the self-force to the order of this expansion.
Denote the fifth and seventh order terms in equation (16) as
$`f_{\mathrm{QL}}^r[5]{\displaystyle \frac{9}{2240}}{\displaystyle \frac{q^2M^2}{R^{15}}}\left(4R11M\right)\left(R2M\right)^5\mathrm{\Delta }\tau ^5,`$ (17)
$`f_{\mathrm{QL}}^r[7]{\displaystyle \frac{1}{3360}}{\displaystyle \frac{q^2M^2}{R^{20}}}\left(20R^3195MR^2+598M^2R585M^3\right)\left(R2M\right)^6\mathrm{\Delta }\tau ^7.`$ (18)
Then the fractional truncation error which is induced by only taking terms in the series expansion of the quasi-local part of the self-force to order $`\mathrm{\Delta }\tau ^7`$ can be estimated as
$$\epsilon \frac{f_{\mathrm{QL}}^r[7]}{f_{\mathrm{QL}}^r[5]+f_{\mathrm{QL}}^r[7]}.$$
(19)
This gives an estimate of the upper limit on the local truncation error (the error in truncating the next term in the series) rather than the more desirable global truncation error (the error in truncating all remaining terms in the series). Nonetheless, it is a standard measure of truncation error for this kind of analysis where neither the exact solution nor a form for the general term in the series is known, and is in any case the best estimate of truncation error we have for this series.
As noted previously, the truncation error can be expected to grow with $`\mathrm{\Delta }\tau `$. We have calculated the error as a function of $`\mathrm{\Delta }\tau `$ for particles located at $`r=6M`$, $`10M`$, $`20M`$ and $`100M`$. The results are presented in figure 2.
We see that, despite our concern that the quasi-local expansion might fail to converge at the boundary of the normal neighbourhood, it does, in fact, appear to converge everywhere within the normal neighbourhood. Furthermore, it converges quite well for particles at all radii, with estimated fractional error less than 0.1, back to almost $`\mathrm{\Delta }\tau =10M`$. For $`R=6M`$, which is the closest particle that we consider and also the case for which the self-force should be most important in calculating templates for LISA, the fractional error is less than $`0.1`$ more than half way out to the boundary of the normal neighbourhood.
We can apply exactly the same method to a particle in a circular geodesic orbit at radius $`R`$ around Schwarzschild. In this case, the coordinates are related to proper time by
$$\mathrm{\Delta }\theta =0,\mathrm{\Delta }r=0,\mathrm{\Delta }\varphi =\frac{1}{R}\sqrt{\frac{M}{R3M}}(\tau \tau ^{}),\mathrm{\Delta }t=\sqrt{\frac{R}{R3M}}(\tau \tau ^{}).$$
(20)
Again, we have set $`\theta =\pi /2`$. In this case, there are three non-vanishing components of the quasi-local part of the self-force
$`f_{\mathrm{QL}}^t={\displaystyle \frac{q^2M^3}{R^{12}}}\sqrt{{\displaystyle \frac{R}{R3M}}}{\displaystyle \frac{1}{\left(R3M\right)^3}}`$
$`\times [{\displaystyle \frac{3}{2240}}R^3(R2M)(R3M)(5R19M)\mathrm{\Delta }\tau ^4`$
$`+{\displaystyle \frac{1}{13440}}\left(56R^4728R^3M+3461R^2M^26990M^3R+5073M^4\right)\mathrm{\Delta }\tau ^6`$
$`+O(\mathrm{\Delta }\tau ^7)]`$ (21)
$`f_{\mathrm{QL}}^r={\displaystyle \frac{q^2M^2}{R^{14}}}{\displaystyle \frac{\left(R2M\right)}{\left(R3M\right)^3}}`$
$`\times [{\displaystyle \frac{9}{11200}}R^3(R3M)(20R^3189MR^2+558M^2R523M^3)\mathrm{\Delta }\tau ^5`$
$`+{\displaystyle \frac{1}{47040}}(\begin{array}{c}280R^54690R^4M+30780R^3M^2\hfill \\ 97302M^3R^2+147777RM^486481M^5)\mathrm{\Delta }\tau ^7\hfill \end{array}`$ (24)
$`+O(\mathrm{\Delta }\tau ^8)]`$ (25)
$`f_{\mathrm{QL}}^\varphi ={\displaystyle \frac{q^2M^2}{R^{13}}}\sqrt{{\displaystyle \frac{M}{R3M}}}{\displaystyle \frac{\left(R4M\right)\left(R2M\right)}{\left(r3M\right)^3}}`$
$`\times [{\displaystyle \frac{9}{2240}}R^3(R3M)(3R7M)\mathrm{\Delta }\tau ^4`$
$`+{\displaystyle \frac{1}{13440}}\left(70R^3685MR^2+1968M^2R1731M^3\right)\mathrm{\Delta }\tau ^6`$
$`+O(\mathrm{\Delta }\tau ^7)]`$ (26)
We compute the estimated fractional error for each of these components in the same manner as we did for the static particle. The results are presented in figures 3, 4 and 5.
Again, we see convergence everywhere within the normal neighbourhood for all components of the quasi-local part of the self-force, and good convergence up to fairly high values of $`\mathrm{\Delta }\tau `$.
## 3 The contribution to the self-force from the distant past
In the previous section, we examined the contribution to the self-force from the portion of the worldline within the normal neighbourhood of the field point at $`z(\tau )`$, i.e. from the first integral in equation (7). In this section, we will examine the contribution from the earlier part of the trajectory. In particular, we will examine the feasibility of computing the second integral in equation (7) using a mode sum expansion for the Green’s function.
### 3.1 The O\[M\] Green’s function
In studying the forces on a freely falling electric charge in a curved spacetime, DeWitt and DeWitt developed some very clever techniques for finding the Green’s function for a spherically symmetric spacetime with mass $`M`$ at the centre. They give an approximate expression – accurate to leading order in the central mass – for the Green’s function for an electric charge in this Schwarzschild-like spacetime. This method was later extended by Wiseman and Pfenning and Poisson to the case of a scalar charge. The Green’s function is found to be
$`G(x,x^{})={\displaystyle \frac{\delta [tt^{}|𝐱𝐱^{}|]}{|𝐱𝐱^{}|}}`$
$`+M{\displaystyle \frac{}{t^{}}}\left\{2{\displaystyle \frac{\delta [tt^{}|𝐱𝐱^{}|]}{|𝐱𝐱^{}|}}\mathrm{ln}\left({\displaystyle \frac{r+r^{}+(tt^{})}{r+r^{}(tt^{})}}\right)4{\displaystyle \frac{\mathrm{\Theta }[(tt^{})(r+r^{})]}{(tt^{})^2|𝐱𝐱^{}|^2}}\right\}`$
$`+O[M^2],`$ (27)
where $`(t,𝐱)`$ is the field-point and $`(t^{},𝐱^{})`$ is the source point. The first term is clearly identifiable as the retarded Green’s function in Minkowski (flat) spacetime. The logarithmic term is a correction to the light cones<sup>2</sup><sup>2</sup>2This term reflects the fact that the light cones are bent in curved spacetime. It was not included in the DeWitt-DeWitt calculation.. The final term is the “tail” of the Green’s function. Notice that this term only contributes prior to $`t^{}=t(r+r^{})`$, the light reflection time. This is depicted in figure 1. This gives a strong indication that the self-force is dependent on the portion of the worldline that is outside the normal neighbourhood.
Multiplying by spherical harmonics and integrating over the solid angle, we can obtain the angular mode decomposition of this Green’s function
$`G(x,x^{})=`$ $`{\displaystyle \underset{l=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(2l+1)}{2rr^{}}}\{[1+2M{\displaystyle \frac{}{t^{}}}\mathrm{ln}\left({\displaystyle \frac{r+r^{}+(tt^{})}{r+r^{}(tt^{})}}\right)]`$ (28)
$`\times P_l(\xi )\mathrm{\Theta }[tt^{}|rr^{}|]\mathrm{\Theta }[r+r^{}(tt^{})]`$
$`+4M{\displaystyle \frac{}{t^{}}}Q_l(\xi )\mathrm{\Theta }[t+t^{}(r+r^{})]\}P_l[\mathrm{cos}(\gamma )]`$
$`+O[M^2],`$
where $`P_l`$ and $`Q_l`$ represent Legendre functions and
$`\xi ={\displaystyle \frac{r^2+r_{}^{}{}_{}{}^{2}(tt^{})^2}{2rr^{}}}=\mathrm{cos}\beta `$ (29)
$`\mathrm{cos}\gamma =\mathrm{cos}\theta cos\theta ^{}+\mathrm{sin}\theta \mathrm{sin}\theta ^{}\mathrm{cos}(\varphi \varphi ^{})`$ (30)
and $`(t,r,\theta ,\varphi )`$ and $`(t^{},r^{},\theta ^{},\varphi ^{})`$ are the spherical coordinates of the field-point and source-point respectively. The angle $`\beta `$ is shown in figure 6. In equation (28), it is understood that the first partial derivative not only operates on the natural log, but also on $`P_l(\xi )`$ and the $`\mathrm{\Theta }`$-functions.
In numerical implementations, we would only be able to include a finite number of multipoles. In this case, we can avoid doing a numerical summation by computing the partial sum analytically, i.e.
$`G(N;t,r,t^{},r^{},\mathrm{cos}\gamma )={\displaystyle \underset{l=0}{\overset{N}{}}}\text{terms in equation (}\text{28}\text{)}`$
$`={\displaystyle \frac{1}{2rr^{}}}\{[1+2M{\displaystyle \frac{}{t^{}}}\mathrm{ln}\left({\displaystyle \frac{r+r^{}+(tt^{})}{r+r^{}(tt^{})}}\right)]`$
$`\times \left[(N+1){\displaystyle \frac{P_{N+1}(\xi )P_N(\mathrm{cos}\gamma )P_N(\xi )P_{N+1}(\mathrm{cos}\gamma )}{\xi \mathrm{cos}\gamma }}\right]\mathrm{\Theta }[tt^{}|rr^{}|]\mathrm{\Theta }[r+r^{}(tt^{})]`$
$`+4M{\displaystyle \frac{}{t^{}}}\left[{\displaystyle \frac{1+(N+1)\left(P_{N+1}(\mathrm{cos}\gamma )Q_N(\xi )P_{N+1}(\mathrm{cos}\gamma )Q_{N+1}(\xi )\right)}{\xi \mathrm{cos}\gamma }}\mathrm{\Theta }[tt^{}(r+r^{})]\right]\}`$
$`+O[M^2].`$ (31)
As above, the first partial derivative acts on the natural log as well as the Legendre functions and the $`\mathrm{\Theta }`$-functions.
This formula can now be used in the second integral in equation (7). As in Section 2, we focus our attention on a static source at radius $`R`$ and $`\mathrm{cos}\gamma =1`$ (see figure 6). Notice that the final term – the tail term – is a total derivative. When this is substituted into equation (7), we can integrate by parts and see immediately that the tail contribution to the force vanishes. This is just a re-confirmation of the well known result that there is no self-force on a static scalar charge in Schwarzschild spacetime (c.f. Wiseman ). By the same argument, the natural log term also gives no contribution to the force integral.
Substituting the remaining term – the flat-space portion of equation (31) – into equation (7), a straightforward calculation reveals that the first $`N`$ multipoles of the source give a radial component of the force
$`f_r(\mathrm{\Delta }\tau ,N)={\displaystyle \frac{q^2}{R^2u^t}}{\displaystyle \frac{(N+1)}{4\sqrt{2}}}{\displaystyle _1^{\mathrm{cos}\beta _{\mathrm{max}}}}{\displaystyle \frac{P_{N+1}(\xi )P_N(\xi )}{(1\xi )^{3/2}}}𝑑\xi ,`$ (32)
where $`u^t=1M/b+O[M^2]`$ is the leading order contribution to the time-component of the four-velocity of the static charge. The upper limit is defined as
$`\mathrm{cos}\beta _{\mathrm{max}}1{\displaystyle \frac{1}{2}}\left({\displaystyle \frac{\mathrm{\Delta }\tau }{u^tb}}\right)^2.`$ (33)
As can be seen from figure 6, when the field point is on the worldline of the charge, $`\mathrm{cos}\beta `$ can be used to parameterise the portion of the worldline from the field point ($`\mathrm{cos}\beta =1`$) backwards to the time when the past null cone of the field point first intersects the cylinder around the central mass on which the worldline lies ($`\mathrm{cos}\beta =1`$). As we will discuss below, this is the only portion of the worldline which can contribute to the self-force in this case. In our case, we do not integrate to the field point, but rather to some time $`\mathrm{\Delta }\tau `$ to the past of it as per figure 1, thus terminating our the integral at $`\mathrm{cos}\beta _{\mathrm{max}}`$. This is indicated on the left hand side of equation (32) by the argument $`\mathrm{\Delta }\tau `$.
The only portion of the Green’s function that contributed to equation (32) was the flat-space term. Since the static point particle would never intersect the past light cone of a field point located on its own worldline, we can ask: why doesn’t $`f_r(\mathrm{\Delta }\tau ,N)`$ vanish identically? The answer lies in the fact that we are using only a finite number of modes to describe the field gradient.
The decomposition of the field into angular modes can also be thought of as a mode decomposition of the source. In the case of the static point particle, we are, in effect, replacing the point charge with a sum of spherical shells of charge where each shell has radius $`R`$ and an angular distribution of charge smeared on it. The point particle is only recovered when a large (infinite) number of shells of charge are summed. As is evident from figure 6, although the point charge doesn’t come into null contact with the source point, the spherical shells of charge do intersect the past light cone of the field point. If we terminate the integral at some $`\mathrm{\Delta }\tau >0`$, but include an infinite number of terms in the summation (i.e. $`N\mathrm{}`$), we will recover the point-particle nature of the source, and therefore there would be no contribution to the force. However, in any practical calculation, we will need to terminate at some finite value of $`N`$, and thus we will be left with an unwanted, and inescapable, contribution to the force. The key question is can we include enough terms in the summation to squeeze this unwanted contribution to an acceptably low level?
Figure 7 shows how the field gradient accumulates as we integrate over $`\tau ^{}`$ for a fixed number of modes. The oscillations are artefacts of the finite number of modes summed, and it is the envelope of the function which reflects the level of accuracy that can be achieved. If we are to accurately evaluate the value for the self-force, we must squeeze the envelope of this function below the value of the self-force. For example, on dimensional grounds, one might expect there to be a radial component of the self-force for a static scalar charge in Schwarzschild spacetime of the form
$`f_r=\lambda {\displaystyle \frac{q^2M}{b^3}},`$ (34)
where $`\lambda `$ is an unknown coefficient we would be trying to dig out from under the unwanted contribution to the force from our finite mode sum (in the electrostatic case, there is a finite contribution to the self-force of exactly the form shown in equation (34) with $`\lambda =1`$). The amplitude of the envelope is therefore an indication of the bound we place on the contribution to the self-force. In the present case, we know the solution vanishes identically, and evaluating equation (16) from section 2 at $`R=6M`$, it can be seen that the first integral in equation (7) gives only a tiny contribution to the self-force of a static scalar charge, even when $`\mathrm{\Delta }\tau `$ is extended out to near the edge of the normal neighbourhood. Therefore, in this case, the amplitude of the envelope is nothing but the approximate accuracy to which we have calculated the self-force.
If we choose $`\mathrm{\Delta }\tau `$ such that $`\mathrm{cos}\beta _{\mathrm{max}}=0`$, which is about half-way out to our upper bound on the edge of the normal neighbourhood, then the we see from table 2 that using 200 multipoles we can constrain the $`\lambda `$ to be less than $`0.1`$.
## 4 Conclusions
Our goal in this paper was to assess the feasibility of the Poisson-Wiseman matched expansion scheme for calculating the self-force. This scheme involves splitting the non-local term in the self-force (which is an integral over the particles worldline $`z(\tau )`$) at some proper time $`\tau \mathrm{\Delta }\tau `$. For the $`\mathrm{}<\tau ^{}<\tau \mathrm{\Delta }\tau `$ (distant) part of the integral, the integrand is expanded as a mode sum. For the $`\tau \mathrm{\Delta }\tau <\tau ^{}<\tau `$ (quasi-local) part, the integrand is expanded in a covariant Taylor series, which requires $`z(\tau \mathrm{\Delta }\tau )`$ be in the normal neighbourhood of $`z(\tau )`$. For the resulting self-force to be accurate, the split must be done in such a way that each expansion reach sufficient accuracy with a finite and feasible number of terms. Until now, it has not been clear that such a split even existed.
For this preliminary investigation, we have studied the simplest scenarios: minimally-coupled massless scalar field, Schwarzschild geometry, static particles and circular geodesic motion. Our results are somewhat surprising. The quasi-local expansion appears to converge well - when the particle is at $`6M`$ and expanding to order $`\mathrm{\Delta }\tau ^7`$, we achieve an estimated truncation error of the order of a few percent for $`\mathrm{\Delta }\tau 6M`$. This would naïvely have seemed to us enough of a buffer between the mode-sum integral and the singularity to allow rapid convergence of the mode sum as well.
However, we have found the convergence of the mode sum to be poor. This is because the mode sum smears the charge of the particle over spheres of finite radius which extend outside of the normal neighbourhood. Thus, even the flat-space term in the Green’s function, which cannot contribute to the self-force for a scalar particle, seems to contribute at every order. For our simple case of a static particle at $`6M`$ and $`\mathrm{\Delta }\tau 6M`$, going from 10 modes to 100 increased our accuracy by only a factor of 3.
Clearly, we need a way to speed convergence of mode sum. In the case presented, one could have achieved infinite accuracy at every order by regularizing the modes, which would have removed the flat-space portion of the Green’s function. It might, therefore, be possible to speed the convergence of the mode sum by regularizing the modes. This would be ironic, since this method is supposed to provide an alternative to mode regularization. Nonetheless, it might be that combining regularization with a two expansion approach will provide much better convergence than either alone. Alternatively, we note in figure 7 that the mode-sum integrand oscillates about its true value. This leads us to speculate that averaging over these oscillations would lead to much quicker convergence of the mode-sum integral.
While we believe our convergence analysis is general, and will apply beyond the simple scenarios we have explored, we would be remiss in not pointing out that there are additional issues for higher spin fields. The foremost of these is the issue of gauge - in order to meaningfully combine the quasi-local and more distant self-force contributions, they will need to be expressed in the same gauge. For instance, the quasi-local contribution for the gravitational field has only be derived in the Lorentz gauge, where no expression for the modes is available. We note, however, that this is a generic difficulty for gravitational self-force calculations because the formal expressions for the self-force are themselves derived in the Lorentz gauge, and we are heartened by current progress in understanding gauge transformations in this context.
In any case, we feel that these preliminary results are promising enough to warrant further investigation. A good figure of merit for calculational schemes like this is accuracy per floating point operation. As is, this approach would seem to lag approaches like mode-sum regularization when measured on this scale. Nonetheless, it might, for instance, provide confirmation for results from other approaches. Further, as mentioned above, it is still possible that this calculational scheme, with refinements, could rival or exceed in computational efficiency other known schemes, especially since the quasi-local expansion, once calculated, can be applied with little further effort to any spacetime or particle motion. In the mean time, this approach remains, in our opinion, in the category of “promising but possessing some technical challenges”.
This work was supported primarily by NSF grants PHY-0140326 and PHY-0200852, with additional support by the Center for Gravitational Wave Physics, which is funded by the National Science Foundation under Cooperative Agreement PHY 0114375, the Center for Gravitational Wave Astronomy, which is funded by the National Aeronautic and Space Administration through the NASA University Research Center Program, and the Yukawa Institute for Theoretical Physics. We would like to thank the Capra Gang, and especially Eric Poisson, for useful conversations. We would also like to thank the organizers of the Capra Meeting series for providing excellent forums for the exchange of ideas and information.
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# Untitled Document
On the non additivity of the trace in derived categories
Daniel Ferrand
In this note we provide an example of an endomorphism of a short exact sequence of perfect complexes, with the trace of the middle map not equal to the sum of the traces of the two other ones. The point is that the squares involved are commutative only up to homotopy. In view of this example I have found in 1968, Deligne immediately created his ”categories spectrales”, and soon afterwards Illusie introduced the ”filtered derived categories” where a satisfactory kind of additivity is restored for the trace.
1. Rappels
Réduisons les rappels au minimum. Soit $`A`$ un anneau commutatif. On ne considère que des complexes bornés de $`A`$-modules projectifs de type fini
$$K=\mathrm{}K^{n1}\stackrel{d}{}K^n\stackrel{d}{}K^{n+1}\mathrm{}$$
Dans la suite, ils seront nommés *complexes parfaits* (En toute rigueur, il faudrait dire : strictement parfaits). Pour un endomorphisme de complexe parfait $`u:KK`$, on pose
$$\mathrm{T}r(u)=(1)^i\mathrm{T}r(u^i).$$
Considérons trois complexes parfaits, $`K,L`$ et $`M`$ et une suite exacte courte $`0K\stackrel{j}{}L\stackrel{q}{}M0`$ (ce qui veut dire qu’en chaque degré $`n`$, la suite $`0K^nL^nM^n0`$ est exacte). Considérons un endomorphisme de cette suite exacte, c’est-à-dire trois endomorphismes de complexes $`u,v`$ et $`w`$ tels que le diagramme
$$\begin{array}{ccccc}K& \stackrel{j}{}& L& \stackrel{q}{}& M\\ u& & v& & w& & \\ K& \underset{j}{}& L& \underset{q}{}& M\end{array}$$
$`(1)`$
soit *commutatif*. Alors on a la formule d’additivité
$`(2)\mathrm{T}r(v)=\mathrm{T}r(u)+\mathrm{T}r(w),`$
puisqu’elle est vraie degré par degré.
On vérifie sans peine que deux endomorphismes homotopes d’un complexe parfait ont la même trace. On a pu espérer que la propriété d’additivité (2) serait conservée lorsque le diagramme (1) ne commute *qu’à homotopie près*. Il n’en est rien.
2.- L’exemple
On suppose que l’anneau $`A`$ contient un élément $`\epsilon `$, non nul et de carré nul, et on considère les complexes (parfaits) suivants :
* $`K`$ est $`A`$ placé en degré $`1`$ ;
* $`L=(A\stackrel{\epsilon }{}A)`$, en degré $`0`$ et $`1`$ ;
* $`M`$ est $`A`$ placé en degré $`0`$.
On a une suite exacte courte de complexes
$$\begin{array}{ccccc}K& & L& & M\\ & & & & & \\ & & A& =& A\\ & & \epsilon & & & & \\ A& =& A& & \end{array}$$
Finalement, on considère les trois endomorphismes $`u=0,w=0`$, et $`v=\left(\genfrac{}{}{0pt}{}{0}{\epsilon }\right):LL`$ (c’est-à-dire $`0`$ en degré $`0`$, et $`\epsilon `$ en degré $`1`$).
L’additivité des traces n’est pas respectée puisque $`\mathrm{T}r(u)=\mathrm{T}r(w)=0`$, et $`\mathrm{T}r(v)=\epsilon `$. Or, dans le diagramme
$$\begin{array}{ccccc}K& & L& & M\\ 0& & v& & 0& & \\ K& & L& & M\end{array}$$
le carré de droite est commutatif, et le carré de gauche est commutatif *à homotopie près*, comme on le voit en prenant pour homotopie $`h`$ l’application identique $`K^1=AL^0=A`$, qui est en tirets sur la figure suivante qui résume la situation :
3.- Commentaires
Au milieu des années 60 je suivais le Séminaire de Géométrie Algébrique piloté par Grothendieck. Au début de 1967 il me demanda de faire un exposé, et de le rédiger, sur le déterminant des complexes parfaits (l’exposé XI, manquant dans SGA 6). Ce devait être, me dit-il, un simple travail de vérification à partir de ses notes, lesquelles annonçaient, entre autre, la multiplicativité du déterminant pour les morphismes de triangles distingués *dans la catégorie dérivée des complexes parfaits* (cette multiplicativité antraîne l’additivité de la trace).
Deux difficultés m’empêchèrent d’avancer.
\- Tenir compte des signes - omniprésents - dont on doit affecter les isomorphismes liés à des permutations des facteurs.
\- Passer à la catégorie dérivée, c’est-à-dire essentiellement envisager des diagrammes qui ne commutent qu’à homotopie près.
Pour surmonter la première difficulté et rendre compréhensible le destin de ces signes, Grothendieck eut, peu après, l’idée lumineuse de faire du déterminant, non un simple module inversible, mais un module inversible *gradué* (concentré en un seul degré si Spec($`A`$) est connexe) : le déterminant du complexe parfait $`\mathrm{}K^nK^{n+1}\mathrm{}`$$`K^n`$ est projectif de rang $`r(n)`$, devient le couple $`(L,r)`$, où
$$L=\underset{n}{}(^{r(n)}K^n)^{(1)^n},\text{et}r=(1)^nr(n).$$
L’isomorphisme de commutativité $`(L,r)(M,s)(M,s)(L,r)`$ doit être affecté du signe « de Koszul » $`(1)^{rs}`$. Voir \[KM\] (il s’agit bien du produit $`rs`$, et non de la somme, comme une coquille p.20 de cet article pourrait le laisser croire).
Que les questions d’homotopie me gênassent à ce point montrait, aux yeux de Grothendieck, mon inaptitude flagrante. Mais, finalement, j’ai osé penser que c’était peut-être faux, ce qui me conduisit alors très vite au contre-exemple trivial signalé plus haut. C’était en juillet 1968. Deligne se mit au travail, et dès septembre, il présenta ses « catégories spectrales » où les triangles, les « vrais » , sont assujéttis à des conditions fortes qui rendent une théorie du déterminant possible.
À ma connaissance, ce texte n’a pas été publié.
Voici un extrait de son introduction : L’idée est la suivante : les triangles qu’on rencontre « en pratique » sont toujours déduits d’un objet plus fin : « un vrai triangle » et, plus important, les endomorphismes de triangles qu’on rencontre « en pratique » proviennent toujours d’endomorphismes de vrais triangles. Pour de tels endomorphismes, la formule \[d’additivité des traces\] redevient correcte.
Malheureusement, les « vrais triangles » de $`D(𝒜)`$ forment une catégorie triangulée, de même que les catégories des « vrais triangles » de « vrais triangles »…Cette machine infernale explique la complication du formalisme de compatibilité requis pour définir les « catégories spectrales »…
Pour mener à bien cette rédaction il m’aurait fallu alors tout reprendre dans ce nouveau contexte, et d’abord le comprendre. Une certaine lassitude, et un intérêt pour des sujets plus directement géométriques m’ont fait abandonner ce travail, et SGA 6 est paru sans l’exposé XI.
Peu après, Illusie introduisit ses « catégories dérivées filtrées » \[Il\] ch. V, où « la construction du cône, non fonctorielle, est remplacée par celle de gradué associé, qui l’est. » Dans ce cadre l’additivité de la trace est restaurée sous la forme suivante (V 3.7.7, p.310) : la trace d’un endomorphisme d’un complexe filtré dont le gradué est parfait, est égale à la trace de l’endomorphisme induit sur le complexe gradué associé.
En 1975, Knudsen et Mumford publièrent un article sur le déterminant des complexes parfaits \[KM\]. Ils établissent sa multiplicativité pour les endomorphismes de suites exactes courtes (p.41), et ils soulignent en note que les carrés en jeu doivent être « vraiment » commmutatifs, et non pas seulement à homotopie près. L’exemple signalé plus haut justifie cette mise en garde. De plus, p.44, ils montrent que si le schéma de base est *réduit*, il suffit de supposer la commutativité à homotopie près. L’emploi d’un nilpotent dans l’exemple est donc inévitable.
Bibliographie
\[D\] DELIGNE P., Catégories spectrales, Manuscrit, (Sept. 1968).
\[Il\] ILLUSIE L.,Complexe cotangent et déformations, I, LNM 239, Springer (1971).
\[KM\] KNUDSEN F. and MUMFORD D., The projectivity of the moduli space of stable curves. I : Preliminaries on ”det” and ”Div”, Math. Scand. 39 (1976) 19-55.
IRMAR, Université de Rennes 1, Campus de Beaulieu, F-35040 Rennes Cedex
E-mail address : daniel.ferrand@univ-rennes1.fr
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# Foundations of Real Analysis and Computability Theory in Non-Aristotelian Finitary Logic
## 1 Introduction to NAFL
The basic description of non-Aristotelian finitary logic (NAFL) given in is outlined below for the sake of completeness. The language, well-formed formulae and rules of inference of NAFL theories are formulated in exactly the same manner as in classical first-order predicate logic with equality (FOPL), where we shall assume, for convenience, that natural deduction is used; however, there are key differences and restrictions imposed by the requirements of the Main Postulate of NAFL, which is defined below. In NAFL, truths for *formal propositions* can exist *only* with respect to axiomatic theories. There are no absolute truths in just the *language* of an NAFL theory, unlike classical/intuitionistic/constructive logics. There do exist absolute (metamathematical, Platonic) truths in NAFL, but these are truths *about* axiomatic theories and their models. As in FOPL, an NAFL theory T is defined to be consistent if and only if T has a model, and a proposition $`P`$ is undecidable in T if and only if neither $`P`$ nor its negation $`\neg P`$ is provable in T. The truth definition given below in the Main Postulate is the heart of NAFL and its vital importance is underscored by Proposition 1, which follows it. The existing mindset that “truth definitions are not part of logic” most emphatically does not hold of NAFL.
### 1.1 The Main Postulate of NAFL
Let $`P`$ be a legitimate proposition of a consistent NAFL theory T. We will see in Sect. 1.2 that legitimacy of a proposition requires it to be in the ‘theory syntax’ of T. If $`P`$ is provable/refutable in T, then $`P`$ is true/false with respect to T (henceforth abbreviated as ‘true/false in T’); i.e., a model for T will assign $`P`$ to be true/false. If $`P`$ is undecidable in a consistent NAFL theory T, then the Main Postulate provides the appropriate truth definition as follows: $`P`$ is true/false in T if and only if $`P`$ is provable/refutable in an *interpretation* T\* of T. Here T\* is an axiomatic NAFL theory that, like T, temporarily resides in the human mind and acts as a ‘truth-maker’ for (a model of) T. The theorems of T\* are precisely those propositions that are assigned ‘true’ in the NAFL model of T, which, unlike its classical counterpart, is not ‘pre-existing’ and is instantaneously generated by T\*. Obviously, the theorems of T must necessarily be included in those of T\*. Note that for a given consistent theory T, T\* could vary in time according to the free will of the human mind that interprets T; for example, T\* could be T+$`P`$ or T+$`\neg P`$ or just T itself at different times for a given human mind, or in the context of quantum mechanics, for a given observer . Further, T\* could vary from one observer to another at any given time; each observer determines T\* by his or her own free will. The essence of the Main Postulate is that $`P`$ is true/false in T if and only if it has been *axiomatically declared* as true/false by virtue of its provability/refutability in T\*. In the absence of any such axiomatic declarations, i.e., if $`P`$ is undecidable in T\* (e.g. take T\*=T), then $`P`$ is ‘neither true nor false’ in T and Proposition 1 shows that consistency of T requires the laws of the excluded middle and non-contradiction to fail in a non-classical model for T in which the superposition $`P\&\neg P`$ is the case.
###### Proposition 1.
Let $`P`$ be undecidable in a consistent *NAFL* theory *T*. Then $`P\neg P`$ and $`\neg (P\&\neg P)`$ are not theorems of *T*. There must exist a non-classical model $``$ for *T* in which $`P\&\neg P`$ is the case.
For a proof of Proposition 1, see or ; this proof, also reproduced in Appendix A, seriously questions the logical/philosophical basis for the law of non-contradiction in both classical and intuitionistic logics. The non-classical model $``$ of Proposition 1 is a superposition of two or more classical models for T, in at least one of which $`P`$ is true and $`\neg P`$ in another. Here ‘classical’ or ‘non-classical’ is used strictly with respect to the status of $`P`$, and ‘superposition’ means that *all* the truths in each of the superposed classical models must hold in the non-classical model $``$. Thus if $`P`$ is true in one classical model and $`\neg P`$ in another, the superposed state $`P\&\neg P`$ will hold when these two classical models are superposed to form $``$. In $``$, ‘$`P`$’ (‘$`\neg P`$’) denotes that ‘$`\neg P`$’ (‘$`P`$’) is not provable in T\*, or in other words, $``$ expresses that neither $`P`$ nor $`\neg P`$ has been axiomatically declared as (classically) true with respect to T; thus $`P`$, $`\neg P`$, and hence $`P\&\neg P`$, are indeed (non-classically) true *in our world*, according to their interpretation in $``$. Note also that $`P`$ and $`\neg P`$ are *classically* ‘neither true nor false’ in $``$, where ‘true’ and ‘false’ have the meanings given in the Main Postulate. The quantum superposition principle is justified by identifying ‘axiomatic declarations’ of truth/falsity of $`P`$ in T (via its provability/refutability in T\*) with ‘measurement’ in the real world .
Proposition 1 is a metatheorem, i.e., it is a theorem *about* axiomatic theories. The concepts in Proposition 1, namely, consistency, undecidability (provability) and the existence of a non-classical model for a theory and hence, quantum superposition and entanglement, are strictly metamathematical (i.e., pertaining to semantics or model theory) and not formalizable in the syntax of NAFL theories. An NAFL theory T is either consistent or inconsistent, and a proposition $`P`$ is either provable or refutable or undecidable in T, i.e., the law of the excluded middle applies to these metamathematical truths. Note that the existence of a non-classical model does not make T inconsistent or paraconsistent, because T does not *prove* $`P\&\neg P`$. However, one could assert that the model theory for T requires the framework of a paraconsistent logic, in which the non-classical models can be analyzed. NAFL is the only logic that correctly embodies the philosophy of formalism ; NAFL truths for formal propositions are axiomatic, mental constructs with strictly no Platonic world required.
### 1.2 Theory syntax and proof syntax
An NAFL theory T requires two levels of syntax, namely the ‘theory syntax’ and the ‘proof syntax’. The theory syntax consists of precisely those propositions that are legitimate, i.e., whose truth in T satisfies the Main Postulate; obviously, the axioms and theorems of T are required to be in the theory syntax. Further, one can only add as axioms to T those propositions that are in its theory syntax. In particular, neither $`P\&\neg P`$ nor its negation $`P\neg P`$ is in the theory syntax when $`P`$ is undecidable in T; this will be clear from the explanation given below. The proof syntax, however, is classical because NAFL has the same rules of inference as FOPL; thus $`\neg (P\&\neg P)`$ is a valid deduction in the proof syntax and may be used to prove theorems of T. For example, if one is able to deduce $`AP\&\neg P`$ in the proof syntax of T where $`P`$ is undecidable in T and $`A`$ is in the theory syntax, then one has proved $`\neg A`$ in T despite the fact that $`\neg (P\&\neg P)`$ is not a theorem (in fact not even a legitimate proposition) of T. This is justified as follows: $`\neg (P\&\neg P)`$ may be needed to prove theorems of T, but it does not follow in NAFL that the theorems of T imply $`\neg (P\&\neg P)`$ if $`P`$ is undecidable in T. Let $`A`$ and $`B`$ be undecidable propositions in the theory syntax of T. Then $`AB`$ (equivalently, $`\neg AB`$) is in the theory syntax of T if and only if $`AB`$ is *not* (classically) deducible in the proof syntax of T. It is easy to check that if $`AB`$ is deducible in the proof syntax of T, then its (illegal) presence in the theory syntax would force it to be a theorem of T, which is not permitted by the Main Postulate. For, in a non-classical model $``$ for T in which both $`A`$ and $`B`$ are in the superposed state (i.e., both $`A\&\neg A`$ and $`B\&\neg B`$ hold), $`A\&\neg B`$ must be non-classically true, but theoremhood of $`AB`$ will prevent the existence of $``$. If one replaces $`B`$ by $`A`$ in this result, one obtains the previous conclusion that $`\neg (A\&\neg A)`$ is not in the theory syntax. For example, take $`\text{T}_0`$ to be the null set of axioms. Then nothing is provable in $`\text{T}_0`$, i.e., every legitimate proposition $`P`$ of $`\text{T}_0`$ is undecidable in $`\text{T}_0`$. Hence $`P\&\neg P`$ must be satisfiable in a non-classical model for $`\text{T}_0`$; this is obviously true, for $`P\&\neg P`$ contradicts only $`P\neg P`$, which is not in the theory syntax (and hence cannot be a theorem) of $`\text{T}_0`$. In particular, the proposition $`(A\&(AB))B`$, which is deducible in the proof syntax of $`\text{T}_0`$ (via the *modus ponens* inference rule), is not in the theory syntax; however, if $`AB`$ is not deducible in the proof syntax of $`\text{T}_0`$, then it is in the theory syntax. Note also that $`\neg \neg AA`$ is not in the theory syntax of $`\text{T}_0`$; nevertheless, the ‘equivalence’ between $`\neg \neg A`$ and $`A`$ holds in the sense that one can be replaced by the other in every model of $`\text{T}_0`$, and hence in all NAFL theories. Indeed, in a non-classical model for $`\text{T}_0`$, this equivalence holds in a non-classical sense and must be expressed by a different notation .
## 2 Infinite Sets Do Not Exist in NAFL Theories
Consider set theory with classes and assume that the only mathematical objects permitted to belong to classes are sets. A class is specified as an extension of a definite property in the language of set theory. Let $`P(x)`$ be one such property of sets, in which all variables (bound and unbound) are restricted to be set variables and in which finitely many independently defined constant class symbols may appear. We denote the class $`C`$ associated with $`P(x)`$ as $`x[xCP(x)]`$ or more concisely, $`C=\{x:P(x)\}`$. Certain classes may be sets; a proper class is a collection that is deemed to be not a set and therefore cannot belong to any class. The universe of a model for a theory is a nonempty class $`U`$ over which all the variables of the theory range. The theory of finite sets F, with classes, may be obtained by removing the axiom of infinity (AI) from either Zermelo’s set theory with classes (Z) \[6, Chaps. 1 and 5\], or from Gödel-Bernays theory (GB) \[7, pp. 73–78\]. That every property expressible in the language determines a class is a theorem scheme of GB \[7, pp. 76–77\] and in particular, of F; we denote this as the theorem of comprehension for classes. The natural numbers in F are defined in the usual set-theoretic sense, with $`0`$ as the empty set $`\mathrm{}`$ and $`n+1=n\{n\}`$. The class $`N`$ of all natural numbers is defined by
$$N=\{x:x\text{is a natural number}\}.$$
Note that we have admitted “$`x`$ is a natural number” as a valid property in the language of F. Consider the class $`D`$, defined as:
$$D=\{x:n(n\text{N}x\mathrm{}^{(n)}(\mathrm{}))\}.$$
(1)
Here $`\mathrm{}^{(n)}(\mathrm{})`$ is the power set operation performed $`n`$ times on the null set $`\mathrm{}`$; we let $`\mathrm{}^{(0)}(\mathrm{})=\mathrm{}`$. Note that $`D`$ exists provably in F by the theorem of comprehension. Consider the following proposition:
$$D=\mathrm{}.$$
(2)
According to the standard interpretation, (2) is undecidable in F, for the following reason. Define a model $`𝔅`$ for F in which only finite sets exist; the universe $`B`$ of $`𝔅`$ is given by
$$B=\{x:n(n\text{N}\&x\mathrm{}^{(n)}(\mathrm{}))\}.$$
(3)
In the interpretation of $`𝔅`$, $`D=\mathrm{}`$ holds and hence (2) is true. But in a model in which infinite sets do exist, (2) is false.
###### Theorem 1.
The proposition $`D=\mathrm{}`$ defined in (2) is required to be provable in *F* by the Main Postulate of *NAFL* and Proposition 1. Consequently, infinite sets do not exist in the *NAFL* version of *F*.
###### Proof.
We will outline the proof given in \[5, Sect. 3\]. Uniqueness of $`D`$ is deducible from the axiom of extensionality for classes (AE), which implies:
$$D=\mathrm{}D\mathrm{}.$$
(4)
The key point is that AE *requires* (4) (equivalent to $`\neg (D=\mathrm{}\&D\mathrm{})`$) to be a theorem of F. But by Proposition 1, (4) cannot be a theorem of F, given that F is consistent and that (2) is undecidable in F. We are forced to the conclusion that one of these two “givens” must be false in NAFL. Either (2) is undecidable in F, which is therefore inconsistent, or else (2) must be decidable in F.
What we have demonstrated is that if F is to be a consistent theory of NAFL, then we need a principle of proof for either $`D=\mathrm{}`$ or for $`D\mathrm{}`$. We will argue the case for $`D=\mathrm{}`$. The axioms of F tacitly presume that a definite class $`U`$ is specified as the universe, because these formulas contain the universal quantifiers $``$ and $``$ which may be thought of as referring to the universe. However, to define $`U`$ we need to consider the formulas of F as meaningful because these assert the existence of sets that must necessarily belong to $`U`$. Thus both $`U`$ and all the sets of F are impredicatively defined. However, finite sets may also be defined predicatively, at least in principle, merely by listing all of their elements; but this is not possible for infinite sets. In conclusion, infinite sets can only be defined impredicatively, i.e., by self-reference via the universal class (which must again be defined impredicatively by reference to infinite sets). The appropriate principle of proof in NAFL for $`D=\mathrm{}`$ is that such an impredicative definition of sets must be banned, for it generates undecidable propositions that contradict the Main Postulate of NAFL. ∎
###### Remark 1.
Note that Theorem 1 is a metatheorem and its proof required Proposition 1. Classically, it is not known whether $`D=\mathrm{}`$ is either provable or refutable or undecidable in F; the undecidability is merely *assumed*. What Theorem 1 *predicts* is that by the classical rules of inference, the theory $`\text{F}+(D\mathrm{})`$ must be inconsistent. In other words, Theorem 1 proves that F proves (2) (assuming consistency of F), but does not provide a specific proof in F of (2).
###### Remark 2.
The proof of Theorem 1 illustrates a very important requirement of NAFL. Whenever an NAFL theory proves that a class or a set must exist uniquely (e.g. via AE), then Proposition 1 requires that such uniqueness must be enforced with respect to that *theory*, rather than just with respect to *models* of that theory. Thus the theory is also required to provide a unique *construction* for the class or set in question. E.g. NAFL requires the theory F to specify the class $`D`$ uniquely; it is not valid in NAFL to assert that either $`D=\mathrm{}`$ or $`D\mathrm{}`$ can hold in models of F, for then uniqueness is lost with respect to F. One can immediately conclude that non-Euclidean geometries are inconsistent in NAFL , for Euclid’s first four postulates prove that there must exist a *unique* straight line between any two distinct points. The uniqueness here must again be enforced with respect to a theory comprising the first four postulates, which implies that such a theory must provide a *construction* for the said straight line via provability of the fifth postulate.
## 3 Infinite proper classes must provably exist in NAFL theories with infinite domains
###### Proposition 2.
If *T* is a consistent *NAFL* theory that proves the existence of infinitely many objects identifiable by a property $`P(x)`$ in the language of *T* (e.g. $`P(x)`$ could be ‘$`x=x`$’ with *T* taken as Peano Arithmetic) then *T* must prove the existence of the infinite proper class $`\{x:P(x)\}`$.
Note that Proposition 2 addresses only *infinite* classes; there is no such corresponding result for finite classes (or sets). In a sense the axioms for the existence of infinite classes in the theory F of Sect. 2 are redundant; by Proposition 2, they must be provable in the NAFL version of F, and the language of F must necessarily include the required properties that specify the infinite classes.
###### Proof.
The simplest proof is as follows. Consider the NAFL version of Peano Arithmetic (NPA; see Appendix B), which is equivalent to F. From Theorem 1, it follows in NAFL that nonstandard models of NPA cannot exist. For the Main Postulate of NAFL (see Sect. 1) requires that models of NPA must necessarily be specified by its interpretation NPA\*, which is also an NAFL theory that denies the existence of infinite sets (by Theorem 1). But without infinite sets, nonstandard models of arithmetic cannot be specified and therefore cannot exist in NAFL. It immediately follows that the only model of NPA is the standard model and hence, by the completeness theorem of first-order logic, NPA must prove the existence of the domain of standard natural numbers specified by the class $`N`$ defined in Sect. 2 (in the language of NPA, $`N`$ may be specified as $`\{x:x=x\}`$). This result is immediately generalized and Proposition 2 follows. ∎
###### Remark 3.
Note that the completeness theorem of classical first-order logic, which speaks *about* axiomatic theories, is taken for granted in NAFL as a metamathematical principle, but is not provable in any NAFL theory because of the infinitary reasoning required; in fact the notion of ‘axiomatic theory’ is not even formalizable in NAFL theories. The nonexistence of nonstandard models of NPA also follows from the arguments given in \[2, Sect. 1\], which raise serious questions about the existence of nonstandard integers even within classical logic. Indeed, undecidability of the Gödel sentence for NPA immediately produces a contradiction similar to that discussed in the proof of Theorem 1. It follows as a corollary that consistency of NPA demands its *completeness*, i.e., every sentence of NPA must be either provable or refutable in NPA. This result, together with Theorem 1, means that one no longer has the freedom to invoke infinite sets to prove purely arithmetical propositions, as is done classically. Obviously, Gödel’s incompleteness theorems and Turing’s argument for the undecidability of the halting problem fail in NAFL, as discussed in and . Another consequence of the nonexistence of nonstandard models is that an ‘arbitrary but fixed’ constant does not exist in NAFL; it is the same as a free variable and assumed to be universally quantified with respect to a *standard* domain (e.g. natural numbers). The Main Postulate requires that the free (or universally quantified) variable is in a superposed state of all possible values and assumes a particular value when and only when it is specified by the human mind that interprets the theory in question.
###### Remark 4.
An intuitive and useful explanation for the validity of Proposition 2 is as follows. In NAFL, the process of counting one or finitely many natural numbers at a time will not exhaust ‘all’ the natural numbers. A simple induction tells us that such a counting process will always leave out infinitely many natural numbers, and so will never be complete. Yet NPA obviously specifies ‘all’ natural numbers; it can only do so by accessing (at some unspecified point in the counting process) infinitely many natural numbers *at the same time*. Thus NPA must prove the existence of the infinite proper class of natural numbers, which may be thought of as the *simultaneous* existence of ‘all’ natural numbers. A class is identified by all of its elements, as required by AE. When a theory ‘constructs’ all the elements of an infinite class, it must unavoidably construct the class itself. Proposition 2 is a metatheorem; it proves that T must prove the existence of the infinite class $`\{x:P(x)\}`$, but does not provide the said proof in T. The proof in T that infinitely many elements exist that satisfy $`P(x)`$ may also be taken as the proof in T of the existence of $`\{x:P(x)\}`$. For example, if $`P(x)`$ expresses that ‘$`x`$ is a natural number’ (say, via $`P(x)\mathrm{`}x=x^{}`$) in the theory NPA, one might take the proof in NPA of the existence of the infinite class $`N=\{x:P(x)\}`$ as that of the proposition $`nm(m>n)`$.
###### Remark 5.
It is clear from Proposition 2 and Remark 4 that an NAFL theory T always identifies an infinite class $`\{x:P(x)\}`$ by all of its elements. In fact $`\{x:P(x)\}`$ can *only* be so identified as an infinite class, i.e., in the extensional sense. Thus the axiom of extensionality for classes, in those instances where infinite classes are involved, is built into the very definition of the word ‘class’ and must also be provable in every NAFL theory T that proves the existence of infinitely many objects.
###### Proposition 3.
‘Arbitrary’ infinite classes and quantification over infinite classes (which are equivalent; see Remark 3) are not permitted in *NAFL* theories. An infinite (proper) class must always be a constant that is specified constructively in *NAFL* theories, via a mapping to the class $`N`$ of all natural numbers (which is assumed to be a constructive specification). This mapping must be shown to exist even when one is unable to specify a general formula for it (e.g. the class of all prime natural numbers).
Proposition 3 emphasizes the difference between an infinite set (which is a mathematical object in classical logic) and an infinite proper class. Quantification over proper classes amounts to treating these as sets.
###### Proof.
From Proposition 2 and Remark 3, it is clear that Gödel’s incompleteness theorems (and Turing’s argument for the unsolvability of the halting problem) must fail in NAFL. To justify such a failure, one must impose a ban on quantification over proper classes. The proofs of Gödel’s theorems (as well as those of all versions of Turing’s argument for the undecidability of the halting problem ) always require quantification over infinitely many infinite proper classes. The said quantification in various versions of the proofs of these theorems appears either directly (e.g. via Cantor’s diagonalization argument) or indirectly, e.g. when one encodes self-referential sentences into Peano Arithmetic using Gödel numbering. In other words, Gödel’s theorems and Turing’s argument serve as reductio ad absurdum proofs of the infinitary nature, and hence, invalidity in NAFL, of quantification over proper classes. Indeed, the very fact that nonstandard models of arithmetic follow as a logical consequence of Gödel’s theorems and Turing’s argument is proof of their infinitary nature. ∎
## 4 Foundations of Real Analysis in NAFL
Extend NPA to the NAFL theory NPAR in which the integers ($`Z`$), rationals ($`Q`$) and sequences of rationals ($`q_n`$), including Cauchy sequences, are defined. See Appendix B for the definitions of NPA and NPAR. By Propositions 2 and 3, NPAR must prove the existence of every Cauchy sequence of rationals definable by formulas (‘properties’) in the language of NPAR. The usual extension of NPAR to treat real numbers as equivalence classes of such Cauchy sequences is not possible, for direct quantification over reals is banned. Neither can we talk of an *arbitrary* real number $`x`$ (in the usual sense, with no construction specified for $`x`$) in NAFL, for that amounts to quantification, via Proposition 3; however, see also Remark 10. Nevertheless, indirect quantification over reals may be accomplished by a metamathematical translation of diagrammatic concepts of Euclidean geometry into operations on sequences of rationals in NPAR. We denote as NRA the appropriate NAFL metatheory (to be defined below) of NPAR in which this translation is performed.
###### Definition 1.
In *NRA*, we associate real numbers with ‘points’ on the real line, which exists a priori as a geometric entity. Euclidean geometry and its diagrammatic concepts are taken for granted as true. In *NRA*, every mention of a constant real number $`r`$ (e.g. $`\pi `$) *must* be accompanied by a ‘construction’ for $`r`$, namely, a specific Cauchy sequence of rationals (that classically converges to $`r`$); ‘non-constructive’ existence of real numbers is not permitted in *NAFL* (see Proposition 3). Any such Cauchy sequence is appropriate in a given instance, and different sequences may be used in different instances to represent the same real number. In *NRA*, a constant real number $`r`$ is specified only *after* a Cauchy sequence representing it has been displayed. Define addition, subtraction, multiplication and division of reals by corresponding termwise operations on the rationals in the Cauchy sequences representing the reals; we do not permit sequences containing zero terms to represent the denominator of a quotient of reals. These definitions, adapted from , are as follows. If $`x=q_n`$ and $`y=q_n^{}`$ are real numbers, we write $`x=_Ry`$ in *NRA* to mean that $`lim_n|q_nq_n^{}|=0`$, i.e., in *NPAR*,
$$ϵ(ϵ>0mn(m<n|q_nq_n^{}|<ϵ)),$$
and we write $`x<_Ry`$ in NRA to mean that
$$ϵ(ϵ>0\&mn(m<nq_n+ϵ<q_n^{}))$$
holds in *NPAR*. Similarly, we also make the following translations in *NRA* to formulas of *NPAR*:
$`x+_Ryq_n+q_n^{}`$
$`x_Ryq_nq_n^{}`$
$`_Rxq_n`$
$`x/_Ryq_n/q_n^{},\text{where }q_n^{}0`$
$`0_R0;1_R1.`$
We have dropped the subscript $`Q`$ from the relations $``$, $`/`$, $`+`$, $``$, $`<`$ and $`>`$ on the rationals, as well as from the symbols $`0`$ and $`1`$ (as defined in Appendix B). We will henceforth drop the subscript $`R`$ on these relations and symbols for the reals as well.
###### Remark 6.
Note that the expression ‘$`\pi /\sqrt{2}`$’ by itself makes no sense in NAFL unless Cauchy sequences representing $`\pi `$ and $`\sqrt{2}`$ are also specified. The result will be again be a specific Cauchy sequence. Of course, in this case we know that any such resulting sequence will always represent the same real number. Consider the expression ‘$`0/0`$’, where the numerator and denominator are again specified as Cauchy sequences representing the real number zero. The answer in this case will be dependent on how we specify these sequences, and may correspond to different real numbers in different instances, or may not converge at all. Whenever the answer converges, $`0/0`$ is a well-defined operation in the NAFL version of real analysis. This is so because we do not require $`0/0`$ to be a uniquely defined real number like $`\pi /\sqrt{2}`$; $`0/0`$ is a (possibly) meaningful real number only *after* Cauchy sequences representing the zeroes in the numerator and denominator are specified. The definitions of the real zeroes in the numerator and denominator that *we* provide are only a means for making an axiomatic assertion of the value of $`0/0`$, in tune with the Main Postulate.
We have defined real numbers as Cauchy sequences of rationals representing points on the real line and also the usual arithmetic operations on these. Next let us consider how one may quantify over specific collections of real numbers on the real line in NAFL. We will denote such collections as ‘super-classes’, since the reals do not constitute a class in NAFL (being themselves proper classes).
###### Definition 2.
A super-class $`I`$ of real numbers that has a geometric representation on the real line is defined by specifying a constant sequence $`S`$ of rationals that has precisely $`I`$ as its (super-class of) limit points. Different sequences $`S`$ may be used to represent $`I`$ in different instances. Here $`I`$ may be a finite collection of reals, an infinite sequence, or an interval (finite or infinite or semi-infinite), or a combination of these.
The key point of Definition 2 is that only entities $`I`$ which have a geometric representation on the real line may be defined as super-classes, as noted below.
###### Proposition 4.
Every super-class $`I`$ of real numbers defined in Definition 2 must necessarily be closed, i.e., include all of its limit points. In particular, open or semi-open intervals of reals, or sequences of reals excluding their limit points, cannot be represented in *NRA*. Note that such entities do not have a geometric representation, e.g. there is no way to geometrically specify an open/semi-open interval of reals by a diagram. Similarly, the rationals, when viewed as real numbers in the above geometrical sense, do not constitute a super-class; any sequence $`S`$ that includes all rational points represents the entire real line.
###### Proof.
This is immediate from Definition 2. Any sequence $`S`$ of rationals that has all the elements of an open/semi-open interval of reals as its limit points *must* also have the end-points of such an interval as its limit points; this is a classically known result. Similarly, if the sequence $`S`$ of rationals has an infinite sequence $`\rho _j`$ of reals as its limit points, it must also have the limit points of $`\rho _j`$ as its limit points. ∎
###### Remark 7.
The sequence $`S`$ of Definition 2 has two roles, namely: (a) As a sequence of rationals specified by a definite property in NPAR, without any reference to real numbers, and (b) As a super-class $`I`$ of reals when viewed metamathematically from ‘outside’ NPAR, i.e., in NRA. In the latter case, one must view $`S`$ as a *simultaneous* specification of *all* of its Cauchy subsequences that represent the super-class $`I`$. Thus $`S`$ satisfies our notion of an infinite class outlined in Remark 4. Indeed, when the theory NPAR *constructs* $`S`$, it has also unavoidably *constructed* every subsequence of $`S`$ in a metamathematical sense. This argument justifies the constructive methods used in the NAFL version of real analysis (NRA); there is no other way, given Proposition 3. Cantor’s diagonalization argument for the uncountability of reals is not legal in NRA (and cannot even be formulated in NPAR), because quantification over reals is not permitted; there is no ‘cardinality’ for a super-class of reals. In fact there is no ‘list’ of reals in NRA, i.e., there is no way to establish a mapping between any super-class of reals and the natural numbers $`N`$ because of the above restriction on quantification.
###### Remark 8.
Note that Definition 2 requires us to extend the meaning of the equality relation $`=_R`$ in NRA (see Definition 1) as follows. Two super-classes of reals represented by sequences $`q_n`$ and $`q_n^{}`$ of rationals are equal if and only if they have precisely the same (super-class of) limit-points. If $`q_n`$ and $`q_n^{}`$ happen to be Cauchy sequences, this modified definition of equality does agree with that in Definition 1. However, the notion of ‘having the same super-class of limit points’ cannot be formalized within an NAFL theory, such as, NPAR, because it would require quantification over reals. This is precisely why NRA must remain a metatheory that appeals to geometric concepts and not a formal NAFL theory of reals (which cannot exist).
###### Remark 9.
The fact that the rationals, when represented metamathematically as reals, do not constitute a super-class in NAFL (see Proposition 4) means that we cannot consider the rationals as having a separate identity from the real numbers, under a *geometric* interpretation. The rationals have a purely arithmetical construction in NPAR as ordered pairs of integers, but they are *not* points on the real line (or solutions of polynomial equations that permit reals as solutions, etc.) when viewed in this sense. Their construction in NPAR as ordered pairs of integers is required only for the purpose of defining the reals and such a construction does *not* imply geometric numberhood for the rationals. When viewed as points on the real line or in any other geometric sense, the rationals must *always* be viewed as reals, and do not constitute a super-class by themselves. E.g. it is not valid in NAFL to talk of all the ‘rational points’ in the interval , for no construction is possible for such a super-class of real numbers. Similarly, integer points on the real line must always be viewed as real numbers. Any super-class that includes all integer points on the real line must also include $`\pm \mathrm{}`$, where one may define $`\pm \mathrm{}`$ as real numbers represented by appropriate unbounded sequences of rationals. E.g., an instance is given below:
$$\pm \mathrm{}=_R\pm _Q((n,0),(1,0)),\text{where }nN.$$
###### Remark 10.
One may wish to represent *arbitrary* real numbers belonging to a super-class in NRA, say, for stating identities like $`x+y=_Ry+x`$ or $`x=_Rx`$. Such a representation is possible if one provides a *construction* for such an arbitrary real number in the super-class. This is accomplished in NAFL by using the Main Postulate; an arbitrary real number $`x`$ in a super-class is taken to be in a superposed state of assuming all possible values in the super-class until the human mind specifies a value for it . Thus $`x`$, when unspecified, has the same construction as the super-class itself, i.e., as a sequence of rationals with the entire super-class as its only limit points. One may verify that this construction works when one performs operations like $`x+y=_Ry+x`$, with the usual arithmetical operations defined similarly to Definition 1. The operations are defined even when we specify $`y`$ in the above example while letting $`x`$ be arbitrary, or let both $`x`$ and $`y`$ be arbitrary; however, note that in these latter instances one needs to interpret $`=_R`$ as defined in Remark 8. It must be emphasized that *arbitrary real numbers can only be defined within legitimate super-classes of NRA*. For example, the following seemingly innocent proposition of classical real analysis is illegitimate in NRA when $`x`$ is left arbitrary:
$$1/x>_R0x>_R0.$$
(5)
One might attempt to represent $`x`$ here as an arbitrary real number in $`[\mathrm{},+\mathrm{}]`$, which is a legitimate super-class of NRA. Indeed whenever *we specify* $`x`$ as a positive or a negative real number (via appropriate Cauchy sequences of rationals), (5) does hold. However, when we leave $`x`$ unspecified, it is in a superposed state of all possible values as required by NAFL and this includes positive, negative and zero real numbers; hence $`x`$ is represented in NRA by the super class $`[\mathrm{},+\mathrm{}]`$. In this instance, NAFL would require that both $`1/x>_R0\&\mathrm{\hspace{0.33em}1}/x_R0`$ and $`x>_R0\&x_R0`$ hold in NRA. But then it follows that $`1/x>_R0\&x_R0`$ also holds when we leave $`x`$ unspecified, and this is formally the negation of (5). This is the reason for the illegality of (5) in NRA. Note also that when we specify $`x=0`$ (by an appropriate Cauchy sequence of rationals) in NRA, one could have $`1/x=+\mathrm{}>_R0`$, which again results in the falsity of the above assertion. Hence restricting $`x`$ to the super-class $`[0,+\mathrm{}]`$ of reals also fails to rescue (5) in NRA. The reason (5) fails to be legitimate in NRA is that it attempts to create a super-class $`(0,+\mathrm{}]`$ that does not exist in NRA. However, restricting $`x`$ to other super-classes of strictly positive reals (such as, $`[1,+\mathrm{}]`$) will make (5) legitimate in NRA, provided $`>_R`$ is appropriately defined when $`x`$ is left arbitrary within this super-class. One may also easily verify that when (5) is modified to the proposition
$$1/x_R0x_R0,$$
it can indeed legitimately be asserted as true in NRA when $`x`$ is restricted to be an arbitrary real number in $`[0,+\mathrm{}]`$. In this case, one must *define* $`_R`$ such that $`x_R0`$ is true when $`x`$ is left unspecified (and hence is represented by the super-class $`[0,+\mathrm{}]`$).
###### Remark 11.
One may extend the above definitions on the real line to the Euclidean plane or to three-dimensional Euclidean space by defining points as Cauchy sequences of ordered pairs or triples of rationals that converge to the real coordinates of the points. One may then define complex numbers, partial derivatives and functions of reals (including discontinuous ones) that have a geometric representation, by an obvious extension of the results given in this paper. The arbitrary variables defined in Remark 10 are also needed to define functions. For example, to represent a function as $`y=f(x)`$, one needs a sequence of pairs of rationals such that it has every pair of values $`(x,f(x))`$ as its (super-class of) limit points. Further given any specific $`x`$ within the super-class representing the domain, the formula $`f(x)`$ generates the value $`y`$ of the function. Finally, if $`x`$ is left unspecified, its value is by definition the super-class representing the domain, and the value of $`f(x)`$ is the super-class representing the range of the function. The domain and range of $`f(x)`$ must be legitimate super-classes of NRA. Discontinuous functions in NRA will take on multiple values at points of discontinuity, as the classical representation using (semi-) open intervals is not possible; hence these are not functions in the classical sense. For example, to represent the step function
$$y=1,\text{for}x<0;y=2,\text{for}x>0;y=[1,2]\text{for}x=0,$$
one needs a sequence $`q_n,q_n^{}`$ of ordered pairs of rationals such that the (classical) limit points of these sequences are pairs of real numbers of the form $`(x,1)`$ for each $`x0`$ and $`(x,2)`$ for each $`x0`$. Further one should also be able to extract the pair $`(0,y)`$ as limit points of the above sequence, for each y satisfying $`1y2`$. The above should, of course, be the complete super-class of limit points of $`q_n,q_n^{}`$. Here we have taken $`y=[1,2]`$ as the super-class representing the ‘value’ of the function at $`x=0`$, to conform with the usual diagrammatic representation of the step function.
###### Remark 12.
Note that *proofs* about operations on real numbers may, in general, be obtained diagrammatically. Diagrammatic concepts from Euclidean geometry are taken to be consistent and *a priori* true in NRA, before translation into NPAR. As an example, see for a suggested diagrammatic proof of Euclid’s fifth postulate from the first four, in plane geometry. The basic idea of this proof is as follows. Given a line $`L`$ and a point $`P`$ of the plane that is not on $`L`$, one can deduce from Euclid’s first four postulates the existence of infinitely many line segments of a line $`M`$, each of which is parallel to $`L`$; further, this collection of line segments includes all the points of $`M`$. In NAFL, this implies the existence of $`M`$ itself. In order to translate this proof in NRA, one is faced with a seeming problem: how does one quantify over line segments, which are themselves super-classes? For example, one needs a representation of infinitely many line segments of the form $`[1,1],[2,2],[3,3],\mathrm{}`$ in NRA. The appropriate representation in this case is achieved by two sequences, one of which is a sequence of rationals that has the entire real line as its limit points and the other is a sequence of pairs of rationals having the limit points $`(1,1),(2,2),(3,3),\mathrm{}`$ (representing the end-points of the line segments). In particular, the second sequence must also include $`(\mathrm{},+\mathrm{})`$ as its limit point; this is the proof of the fact that the entire line must necessarily be included in any representation of its line segments that includes all the points of the line.
###### Remark 13.
In several previous papers we have highlighted the fact that special relativity theory (SRT) cannot be formalized as a consistent theory of NAFL. The present paper provides one more such argument, as follows. SRT requires that a particle can, in principle, acquire *all* velocities strictly less than $`c`$, the velocity of light, but the particle velocity can never equal $`c`$. Formalization of this result requires the existence of the semi-open interval of reals $`[0,c)`$. By Proposition 4, such a semi-open interval of reals does not exist in NRA.
###### Definition 3.
The derivative dy/dx in *NRA* may be reduced to a particular means of defining 0/0 as outlined in Remark 6. The classical limit of a sequence of reals of the type $`\{\mathrm{\Delta }y/\mathrm{\Delta }x\}`$ cannot be separated from the sequence itself in *NRA* (see Proposition 4); the limit in fact represents $`0/0`$. Given $`y=f(x)`$ and a particular value of $`x`$, one defines $`\mathrm{\Delta }x=0`$ as a specific Cauchy sequence of rationals, and then computes $`\mathrm{\Delta }y=f(x+\mathrm{\Delta }x)f(x)`$ as another Cauchy sequence representing zero. Performing the division $`\mathrm{\Delta }y/\mathrm{\Delta }x`$ yields $`dy/dx`$ at that value of $`x`$. Similarly, integration in *NRA* is an operation of the type $`0\times \mathrm{}`$.
###### Remark 14.
A detailed exposition of the calculus and the theory of functions of a real variable in NRA will be dealt with in future work. The derivative, of course, has the geometric interpretation of the slope of a curve. It is easily seen that essentially every object of Euclidean geometry (whether plane or three-dimensional) can be translated into NPAR via NRA. As an example of the derivative, consider $`y=f(x)=x^2`$. Performing the operation noted in Definition 3, one obtains
$`\mathrm{\Delta }y=_Rf(x+\mathrm{\Delta }x)f(x)=_R2x\mathrm{\Delta }x+(\mathrm{\Delta }x)^2`$
$`{\displaystyle \frac{dy}{dx}}=_R{\displaystyle \frac{\mathrm{\Delta }y}{\mathrm{\Delta }x}}=_R2x+\mathrm{\Delta }x=_R2x,`$
where $`\mathrm{\Delta }x`$ and $`x`$ are to be replaced by specific Cauchy sequences of rationals representing the real numbers zero and $`x`$ respectively. Observe that $`x`$ may also be left arbitrary within a legitimate super-class of NRA (see Remark 10). It is extremely important to note that the functional representation $`y=f(x)`$ must output a specific Cauchy sequence for $`y`$ when a Cauchy sequence for $`x`$ is substituted into $`f(x)`$. This does happen in the above example. Note also that the cancellation of $`\mathrm{\Delta }x`$, representing zero, from the numerator and denominator of $`\mathrm{\Delta }y/\mathrm{\Delta }x`$ has no bearing on the final outcome; it is a legal operation in NRA.
###### Remark 15.
The Weierstrass $`ϵ\delta `$ argument for the derivative $`dy/dx`$ requires the existence of open intervals of reals and therefore fails in NRA. The NAFL version of real analysis excludes the paradoxes of classical real analysis, all of which arise from the assumption that non-closed super-classes of reals, in particular, open/semi-open intervals, exist. As an example, consider one of Zeno’s paradoxes of motion, in which Achilles chases the tortoise and seemingly never catches up with it despite running at a higher velocity. This paradox only arises when one considers Achilles as reaching, and being confined to, infinitely many locations (points) at all of which the tortoise must always be ahead. But such a super-class of points does not exist in NRA; every attempt to represent such a super-class by the NAFL definition will result in its classical limit point to also be included in it (by Proposition 4), and Achilles does catch up with the tortoise at this limit point. Similarly the paradoxes of classical measure theory, such as, the Banach-Tarski paradox, can also be attributed to requiring non-closed super-classes of reals that do not exist in NRA.
###### Remark 16.
It is worth highlighting the fundamental reasons for the paradoxes of classical real analysis and precisely how they are resolved in NRA. Let us consider again Zeno’s paradoxes of motion. The essential paradox here arises from the fact that infinitely many finite, non-zero intervals of real numbers sum to a finite interval in classical real analysis. For example, consider the following infinite sequence of real intervals (in either space or time):
$$[0,1/2],[1/2,3/4],[3/4,7/8],\mathrm{}$$
If the above intervals are ‘stacked’ side by side, one gets the real interval $`[0,1]`$, reflecting the fact that $`1/2+1/4+1/8+\mathrm{}=1`$. But note that *each* of the above intervals is finite, non-zero and non-infinitesimal (infinitesimals do not exist in NRA because NAFL does not permit nonstandard models of NPA). If there are indeed infinitely many such intervals (as classical real analysis asserts, by mapping these to the class of natural numbers $`N`$), then how can their sum be an interval of finite length? For a related paradox, consider the sequence of nested intervals
$$[1,1],[1/2,1/2],[1/4,1/4],\mathrm{}$$
Note that classically, *each* of these nested intervals contains infinitely many (in fact, uncountably many) points. Then how can the intersection of these intervals contain only a single point, namely, zero? NAFL answers both of these paradoxes as follows. In each instance, the mapping of the real intervals to $`N`$ is not legal in NRA as it amounts to quantification over intervals of reals. When one attempts to construct either of the above sequences of intervals in NRA by the method noted in Remark 12, one finds that the interval $`[1,1]`$ or $`[0,0]`$ respectively of zero length must necessarily get included. This immediately explains these paradoxes, for it is not possible in NRA to assert that there exist infinitely many finite, non-zero intervals that sum to a finite interval, as in the above examples. Secondly, it is not even legal in NRA to ask *how many* intervals exist (including $`[1,1]`$ or $`[0,0]`$) in each of the above sequences of intervals, when one constructs them legally in NRA. Such a question assumes that it is legal to quantify over intervals of reals, whereas it is illegal to even quantify over real numbers in NRA.
Chaitin cites the paradoxes of classical real analysis and the incompleteness results of Gödel, Turing and himself as reasons why the classical real number system must be abandoned. Lynds asserts that there are no ‘instants’ of time in the real world, which means that Zeno’s paradoxes cannot even be formulated. These authors have essentially rejected the continuum as aphysical; it is interesting to compare their arguments with the NAFL approach to real analysis presented in this paper.
###### Remark 17.
Exponentiation of reals (or rationals) with rational exponents is relatively straightforward to define. But exponentiation with real exponents seems problematic in NPAR and will be analyzed in future work.
###### Remark 18.
As an interesting application of real analysis in NAFL, consider Fermat’s Last Theorem (FLT). This may be put in the form that $`r^n+s^n=1`$ has no solutions in the positive rationals $`r`$ and $`s`$ for integers $`n3`$. Let us fix $`n`$ and $`s`$ to be specific constants, and consider a sequence of rationals $`\{r_j\}`$, such that the limit of this sequence is a real number $`\varphi `$ that solves $`\varphi ^n+s^n=1`$; FLT requires that no member of the sequence $`\{r_j\}`$ should solve this equation. But we have seen from Remark 9 that the rationals, when substituted into polynomial equations that permit reals as solutions, *must* be considered as reals. So FLT requires that a super-class of reals $`\{r_j\}`$ must exist that fails to solve the above equation, while its limit point does solve it. This amounts to requiring that the limit point of this purported super-class must be excluded from it, but this is not possible by Proposition 4. One concludes that the truth of FLT cannot be meaningfully represented if the geometric interpretation of real numbers is imposed upon Peano Arithmetic (NPA). Therefore NAFL requires that any proof (or refutation) of FLT must be carried out entirely within NPA, without invoking real numbers. This reinforces our earlier conclusion of Remark 3 that consistency of NPA demands its completeness.
###### Remark 19.
Simpson (with his subsystems of second-order arithmetic), Feferman (predicativism) and Weaver (mathematical conceptualism) have all attempted to restrict classical infinitary reasoning so as to eliminate the set-theoretic paradoxes. The fundamental difference between the NAFL approach of this paper and those of the above authors is that NAFL does not accept the existence of infinite sets and quantification over (infinite) proper classes. Neither does NAFL accept the notion of an arbitrary infinite class (as a free variable), which amounts to quantification in NAFL. The point we wish to make here is that quantification over proper classes is fundamentally and unavoidably part of infinitary reasoning and so stands rejected by NAFL, as noted in the proof of Proposition 3. From the NAFL point of view, we dispute Simpson’s claim that he has achieved partial realizations of Hilbert’s program (to justify classical infinitary reasoning from the finitary standpoint) because Simpson has really used infinitary methods. In fact NAFL shows that Hilbert’s program is decisively settled negatively (see also ); classical infinitary reasoning stands refuted from the strictly finitary standpoint developed in NAFL. We also mention in passing that Weaver has severely criticized Feferman’s approach to predicativism. We honestly believe though, that neither Weaver’s nor Feferman’s approach is predicativism in the strict sense; they accept the existence of infinite sets of natural numbers, which are essentially impredicative objects from the NAFL standpoint (see Sect. 2). In fairness to these authors, they have characterized their approach as predicativism ‘given the natural numbers’.
## 5 Foundations of Computability Theory in NAFL
Here we will outline the basic arguments in \[4, Sects. 6 and 7\], and establish the connection with real analysis in NAFL. A Turing machine (TM), *by definition*, must either halt or not halt. If $`P`$ is the proposition that a given TM halts, then $`P\neg P`$ is unavoidably built into the definition of that TM. It follows that $`P`$ cannot be undecidable in any consistent NAFL theory T in which the existence of that TM is formalized; the non-classical model for T required by Proposition 1 in which $`P\&\neg P`$ is the case cannot exist, for T must prove $`P\neg P`$ (and hence, either $`P`$ or $`\neg P`$). Any infinite (proper) class in an NAFL theory must be recursive; whether an object belongs to or does not belong to that class cannot be undecidable in a consistent NAFL theory because such undecidability will violate the uniqueness required by the axiom for extensionality for classes, which is an essential ingredient of NAFL theories in which infinite classes exist (e.g. see Sect. 2, as well as Remarks 4 and 5). Note that the above arguments require the non-classical features of NAFL, namely, Proposition 1 and the Main Postulate. It is clear that a new paradigm for computability theory is required in NAFL.
###### Proposition 5 (NAFL Computability Thesis, NCT).
Every infinite (recursive, proper) class that exists in *NAFL* theories must be effectively computable, i.e., there exists an algorithm that computes it. Here ‘recursive’ or ‘computable’ does not necessarily mean ‘computable by a classical Turing machine’.
###### Proof.
It is clear that for every element of an infinite class, which is a finite set, there exists an algorithm that computes it, i.e., it is effectively computable. Note that nonstandard models of arithmetic do not exist in NAFL; see the proof of Proposition 2 and Remark 3. Therefore the existence of infinitely many algorithms that compute every (standard finite) element of an infinite class implies that the entire infinite class has been computed by these algorithms. But by Remark 4, these algorithms cannot operate one or finitely many at a time and complete the computation process; by induction, there will always exist infinitely many elements remaining to be computed. One concludes that in NAFL, the existence of infinitely many algorithms as noted above necessarily implies the existence of *an* algorithm $`𝒜`$ that computes the infinite class. One may think of $`𝒜`$ as operating on the infinitely many algorithms corresponding to all elements and ‘parallelizing’ them, so that the infinite class gets computed *simultaneously*, rather than finitely many at a time, as required by the NAFL interpretation of an infinite class (see Remark 4). Note that the nonexistence of $`𝒜`$ would leave us with a paradox in NAFL, namely, that a class of algorithms collectively managed to compute an infinite class, although each of them did only a finite computation. The above proof of the existence of $`𝒜`$ required the non-classical concepts of NAFL, namely, the Main Postulate, Proposition 1 and the nonexistence of nonstandard integers. The proof that every infinite class must be recursive also required these concepts. Hence $`𝒜`$ need not be a classical algorithm, i.e., the infinite class in question need not be computable by a classical Turing machine. ∎
###### Remark 20.
It is a theorem of classical recursion theory that every recursive class must be computable by a Turing machine. But many of these standard results fail in NAFL, whose concepts of ‘computability’ and ‘algorithm’ must necessarily be different from the classical concepts. Proposition 5 implies that Turing’s halting routine $`H`$ must exist in NAFL. We believe that such an algorithm must be non-classical, not because of Turing’s argument (which fails in NAFL), but because requiring $`H`$ to be classical would make it a self-referential entity that would lead to contradictions in NAFL. Thus NAFL permits hypercomputation, and there is a case for believing that $`H`$ must be a quantum algorithm, for the following reasons. NAFL also justifies quantum superposition and entanglement , which are non-classical and metamathematical phenomena, i.e., they are confined to the metatheory (semantics) of NAFL theories. Hence a purported quantum algorithm for $`H`$, which only computes the halting decisions for classical Turing machines, would reside in the metatheory and will not be a self-referential entity. Secondly, there is evidence that quantum algorithms permit hypercomputation and infinite parallelism . The latter feature may also be justified by noting that in NAFL, a quantum algorithm is permitted to access and compute a truly random element of an infinite class. When such an element is not specified, it must be in a superposed state of assuming all values in the infinite class (see Remark 10) and this corresponds to infinite parallelism.
###### Remark 21.
An interesting analogy is suggested between the rational/real numbers on the one hand, and classical/quantum algorithms on the other. It is known that for every classical algorithm, there is a quantum algorithm that achieves the same result; the converse of this assertion is controversial and not necessarily true. The classical and quantum representations of an algorithm (when they both exist) may be thought of as corresponding to the dual representation of a number as a rational and real respectively. As noted in Remark 9, the rational representation does not have the geometrical significance of a point on the real line; in a similar sense, the classical representation of an algorithm merely encodes a meaningless finite string of symbols, rather than an algorithm that executes. When thought of as an algorithm, the quantum representation *must* always be used. This is justified by the fact that the infinitely many classical algorithms (in the execution mode) by themselves do not constitute a class because they cannot be separated from the quantum algorithms that ‘parallelize’ them (see Proposition 5, its proof and Remark 20). Quantum algorithms may be thought of as infinite sequences of the classical representations; these sequences are only needed to define the quantum algorithms (in the same sense that reals are infinite sequences of rational representations that are only needed to define the reals). Such sequences of classical representations may also be used to define ‘super-classes’ of quantum algorithms, in the same manner that NRA defines super-classes of reals. An executing algorithm, even when it halts and has a classical representation, is to be thought of as an infinite object corresponding to its quantum representation. This may be because the machine that executes the algorithm has an infinite tape, for example; the instruction ‘halt’ may be interpreted as a tacit requirement not to operate further on this infinite tape. Secondly, the classical model of computation is not valid in NAFL and the algorithm may at some point access infinitely many execution states at the same time. A quantum algorithm that does not have a classical representation may be thought of as corresponding to an irrational real number. In conclusion, we observe that important potential applications of the NAFL paradigms for real analysis and computability theory are in the areas of quantum and autonomic computing .
## Appendix A. Proof of Proposition 1
###### Proof.
By the Main Postulate of NAFL, $`P`$ ($`\neg P`$) can be the case in T if and only if $`P`$ ($`\neg P`$) has been asserted *axiomatically*, by virtue of its provability in T\*. In the absence of any such axiomatic assertions (*e.g.* if T\*=T), it follows that neither $`P`$ nor $`\neg P`$ can be the case in T and hence $`P\neg P`$ cannot be a theorem of T. The classical refutation of $`P\&\neg P`$ in T proceeds as follows: ‘If $`P`$ ($`\neg P`$) is the case, then $`\neg P`$ ($`P`$) cannot be the case’, or equivalently, ‘$`\neg P`$ ($`P`$) contradicts $`P`$ ($`\neg P`$)’. But, by the Main Postulate, this argument fails in NAFL and amounts to a refutation of $`P\&\neg P`$ in T\*=T+$`P`$ (T+$`\neg P`$), and *not* in T as required. Careful thought will show that the classical refutation of $`P\&\neg P`$ in T is the *only possible* reason for $`\neg (P\&\neg P)`$ to be a theorem of T, and it fails in NAFL. The intuitionistic refutation of $`P\&\neg P`$ in T is flawed and also fails in NAFL, as will be shown in the ensuing paragraph. By the completeness theorem of FOPL (which, as noted in Remark 3, is taken for granted as a metamathematical principle in NAFL) it follows that there must exist a non-classical model for T in which $`P\&\neg P`$ is satisfiable. ∎
Consider the law of non-contradiction as stated in a standard system of intuitionistic first-order predicate logic due to S. C. Kleene, namely,
$`\neg P(PQ)`$. This formula asserts that from contradictory premises $`P`$ and $`\neg P`$, an *arbitrary* proposition $`Q`$ can be deduced, which is absurd. Hence $`\neg (P\&\neg P)`$ seemingly follows. However, note that in intuitionism, truth is provability (not necessarily in a specific theory T); together with the intuitionistic concept of negation, it follows that an assertion of $`\neg (P\&\neg P)`$ is the same as deducing an absurdity from $`P\&\neg P`$, or equivalently, from contradictory premises $`P`$ and $`\neg P`$. But we have seen that the ‘absurdity’ referred to here is precisely the fact that *any proposition can be deduced*, given contradictory premises! The above ‘proof’ of $`\neg (P\&\neg P)`$ from contradictory premises, mandated by the intuitionistic concepts of truth and negation, is flawed because *any proposition can be so deduced*. Note that this ‘proof’ is formally indistinguishable from one in which $`\neg (P\&\neg P)`$ is *substituted* for the deduced arbitrary proposition $`Q`$. In NAFL, it is not possible to deduce an arbitrary proposition from contradictory premises in a non-classical model, and so the flawed intuitionistic argument for $`\neg (P\&\neg P)`$ fails in any case. Indeed, as explained in , the argument for deducing an arbitrary proposition would normally proceed as follows:
$`P\&\neg PP,`$
$`PPQ,`$
$`P\&\neg P\neg P,`$
$`\neg P\&(PQ)Q.`$
The final step fails in a non-classical NAFL model for a theory T (in which $`P\&\neg P`$ is the case) because this step presumes $`\neg (P\&\neg P)`$.
## Appendix B. The formal system NPA and its extension NPAR
We will outline the formal system NPA (i.e., the NAFL version of first-order Peano Arithmetic) along the lines of the description given in Chapter 1 of , suitably modified for our purposes. Throughout, a natural deduction system of classical first order predicate logic with equality (FOPL) is assumed. The language of NPA admits number variables, denoted by $`i,j,k,m,n,\mathrm{},`$ which are intended to range over the class $`N`$ of all natural numbers $`\{0,1,2,\mathrm{}\}`$. Capital letters $`L,M,N,\mathrm{},`$ are reserved for constant (infinite) classes; there are no class variables, and we do not need finite classes for our purposes (though these can be added if required). These infinite classes are defined by the notation $`\{n:\varphi (n)\}`$, for each L-formula $`\varphi (n)`$ that holds for infinitely many values of $`n`$.
The terms and formulas of the language (L) of NPA are as follows. Numerical terms are number variables, the constant symbols $`0`$ and $`1`$, and $`t_1+t_2`$ and $`t_1t_2`$ whenever $`t_1`$ and $`t_2`$ are numerical terms. Here $`+`$ and . are binary operation symbols intended to denote addition and multiplication of natural numbers. Numerical terms are intended to denote natural numbers. Atomic formulas are $`t_1=t_2`$, $`t_1<t_2`$ and $`t_1X`$ where $`t_1`$ and $`t_2`$ are numerical terms and $`X`$ is any (constant) infinite class. The intended meanings of these atomic formulas are that $`t_1`$ equals $`t_2`$, $`t_1`$ is less than $`t_2`$, and $`t_1`$ is an element of $`X`$. Formulas are built up from atomic formulas by means of propositional connectives $`\&`$, $``$, $`\neg `$, $``$, $``$ (and, or, not, implies, if and only if) and universal (number) quantifiers $`n`$, $`n`$ (for all $`n`$, there exists $`n`$). There are no class quantifiers. A sentence is a formula with no free variables. In writing terms and formulas of L, parentheses and brackets will be used to indicate grouping and some obvious abbreviations including the symbols $``$ and $``$ and the numbers $`2,3,\mathrm{}`$ will also be used for convenience.
The axioms of NPA are the following L-formulas (universal quantification of number variables is assumed):
(i) basic axioms
$`n+10`$
$`m+1=n+1m=n`$
$`m+0=m`$
$`m+(n+1)=(m+n)+1`$
$`m0=0`$
$`m(n+1)=(mn)+m`$
$`\neg m<0`$
$`m<n+1(m<nm=n)`$
(ii) induction axiom scheme (for the L-formulas $`\varphi (n)`$)
$`(\varphi (0)\&n(\varphi (n)\varphi (n+1)))n\varphi (n)`$
Note that the L-formulas $`\varphi (n)`$ may contain independently defined constant class symbols. Further, by Remark 3, $`\varphi (n)`$ must be a sentence; the free (number) variables are automatically assumed to be universally quantified in NAFL. By Proposition 2, NPA must prove the existence of every infinite class $`𝒞`$ generated by L-formulas, as indicated by the following theorem scheme, which we will denote as the theorem scheme of comprehension:
$$n(n𝒞\varphi (n)).$$
Here $`𝒞`$ is the generic notation that stands for $`\{n:\varphi (n)\}`$. The theorem scheme applies for each L-formula $`\varphi (n)`$ that holds for infinitely many values of $`n`$. Note that $`\varphi (n)`$ may not contain $`𝒞`$, but could possibly contain other (constant) class symbols that have been defined independently of $`𝒞`$. The given instance of the theorem scheme says that there exists an infinite class $`𝒞=\{n:\varphi (n)\}`$, which is the class of all $`n`$ such that $`\varphi (n)`$ holds. As noted in Remarks 4 and 5, the ‘infinite class’ referred to here is always identified by all of its elements, i.e., it is always a ‘class’ in the extensional sense and NPA must also prove the axiom of extensionality for classes in those instances where infinite classes are involved.
The above completes the description of the language L and the axioms of NPA, whose objects are the natural numbers. The theorems of NPA are deduced from the axioms using the classical rules of inference. The restrictions required by NAFL as noted in Sect. 1, and in particular, Sects. 1.1 and 1.2, apply to NPA.
Next consider the extension NPAR of NPA in order to handle integers, rationals and sequences of rationals. The language LR of NPAR augments L, firstly by admitting ordered pairs of natural numbers in the form $`(m,n)`$, which are the integers, belonging to the infinite class $`Z`$. The axioms for the integers are as follows (universal quantification over the number variables is assumed):
(i) equality
$`(m,n)=_Z(i,j)m+j=n+i`$
(ii) addition
$`(m,n)+_Z(i,j)=_Z(m+i,n+j)`$
(iii) negative integers
$`_Z(m,n)=_Z(n,m)`$
(iv) multiplication
$`(m,n)_Z(i,j)=_Z(mi+nj,mj+ni)`$
(v) zero and one
$`0_Z=_Z(0,0);1_Z=_Z(1,0)`$
(vi) order
$`(m,n)<_Z(i,j)m+j<n+i`$
Note that we have avoided mentioning equivalence classes while defining the integers. Although equivalence classes of pairs of naturals defining specific integers do exist in NPAR (via the theorem of comprehension), one cannot quantify over these in NAFL because they are infinite classes. Hence it is not possible to use equivalence classes to define infinitely many integers in NPAR. Following Simpson \[8, Chap. 1\], one may also define $`Z`$ to be a class of *representatives* of the appropriate equivalence classes, e.g.:
$$Z=\{(0,0),(1,0),(0,1),(2,0),(0,2),(3,0),(0,3),\mathrm{}\}.$$
But note that *formally*, the above definition of $`Z`$ does not require any mention of equivalence classes. Thus in NPAR, whenever we formulate a proposition to hold for all integers, it is to be understood that the proposition is true of all members of $`Z`$ in the above form, or alternatively, for all ordered pairs $`(m,n)`$ as noted in the above axioms.
The rationals (denoted by the infinite class $`Q`$) are defined in NPAR as ordered pairs $`(a,b)`$ of integers, where $`b`$ is restricted to be a positive integer. The axioms for rationals are as follows (universal quantification over free variables is assumed).
(i) equality
$`(a,b)=_Q(c,d)a_Zd=_Zb_Zc`$
(ii) addition
$`(a,b)+_Q(c,d)=_Q(a_Zd+_Zb_Zc,b_Zd)`$
(iii) negative rationals
$`_Q(a,b)=_Q(_Za,b)`$
(iv) multiplication
$`(a,b)_Q(c,d)=_Q(a_Zc,b_Zd)`$
(v) division
$`(a,b)/_Q(c,d)=_Q(a_Zd,b_Zc),\text{where }c0_Z`$
(vi) zero and one
$`0_Q=_Q(0_Z,1_Z);1_Q=_Q(1_Z,1_Z)`$
(vii) order
$`(a,b)<_Q(c,d)a_Zd<_Zb_Zc`$
Note that the axiom for division is redundant, but it will be useful for our purposes in defining division of real numbers. Following Simpson \[8, Chap 1\], one may also define $`Q`$ in NPAR as a class of representatives of equivalence classes; these representatives may be chosen to be minimal in the sense that all common factors are cancelled out. Again, $`Q`$ must be constructively defined in NPAR without any reference to equivalence classes in order to avoid quantification over these.
The idea behind the above definitions of $`Z`$ and $`Q`$ within NPAR is that $`(m,n)`$ corresponds to the integer $`mn`$, while $`(a,b)`$ (with $`b`$ restricted to be a positive integer) corresponds to the rational $`a/b`$.
A sequence of rational numbers is defined to be a function $`f:NQ`$. We denote such a sequence as $`q_n:nN`$ or simply $`q_n`$, where $`q_n=f(n)`$. A Cauchy sequence of rational numbers is a sequence $`q_n:nN`$ such that
$$ϵ(ϵ>0mn(m<n|q_mq_n|<ϵ)).$$
Here $`ϵ`$ ranges over $`Q`$ and we have dropped the obvious subscripts on the operations $`>`$, $`<`$, $``$, etc. Propositions 2 and 3 require that NPAR must prove the existence of every sequence of rationals (including Cauchy sequences) definable in the language of NPAR. Remarks 4 and 5 apply. Of course, the language of NPAR must also admit sequences of integers definable by its formulas, and NPAR must prove the existence of these as well.
IBM and Autonomic Computing are trademarks of the International Business Machines Corporation in the United States, other countries, or both.
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# Secular Evolution of Galaxies
## 1 The three main processes/scenarios of galaxy formation
### 1.1 Monolithic Collapse (MC)
The scenario of monolithic collapse for the formation of the Galaxy was suggested by Eggen, Lynden-Bell and Sandage (1962, ELS), on the basis of the observed correlation between \[Fe/H\] and orbital excentricity $`e`$ for old stars. Since stars with low metallicity had very low angular momentum L<sub>z</sub>, they suggested that the old stars were formed out of gas falling towards the center in radial orbits, collapsing quickly from a halo to a thin rotating disk plane enriched in heavy elements by star formation. In this picture, the gravitational potential is varying slowly and the stellar parameters $`e`$ and L<sub>z</sub> can be considered as invariants.
In an opposite view, Searle & Zinn (1978) consider that the ELS results are simply the consequence of the selection bias on high proper motion stars. They remark that halo globular clusters have a large range of metallicity, with \[Fe/H\] uncorrelated with distance, and favor a formation through merging of small protogalaxies. The gravitational potential is varying rapidly in this view, and there is also substantial gas accreted later.
The monolithic collapse scenario for our Galaxy is not viable anymore, or only for a restricted component, the central bulge. The stellar halo is smaller than the disk, and could not collapse into the disk. The only galactic component to collapse into is the old bulge (Gilmore 1996). Note however that, although the bulge stars are old, they have high metallicity, as the disk, so this is not easily explained by this scenario.
Rapid monolithic collapse is not likely to be the main formation process of most normal galaxies, since huge ULIRGs at high redshift are very rare. In any case, departures from MC are larger in low-mass galaxies (Ferreras & Silk 2000).
### 1.2 Hierarchical scenario (HS)
In this frame, a system like our own Galaxy is the result of the hierarchical assembly of dark halo building blocks. Accretion of baryonic gas occurs later, in the assembled structure, to form the bulge, and progressively the thin disk, which forms last. The thick disk could be due to the past heating of the thin disk by companions.
This scenario is supported by the existence of stellar streams in the stellar halo (Helmi 2002, Ibata et al 2002): they are tidal debris from the accretion of small companions. The halo could be entirely built from these minor mergers.
The observed large increase with redshift of the merger frequency comes in support of the hierarchical scenario. This is particularly spectacular in galaxy clusters, for massive galaxy formation (e.g. van Dokkum et al 1999). In some massive ellipticals, two populations of globular clusters have been detected, hinting towards two major merging episodes, triggering their formation (Zepf & Ashman 1993).
### 1.3 Secular Evolution, with slow and continuous external matter accretion (SE)
In secular evolution, the bulge component is formed slowly from the disk through the bar action, and the disk can be replenished through continuous external gas accretion. The Galaxy can be considered an open system, with slow mass growth through time, from external accretion of gas progressively transformed into stars, and the various components interacting with each other.
According to a statistical analysis of 257 spiral galaxies, the radial color gradients observed and the relation between disk and bulge colors favor the formation of bulges through SE (Gadotti and dos Anjos 2001). SE is also supported by the relation between bulge and disk masses and radii (Courteau et al 1996).
In a large sample of early-type galaxies, moderate radial gradients of metallicity are observed, which are independent of luminosity (De Propris et al 2005). Ellipticals and bulges of lenticular galaxies reveal similar behaviour. These observations are contrary to what is expected with MC (large gradients which should increase with mass and luminosity) and contrary to what is predicted by HS (which smears out gradients).
Cosmological simulations based on a $`\mathrm{\Lambda }`$CDM universe may include all processes with different emphasis, according to the treatment of baryon physics. In addition to hierarchical merging of dark haloes, the hot gas is cooling in dark haloes within the cooling radius, and contracts monolithically to form galaxies. However, the gas infall might also occur in fragmented and dense colder gas clouds (Maller & Bullock 2004), or the gas accretion could be continuous, and favor secular evolution.
Successive episodes of dissipational collapse, followed by mergers, lead in numerical simulations to elliptical-like objects, which are in better agreement with observations than the results of stellar disks mergers only (Meza et al 2003). In this scenario, the final merger episode forms hot gas, the absence of cooling prevents star formation, and the final object is an elliptical, with no recent accretion, and consequently an old stellar population with no disk morphology.
## 2 Successes and difficulties of the different processes
### 2.1 Massive galaxies early in the universe
Very massive elliptical systems (10$`{}_{}{}^{11}M_{}^{}`$), possessing a large fraction of old stars, have been observed in deep optical and near-infrared surveys, at redshifts between 1 and 2 (Daddi et al 2004, Cimatti et al 2004). The presence of these rather evolved systems, when the universe was only 25% of its age was a surprise, since it is expected that massive ellipticals are the end result of the hierarchical merging process.
However, these massive objects are highly clustered and appear to correspond to the early formation of groups in rich environments. It is predicted that evolution proceeds much more rapidly in those environments with respect to the field, and the presence of some fraction of massive objects is thus expected.
Such massive objects can be traced by their star formation, either by optical techniques, when their light is not highly obscured, as in Lyman Break Galaxies (LBG), or through their dust emission in the far-infrared, redshifted in the submillimeter range. The latter are then called SMGs (SubMillimeter Galaxies), and are found preferentially between z=1 and 4. The physical nature of these objects can be determined more precisely when CO emission is detected. Given their gas content, and star formation rate (larger than 700 M$``$/yr), their life-time is between 40 Myr and 200 Myr (Greve et al 2005). The correlation between the Far-Infrared and CO emission suggests that they are the equivalent of local ULIRGs, with an even more efficient star formation rate. These starbursts are associated to galaxy mergers, which could be more violent than locally, since bulges are less prominent at high redshifts.
The number of SMGs detected can be used to estimate the density of massive galaxies at these redshifts (Greve et al 2005). The estimate is a lower limit, since the life-time of the starburst phase is not precisely known (and the correction is made with the conservative value of 200 Myr). Within the error bars, the derived number density of SMGs is compatible with the predictions of the $`\mathrm{\Lambda }`$CDM simulations and semi-analytical calculations, if 10% of baryons are rapidly converted into galaxies (cf Figure 1).
This agreement with observations is obtained in the new model from Baugh et al (2005) through a few changes in the assumptions of the semi-analytical simulations, in particular a much larger influence of bursts in the star formation history, a longer time-scale for star formation in disks, and most important, the assumption of a top-heavy IMF in starbursts at high redshift. The consequences of the new model are that the star formation is dominated by bursts at redshifts larger than 4, and later by quiescent mode (while it was dominated by the quiescent mode at all redshifts before). Integrated over the Hubble time, bursts are now responsible for 30% of all star formation (Baugh et al 2005). Massive galaxies today have only a few percent of their stars formed in bursts at z $``$ 2. The galaxy formation is still quite far from the monolithic collapse view, since the bulk of the stars in massive ellipticals is formed in disks, and then re-arranged in spheroids during mergers. Note that the estimated baryonic mass of SMGs has now been re-scaled to 0.6 M in average, while it was 4 times larger before (Greve et al 2005).
### 2.2 Star formation history versus mass
The star formation history (SFH) has been derived from the direct observations of young stars in distant galaxies at different redshifts, which reveal a peak in the star formation rate about 8 Gyr ago, and then a decline by a factor of $``$ 10. Alternatively, the study of the fossil record of stellar populations, in the local large sample of galaxies in the SDSS has also led to a similar star formation history, with a peak occuring slightly later, about 5 Gyr ago (Heavens et al 2004, Jimenez et al 2004). This study also shows that the SFH is different according to the mass. Stars in massive galaxies appear to have formed at early times, and these galaxies are not actively forming stars now, while dwarf galaxies appear to experience starburts. Only intermediate masses have in average maintained their star formation rate over a Hubble time.
It has been known for a long time that galaxies in the middle of the Hubble sequence had about constant SFR across the Hubble time (Kennicutt 1983, Kennicutt et al 1994). Even taking into account the stellar mass loss, an isolated galaxy should have an exponentially decreasing star formation history. To account for observations, a source of gas should be found to replenish the star formation fuel in such galaxies. This cannot come from major mergers, which destroy disks, and lead to early-type or elliptical galaxies. On the other hand, small companions are not sufficient, for instance systems falling now on the Milky Way (Sag dw, Canis major, etc..) are of the order of 1/400th of the mass of the Galaxy (Ibata et al 2001, 2003). Accretion of gas from the cosmic filaments should provide the required fuel, which corresponds to a slow and continuous secular evolution.
### 2.3 Quiescent star formation versus bursts
Below a redshift of z =0.7, which corresponds to about half of the Hubble time, the star formation rate of the universe is dominated by a rather quiescent mode in normal galaxies. But at higher redshift, starbursts begin to dominate, first moderate bursts, giving rise to LIRGs (Luminous Infra-Red Galaxies, with L $`>10^{11}L_{}`$), and after z=4, ULIRGs (Ultra-Luminous Infra-Red galaxies, with L $`>10^{12}L_{}`$). Given the shape of the star formation history, most of the stars today have been born in LIRGs (63%), the rest being in ULIRGs (21%) or normal galaxies (16%); this is consistent with the cosmic infrared backgound (CIRB), which is dominated at least within a proportion of two thirds by LIRGs (Elbaz & Moy, 2004).
The fact that LIRGs are more numerous at z $`>`$ 0.4 (15%) than today ($``$ 1%) has been interpreted in terms of recent formation (in the second half of the universe) of the bulk of stars in intermediate-mass galaxies (Hammer et al 2005). At this epoch, LIRGs can account for 38% of all star formation. The high frequency of LIRGs implies episodic star formation bursts, maybe triggered by galaxy interactions. The starburst phase is of the order of 100 Myr, while an interaction is of the order of 10<sup>9</sup> yrs. This means that all galaxies should be interacting at z$`>`$ 0.4, which is not what is observed.
The frequency of interactions must be lower, but another process could compensate: external gas accretion from the cosmic filaments. Secular evolution could then be responsible for these episodes of star formation. In any scenario, gas must be provided to replenish spiral disks, whatever is the final dynamical trigger, either a passing companion or internal evolution. In both cases, the amount of fuel is the same, the gas from the disk of one (or two) galaxy, and a bar drives the gas towards the center. The gravity torques from the bar can fuel the central disk and star formation intermittently, through self-regulated mechanisms (see below).
### 2.4 Late bulge formation
Contrary to massive spheroids, consisting of old and red stellar populations, many bulges of spiral galaxies appear to have formed later than their disk. They host young populations, that must have been formed from recent gas inflow followed by star formation, due either to secular evolution or to galaxy interactions. Using the criterion that the color inside the half-mass radius is bluer than the outer parts, Kannappan et al (2004) found that 10% of galaxies in the NFGS (Nearby Field Galaxy Survey) are experiencing bulge growth at the present time.
To distinguish the actual source of gas inflow in those galaxies, a correlation was searched between the blue-center disk characteristic and the presence of companions (see Figure 2). Although there is not a large correlation with the presence of visible companions, there is a strong one with perturbed morphology. The authors then conclude that external drivers, minor mergers or interactions, are more likely the origin of bulge growth today, than internal secular evolution. In fact, as will be developped later, secular evolution can also be driven by external gas accretion, which if asymmetrical, can lead to perturbed morphology as well.
### 2.5 Respective role of MC, HS and SE
In summary, a certain amount of all three processes are certainly present in galaxy formation, and it is quite difficult to distinguish the relative role of each. For instance, although MC is likely more effective at high redshifts and in high-density environments, it is possible that old stars have been accumulating in massive systems at various redshifts.
Galaxy interactions and mergers increase with $`(1+z)^m`$, with $`m`$ 4. It has been thus suggested that hierarchical merging dominate in the past and secular evolution will in the future (Kormendy & Kennicutt 2004). But the availability of gas accretion also decreases with time, and it is quite possible that today both processes are occuring with comparable importance, depending strongly on environment. Both the effects of interactions and gas accretion are quenched in clusters. Even inside a given group, the past evolution appears to depend on morphological type: the Milky Way reveals more signs of secular evolution (pseudo-bulge, old globular clusters..), while Andromeda has experienced a major merger more recently.
The perturbed morphology of galaxies should not always be interpreted in terms of galaxy interactions, it could also be the result of external gas accretion, which is most often asymmetric: lopsided systems, warps, polar accretion… (cf Figure 3). External gas accretion, followed by secular evolution from the bar, can also provide starbursts. The gravity torques of the bar either prevent or favor the gas infall towards the center, as will be described now. The self-regulated bar action results in intermittent periods of activity, starbursts ad well as nuclear activity.
## 3 Bars and secular evolution
Dynamical instabilities in spiral disks are responsible for its evolution. Detailed numerical simulations since several years have unveiled a self-regulated cycle in this secular evolution. The cycle begins by a first bar instability in a cold spiral disk with stars and gas. When the bar is strong enough, it produces gravity torques on the gas component, due to the phase shift of the gas response (the gas being dissipative). These torques drive gas inflow, from corotation to the center (or ILRs, when they exist). The gas is losing its angular momentum to the benefit of the bar, which is then weakening. The destruction of the bar is due essentially to the gas inflow itself, and also to a lesser degree to the presence of a central mass concentration (CMC), built in the process. Only if enough external gas accretion occurs to replenish the disk mass and trigger another bar instability, can the cycle loop again.
There have been debates about this cycle, in particular about the efficiency of bar destruction, about their ability to reform, and about the required central mass (CMC) to weaken the bar.
### 3.1 Role of gas in bar destruction
Self-consistent simulations of spiral galaxies with gas have shown that only the infall of 1-2% of mass in gas is sufficient to destroy the stellar bar or to transform it into a lens (Friedli 1994, Berentzen et al 1998, Bournaud & Combes 2002, 2005). But the mechanism of destruction was attributed to the formation of a central mass concentration (CMC), and simulations with the growth of an artificial CMC in the center of a few percent is not enough to destroy the bar (Shen & Sellwood 2004). Now, several simulations with and without gas or CMC have clearly shown the role of gas in bar destruction: gas is driven in by the bar gravity torques, and its angular momentum is taken up by the bar wave. In other words, the reciprocal torques from the gas provides angular momentum to the bar, weakening it, since the bar wave has negative angular momentum inside corotation (Bournaud & Combes 2005). Since the formed CMC is not enough to destroy the bar by itself, it is then more easy to reform a bar, when the disk has become unstable again, through external gas accretion.
In this cycle, it is interesting to note that the galaxy accretes gas towards its center by intermittence. Indeed, while the bar is strong in the disk, the gas from corotation to OLR (Outer Lindblad Resonance) is driven outwards by the positive gravity torques from the bar. Gas is then stalled at the OLR in the border of the disk, prevented to enter. The external gas remains there, while the gas inside corotation is driven inwards, until the bar weakens. After the bar destruction, the external gas can enter and replenish the disk, to make it unstable again to bar formation (see Figure 4). This scenario explains why nuclear activity is not well correlated with the presence of bars in galaxies, even if the gas driven by bars fuels the AGN (e.g. Garcia-Burillo et al 2005). In this frame, AGN activity can be triggered in the weakening phases of the bar.
### 3.2 Bar frequency
Recent near-infrared surveys have allowed a statistical estimation of bar strength in local galaxies (Block et al 2002, Laurikainen et al 2004). The observed bar frequency is much larger than what is expected for isolated galaxies, where bars should have been destroyed. These statistics can be used to quantify the gas accretion rate, required to reform bars in the right proportions. Note that the bar frequency has been observed comparable at high redshift (Jogee et al 2004), and therefore gas accretion must have played a role all along the Hubble time.
The gas is required to replenish the disk, and maintain it cold and unstable. Not more than 10% of the accretion can be provided by dwarf companions, since interactions between galaxies heat the disk (Toth & Ostriker 1992), and massive interactions develop the spheroids. The fitting of bar frequency between simulations and observations implies that a galaxy doubles its mass through gas accretion in about 10 Gyr (Bournaud & Combes 2002, Block et al 2002). The source of continuous cold gas accretion can come from the cosmic filamentary structure in the near environment of galaxies. Cosmological accretion in cosmological simulations confirm this gas accretion rate (e.g. Semelin & Combes 2005), and therefore bar reformation.
## 4 Warps and polar rings
The presence of warps in almost all galaxies (conspicuous in HI-21cm) can only be explained through misaligned external gas accretion (e.g. Binney 1992). The amount required can lead to the reorientation of the angular momentum of the galaxy in 7-10 Gyr (Jiang & Binney 1999). This would correspond to the same amount of gas accretion that can explain bar frequency.
Polar Ring Galaxies (PRG) are composed of an early-type host surrounded by a perpendicular ring of young stars and gas, akin to late-type galaxy disks. Stars in the polar ring have formed after the interaction/accretion event, from the gas settled afterwards in the polar plane. The frequency of PRGs deduced from observations is about 5% (Whitmore et al 1990).
The formation of polar rings can be explained either by gas accretion (Schweizer et al 1983, Reshetnikov et al 1997), or by a galaxy merger (Bekki 1997, 1998). External gas can be accreted either from a passing by companion, or from the cosmic web filaments. Numerical simulations reveal that it is about 5 times more probable to form a PRG by gas accretion (Bournaud & Combes 2003).
## 5 Lopsided galaxies
The frequency of asymmetries in galaxies can also help to constrain the gas accretion rate. Peculiar galaxies without any companion are quite frequent, about 50% out of a sample of 1700 galaxies have an HI asymmetric profile (Richter & Sancisi 1994). The asymmetry is even more frequent for late-type galaxies, about 77% (Matthews et al 1998). The asymmetry is also observed in the stellar disk in the optical light (Zaritsky & Rix 1997).
Recently, the amplitude of lopsidedness has been quantified precisely in the NIR distribution (from 150 galaxies of the OSU sample) and compared to what is expected from numerical simulations (Bournaud et al 2005). The $`m=1`$ perturbations could be excited by a companion, or during the formation of the galaxy, since they can persist under the form of kinematic waves for quite a long time, but still not sufficient to explain the observed frequency now (Baldwin et al 1980).
It is difficult to explain all lopsidedness through the tidal interaction from companions, since most galaxies are isolated (Wilcots & Prescott 2004), and for those with visible companions, the amplitude of the $`m=1`$ perturbation does not correlate with the tidal index (proportional to the mass of the companion and inversely to the cube of its distance). In addition, the amplitude of the $`m=1`$ and $`m=2`$ perturbations are correlated, and also the $`m=1`$ intensity correlates with type, which is not predicted for galaxy interactions.
The cause of the $`m=1`$ could then be either a minor merger (no companion is seen), or external gas accretion. Both processes have been simulated, and the amplitude of the $`m=1`$ perturbation then quantified, as a function of time (Bournaud et al 2005). It is found that only asymmetric gas accretion is able to reproduce the high frequency of the asymmetries observed, since the result of a minor merger becomes symmetric quite early (see Figure 5). Only gas accretion (here with 4 M/yr) can explain the observed frequency of $`m=1`$ and the long life-time of the perturbation in NGC 1637. A large number, at least two thirds, of strong $`m=1`$ galaxies require external gas accretion.
## 6 Environmental effects in clusters
The relative importance of the hierarchical scenario and secular evolution must depend strongly on environment. In particular, in rich clusters, mergers have considerably increased the fraction of spheroids and ellipticals, and tidal interactions and ram pressure have stripped and heated the cold gas around galaxies, so that external accretion will be reduced or suppressed.
Although massive galaxies in rich clusters are passively evolving ellipticals, devoid of any star formation, this has not always been the case. Clusters have evolved in a recent past. High resolution images with the HST, followed by spectroscopic surveys, have shown that there exist signs of tidal interaction/mergers in z=0.4 clusters. There is a much larger fraction of perturbed galaxies, a larger fraction of late-type and starbursting objects. Rings of star formation are much more frequent than 2-arm spirals (Oemler et al 1997). And it is well known that there is an excess of blue galaxies as a function of z (Butcher, Oemler 1978, 1984).
Moreover $`z0.5`$ clusters possess a large fraction of peculiar galaxies, likely post-starburst, called E+A (or k+a), devoid of emission lines (and therefore with no current star formation), but very strong Balmer absorption lines (Dressler et al 1999, Poggianti et al 1999). This means that they have a large fraction of A stars, implying that the galaxy was experiencing a strong starburst that has just been abruptly interrupted. Star formation was quenched, in these galaxies in majority disk-dominated. Their fraction is about 20%, much larger than in the field, and in the clusters at $`z=0`$ (Dressler et al 1999).
Star formation and morphological evolution appear decoupled in cluster galaxies (Couch et al 2001), while star formation is still occuring actively now in groups (Balogh et al 2004). The star formation rate is strongly dependent on the local projected density, star formation is quenched as soon as the density is above $`\mathrm{\Sigma }`$ = 1 galaxy/Mpc<sup>2</sup>, independent of the size of the structure. It is likely that the bulk of the stars now in massive galaxies have been formed in spiral galaxies in groups, before merging into ellipticals.
## 7 Conclusion
Three main processes have been invoked for the mass assembly of galaxies: monolithic collapse (MC), hierarchical merging (HS) and external gas accretion associated with secular evolution (SE). The MC scenario appears limited at high redshift, involves only sub-components of galaxies and has a limited role, while the two others (HS and SE) compete with comparable weights to galaxy formation.
Observations constrain the relative importance of these processes. The star formation history across the Hubble time, either through the fossil record of present stellar populations, or direct observatons of galaxies at various redshifts, reveals that massive galaxies have now stopped their star forming activity, while intermediate-mass spirals are still continuing with an almost constant rate, and dwarf galaxies experience starbursts. The observation of massive galaxies at high redshift is compatible with the predictions of the $`\mathrm{\Lambda }`$CDM model of galaxy formation. The observation of the dynamical state of galaxies (bars, spirals, warps, polar rings, $`m=1`$ asymmetries) also constrain the importance of external gas accretion, and secular evolution.
The different processes play a different role, according to the environment. In the field, gas accretion is dominant, to reform bars and spirals, to explain warps, polar rings or lopsidedness. In rich environments, the evolution proceeds at a much faster pace, hierarchical merging had much more importance, and secular evolution of galaxies is halted at z $``$ 1, since galaxies are stripped from their gas, and from their cold filamentary structure, acting as a gas reservoir.
## Acknowledgments
Many thanks to the organisers, and in particular the chairmen D. ELbaz and H. Aussel, for inviting me to such a nice and fruitful conference.
## References
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# 1 Deformation
## 1 Deformation
Let $`X=G/P`$ be a homogeneous space of a Lie group $`G`$, and consider a submanifold $`MX`$ defined by an immersion
$$f:MX.$$
(2)
In many geometric situations, the standard reduction process of moving frame method gives rise to an adapted subbundle $`E_f:B_fG`$ with a structure group $`HP`$ in a canonical way \[Ga\]\[Ch\].
Such $`B_f`$ captures the geometry of $`f`$ in the following sense. Let $`𝔤`$ be the Lie algebra of $`G`$, and let $`\varphi `$ be the $`𝔤`$-valued left invariant Maurer-Cartan form of $`G`$: Let $`f_1`$, $`f_2`$ be two immersions with the induced $`H`$-bundles $`B_{f_1}`$, $`B_{f_2}`$ respectively. $`f_1`$ and $`f_2`$ are congruent up to a left motion by an element of $`G`$ iff there exists a $`H`$-bundle isomorphism of the pair $`(B_{f_1},E_{f_1}^{}\varphi )`$ and $`(B_{f_2},E_{f_2}^{}\varphi )`$ \[Br1\]\[Gr\].
We wish to describe in this section the geometry of the deformation of a submanifold in this setting. Assume for definiteness $`g:G\text{GL}(N,)`$ or $`\text{GL}(N,)`$ for an integer $`N`$. Then in particular $`\varphi =g^1dg`$, and $`\varphi `$ satisfies Maurer-Cartan structure equation
$$d\varphi =\varphi \varphi .$$
Consider a deformation $`f_t`$ of $`f_0=f`$ for $`t(ϵ,ϵ)`$, and let $`E_t:B_tG`$ be the associated adapted $`H`$-bundle.
A general element $`E_t=(e_1(t),e_2(t),\mathrm{},e_N(t))`$ may be written as
$$e_A(t)=e_B(0)g_A^B(t),A,B=1,\mathrm{},N.$$
Equivalently in matrix form
$$E_t=E_0g_t,$$
where $`g_t=(g_A^B(t))`$ is the unique $`G`$-valued function on $`B_f\times (ϵ,ϵ)`$ that describes the deformation of the bundle $`B_f`$. Let us expand $`g_t`$ (formally) with respect to $`t`$
$$g_t=\mathrm{exp}(\underset{k=1}{\overset{\mathrm{}}{}}\frac{t^k}{k!}U_k)$$
for a sequence of $`𝔤`$-valued functions $`\{U_k\}`$ on $`B_f`$. Then
$`dE_t`$ $`=dE_0g_t+E_0dg_t`$
$`=E_0\pi _0g_t+E_0dg_t`$
$`=E_t(g_t^1\pi _0g_t+g_t^1dg_t)`$
$`=E_t\pi _t`$
where $`\pi _t=E_t^1dE_t=E_t^{}\varphi `$. The expression for $`\pi _t=_{k=1}^{\mathrm{}}\frac{t^k}{k!}\pi _k`$ can be computed as follows.
$`\pi _1`$ $`=dU_1+\pi _0U_1U_1\pi _0`$ (3)
$`\pi _2`$ $`=dU_2+\pi _0U_2U_2\pi _0+\pi _1U_1U_1\pi _1`$
$`\mathrm{}`$
$`\pi _k`$ $`dU_k+\pi _0U_kU_k\pi _0mod\{\pi _{k1},\pi _{k2},\mathrm{}\pi _1\}`$
$`\mathrm{}`$
On the other hand, the deformations we are interested in are not arbitrary ones but those deformations that satisfy certain geometric constraint. Since the induced Maurer-Cartan form $`\pi _t`$ captures the geometry of the immersion $`f_t`$, it is reasonable to define the condition of admissible deformation in terms of *a finite set of ordinary differential relations among the coefficients of $`\pi _t`$*. These relations are in turn expressed as a sequence of (algebraic) relations among the coefficients of $`\{\pi _1,\pi _2,\mathrm{}\}`$.
*Remark.* There are situations where the condition of admissible deformation is given in a form which is more general than a sequence of algebraic relations among the coefficients of $`\{\pi _1,\pi _2,\mathrm{}\}`$, e.g., the jet rigidity of homogeneous varieties in projective spaces \[LM\]. The condition of admissible deformation in this case may involve a sequence of differential relations among the coefficients of $`\{\pi _1,\pi _2,\mathrm{}\}`$.
###### Definition 1.1
Let $`f:MX`$ be a submanifold in a homogeneous space $`X=G/P`$. Suppose a condition of admissible deformation is given as a finite set of ordinary differential relations among the coefficients of $`\pi _t`$. The *equation of deformation* at $`f_0=f`$ is the sequence of first order equations (3) for a sequence of $`𝔤`$-valued functions $`\{U_1,U_2,\mathrm{}\}`$ on the adapted $`H`$-bundle $`B_f`$ with $`\{\pi _1,\pi _2,\mathrm{}\}`$ satisfying the sequence of relations that correspond to the given set of differential relations on $`\pi _t`$.
The idea of this set up is to lift the problem of deformation to Lie group $`G`$ where a uniform treatment is available through Maurer-Cartan form and its structure equation.
Note by expanding the structure equation
$$d\pi _t=\pi _t\pi _t$$
with respect to $`t`$, we obtain a sequence of differential equations for $`\{\pi _1,\pi _2,\mathrm{}\}`$.
$`d\pi _1`$ $`=\pi _0\pi _1+\pi _1\pi _0,`$ (4)
$`d\pi _2`$ $`=\pi _0\pi _2+\pi _2\pi _0+2\pi _1\pi _1,`$
$`\mathrm{}`$
$`d\pi _k`$ $`={\displaystyle \underset{i+j=k}{}}{\displaystyle \frac{k!}{i!j!}}\pi _i\pi _j,`$
$`\mathrm{}`$
These are in fact the set of compatibility conditions obtained by differentiating (3) once.
We now introduce a certain *group action* on the space of deformations that is tangent to Cauchy characteristics of the deformation equation, which plays an important role in our treatment of rigidity questions.
Assume a deformation $`E_t`$ is given. Suppose $`h_t`$ is a $`H`$-valued function on $`B_f\times (ϵ,ϵ)`$, and consider a new deformation
$$\widehat{E_t}=E_th_t.$$
Let $`\Pi :GX=G/P`$ be the projection map. Since $`HP`$, $`\Pi (\widehat{E_t})=\Pi (E_t)=f_t`$ describes the same deformation of the immersion $`f`$. This action of $`C^{\mathrm{}}(B_f\times (ϵ,ϵ),H)`$ simply rotates along the fiber of $`B_tM`$.
The nature of this ambiguity is similar to the one in the reduction process of moving frame method. In order to see this group action explicitly, let us write $`h_t=\mathrm{exp}(_{k=1}^{\mathrm{}}\frac{t^k}{k!}V_k)`$ for a sequence of $`𝔥`$-valued functions $`\{V_1,V_2,\mathrm{}\}`$ on $`B_f`$, where $`𝔥`$ is the Lie algebra of $`H`$. Then $`V_k`$’s *acts* on $`\pi _k`$’s by translation as follows.
$`\mathrm{\Delta }\pi _1`$ $`=\pi _0V_1V_1\pi _0`$ (5)
$`\mathrm{\Delta }\pi _2`$ $`\pi _0V_2V_2\pi _0mod\{V_1\}`$
$`\mathrm{}`$
$`\mathrm{\Delta }\pi _k`$ $`\pi _0V_kV_k\pi _0mod\{V_{k1},V_{k2},\mathrm{}V_1\}`$
$`\mathrm{}`$
Here $`\mathrm{\Delta }\pi _k`$ stands for the contribution from $`\{V_1,V_2,\mathrm{}\}`$. We will utilize this extra degree of freedom to normalize certain coefficients in the rigidity related computations. This is analogous to the process of *absorption* in the method of equivalence \[Ga\].
## 2 1-rigidity
The notion of 1-rigidity under admissible deformations of a submanifold in a homogenous space is introduced in this section. 1-rigidity can be considered as a geometric and unifying definition of infinitesimal rigidity in classical differential geometry. We continue to use the notations adopted in Section 1.
Suppose in the deformation equation (3)
$$\pi _i=0\text{for}\mathrm{\hspace{0.33em}1}ik,$$
or equivalently
$$dU_i+\pi _0U_iU_i\pi _0=0\text{for}\mathrm{\hspace{0.33em}1}ik.$$
Then it is easy to check $`U_i=E_0^1a_iE_0`$ for a constant $`a_i𝔤`$, $`1ik`$, and
$`E_t`$ $`=E_0g_t=E_0\mathrm{exp}({\displaystyle \frac{t^i}{i!}}U_i)`$
$`\mathrm{exp}({\displaystyle \underset{i=1}{\overset{k}{}}}{\displaystyle \frac{t^i}{i!}}a_i)E_0modt^{k+1}.`$
The vanishing of the first $`k`$ $`\pi _i`$’s thus means that the deformation osculates a one parameter family of motions by $`G`$ up to order $`k`$.
###### Definition 2.1
Suppose a condition of admissible deformations is given as a finite set of ordinary differential relations among the coefficients of $`\pi _t`$. A submanifold in a homogeneous space is *1-rigid* under admissible deformations if every solution to the deformation equation (3) necessarily has $`\pi _1=0`$ modulo the group action (5).
The original deformation equation (3) is however rather difficult to use in practice to verify 1-rigidity of a submanifold. Instead, note $`\pi _1`$ satisfies the compatibility equation (4),
$$d\pi _1=\pi _0\pi _1+\pi _1\pi _0.$$
For all practical purposes, it is this equation modulo the group action (5) that we will use to test for 1-rigidity. The process of simplification of this sort, that is removing nongeometric terms by differentiating once, has been capitalized in \[BG\] in their study of characteristic cohomology.
As it is the case with the method of moving frames, this definition is best explained and understood in practice. We consider in this section as an example a submanifold in the homogeneous space of conformal sphere $`S^m`$, and its conformal deformation. In the course of computation, a set of criteria for conformal 1-rigidity of a submanifold naturally emerges. This example also serves as a guide through the computation for CR 1-rigidity in Section 4. For general reference for conformal geometry, \[Br2\].
Let $`^{m+1,1}`$ be given a metric of signature $`(m+1,1)`$. Let $`G=`$ SO$`{}_{0}{}^{}(m+1,1)`$ be the identity component of the group of linear transformations that preserve the metric. Let $`X=S^m`$ be the set of *positive* null rays through the origin. It is well known that SO$`{}_{0}{}^{}(m+1,1)`$ acts transitively on $`S^m`$, and that $`S^m`$ inherits an SO$`{}_{0}{}^{}(m+1,1)`$ invariant conformal structure.
Let $`f:MX`$ be an $`n`$-dimensional submanifold. Upon an appropriate choice of basis for $`^{m+1,1}`$, we may arrange so that the induced Maurer-Cartan form $`\pi `$ on a *1-adapted* $`H`$-bundle $`B_f`$ takes the following form, \[Br2\] for details on 1-adapted bundle $`B_f`$.
$$\pi =\left(\begin{array}{cccc}\pi _0^0& \pi _j^0& \pi _b^0& 0\\ \pi _0^i& \pi _j^i& \pi _b^i& \pi _{m+1}^i\\ 0& \pi _j^a& \pi _b^a& \pi _{m+1}^a\\ 0& \pi _j^{m+1}& 0& \pi _{m+1}^{m+1}\end{array}\right)$$
(6)
where $`\pi _j^i=\pi _i^j`$, $`\pi _b^a=\pi _a^b`$, $`\pi _j^{m+1}=\pi _0^j`$, $`\pi _j^0=\pi _{m+1}^j`$, $`\pi _b^0=\pi _{m+1}^b`$, $`\pi _{m+1}^{m+1}=\pi _0^0`$, and $`\mathrm{\hspace{0.17em}1}i,jn`$, $`n+1a,bm`$. $`\pi _0^i`$’s are semibasic 1-forms, and the quadratic form $`\pi _0^i\pi _0^i`$ represents the induced conformal structure on $`M`$.
$`\pi `$ satisfies the structure equations $`d\pi =\pi \pi `$. Differentiating $`\pi _0^a=0`$, we get $`\pi _i^a\pi _0^i=0`$, and by Cartan’s Lemma
$$\pi _i^a=h_{ij}^a\pi _0^j$$
for a set of coefficients $`h_{ij}^a=h_{ji}^a`$. Using the group action by $`\pi _{m+1}^a`$, we normalize so that tr $`h_{ij}^a=0`$. The trace free parts of the second fundamental forms $`H^a=h_{ij}^a\pi _0^i\pi _0^j`$ are conformal invariant of the submanifold $`M`$. Note that an element of the Lie algebra $`𝔥`$ of $`H`$ is of the following form.
$$\left(\begin{array}{cccc}V_0^0& V_j^0& 0& 0\\ 0& V_j^i& 0& V_{m+1}^i\\ 0& 0& V_b^a& 0\\ 0& 0& 0& V_{m+1}^{m+1}\end{array}\right)$$
Suppose we are interested in 1-rigidity under the deformations that preserve the induced conformal structure on $`M`$. From (6), this is written as an ordinary differential equation
$$\frac{d}{dt}\pi _0^i(t)\pi _0^i(t)|_{t=0}=\lambda \pi _0^i(0)\pi _0^i(0)$$
for a scaling factor $`\lambda `$. Using the group action by $`V_0^0,V_j^i`$ components, however, we may in fact normalize so that
$$\frac{d}{dt}\pi _0^i(t)|_{t=0}=0.$$
Hence $`\pi _1=\delta \pi `$ is of the following form.
$$\delta \pi =\left(\begin{array}{cccc}\delta \pi _0^0& \delta \pi _j^0& \delta \pi _b^0& 0\\ 0& \delta \pi _j^i& \delta \pi _b^i& \delta \pi _{m+1}^i\\ 0& \delta \pi _j^a& \delta \pi _b^a& \delta \pi _{m+1}^a\\ 0& 0& 0& \delta \pi _{m+1}^{m+1}\end{array}\right)$$
(7)
The remaining *group variables* at this stage are $`V_b^a`$ and $`V_i^0=V_{m+1}^i`$.
Differentiating $`\delta \pi _0^i=0,\delta \pi _0^a=0`$ in (7) using (3), and by Cartan’s lemma, it is easy to check there exists a set of coefficients $`p_i,p_{ij}^a=p_{ji}^a`$ such that
$`\delta \pi _0^0`$ $`=p_i\pi _0^i,`$
$`\delta \pi _j^i`$ $`=p_i\pi _0^jp_j\pi _0^i,`$
$`\delta \pi _i^a`$ $`=p_{ij}^a\pi _0^j.`$
Using the group action by $`V_j^0`$ component as above, we may translate $`p_j=0`$. Now differentiating $`\delta \pi _0^0=0,\delta \pi _j^i=0`$ with these relations, we obtain the following compatibility conditions for deformation.
$`\delta \pi _i^0`$ $`=p_{ij}\pi _0^j,p_{ij}=p_{ji},`$
$`\pi _i^a`$ $`\delta \pi _j^a+\delta \pi _i^a\pi _j^a\pi _0^i\delta \pi _j^0\delta \pi _i^0\pi _0^j=0.`$ (8)
Fix a point $`pM`$, and let $`V^n=T_pM`$, $`W^{mn}=N_pM`$ represent the tangent space and the normal space respectively. Take a conformal basis $`\{x^i\}`$ for $`V^{}`$ and $`\{w_a\}`$ for $`W`$. Set
$`H=h_{ij}^ax^ix^jw_aS^2V^{}W,`$
$`P=p_{ij}^ax^ix^jw_aS^2V^{}W.`$
Let $`K(V)^2V^{}^2V^{}`$ be the space of curvature like tensors, which decomposes into $`K(V)=𝒲Ric`$, Weyl curvature tensor and Ricci tensor. Then the equation (8) is equivalent to
$`\gamma (H,P)^𝒲=0,`$ (9)
$`\gamma (H,P)^{Ric}=p_{ij}x^ix^j`$ (10)
where $`\gamma (H,P)`$ is the linearized Gauss map \[BBG\]. An element $`HS^2V^{}W`$ is called *Weyl rigid* when the space of solutions to (9) is $`\{P=v_b^ah_{ij}^bx^ix^jw_a|v_b^a=v_a^b\}`$. Note for any $`QS^2V^{}W`$, $`\gamma (Q,P)^𝒲=\gamma (Q_0,P)^𝒲`$, where $`Q_0`$ is the trace free part of $`Q`$. We mention in passing that when dim$`W=1`$, a traceless quadratic form $`H`$ is Weyl rigid if it has 3 nonzero eigenvalues.
Assume the trace free part of the second fundamental form $`H`$ of $`M`$ is Weyl rigid. Then using the group action by $`V_b^a`$ component, we may translate $`\delta \pi _i^a=0`$. It follows $`\delta \pi _i^0=0`$ by (10), and the only remaining nonzero elements of $`\delta \pi `$ are $`\delta \pi _b^0`$ and $`\delta \pi _b^a`$. Differentiating $`\delta \pi _i^a=0,\delta \pi _i^0=0`$, we finally get
$`\delta \pi _b^0\pi _i^b=0,`$ (11)
$`\delta \pi _b^a\pi _i^b=0.`$ (12)
An element $`HS^2V^{}W`$ is nondegenerate if the associated map $`H:S^2VW`$ is surjective. We first show that for a nondegenerate $`H`$, (12) implies $`\delta \pi _b^a=0`$. Let $`\delta \pi _b^a=q_{bi}^a\pi ^i`$, $`q_{bi}^a=q_{ai}^b`$. Then (12) is equivalent to $`q_{bi}^ah_{jk}^b=q_{bk}^ah_{ji}^b=q_{ijk}^a`$ with $`q_{ijk}^a`$ fully symmetric in lower indices. Multiplying both sides by $`h_{sp}^a`$ and summing over $`a`$,
$`q_{ijk}^ah_{sp}^a`$ $`=q_{bi}^ah_{jk}^bh_{sp}^a`$
$`=q_{ai}^bh_{jk}^bh_{sp}^a`$
$`=q_{isp}^bh_{jk}^b.`$
But $`q_{ijk}^a`$ is symmetric in lower indices and $`h_{ij}^a`$ is nondegenerate, hence $`q_{ijk}^a=0`$ and consequently $`q_{bi}^a=0`$.
We next show if $`H`$ is both nondegenerate and Weyl rigid, then (11) implies $`\delta \pi _b^0=0`$. Let $`\delta \pi _b^0=q_{bi}\pi ^i`$. Then (11) is equivalent to $`q_{bi}h_{jk}^b=q_{bk}h_{ji}^b=q_{ijk}`$ with $`q_{ijk}`$ fully symmetric in lower indices. For an arbitrary set of coefficients $`\{f_i\}`$, set
$$Q=Q^ae_a=(q_{ai}f_j+q_{aj}f_i)x^ix^je_aS^2V^{}W.$$
Then one easily sees $`\gamma (H,Q)^𝒲=0`$. Since $`H`$ is Weyl rigid, $`q_{ai}f_j+q_{aj}f_i=v_b^ah_{ij}^b`$ for skew symmetric $`v_b^a=v_a^b`$. Multiplying both sides by $`h_{sp}^ax^ix^jx^sx^p`$ and summing over $`a,i,j,s,p`$, we get
$$(f_ix^i)(q_{jsp}x^jx^sx^p)=0.$$
Since $`\{f_i\}`$ is arbitrary, $`q_{ijk}=0`$, and consequently $`q_{bi}=0`$ for $`h_{ij}^a`$ is nondegenerate.
Thus *a submanifold $`MS^m`$ is conformally 1-rigid if it has a nondegenerate Weyl rigid second fundamental form*. One can apply the results of \[KT\] and show that the canonical isometric embeddings of irreducible compact Hermitian symmetric spaces into spheres satisfy this criteria, and hence they are conformally 1-rigid.
## 3 CR submanifolds
### 3.1 Fundamental forms
In this section, we set up the basic structure equations for CR submanifolds in spheres. A sequence of invariants called *fundamental forms* are derived from the structure equations. For general reference in CR geometry, \[ChM\]\[Ja\].
Let $`^{m+1,1}`$ be the complex vector space with coordinates $`z=(z^0,z^A,z^{m+1})`$, $`\mathrm{\hspace{0.17em}1}Am`$, and a Hermitian scalar product
$$z,\overline{z}=z^A\overline{z}^A+\text{i}(z^0\overline{z}^mz^m\overline{z}^0).$$
Let $`\mathrm{\Sigma }^m`$ be the set of equivalence classes up to scale of null vectors with respect to this product. Let SU$`(m+1,1)`$ be the group of unimodular linear transformations that leave the form $`z,\overline{z}`$ invariant. Then SU$`(m+1,1)`$ acts transitively on $`\mathrm{\Sigma }^m`$, and
$$p:\text{SU}(m+1,1)\mathrm{\Sigma }^m=\text{SU}(m+1,1)/P$$
for an appropriate subgroup $`P`$ \[ChM\].
Explicitly, consider an element $`Z=(Z_0,Z_A,Z_{m+1})\text{SU}(m+1,1)`$ as an ordered set of $`(m+2)`$-column vectors in $`^{m+1,1}`$ such that det$`(Z)=1`$, and that
$$Z_A,\overline{Z}_B=\delta _{AB},Z_0,\overline{Z}_{m+1}=Z_{m+1},\overline{Z}_0=\text{i},$$
(13)
while all other scalar products are zero. We define $`p(Z)=[Z_0]`$, where $`[Z_0]`$ is the equivalence class of null vectors represented by $`Z_0`$. The left invariant Maurer-Cartan form $`\varphi `$ of SU$`(m+1,1)`$ is defined by the equation
$$dZ=Z\varphi ,$$
which is in coordinates
$$d(Z_0,Z_A,Z_{m+1})=(Z_0,Z_B,Z_{m+1})\left(\begin{array}{ccc}\varphi _0^0& \varphi _A^0& \varphi _{m+1}^0\\ \varphi _0^B& \varphi _A^B& \varphi _{m+1}^B\\ \varphi _0^{m+1}& \varphi _A^{m+1}& \varphi _{m+1}^{m+1}\end{array}\right).$$
(14)
Coefficients of $`\varphi `$ are subject to the relations obtained from differentiating (13) which are
$`\varphi _0^0+\overline{\varphi }_{m+1}^{m+1}`$ $`=0`$
$`\varphi _0^{m+1}`$ $`=\overline{\varphi }_0^{m+1},\varphi _{m+1}^0=\overline{\varphi }_{m+1}^0`$
$`\varphi _A^{m+1}`$ $`=\text{i}\overline{\varphi }_0^A,\varphi _{m+1}^A=\text{i}\overline{\varphi }_A^0`$
$`\varphi _B^A+\overline{\varphi }_A^B`$ $`=0`$
$`\text{tr}\varphi `$ $`=0,`$
and $`\varphi `$ satisfies the structure equation
$$d\varphi =\varphi \varphi .$$
(15)
It is well known that the SU$`(m+1,1)`$-invariant CR structure on $`\mathrm{\Sigma }^mP^{m+1}`$ as a real hypersurface is biholomorphically equivalent to the standard CR structure on $`S^{2m+1}=𝔹^{m+1}`$, where $`𝔹^{m+1}^{m+1}`$ is the unit ball. The structure equation (14) shows that for any local section $`s:\mathrm{\Sigma }^m`$ SU$`(m+1,1)`$, this CR structure is defined by the hyperplane fields $`(s^{}\varphi _0^{m+1})^{}=`$ and the set of (1,0)-forms $`\{s^{}\varphi _0^A\}`$.
###### Definition 3.1
Let $`M`$ be a manifold of dimension $`\mathrm{\hspace{0.17em}2}n+1`$. A submanifold defined by an immersion $`f:M\mathrm{\Sigma }^m`$ is a *CR submanifold* if $`f_{}T_pM_{f(p)}`$ is a complex subspace of $`_{f(p)}`$ of dimension $`n`$ for each $`pM`$.
Note the induced hyperplane fields $`f_{}^1(f_{}(TM))`$ is necessarily a contact structure on $`M`$, and thus $`M`$ has an induced nondegenerate CR structure.
Let $`f:M\mathrm{\Sigma }^m`$ be a CR submanifold. For each $`pM`$, we wish to define an associated decomposition of $`^{m+1,1}`$, which is a part of the reduction process of moving frame method applied to CR submanifold. Let $`\mathrm{\Pi }:^{m+1,1}\{0\}P^{m+1}`$ denote the projection map.
Set $`V_0=\mathrm{\Pi }^1(f(p))`$ and $`V_0+V_1=\mathrm{\Pi }^1(f(p)),\mathrm{\Pi }^1(f)`$, where $`f`$ stands for the holomorphic or (1,0) derivatives of $`f`$. Successively define the sequence of subspaces $`W_2,W_3,\mathrm{},W_\tau `$ with dim $`W_l=r_l`$ so that
$$(V_0+V_1)W_2\mathrm{}W_l=\mathrm{\Pi }^1(f(p)),\mathrm{\Pi }^1(f),\mathrm{\Pi }^1(^2f),\mathrm{}\mathrm{\Pi }^1(^lf)$$
(16)
and $`n+r_2+\mathrm{}r_\tau =m`$, where $`^lf`$ stands for the $`l`$-th order (1,0) derivatives of $`f`$. The set of numbers $`(r_2,r_3,\mathrm{}r_\tau )`$ are called the *type numbers* of $`f`$ at $`p`$, and $`\tau `$ is called the *height*.
*Remark.* The structure equation (14), (19), and $`\pi _A^0=\text{i}\overline{\pi }_{m+1}^A`$ show this orthogonal decomposition, rather than a filtration, is well defined.
###### Definition 3.2
A CR submanifold is of *constant type* if the type numbers are constant.
We assume $`M\mathrm{\Sigma }^m`$ is a CR submanifold of constant type from now on.
In terms of the structure equation (14), we may arrange so that
$`V_0`$ $`=Z_0`$ (17)
$`(V_0+V_1)`$ $`=Z_0+Z_i_{i=1}^n`$
$`(V_0+V_1)W_2`$ $`=\left(Z_0+Z_i_{i=1}^n\right)Z_{i_2}_{i_2=n+1}^{n_2}`$
$`\mathrm{}`$
$`(V_0+V_1)W_2\mathrm{}W_l`$ $`=\left(Z_0+Z_i_{i=1}^n\right)Z_{i_2}_{i_2=n+1}^{n_2}\mathrm{}Z_{i_l}_{i_l=n_{l1}+1}^{n_l}`$
$`\mathrm{}`$
$`(V_0+V_1)W_2\mathrm{}W_\tau `$ $`=\left(Z_0+Z_i_{i=1}^n\right)Z_{i_2}_{i_2=n+1}^{n_2}\mathrm{}Z_{i_\tau }_{i_\tau =n_{\tau 1}+1}^{n_\tau },`$
where $`n_l=n_{l1}+r_l`$ and $`n_1=n`$. This decomposition then induces an associated adapted subbundle $`B`$
$$E:f^{}(\text{SU}(m+1,1))B\text{SU}(m+1,1)$$
on which the following weak structure equations hold. Let $`\pi =E^{}\varphi `$.
$`\pi _0^A`$ $`=0\text{for}n+1Am,`$ (18)
$`\pi _i^A`$ $`0mod\theta ,\pi ^{0,1}\text{for}n_2+1Am,`$
$`\pi _{i_l}^A`$ $`0mod\theta ,\pi ^{0,1}\text{for}n_{l+1}+1Am,2l\tau 2,`$
where $`n_{l1}+1i_ln_l`$ for $`\mathrm{\hspace{0.17em}2}l\tau `$, and$`mod\pi ^{0,1}`$ means$`mod\overline{\pi }_0^1,\overline{\pi }_0^2,\mathrm{}\overline{\pi }_0^n`$. (18) implies the following equation in matrix form.
$$\pi \left(\begin{array}{cccccc}& & & & & \\ \pi _0^i& & & & \mathrm{}& \\ 0& \pi _{i_1}^{i_2}& & & & \\ 0& 0& \pi _{i_2}^{i_3}& & & \\ & & \mathrm{}& & \mathrm{}& \\ 0& 0& 0& 0& & \end{array}\right)mod\pi _0^{m+1},\pi ^{0,1}.$$
Denote $`\pi _0^{m+1}=\theta `$, $`\pi ^i=\pi _0^i`$. Differentiating $`\pi _0^A=0`$ $`\text{for}n+1Am`$, we get $`\pi _i^A\pi ^i+\pi _{m+1}^A\theta =0`$. By Cartan’s lemma, there exists a set of coefficients $`h_{ij}^A=h_{ji}^A,h_i^A,h^A`$ such that
$$\left(\begin{array}{c}\pi _i^A\\ \pi _{m+1}^A\end{array}\right)=\left(\begin{array}{cc}h_{ij}^A& h_i^A\\ h_j^A& h^A\end{array}\right)\left(\begin{array}{c}\pi ^j\\ \theta \end{array}\right).$$
(19)
Note by definition of $`W_2`$, $`h_{ij}^A=0`$ for $`An_2+1`$.
Since the CR structure on $`\mathrm{\Sigma }^m`$ is integrable, the differential ideal generated by $`\{\theta ,\pi ^{0,1}\}`$ is closed. By successively differentiating $`\pi _{i_{l1}}^{i_{l+1}}0`$ $`mod\theta ,\pi ^{0,1}`$ for $`l=2,\mathrm{}\tau 1`$, we obtain the following structure equations.
$$\pi _{i_l}^{i_{l+1}}\pi _{i_{l1}}^{i_l}0mod\theta ,\pi ^{0,1},\text{for}\mathrm{\hspace{0.33em}\hspace{0.33em}2}l\tau 1.$$
(20)
Set $`\pi _{i_l}^{i_{l+1}}h_{i_lj}^{i_{l+1}}\pi ^jmod\theta ,\pi ^{0,1}`$. The sequence of equations (20) then gives
$$h_{i_lj}^{i_{l+1}}h_{i_{l1}k}^{i_l}=h_{i_lk}^{i_{l+1}}h_{i_{l1}j}^{i_l},$$
and thus
$$h_{i_{l1}i}^{i_l}h_{i_{l2}j}^{i_{l1}}\mathrm{}h_{ks}^{i_2}=h_{ij\mathrm{}ks}^{i_l}$$
is fully symmetric in lower indices.
###### Definition 3.3
Let $`M\mathrm{\Sigma }^m`$ be a CR submanifold of constant type with the associated decomposition (17) and the Maurer Cartan form (18). The $`l`$-th fundamental form $`F^l`$ is defined by
$`F^l`$ $`=h_{i_{l1}i}^{i_l}h_{i_{l2}j}^{i_{l1}}\mathrm{}h_{ks}^{i_2}\pi ^i\pi ^j\mathrm{}\pi ^k\pi ^sZ_{i_l}`$ (21)
$`\pi _{i_{l1}}^{i_l}\pi _{i_{l2}}^{i_{l1}}\mathrm{}\pi _i^{i_2}\pi ^iZ_{i_l}mod\theta ,\pi ^{0,1}.`$
A computation with the structure equations (14) shows $`F^l`$ is up to scale a well defined section of the bundle $`S^{l,0}(V^{})W_l`$ over $`M`$, where $`V`$ is the holomorphic or $`(1,0)`$ tangent space of $`M`$. $`F^l`$ is one of the simplest $`l`$-th order extrinsic invariant of the CR submanifold $`M`$.
### 3.2 Bochner rigid submanifolds
The weak structure equations (18) is refined in this section for a class of CR submanifolds whose fundamental forms are algebraically rigid.
We start with an algebraic definition. Let $`V=^n,W=^r`$ with the standard Hermitian metric. Let $`S^{k,0}(V^{})`$ ($`S^{0,k}(V^{}))`$ be the space of holomorphic(anti-holomorphic) homogeneous polynomials of degree $`k`$, and denote $`S^{k,l}=S^{k,0}(V^{})S^{0,l}(V^{})`$. Let $`\{x^i\}`$ be a unitary basis of $`S^{1,0}(V^{})`$. For $`k,l1`$, we define the subspace
$$S_1^{k1,l1}=\{x^k\overline{x}^kf|fS^{k1,l1}\}S^{k,l}.$$
Let $`\{e_a\}`$ be a unitary basis of $`W`$. For $`H=H^ae_aS^{k,0}(V^{})W`$ and $`P=P^ae_aS^{l,0}(V^{})W`$, we define
$$H,\overline{P}=H^a\overline{P}^aS^{k,l}(V^{}).$$
###### Definition 3.4
An element $`HS^{k,0}(V^{})W`$ is *Bochner rigid* if for any $`PS^{k,0}(V^{})W`$, the equation
$$\gamma (H,P)=H,\overline{P}+P,\overline{H}S_1^{k1,k1}S^{k,k}$$
implies
$$\gamma (H,P)=0.$$
(22)
Suppose a Bochner rigid $`H=H^ae_aS^{k,0}(V^{})W`$ is nondegenerate, or equivalently the associated map $`H:S^{k,0}(V)W`$ is surjective. Then by Cartan’s Lemma, (22) implies
$$P^a=u_b^aH^b$$
(23)
for a skew Hermitian matrix $`u_b^a=\overline{u}_a^b`$.
It is known that when dim $`W\frac{1}{2}`$ dim $`V`$, every $`HS^{k,0}(V^{})W`$ is Bochner rigid \[Hu\]. We record the following for later application.
###### Lemma 3.1
Let $`HS^{k,0}(V^{})W`$ be a Bochner rigid polynomial. Suppose for $`BS^{1,0}(V^{})W`$,
$$H,\overline{B}S_1^{k1,0}S^{k,1}.$$
Then we necessarily have
$$H,\overline{B}=0.$$
If $`H`$ is moreover nondegenerate, then $`B=0`$.
*Proof*. Let $`\{x^i\}`$ be a unitary (1,0)-basis of $`V^{}`$, and denote $`H,\overline{B}=(x^i\overline{x}^i)g_B`$ for $`g_BS^{k1,0}`$. Take any $`gS^{k1,0}`$. Then $`\gamma (H,gB)=(x^i\overline{x}^i)(g_B\overline{g}+g\overline{g}_B)=0`$ by Bochner rigidity of $`H`$. Since $`g`$ is arbitrary, $`g_B=0`$. The rest follows from the definition of a nondegenerate form. $`\mathrm{}`$
###### Definition 3.5
Let $`M\mathrm{\Sigma }^m`$ be a CR submanifold of constant type with the associated decomposition (17) and the fundamental forms (21). $`M`$ is a *Bochner rigid submanifold* if each of its fundamental forms is Bochner rigid.
###### Example 3.1
Let $`M_{n,p}P^N=P(^n^{n+p})`$ be the Plücker embedding of the Grassmannian $`Gr(n,^{n+p})`$, $`pn`$. Let $`\widehat{M}_{n,p}S^{2N+1}`$ be the inverse image of $`M_{n,p}`$ under Hopf map, which is an $`S^1`$-bundle over $`M_{n,p}`$. When $`p=2,\mathrm{\hspace{0.17em}3}`$, $`\widehat{M}_{n,p}`$ is a Bochner rigid CR submanifold of $`S^{2N+1}`$.
We postpone the computation of this example to Section 6. It is likely the case that $`\widehat{M}_{n,p}S^{2N+1}`$ is Bochner rigid for all $`p`$.
The main result of this section is the following refined structure equations for Bochner rigid submanifolds. We continue to use the notations of Section 3.1.
###### Theorem 3.1
Let $`M\mathrm{\Sigma }^m`$ be a Bochner rigid submanifold. Then the weak structure equations (18) can be refined as follows.
$`\pi _0^A`$ $`=0\text{for}n+1Am`$ (24)
$`\pi _i^A`$ $`=0\text{for}n_2+1Am`$
$`\pi _{i_l}^A`$ $`=0\text{for}n_{l+1}+1Am,2l\tau 2`$
These equations furthermore imply
$`\pi _{i_l}^{i_{l+1}}`$ $`0mod\theta ,\pi ^{1,0}\text{for}\mathrm{\hspace{0.33em}\hspace{0.33em}1}l\tau 1,`$ (25)
$`\pi _{m+1}^A`$ $`=0\text{for}n_2+1Am.`$
*Proof*. By definition of $`W_2`$ and the weak structure equations (18), we in fact have
$`\pi _i^A`$ $`=h_i^A\theta ,`$ (26)
$`\pi _{m+1}^A`$ $`h_i^A\pi ^imod\theta ,\text{for}n_2+1Am`$
for a set of coefficients $`h_i^A`$. Set $`\pi _{i_2}^AB_{i_2j}^A\overline{\pi }^j`$ mod $`\theta ,\pi ^{1,0}`$ $`\text{for}n_2+1Am`$. Differentiating (26) mod $`\theta `$,
$$\text{i}h_i^A\varpi \pi _{i_2}^A\pi _i^{i_2}+\pi _{m+1}^A\pi _i^{m+1}mod\theta ,\text{for}n_2+1Am,$$
where $`\varpi =\pi ^k\overline{\pi }^k`$, and this gives
$$B_{i_2p}^Ah_{ij}^{i_2}=\text{i}(h_i^A\delta _{jp}+h_j^A\delta _{ip}).$$
But $`h_{ij}^{i_2}`$ represents the second fundamental form $`F^2`$ which is Bochner rigid. By Lemma 3.1, $`B_{i_2p}^A=0,h_j^A=0`$ for $`n_2+1Am`$. We thus have, since $`\pi _{i_2}^A0mod\theta ,\pi ^{0,1}`$ for $`n_3+1Am`$ from (18),
$`\pi _i^A`$ $`=0\text{for}n_2+1Am,`$
$`\pi _{i_2}^A`$ $`0mod\theta ,\pi ^{1,0}\text{for}n_2+1An_3,`$
$`\pi _{i_2}^A`$ $`=h_{i_2}^A\theta \text{for}n_3+1Am.`$ (27)
Set $`\pi _{i_3}^AB_{i_3j}^A\overline{\pi }^j`$ mod $`\theta ,\pi ^{1,0}`$ for $`n_3+1Am`$. Differentiating (27) mod $`\theta `$,
$$h_{i_2}^A\varpi \pi _{i_3}^A\pi _{i_2}^{i_3}mod\theta \text{for}n_3+1Am,$$
for $`\pi _{m+1}^A0mod\theta `$ for $`n_2+1Am`$ from (26), and this gives
$$B_{i_3p}^Ah_{i_2k}^{i_3}=\text{i}h_{i_2}^A\delta _{kp}.$$
Multiplying $`h_{ij}^{i_2}`$ and summing over $`i_2`$ we get
$$B_{i_3p}^Ah_{i_2k}^{i_3}h_{ij}^{i_2}=\text{i}h_{i_2}^Ah_{ij}^{i_2}\delta _{kp}\text{for}n_3+1Am.$$
But $`h_{i_2k}^{i_3}h_{ij}^{i_2}=h_{ijk}^{i_3}`$ represents the third fundamental form $`F^3`$ which is Bochner rigid, and both $`F^2,F^3`$ are nondegenerate by definition. By Lemma 3.1, $`B_{i_3p}^A=0,h_{i_2}^A=0`$ for $`n_3+1Am`$. We thus have by similar argument as before
$`\pi _{i_2}^A`$ $`=0\text{for}n_3+1Am,`$
$`\pi _{i_3}^A`$ $`0mod\theta ,\pi ^{1,0}\text{for}n_3+1An_4,`$
$`\pi _{i_3}^A`$ $`=h_{i_3}^A\theta \text{for}n_4+1Am.`$
Continuing in this manner, it is straightforward to verify (24).
The first equation of (25) is already proved. Once (24) is true, we have from (19) $`\pi _{m+1}^A=p^A\theta `$ $`\text{for}n_2+1Am`$. Differentiating this equation mod $`\theta `$ and collecting (1,1)-terms, we get $`p^A\varpi 0`$ for $`n_2+1Am.\mathrm{}`$
## 4 Rigidity of CR submanifolds
### 4.1 CR 1-rigidity
Let $`f:M\mathrm{\Sigma }^m`$ be a CR submanifold of constant type. A *type preserving CR deformation* of $`f`$ is a deformation of $`f`$ through CR immersions with the same type numbers and the equivalent induced CR structures as $`f`$.
###### Theorem 4.1
Let $`f:M\mathrm{\Sigma }^m`$ be a Bochner rigid CR submanifold of CR dimension $`2`$. Then $`f`$ is 1-rigid under type preserving CR deformations.
The idea is to explore the structure of the compatibility equation (4) in the presence of the associated decomposition (17) and the corresponding adapted structure equation (24), (25). The proof is carried out by a sequence of over-determined pde computations. Simply put, Bochner rigidity of the fundamental forms implies that each adjacent sequence of compatibility conditions of the over-determined pde’s are tightly related. It thus plays the role of connecting the propagation of the associated sequence of compatibility conditions.
It would be interesting and more desirable to remove the type preserving hypothesis in the theorem, which is the case when the height $`\tau =2`$. The computation involved without the hypothesis is, however, more differential algebraic than algebraic which we are not able to resolve.
###### Corollary 4.1
\[Hu\]\[EHZ\] Let $`f:M^n\mathrm{\Sigma }^m`$ be a CR submanifold of constant type and of CR dimension $`n2`$. Suppose each of the type numbers $`(r_2,r_3,\mathrm{}r_\tau )`$ are bounded by $`r_l\frac{n}{2}`$, $`\mathrm{\hspace{0.17em}2}l\tau `$. Then $`f`$ is 1-rigid under type preserving CR deformations.
###### Corollary 4.2
Let $`f_{n,p}:\widehat{M}_{n,p}\mathrm{\Sigma }^m`$ be the canonical $`S^1`$-bundle over the Plücker embedding of the Grassmannian $`Gr(n,^{n+p})`$, $`pn`$. When $`p=2,\mathrm{\hspace{0.17em}3}`$, $`f_{n,p}`$ is 1-rigid under type preserving CR deformations.
We present the proof of the theorem for the case $`f`$ has height $`\tau =3`$, as the proof for the general case can be easily deduced form this. We shall agree on the index range
$`1i,j,k,`$ $`l,p,sn`$
$`n+1a,`$ $`bn_2`$
$`n_2+1\nu ,`$ $`\mu n_3,`$
and we continue to use the notations in Section 3.
*Proof of Theorem 4.1.* Since $`M`$ is a Bochner rigid submanifold, the induced Maurer-Cartan form $`\pi `$ on the associated $`H`$-bundle $`B`$ takes the following form.
$$\pi =\left(\begin{array}{ccccc}\pi _0^0& \pi _j^0& \pi _b^0& 0& \pi _{m+1}^0\\ \pi ^i& \pi _j^i& \pi _b^i& 0& \pi _{m+1}^i\\ 0& \pi _j^a& \pi _b^a& \pi _\mu ^a& \pi _{m+1}^a\\ 0& 0& \pi _b^\nu & \pi _\mu ^\nu & 0\\ \theta & \pi _j^{m+1}& 0& 0& \pi _{m+1}^{m+1}\end{array}\right)$$
(28)
with
$`\pi _j^a`$ $`=h_{jk}^a\pi ^k+h_j^a\theta ,`$
$`\pi _{m+1}^a`$ $`h_j^a\pi ^jmod\theta ,`$
$`\pi _b^\nu `$ $`h_{bk}^\nu \pi ^kmod\theta .`$
Here $`h_{jk}^a`$ and $`h_{bk}^\nu h_{ij}^b=h_{ijk}^\nu `$ are fully symmetric in lower indices representing the second and the third fundamental forms. The Lie algebra of $`H`$-valued group variable $`V`$ in the group action (5) is of the following form at this stage.
$$V=\left(\begin{array}{ccccc}V_0^0& V_j^0& 0& 0& V_{m+1}\\ 0& V_j^i& 0& 0& V_{m+1}^i\\ 0& 0& V_b^a& 0& 0\\ 0& 0& 0& V_\mu ^\nu & 0\\ 0& 0& 0& 0& V_{m+1}^{m+1}\end{array}\right).$$
For notational purpose, let $`\delta \pi `$ denote $`\pi _1`$ in the deformation equation (4). Since the type preserving CR deformation is considered, it is clear that $`\delta \pi `$ should take the following form after some absorption by group actions (5).
$$\delta \pi =\left(\begin{array}{ccccc}\delta \pi _0^0& \delta \pi _j^0& \delta \pi _b^0& \delta \pi _\mu ^0& \delta \pi _{m+1}^0\\ 0& \delta \pi _j^i& \delta \pi _b^i& \delta \pi _\mu ^i& \delta \pi _{m+1}^i\\ 0& \delta \pi _j^a& \delta \pi _b^a& \delta \pi _\mu ^a& \delta \pi _{m+1}^a\\ 0& \delta \pi _j^\nu & \delta \pi _b^\nu & \delta \pi _\mu ^\nu & \delta \pi _{m+1}^\nu \\ 0& 0& 0& 0& \delta \pi _{m+1}^{m+1}\end{array}\right)$$
(29)
Here $`\delta \pi _0^i`$ is translated to $`\mathrm{\hspace{0.17em}0}`$ by $`V_j^i\delta _{ij}V_0^0`$ and $`V_{m+1}^i`$ component, and $`\delta \pi _0^{m+1}`$ is translated to $`\mathrm{\hspace{0.17em}0}`$ by Re$`V_0^0`$ component for it is a CR deformation, and $`\delta \pi _j^\nu 0`$ mod $`\theta ,\pi ^{0,1}`$ for it is type preserving. The remaining group variables at this stage are $`V_0^0,V_b^a,V_{m+1}^0`$, and $`V_\mu ^\nu `$ with $`V_0^0=\overline{V}_0^0`$.
*Step 1.* Differentiating $`\delta \pi _0^i=0,\delta \pi _0^a=0,\delta \pi _0^\nu =0,\delta \pi _0^{m+1}=0`$, we get
$`\mathrm{\Delta }_j^i\pi ^j+\delta \pi _{m+1}^i\theta `$ $`=0`$ (30)
$`\delta \pi _i^a\pi ^i+\delta \pi _{m+1}^a\theta `$ $`=0`$
$`\delta \pi _i^\nu \pi ^i+\delta \pi _{m+1}^\nu \theta `$ $`=0`$
$`\text{Re}\delta \pi _0^0\theta `$ $`=0,`$
where $`\mathrm{\Delta }_j^i=\delta \pi _j^i\delta \pi _0^0\delta _{ij}`$. Using the group action (5) by $`V_{m+1}^0`$ component, we can absorb $`\text{Re}\delta \pi _0^0`$ to $`\mathrm{\hspace{0.17em}0}`$. Then $`(\mathrm{\Delta }_j^i)`$ is a skew Hermitian matrix valued 1-form such that
$$\mathrm{\Delta }_j^i\pi ^j0mod\theta ,$$
and Cartan’s lemma implies $`\mathrm{\Delta }_j^i0`$ mod $`\theta `$. $`\delta \pi _j^\nu 0`$ mod $`\theta ,\pi ^{0,1}`$ and (30) implies $`\delta \pi _j^\nu 0`$ mod $`\theta `$. We thus have
$`\left(\begin{array}{c}\mathrm{\Delta }_j^i\\ \delta \pi _{m+1}^i\end{array}\right)`$ $`=\left(\begin{array}{cc}0& q_j^i\\ q_j^i& q^i\end{array}\right)\left(\begin{array}{c}\pi ^j\\ \theta \end{array}\right),`$
$`\left(\begin{array}{c}\delta \pi _i^a\\ \delta \pi _{m+1}^a\end{array}\right)`$ $`=\left(\begin{array}{cc}p_{ij}^a& p_i^a\\ p_i^a& p^a\end{array}\right)\left(\begin{array}{c}\pi ^j\\ \theta \end{array}\right),`$
$`\left(\begin{array}{c}\delta \pi _i^\nu \\ \delta \pi _{m+1}^\nu \end{array}\right)`$ $`=\left(\begin{array}{cc}0& p_i^\nu \\ p_i^\nu & p^\nu \end{array}\right)\left(\begin{array}{c}\pi ^j\\ \theta \end{array}\right),`$
$`\text{Re}\delta \pi _0^0`$ $`=0,`$
where $`q_i^j=\overline{q}_j^i`$, $`p_{ij}^a=p_{ji}^a`$. Differentiating $`\text{Re}\delta \pi _0^0=0`$ with these relations, we get
$$\delta \pi _{m+1}^0=\frac{\text{i}}{2}(\overline{q}^i\pi ^iq^i\overline{\pi }^i)+q^0\theta ,$$
for a real coefficient $`q^0`$. Note $`\delta \pi _{m+1}^{m+1}=\delta \overline{\pi }_0^0=\delta \pi _0^0`$. The remaining group variables at this stage are $`V_0^0,V_b^a`$, and $`V_\mu ^\nu `$.
*Step 2.* Differentiating $`\mathrm{\Delta }_j^i=q_j^i\theta `$ mod $`\theta `$ and collecting terms, we obtain an equation which is equivalent to
$$h_{jk}^a\overline{p}_{il}^a+p_{jk}^a\overline{h}_{il}^a=\text{i}(q_j^i\delta _{kl}+q_k^i\delta _{jl}+q_j^l\delta _{ik}+q_k^l\delta _{ij}).$$
(31)
Since $`h_{ij}^a`$ is Bochner rigid, this equation implies
$`q_j^i`$ $`=0,`$
$`h_{jk}^a\overline{p}_{il}^a+p_{jk}^a\overline{h}_{il}^a`$ $`=0.`$
From (23), $`p_{ij}^a=u_b^ah_{ij}^b`$ for a skew Hermitian matrix $`(u_b^a)`$, and thus can be absorbed to $`\mathrm{\hspace{0.17em}0}`$ by the group action by $`V_b^a\delta _{ab}V_0^0`$ component. Differentiating $`\delta \pi _{m+1}^i=q^i\theta `$ mod $`\theta `$,
$`q^i\varpi `$ $`\pi ^i\delta \pi _{m+1}^0+\pi _a^i\delta \pi _{m+1}^amod\theta `$
$`\pi ^i({\displaystyle \frac{\text{i}}{2}}(\overline{q}^j\pi ^jq^j\overline{\pi }^j))\overline{h}_{ik}^ap_j^a\overline{\pi }^k\pi ^jmod\theta ,`$
where $`\varpi =\pi ^k\overline{\pi }^k`$. Collecting $`(2,0)`$-terms, we get $`q^i=0`$. Since $`h_{ij}^a`$ represents the nondegenerate second fundamental form, the remaining equation $`\overline{h}_{ik}^ap_j^a=0`$ implies $`p_i^a`$ = 0. Now differentiating $`\pi _{m+1}^a=p^a\theta `$ mod $`\theta `$, we get $`p^a=0`$. $`\delta \pi `$ is now reduced to the following form.
$$\delta \pi =\left(\begin{array}{ccccc}\delta \pi _0^0& 0& 0& \delta \pi _\mu ^0& q^0\theta \\ 0& \delta \pi _0^0\delta _{ij}& 0& \delta \pi _\mu ^i& 0\\ 0& 0& \delta \pi _b^a& \delta \pi _\mu ^a& 0\\ 0& \delta \pi _j^\nu & \delta \pi _b^\nu & \delta \pi _\mu ^\nu & \delta \pi _{m+1}^\nu \\ 0& 0& 0& 0& \delta \pi _0^0\end{array}\right)$$
The remaining group variables at this stage are $`V_0^0`$ and $`V_\mu ^\nu `$.
*Step 3.* Differentiating $`\delta \pi _i^\nu =p_i^\nu \theta `$ mod $`\theta `$, we get
$$p_i^\nu \varpi \delta \pi _a^\nu \pi _i^amod\theta .$$
(32)
Set $`\delta \pi _a^\nu p_{ai}^\nu \pi ^i+C_{ai}^\nu \overline{\pi }^imod\theta `$. Collecting $`(2,0)`$, and $`(1,1)`$ terms in (32) we get
$`p_{ai}^\nu h_{jk}^a`$ $`=p_{aj}^\nu h_{ik}^a=p_{ijk}^\nu \text{fully symmetric in lower indices}`$
$`\text{i}p_i^\nu \delta _{jk}`$ $`=C_{ak}^\nu h_{ij}^a.`$ (33)
By Bochner rigidity of $`h_{ij}^a`$, (33) implies $`C_{ak}^\nu =0`$, $`p_i^\nu =0`$. Differentiating $`\delta \pi _{m+1}^\nu =p^\nu \theta `$ mod $`\theta `$ with these relations gives $`p^\nu =0`$. Differentiating $`\delta \pi _{m+1}^0=q^0\theta `$ mod $`\theta `$ at this stage gives $`q^0=0`$. $`\delta \pi `$ is now reduced to the following form.
$$\delta \pi =\left(\begin{array}{ccccc}\delta \pi _0^0& 0& 0& 0& 0\\ 0& \delta \pi _0^0\delta _{ij}& 0& 0& 0\\ 0& 0& \delta \pi _b^a& \delta \pi _\mu ^a& 0\\ 0& 0& \delta \pi _b^\nu & \delta \pi _\mu ^\nu & 0\\ 0& 0& 0& 0& \delta \pi _0^0\end{array}\right)$$
*Step 4.* Differentiating $`\delta \pi _i^a=0`$ mod $`\theta `$,
$$\mathrm{\Delta }_b^a\pi _i^b0mod\theta $$
where $`\mathrm{\Delta }_b^a=\delta \pi _b^a\delta \pi _0^0\delta _{ab}`$. Since the coefficients of $`\pi _i^ah_{ij}^a\pi ^jmod\theta `$ represent the second fundamental form which is nondegenerate by definition, a variant of Cartan’s lemma, \[GH\], implies $`\mathrm{\Delta }_b^amod\theta `$ is in the span of $`\{\pi _i^a\}`$. But $`\pi _i^a`$ consists of $`(1,0)`$-forms mod $`\theta `$ and $`\mathrm{\Delta }_b^a`$ is skew Hermitian. Hence $`\mathrm{\Delta }_b^a0`$ mod $`\theta `$, and we write
$$\mathrm{\Delta }_b^a=q_b^a\theta $$
where $`q_a^b=\overline{q}_b^a`$. Differentiating this equation, we get $`\text{i}q_b^a\varpi \overline{\pi }_a^\nu \delta \pi _b^\nu +\delta \overline{\pi }_a^\nu \pi _b^\nu mod\theta `$, or equivalently
$$\text{i}q_b^a\delta _{sk}=\overline{h}_{as}^\nu p_{bk}^\nu +\overline{p}_{as}^\nu h_{bk}^\nu .$$
Multiplying both sides by $`\overline{h}_{lp}^ah_{ij}^b`$ and summing with respect to $`a,b`$,
$$\text{i}q_b^a\overline{h}_{lp}^ah_{ij}^b\delta _{sk}=\overline{h}_{lps}^\nu p_{ijk}^\nu +\overline{p}_{lps}^\nu h_{ijk}^\nu .$$
By Bochner rigidity of $`h_{ijk}^\nu `$ representing the third fundamental form, this implies
$`q_b^a`$ $`=0,`$
$`p_{ijk}^\nu `$ $`=p_{ai}^\nu h_{jk}^a=u_\mu ^\nu h_{ijk}^\mu `$
$`=u_\mu ^\nu h_{ai}^\mu h_{jk}^a`$
for a skew Hermitian $`u_\nu ^\mu =\overline{u}_\mu ^\nu `$. Since $`h_{ij}^a`$ is nondegenerate, this implies
$$p_{ai}^\nu =u_\mu ^\nu h_{ai}^\mu ,$$
and we may absorb $`p_{ai}^\nu `$ to $`\mathrm{\hspace{0.17em}0}`$ by group action by $`V_\mu ^\nu \delta _{\nu \mu }V_0^0`$ component as before. Put $`\delta \pi _a^\nu =p_a^\nu \theta `$. Differentiating $`\mathrm{\Delta }_b^a=0`$ and collecting $`\theta \pi ^{1,0}`$-terms, we get $`h_{ai}^\nu \overline{p}_b^\nu =0`$. Multiplying $`h_{jk}^a`$ and summing with respect to $`a`$, we get
$$h_{ijk}^\nu \overline{p}_b^\nu =0.$$
Since $`h_{ijk}^\nu `$ represents the third fundamental form which is nondegenerate by definition, this equation implies $`p_b^\nu =0`$. $`\delta \pi `$ is now reduced to the following form.
$$\delta \pi =\left(\begin{array}{ccccc}\delta \pi _0^0& 0& 0& 0& 0\\ 0& \delta \pi _0^0\delta _{ij}& 0& 0& 0\\ 0& 0& \delta \pi _0^0\delta _{ab}& 0& 0\\ 0& 0& 0& \delta \pi _\mu ^\nu & 0\\ 0& 0& 0& 0& \delta \pi _0^0\end{array}\right).$$
Note there are no remaining group variables at this stage.
*Step 5.* Differentiating $`\delta \pi _a^\nu =0`$ mod $`\theta `$,
$$\mathrm{\Delta }_\mu ^\nu \pi _a^\mu 0mod\theta $$
where $`\mathrm{\Delta }_\mu ^\nu =\delta \pi _\mu ^\nu \delta \pi _0^0\delta _{\nu \mu }`$. Since the coefficients of $`\pi _a^\mu h_{aj}^\mu \pi ^j`$ mod $`\theta `$ represents the nondegenerate third fundamental form, a variant of Cartan’s lemma as before implies
$$\mathrm{\Delta }_\mu ^\nu =\delta \pi _\mu ^\nu \delta \pi _0^0\delta _{\nu \mu }0mod\theta .$$
Now $`\delta \pi `$ takes values in $`𝔰𝔲(m+1,1)`$, and in particular tr$`\delta \pi =0`$. Hence $`\delta \pi _0^0`$ is a multiple of $`\theta `$, and consequently
$$\delta \pi =\text{P}\theta $$
for a scalar $`𝔰𝔲(m+1,1)`$-valued coefficient P. But differentiating this equation mod $`\theta `$ gives $`\text{P}\varpi 0`$, hence $`\text{P}=0`$. $`\mathrm{}`$
Here we summarize the process of computation for 1-rigidity in analogy with the game of dominoes. Imagine $`\delta \pi `$ is a $`(\tau +2)`$-by-$`(\tau +2)`$ standing blocks of dominoes. Our objective is to lay down all of these domino blocks. There are three main ingredients in the process.
A. The initial push: This corresponds to imposing the condition of admissible deformations as in (29).
B. When a block gets hit, it falls toward the nearby standing blocks: This is expressed by the compatibility conditions obtained by differentiation using the deformation equation (4).
C. The nearby standing blocks in process B are within hitting distance, and gets hit by at least one domino block, and falls: This in CR case is guaranteed by the Bochner rigid conditions.
Once there was the initial push A, the processes B and C alternate just as it is the case with real dominoes. We finally mention that in process C, it is conceivable that a standing block gets hit but the force is not enough to lay it down, or it is equally conceivable that a standing block gets hit by more than one blocks and the hitting forces cancel. In these cases, one has to analyze the prolongation of the deformation equation. This phenomena occurs when trying to remove the type preserving hypothesis from Theorem 4.1.
### 4.2 Local rigidity
Although it is formulated based on a natural geometric consideration, 1-rigidity introduced in Section 2 is an abstract definition. It is not clear whether 1-rigidity has any bearing on the actual rigidity property of submanifolds under admissible deformation. The main result of this section is that 1-rigidity can be extended to actual local rigidity at least for Bochner rigid CR submanifolds in spheres. The idea of proof is a certain averaging trick, which might be of use in rigidity problems in other homogeneous spaces as well.
###### Theorem 4.2
Let $`f_0:M\mathrm{\Sigma }^m`$ be a Bochner rigid CR submanifold of height $`\tau `$. Then there exists a $`C^\tau `$ neighborhood $`𝒰C^\tau (M,\mathrm{\Sigma }^m)`$ of $`f_0`$ with the following property. Suppose $`f,g𝒰`$ are CR immersions such that their induced CR structures on $`M`$ are mutually equivalent and that they are of the same type as $`f_0`$. Then $`f`$ and $`g`$ are congruent up to an automorphism of $`\mathrm{\Sigma }^m`$.
Note when $`\tau =2`$, the type preserving hypothesis can be dropped.
We present the proof for the case $`\tau =2`$, for the proof for the general case is essentially the same. We follow the notations in Section 4.1.
*Proof of Theorem 4.2.* Given two CR immersions $`f,g`$, let $`F:B_ff^{}`$SU$`(m+1,1)`$, $`G:B_gg^{}`$SU$`(m+1,1)`$ be the associated adapted $`H`$-bundles over $`M`$. Let
$`s_f`$ $`=(F_0,F_1,\mathrm{}F_{m+1}),`$
$`s_g`$ $`=(G_0,G_1,\mathrm{}G_{m+1})`$
be any sections of these bundles, and denote
$`s_f^{}(F^{}\varphi )`$ $`=\alpha ,`$
$`s_g^{}(G^{}\varphi )`$ $`=\beta `$
where $`\varphi `$ is the Maurer-Cartan form of SU$`(m+1,1)`$. Set
$`\pi `$ $`={\displaystyle \frac{1}{2}}(\alpha +\beta ),`$
$`\delta \pi `$ $`=\alpha \beta ,`$
and observe
$`d(\delta \pi )`$ $`=\alpha \alpha \beta \beta `$
$`=\delta \pi \pi +\pi \delta \pi .`$ (34)
For a fixed section $`s_g`$, we wish to show there exists a section $`s_f`$ such that $`\delta \pi =0`$.
Since $`f`$ and $`g`$ induce the equivalent CR structures on $`M`$, we may take a section $`s_f`$ such that
$`\alpha _0^i`$ $`=\beta _0^i=\pi ^i,`$ (35)
$`\alpha _0^{m+1}`$ $`=\beta _0^{m+1}=\theta `$
and $`\delta \pi `$ takes the following form.
$$\delta \pi =\left(\begin{array}{cccc}\delta \pi _0^0& & & \delta \pi _{m+1}^0\\ 0& & & \\ 0& & & \\ 0& 0& 0& \end{array}\right).$$
Differentiating $`\delta \pi _0^{m+1}=0`$ using (34), we get Re$`\delta \pi _0^0\theta =0`$. Modifying $`F_{m+1}`$ by a multiple of $`F_0`$ if necessary, which is still a legitimate section of $`B_f`$ preserving the relations (35), we may in fact have Re$`\delta \pi _0^0=0`$(this requires a computation, which is rather similar to the absorption by group action as in CR 1-rigidity and we omit). Now differentiating $`\delta \pi _0^i=0`$ we get
$$\left(\begin{array}{c}\mathrm{\Delta }_j^i\\ \delta \pi _{m+1}^i\end{array}\right)=\left(\begin{array}{cc}0& q_j^i\\ q_j^i& q^i\end{array}\right)\left(\begin{array}{c}\pi ^j\\ \theta \end{array}\right)$$
(36)
where $`\mathrm{\Delta }_j^i=\delta \pi _j^i\delta \pi _0^0\delta _{ij}`$, $`q_i^j=\overline{q}_j^i`$.
Let $`f^a=f_{ij}^a\pi ^i\pi ^j`$ and $`g^a=g_{ij}^a\pi ^i\pi ^j`$ represent the second fundamental forms of $`f`$ and $`g`$. If $`f`$ and $`g`$ are sufficiently $`C^2`$ close to $`f_0`$ which is Bochner rigid, then $`f`$ and $`g`$ would be Bochner rigid. Thus the average $`\frac{1}{2}(f^a+g^a)`$ would be Bochner rigid too for a section $`s_f`$ with appropriate $`(F_{n+1},F_{n+2},\mathrm{}F_{n+r})`$-part. Differentiating (36) and proceeding as in (31), we get
$`q_j^i`$ $`=0,`$
$`f^ag^a`$ $`=u_b^a(f^a+g^a)`$ (37)
for a skew Hermitian matrix $`(u_b^a)=u=u^{}`$.
Set $`\stackrel{}{f}=(f^{n+1},f^{n+2},\mathrm{}f^{n+r})^t`$, $`\stackrel{}{g}=(g^{n+1},g^{n+2},\mathrm{}g^{n+r})^t`$, and write (37) in matrix form
$$(Iu)\stackrel{}{f}=(I+u)\stackrel{}{g}.$$
When $`f`$ and $`g`$ are sufficiently $`C^2`$ close to $`f_0`$, $`u`$ is small and $`Iu`$ would be invertible. Thus we may write
$$\stackrel{}{f}=(Iu)^1(I+u)\stackrel{}{g}.$$
But
$`(Iu)^1(I+u)((Iu)^1(I+u))^{}`$ $`=(Iu)^1(I+u)((Iu)(I+u)^1`$
$`=(Iu)^1(Iu)((I+u)(I+u)^1`$
$`=I.`$
Thus we may modify $`(F_{n+1},F_{n+2},\mathrm{}F_{n+r})`$ part of the section $`s_f`$ by unitary matrices close to identity to have
$$\delta \pi _i^a=0.$$
$`\delta \pi `$ now becomes
$$\delta \pi =\left(\begin{array}{cccc}\delta \pi _0^0& & & \delta \pi _{m+1}^0\\ 0& \delta \pi _0^0\delta _{ij}& 0& \\ 0& 0& & \\ 0& 0& 0& \delta \pi _0^0\end{array}\right).$$
The rest of the computation proceeds the same as in the proof of Theorem 4.1, and we conclude for this section $`s_f`$,
$$\delta \pi =\alpha \beta =0.$$
By the fundamental theorem \[Gr, p 780\], the sections $`s_f`$ and $`s_g`$ are congruent by an element of SU($`m+1,1`$), and hence $`f`$ and $`g`$ are congruent by an automorphism of $`\mathrm{\Sigma }^m`$. $`\mathrm{}`$
## 5 Whitney submanifold
The main result of this section is the local characterization that every nonlinear CR-flat submanifold $`M^n\mathrm{\Sigma }^{2n}`$, $`n2`$, is a part of a *Whitney submanifold*. Whitney submanifold is an example of a CR submanifold which is not 1-rigid. The result of this section in fact implies it is CR deformable in exactly 1 direction, see the end of this section.
Let $`V^{n+2}=^{n+1,1}`$ be the complex vector space with coordinates $`\xi =(\xi ^0,\xi ^i,\xi ^{n+1})`$, $`\mathrm{\hspace{0.17em}1}in`$, and a Hermitian scalar product
$$Q_n(\xi ,\overline{\xi })=\xi ^i\overline{\xi }^i+\text{i}(\xi ^0\overline{\xi }^{n+1}\xi ^{n+1}\overline{\xi }^0).$$
Let $`\mathrm{\Sigma }^nS^{2n+1}`$ be the set of equivalence classes up to scale of null vectors.
Let $`\mu =(\mu ^0,\mu ^i,\mu ^{n+i}\mu ^{2n+1})`$, $`1in`$, be the coordinates of $`V^{2n+2}`$. *Whitney submanifold* $`\mathrm{\Gamma }_n:\mathrm{\Sigma }^n\mathrm{\Sigma }^{2n}`$ is an immersion induced by the quadratic map $`\widehat{\mathrm{\Gamma }}_n:V^{n+2}V^{2n+2}`$ defined as
$`\mu ^0`$ $`=2\xi ^0\xi ^{n+1},`$
$`\mu ^{2n+1}`$ $`=(\xi ^{n+1})^2(\xi ^0)^2,`$
$`\mu ^i`$ $`=\xi ^i(\text{i}\xi ^0+\xi ^{n+1}),`$
$`\mu ^{n+i}`$ $`=\xi ^i(\text{i}\xi ^0+\xi ^{n+1}).`$
$`\widehat{\mathrm{\Gamma }}_n^{}Q_{2n}=2(\xi ^0\overline{\xi }^0+\xi ^{n+1}\overline{\xi }^{n+1})Q_n`$, and the induced map $`\mathrm{\Gamma }_n`$ on $`\mathrm{\Sigma }^n`$ is well defined. It is easy to check $`\mathrm{\Gamma }_n`$ is CR-equivalent to the boundary map $`W_n:S^{2n+1}S^{4n+1}`$ of the following *Whitney map* $`W_n:𝔹^{n+1}𝔹^{2n+1}`$, where $`𝔹^{}^{}`$ is the unit ball and $`(z^0,z^i)`$, $`\mathrm{\hspace{0.17em}1}in`$, is a coordinate of $`^{n+1}`$.
$$W_n(z^0,z^i)=((z^0)^2,z^0z^i,z^i).$$
(38)
This equivalence is via the isomorphism $`\mathrm{\Sigma }^nS^{2n+1}`$ given in coordinates
$`z^0`$ $`={\displaystyle \frac{\text{i}\xi ^0+\xi ^{n+1}}{\text{i}\xi ^0+\xi ^{n+1}}}`$
$`z^i`$ $`={\displaystyle \frac{\sqrt{2}\xi ^i}{\text{i}\xi ^0+\xi ^{n+1}}}`$
Set $`\mathrm{\Sigma }_0^n=\{[\xi ]\mathrm{\Sigma }^n|\xi ^i=0,i\}`$ and $`\mathrm{\Sigma }_s^n=\{[\xi ]\mathrm{\Sigma }^n|\text{i}\xi ^0+\xi ^{n+1}=0\}`$. Then $`\mathrm{\Gamma }_n`$ is an immersion which is 1 to 1 on $`\mathrm{\Sigma }^n\mathrm{\Sigma }_0^n`$, 2 to 1 on $`\mathrm{\Sigma }_0^n`$, and the second fundamental form vanishes along $`\mathrm{\Sigma }_s^n`$.
###### Theorem 5.1
Let $`M^n\mathrm{\Sigma }^{2n}`$ be a $`C^3`$ nonlinear CR-flat submanifold of CR dimension and codimension $`n2`$. Then $`M`$ is congruent to a part of the Whitney submanifold up to an automorphism of $`\mathrm{\Sigma }^{2n}`$.
CR-flat submanifold $`M^1\mathrm{\Sigma }^2`$ has been classified by Faran \[Fa\]. In contrast to $`n2`$ cases, there are four inequivalent CR flat submanifolds when $`n=1`$.
As a corollary, we have a simple characterization of the proper holomorphic maps from a unit ball $`𝔹^{n+1}`$ to $`𝔹^{2n+1}`$ \[HJ\].
###### Corollary 5.1
Let $`F:𝔹^{n+1}𝔹^{2n+1}`$, $`n2`$, be a nonlinear proper holomorphic map which is $`C^{n+1}`$ up to the boundary. Then $`F`$ is equivalent to the Whitney map (38) up to automorphisms of the unit balls.
Theorem 5.1 is based on the following algebraic lemma due to Iwatani on the normal form of the second fundamental form of a Bochner-Kähler submanifold \[Iw\]\[Br3\]. Let $`V=^n`$, $`W=^n`$ with the standard Hermitian scalar product. Let $`\{z^i\}`$ be a unitary (1,0)-basis for $`V^{}`$, and let $`\{w_a\}`$ be a unitary basis for $`W`$.
###### Lemma 5.1
\[Iw\] Suppose $`H=h_{ij}^az^iz^jw_aS^{2,0}(V^{})W`$ satisfies
$$\gamma (H,H)=h_{ij}^a\overline{h}_{kl}^az^iz^j\overline{z}^k\overline{z}^lS_1^{1,1}S^{2,2},$$
or simply $`\gamma (H,H)`$ is Bochner-flat \[Br3\]. Then up to a unitary transformation on $`V`$,
$$H=h_{in}^az^iz^nw_a.$$
Set $`\nu _i=h_{in}^aw_aW`$. A computation shows $`\gamma (H,H)`$ is Bochner-flat whenever
$`\nu _i,\nu _j`$ $`=0\text{for}ij,`$
$`\nu _n,\nu _n`$ $`=4<\nu _q,\nu _q>\text{for}q<n.`$
Thus up to a unitary transformation on $`W`$, we may assume
$`\nu _q`$ $`=rw_q\text{for}q<n,`$
$`\nu _n`$ $`=2rw_n,`$
for some $`r0`$.
Let $`M^n\mathrm{\Sigma }^{2n}`$ be a CR-flat submanifold. The second fundamental form of $`M`$ is Bochner-flat \[EHZ\]. In the notation of Section 3, Lemma 5.1 then implies we may write
$`\pi _q^{n+i}`$ $`r\delta _{iq}\omega ^nmod\theta ,\text{for}q<n,`$
$`\pi _n^{n+i}`$ $`r(1+\delta _{in})\omega ^imod\theta ,`$
where $`\theta `$ is the dual to the contact hyperplane fields. Assume $`M`$ is not linear, $`H0`$, and we scale $`r=1`$ using the group action by Re$`\pi _0^0`$. We thus obtain the following local normal form of the second fundamental form of a nonlinear CR-flat submanifold $`M^n\mathrm{\Sigma }^{2n}`$.
$`\pi _q^{n+i}`$ $`=\delta _{iq}\omega ^n+h_q^i\theta \text{for}q<n,`$ (39)
$`\pi _n^{n+i}`$ $`=(1+\delta _{in})\omega ^i+h_n^i\theta `$
for some coefficients $`h_j^i`$.
Theorem 5.1 is now obtained by successive differentiation of this normalized structure equation. We assume $`n3`$ for simplicity for the rest of this section, as $`n=2`$ case can be treated in a similar way. We shall agree on the index range $`\mathrm{\hspace{0.17em}1}p,q,s,tn1`$, and denote $`p^{}=n+p`$, $`n^{}=n+n`$. Recall our convention $`d\theta \text{i}\pi ^k\overline{\pi }^k\text{i}\varpi mod\theta `$, and we denote $`\pi ^i=\omega ^i`$ for the sake of notation.
*Step 1.* Differentiating $`\pi _s^{n+n}=h_s^n\theta `$ $`mod\theta `$, we get
$$\text{i}h_s^n\varpi (\pi _s^{}^n^{}2\pi _s^n)\omega ^n+\pi _{2n+1}^n^{}(\text{i}\overline{\omega }^s)mod\theta .$$
Since $`n12`$, this implies $`h_s^n=0`$, and by Cartan’s lemma
$$\left(\begin{array}{c}\pi _s^{}^n^{}2\pi _s^n\\ \pi _{2n+1}^n^{}\end{array}\right)\left(\begin{array}{cc}c_s& u\\ u& 0\end{array}\right)\left(\begin{array}{c}\omega ^n\\ \text{i}\overline{\omega }^s\end{array}\right)mod\theta $$
for coefficients $`c_s,u`$.
*Step 2.* Differentiating $`\pi _s^t^{}=h_s^t\theta `$ $`mod\theta `$ for $`ts`$, we get
$$\text{i}h_s^t\varpi (\pi _s^{}^t^{}\pi _s^t)\omega ^n\pi _s^n\omega ^t+\pi _{2n+1}^t^{}(\text{i}\overline{\omega }^s)mod\theta .$$
Since $`n12`$, this implies $`h_s^t=0`$ for $`ts`$, and by Cartan’s lemma
$$\left(\begin{array}{c}\pi _s^{}^t^{}\pi _s^t\\ \pi _s^n\\ \pi _{2n+1}^t^{}\end{array}\right)\left(\begin{array}{ccc}0& b_s& \text{i}\overline{b}_t\\ b_s& 0& a\\ \text{i}\overline{b}_t& a& 0\end{array}\right)\left(\begin{array}{c}\omega ^n\\ \omega ^t\\ \text{i}\overline{\omega }^s\end{array}\right)mod\theta $$
for coefficients $`c_s,u`$. Since $`\pi _s^{}^t^{}\pi _s^t`$ is skew Hermitian, it cannot have any $`\omega ^n`$-term.
*Step 3.* Differentiating $`\pi _t^t^{}=\omega ^n+\text{i}h_t^t\theta `$ $`mod\theta `$ and collecting terms, we get
$$h_t^t\varpi (\pi _t^{}^t^{}\pi _t^t+\pi _0^0\pi _n^n)\omega ^n+(b_p\omega ^n\text{i}a\overline{\omega }^p)\omega ^p+(b_t\omega ^t+\overline{b}_t\overline{\omega }^t)\omega ^n$$
$`mod\theta `$. Since $`n12`$, this implies $`h_t^t=a`$, and
$`\mathrm{\Delta }_t`$ $`=\pi _t^{}^t^{}\pi _t^t+\pi _0^0\pi _n^n`$
$`=a_t\omega ^n\text{i}a\overline{\omega }^n+(b_t\omega ^t\overline{b}_t\overline{\omega }^t)+{\displaystyle \underset{p}{}}b_p\omega ^pA_t\theta `$
for coefficients $`a_t,A_t`$.
*Step 4.* From *Step 2*, we may use the group action by $`\pi _{2n+1}^n`$ to translate $`a=0`$, which we assume from no on. We also translate $`h_n^t=0`$ by $`\pi _{2n+1}^t`$. Differentiating $`\pi _n^t^{}=\omega ^t`$ $`mod\theta `$ with these relations, we get
$$0\overline{b}_t(\underset{p}{}\omega ^p\overline{\omega }^p)+(2\overline{c}_t3\overline{b}_t)\omega ^n\overline{\omega }^n+\omega ^n(a_t+2\text{i}\overline{u})\omega ^tmod\theta .$$
Thus $`b_t=c_t=0`$, $`a_t=2\text{i}\overline{u}`$.
*Step 5.* Differentiating $`\pi _n^n^{}=2\omega ^n+h_n^n\theta `$ $`mod\theta `$ and collecting terms, we get $`A_t=A`$ for a single variable, and
$$\text{i}h_n^n\varpi 2(\pi _n^{}^n^{}\pi _n^n+\pi _0^0\pi _n^n)\omega ^n+\text{i}u\omega ^p\overline{\omega }^p\text{i}u\omega ^n\overline{\omega }^nmod\theta .$$
This implies $`h_n^n=u`$, and
$`\mathrm{\Delta }_n`$ $`=\pi _n^{}^n^{}\pi _n^n+\pi _0^0\pi _n^n`$
$`=a_n\omega ^n\text{i}u\overline{\omega ^n}A_n\theta `$
for coefficients $`a_n,A_n`$. But $`\mathrm{\Delta }_t\mathrm{\Delta }_n`$ is purely imaginary, and comparing with *Step 3*, $`a_n=3\text{i}\overline{u}`$.
*Step 6.* Now by considering $`\theta `$-terms in *Step 1, 2, 3, 4, 5* and the fact $`\pi _i^a\omega ^i+\pi _{2n+1}^a\theta =0`$, we obtain the following simple structure equations. We omit the details of computations.
$`\left(\begin{array}{c}\pi _s^{}^n^{}2\pi _s^n\\ \pi _{2n+1}^n^{}\end{array}\right)`$ $`=\left(\begin{array}{cc}0& u\\ u& 0\end{array}\right)\left(\begin{array}{c}\omega ^n\\ \text{i}\overline{\omega }^s\end{array}\right),`$
$`\left(\begin{array}{c}\pi _s^{}^t^{}\pi _s^t\\ \pi _s^n\\ \pi _{2n+1}^t^{}\end{array}\right)`$ $`=\left(\begin{array}{c}0\\ 0\\ 0\end{array}\right),`$
$`\pi _{2n+1}^t`$ $`=(A\text{i}u\overline{u})\omega ^t+B_t\theta `$
$`\pi _{2n+1}^n`$ $`=A\omega ^n+B_n\theta `$
*Step 7.* Differentiating $`\pi _s^{}^t^{}\pi _s^t=0,\pi _s^n=0`$, we get first $`B_s=0,B=0`$, and
$$A\overline{A}=\text{i}(u\overline{u}1).$$
(40)
*Step 8.* Differentiating $`\pi _n^n^{}=2\omega ^n+u\theta `$ and $`\pi _{2n+1}^n^{}=u\omega ^n`$,
$$du=u(\pi _n^n+\pi _{2n+1}^{2n+1})+2(AA_n)\omega ^n+\text{i}u(\mathrm{\hspace{0.17em}3}\overline{u}\omega ^n+u\overline{\omega ^n})+u(2AA_n)\theta .$$
(41)
*Step 9.* Differentiating $`\pi _s^{}^n^{}=\text{i}u\overline{\omega }^s`$ using (41) and collecting terms in $`\theta \overline{\omega }^s`$,
$$A_n=2A\overline{A}.$$
*Step 10.* Differentiating $`\pi _{2n+1}^t=(\overline{A}\text{i})\omega ^t`$, $`\pi _{2n+1}^n=A\omega ^n`$, we get
$$dA=\pi _{2n+1}^0+(A+\text{i})(\pi _{2n+1}^{2n+1}\pi _0^0)(A+\text{i})^2\theta .$$
(42)
We normalize $`A=\text{i}\alpha `$ for a real number $`\alpha `$ using group action by $`\pi _{2n+1}^0`$. Since $`\pi _i^i+\overline{\pi }_i^i=0`$, $`\mathrm{\Delta }_t+\overline{\mathrm{\Delta }}_t=\pi _0^0+\overline{\pi }_0^0`$ ana (42) is now reduced to
$`d\alpha `$ $`=2\text{i}(\alpha +1)(\overline{u}\omega ^nu\overline{\omega }^n)`$ (43)
$`\pi _{2n+1}^0`$ $`=(\alpha +1)^2\theta .`$
When $`u0`$, we may also rotate $`u`$ to be a positive number, in which case it is determined by (40)
$$2\alpha +1=u\overline{u}.$$
(44)
At this stage, note that the only independent coefficients in the structure equations are $`\alpha ,u`$, and that the expression for their derivatives does not involve any new variables. The structure equations for the CR-flat submanifold $`M^n\mathrm{\Sigma }^{2n}`$ thus *close up* as follows.
$`\left(\begin{array}{cc}\pi _q^p^{}& \pi _n^p^{}\\ \pi _q^n^{}& \pi _n^n^{}\end{array}\right)`$ $`=\left(\begin{array}{cc}\delta _{pq}\omega ^n& \omega ^p\\ 0& 2\omega ^n+u\theta \end{array}\right)`$ (45)
$`\left(\begin{array}{c}\pi _s^{}^n^{}2\pi _s^n\\ \pi _{2n+1}^n^{}\end{array}\right)`$ $`=\left(\begin{array}{cc}0& u\\ u& 0\end{array}\right)\left(\begin{array}{c}\omega ^n\\ \text{i}\overline{\omega }^s\end{array}\right)`$
$`\left(\begin{array}{c}\pi _s^{}^t^{}\pi _s^t\\ \pi _s^n\\ \pi _{2n+1}^t^{}\end{array}\right)`$ $`=\left(\begin{array}{c}0\\ 0\\ 0\end{array}\right)`$
$`\pi _{2n+1}^t`$ $`=(A\text{i}u\overline{u})\omega ^t`$
$`\pi _{2n+1}^n`$ $`=A\omega ^n`$
$`\mathrm{\Delta }_t`$ $`=2\text{i}\overline{u}\omega ^nA\theta `$
$`\mathrm{\Delta }_n`$ $`=3\text{i}\overline{u}\omega ^n\text{i}u\overline{\omega ^n}A_n\theta `$
$`du`$ $`=u(\pi _n^n+\pi _{2n+1}^{2n+1})+\text{i}u(\mathrm{\hspace{0.17em}3}\overline{u}\omega ^n+u\overline{\omega ^n})`$
$`+2(AA_n)\omega ^n+u(2AA_n)\theta `$
$`d\alpha `$ $`=2\text{i}(\alpha +1)(\overline{u}\omega ^nu\overline{\omega }^n)`$
$`\pi _{2n+1}^0`$ $`=(\alpha +1)^2\theta `$
where $`A=\text{i}\alpha ,A_n=3\text{i}\alpha `$, and $`\alpha ,u`$ satisfy the relation (44). Moreover, a long but direct computation shows that these structure equations are compatible, that is $`d^2=0`$ is a formal identity of the structure equations.
*Remark.* The structure equation closes up at order 3 \[Ha\]. Since both $`\mathrm{\Sigma }^n`$ and $`\mathrm{\Sigma }^{2n}`$ are real analytic, this implies a $`C^3`$ nonlinear CR-flat submanifold $`f:\mathrm{\Sigma }^n\mathrm{\Sigma }^{2n}`$ is real analytic.
Let $`\mathrm{\Sigma }^{}=\mathrm{\Sigma }^n\mathrm{\Sigma }_s^n`$, and note that it is a connected set. Note also the structure equation (45) implies that the set of points where $`u=0`$ or equivalently $`\alpha =\frac{1}{2}`$ cannot have any interior on $`\mathrm{\Sigma }^{}`$. We claim the invariant $`\alpha `$ takes any value $`>\frac{1}{2}`$ on $`\mathrm{\Sigma }^{}`$.
Suppose $`\alpha _+=sup_\mathrm{\Sigma }^{}\alpha >\frac{1}{2}`$ is finite. Applying the existence part of Cartan’s generalization of Lie’s third fundamental theorem on closed structure equations, \[Br3\], there exists for any $`p_0\mathrm{\Sigma }^{}`$ a neighborhood $`U\mathrm{\Sigma }^{}`$ and a CR immersion $`g:U\mathrm{\Sigma }^{2n}`$ with invariant $`\alpha |_{p_0}=\alpha _+`$, hence necessarily $`u_{p_0}=\sqrt{2\alpha _0+1}`$. From (40), $`d\alpha |_{p_0}0`$ and let $`p_{}U`$ be a point with $`\frac{1}{2}<\alpha |_p_{}<\alpha _+`$. Then by uniqueness part of Cartan’s theorem, there exists a neighborhood $`U^{}U`$ of $`p_{}`$ on which $`g`$ agrees with the Whitney map $`\mathrm{\Gamma }_n`$ up to automorphisms of $`\mathrm{\Sigma }^n`$ and $`\mathrm{\Sigma }^{2n}`$. Since $`g`$ and $`\mathrm{\Gamma }_n`$ satisfy the closed set of structure equations, they are real analytic. Thus $`g`$ is a part of $`\mathrm{\Gamma }_n`$. But $`du|_{p_0}0`$, and there exists a point $`p_+U\mathrm{\Sigma }^{}`$ such that $`\alpha |_{p_+}>\alpha _+`$, a contradiction. By similar argument, $`inf_\mathrm{\Sigma }^{}\alpha =\frac{1}{2}`$, and the claim follows for $`\mathrm{\Sigma }^{}`$ is connected.
*Proof of Theorem 5.1.* Since the set of points $`\alpha =\frac{1}{2}`$ cannot have any interior, let $`pM`$ be a point with $`\alpha |_p>\frac{1}{2}`$. Then from the results above, there exists a point $`q\mathrm{\Sigma }^n`$ such that $`\alpha |_q=\alpha |_p`$. By similar argument as above after identifying $`pM`$ with $`q\mathrm{\Sigma }^n`$, there exists an automorphisms $`\tau _{2n}`$ of $`\mathrm{\Sigma }^{2n}`$ such that $`f=\tau _{2n}\mathrm{\Gamma }_n`$ on a neighborhood $`U`$ of $`p`$. The theorem follows for both $`f`$ and $`\mathrm{\Gamma }_n`$ are real analytic. $`\mathrm{}`$
*Proof of Corollary 5.1.* By the regularity theorem \[Mi\], $`F`$ is real analytic up to $`F`$. Since the CR structure on $`S^{2n+1}=\mathrm{\Sigma }^n`$ is definite, the set of points where $`F`$ has holomorphic rank $`n`$ is a dense open subset. Since $`F`$ is not linear, there exists a point $`p\mathrm{\Sigma }^n`$ where the second fundamental form does not vanish either. By Theorem 5.1, $`F`$ agrees with the Whitney map $`\mathrm{\Gamma }_n`$ in a neighborhood of $`p`$ up to automorphisms of the unit balls. The real analyticity then implies $`F=\mathrm{\Gamma }_n`$ on $`\mathrm{\Sigma }^n`$, and hence $`F=W_n`$ on $`𝔹^{n+1}`$. $`\mathrm{}`$
We may apply Cartan’s generalization of Lie’s third fundamental theorem and show that the Whitney submanifold provides an example of a deformable CR-submanifold, and in particular that it is not 1-rigid. Take a point $`p\mathrm{\Sigma }^n`$ and an analytic one parameter family of real numbers $`\alpha _t>\frac{1}{2}`$, and set $`u_t=\sqrt{2\alpha _t+1}`$. Then by the existence part of Cartan’s theorem, there exists a neighborhood $`U`$ of $`p`$ and a one parameter family of CR immersions $`f_t:U\mathrm{\Sigma }^{2n}`$ with the induced structure equations (45) such that the invariants $`\alpha ,u`$ have the prescribed values $`\alpha _t`$, $`u_t`$ at $`p`$. Of course this deformation is *tangential* and does not actually deform the submanifold. It is due to an intrinsic CR symmetry of $`\mathrm{\Sigma }^n`$ that cannot be extended to a symmetry of the ambient $`\mathrm{\Sigma }^{2n}`$ along the Whitney submanifold.
## 6 Proof of Example 3.1
Example 3.1 *Let $`M_{n,p}P^N=P(^n^{n+p})`$ be the Plücker embedding of the Grassmannian $`Gr(n,^{n+p})`$, $`pn`$. Let $`\widehat{M}_{n,p}S^{2N+1}`$ be the inverse image of $`M_{n,p}`$ under Hopf map, which is an $`S^1`$-bundle over $`M_{n,p}`$. When $`p=2,\mathrm{\hspace{0.17em}3}`$, $`\widehat{M}_{n,p}`$ is a Bochner rigid CR submanifold of $`S^{2N+1}`$.*
It is easy to check that the canonical $`S^1`$-bundle $`\widehat{M}S^{2N+1}`$ over any complex submanifold $`MP^N`$ is a CR submanifold. Moreover, the fundamental forms of $`\widehat{M}`$ as a CR submanifold are simply the pull back, in an appropriate sense, of the usual fundamental forms of $`M`$ as a projective subvariety \[HY\]. It thus suffices to show that the fundamental forms of the Plücker embeddings of the complex Grassmannian manifolds are Bochner rigid when $`p=2,\mathrm{\hspace{0.17em}3}`$. Note the height $`\tau `$ of $`\widehat{M}_{n,p}`$ is $`p`$.
Let $`V`$ denote the tangent space of $`M_{n,p}`$, and we follow the notations of Section 3.
*Case $`p=2`$*. There exists a unitary basis $`\{x^i,y^i\}`$, $`i=1,\mathrm{}n`$, of $`V^{}`$ such that the second fundamental form is given by $`\frac{n(n1)}{2}`$ quadratic forms
$$F^{ij}=x^iy^jx^jy^i,i<j.$$
Let $`P^{ij}S^{2,0}(V^{})`$, $`i<j`$, be such that
$$F^{ij}\overline{P}^{ij}+P^{ij}\overline{F}^{ij}=(x^i\overline{x}^i+y^i\overline{y}^i)Q$$
for some $`QS^{1,1}(V^{})`$. Since $`F^{ij}`$ has only $`x,y`$ terms, $`Q`$ cannot have $`x\overline{x}`$ terms nor $`y\overline{y}`$ terms, and thus $`Q`$ only has $`x\overline{y}`$, $`y\overline{x}`$ terms. Consider $`x^1\overline{x}^1Q`$ term. Since it contains three $`x`$’s, the only possible contribution comes from $`x^1y^j\overline{P}_{xx}^{1j}+P_{xx}^{1j}\overline{x}^1\overline{y}^j`$, where $`P_{xx}^{1j}`$ denotes the $`xx`$ component of $`P^{1j}`$. But
$$\frac{}{y^1}(x^1y^j\overline{P}_{xx}^{1j}+P_{xx}^{1j}\overline{x}^1\overline{y}^j)=0.$$
Hence $`\frac{}{y^1}Q=0`$. By permutation symmetry in $`x,y`$, and in $`i`$, and since $`Q=\overline{Q}`$, $`Q=0`$.
*Case $`p=3`$*. There exists a unitary basis $`\{x^i,y^i,z^i\}`$, $`i=1,\mathrm{}n`$, of $`V^{}`$ such that the second fundamental form is given by the quadratic forms
$`F_1^{ij}`$ $`=y^iz^jy^jz^i,i<j`$
$`F_2^{ij}`$ $`=z^ix^jz^jx^i,i<j`$
$`F_3^{ij}`$ $`=x^iy^jx^jy^i,i<j.`$
Let $`P_1^{ij},P_2^{ij},P_3^{ij}S^{2,0}(V^{})`$, $`i<j`$, be such that
$$F_a^{ij}\overline{P}_a^{ij}+P_a^{ij}\overline{F}_a^{ij}=(x^i\overline{x}^i+y^i\overline{y}^i+z^i\overline{z}^i)Q$$
for some $`QS^{1,1}(V^{})`$. Consider this equation$`modx,y`$, and $`z`$ in turn. Since $`Q`$ is quadratic, the result for the case $`p=2`$ implies $`Q=0`$.
The third fundamental form is given by the cubic forms
$$F^{ijk}=x^iy^jz^k+x^jy^kz^i+x^ky^iz^jx^iy^kz^jx^jy^iz^kx^ky^jz^i,i<j<k.$$
Let $`P^{ijk}S^{3,0}(V^{})`$, $`i<j<k`$, be such that
$$F^{ijk}\overline{P}^{ijk}+P^{ijk}\overline{F}^{ijk}=(x^i\overline{x}^i+y^i\overline{y}^i+z^i\overline{z}^i)Q$$
for some $`QS^{2,2}(V^{})`$. Considering this equation$`modx,y`$, and $`z`$ in turn, every monomials of $`Q`$ is either a $`xxyz`$, $`yxyz`$, or $`zxyz`$ term(ignoring the conjugation). Consider $`xxyz`$ component $`Q_{xxyz}`$ and the $`x^1\overline{x}^1Q_{xxyz}`$ term in the above equation. Since there are 4 $`x`$-terms, the only possible contribution comes from $`x^1y^jz^k\overline{P}_{xxx}^{1jk}+P_{xxx}^{1jk}\overline{x}^1\overline{y}^j\overline{z}^k`$, where $`P_{xxx}^{1jk}`$ denotes the $`xxx`$ component of $`P^{1jk}`$. But
$$\frac{}{y^1}(x^1y^jz^k\overline{P}_{xxx}^{1jk}+P_{xxx}^{1jk}\overline{x}^1\overline{y}^j\overline{z}^k)=0.$$
Hence $`\frac{}{y^1}Q_{xxyz}=0`$. By permutation symmetry in $`x,y,z`$, and in $`i`$, and since $`Q=\overline{Q}`$, $`Q=0`$.
References
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Sung Ho Wang
Department of Mathematics
Kias
Seoul, Corea 130-722
shw@kias.re.kr
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# Relaxation of the distribution function tails for systems described by Fokker-Planck equations
## 1 Introduction
The study of Fokker-Planck equations is an important problem in statistical mechanics and kinetic theory . In the simplest models, the diffusion coefficient is constant. However, Fokker-Planck equations with a diffusion coefficient depending on the velocity of the particles have also been introduced in physics. These equations usually describe the relaxation of a “test particle” evolving in a bath of “field particles” at statistical equilibrium when the particles interact via weak long-range forces. In that case, the diffusion coefficient is a function of the velocity of the test particle. For example, in his Brownian theory of stellar dynamics, Chandrasekhar describes the evolution of the velocity distribution of a star in a cluster by a Fokker-Planck equation involving a diffusion and a friction. The coefficients of diffusion and friction are related to each other by an Einstein relation and the diffusion coefficient decreases as $`v^3`$ for large velocities. These results are similar to those obtained in plasma physics for the Coulombian interaction . By using an analogy with stellar dynamics, Chavanis describes the relaxation of a test vortex in a thermal bath of field vortices by a Fokker-Planck equation (in position space) involving a diffusion and a drift along the vorticity gradient. The coefficients of drift and diffusion are related to each other by a form of Einstein relation involving a negative temperature and the diffusion coefficient is inversely proportional to the local shear created by the vortex cloud. For a Gaussian distribution of field vortices, the diffusion coefficient of the test vortex decreases with the distance as $`r^2e^{\lambda r^2}`$ . Similarly, for the Hamiltonian Mean Field (HMF) model, Bouchet & Dauxois and Chavanis et al. find that the velocity distribution of a test particle satisfies a Fokker-Planck equation with a diffusion coefficient decreasing as $`e^{\beta v^2/2}`$ for large velocities. More generally, using the theory developed by Landau, Lenard and Balescu in plasma physics and implementing a thermal bath approximation, one can obtain a general Fokker-Planck equation involving an anisotropic diffusion coefficient depending on the velocity of the test particle. This Fokker-Planck equation is valid for systems with weak long-range potentials of interaction. The preceding kinetic equations can be recovered as particular cases of this general Fokker-Planck equation .
For Fokker-Planck equations with a variable diffusion coefficient, the relaxation towards the Boltzmann distribution is slowed down, especially if the diffusion coefficient decreases rapidly with the velocity. One consequence is that velocity correlation functions can decrease algebraically rapidly with time (instead of exponentially) as investigated by Bouchet & Dauxois in relation with the HMF model <sup>1</sup><sup>1</sup>1In that case, the Fokker-Planck equation with Gaussian diffusion coefficient and linear friction can be transformed into a Fokker-Planck equation with constant diffusion coefficient and logarithmic potential which is known to exhibit power-law correlations (see, in particular, Appendix B of and ).. They therefore explain the observed algebraic tails of the velocity correlations functions in terms of classical kinetic theory without invocating a notion of “generalized thermodynamics”. Here, we consider the relaxation of the system towards equilibrium from another point of view. We focus on the distribution function $`f(v,t)`$ and study the structure and the evolution of the front formed in the high velocity tail. Our study is based on the approach of Potapenko et al. who studied this problem in the case where the diffusion coefficient decreases algebraically with the velocity. In the present paper, we generalize this approach to an arbitrary form of diffusion coefficient and study how the front position $`v_f(t)`$ and the front structure evolve with time. In the case of a Gaussian or exponential decay of the diffusion coefficient with the velocity, we find that the progression of the front is extremely slow (logarithmic). This can leads to a sort of “kinetic blocking” or, at least, a “slowing down” of the relaxation. A similar confining effect was noted in the case where the diffusion coefficient depends on the density , but this situation (related to the theory of violent relaxation) is more difficult to analyze since the kinetic equation is then nonlinear.
The paper is organized as follows. In Sec. 2, we discuss different systems with long-range interactions (stellar systems, Coulombian plasmas, two-dimensional vortices, HMF model,…) which are described in the thermal bath approximation by Fokker-Planck equations with a variable diffusion coefficient. We also consider systems like optical lattices described by Fokker-Planck equations with constant diffusion coefficient but variable friction coefficient. In Sec. 3, we generalize the theory of Potapenko et al. for an arbitrary form of diffusion coefficient and friction force. We provide general equations characterizing the structure and the evolution of the front for large times. In Sec. 4, we address the validity of our approach. In Sec. 5, we consider particular applications of our general formalism for physically motivated Fokker-Planck equations. Analytical results are compared with direct numerical simulations of the Fokker-Planck equation. Particular emphasis is given to the case where the diffusion coefficient decreases with the velocity like an exponential or a Gaussian distribution. This is the situation relevant for one-dimensional systems like the HMF model and for two-dimensional point vortices. Finally, in Sec. 6 we investigate a class of Fokker-Planck equations for which our approach is exact for all times. The Appendices provide technical details and extensions of our main results.
## 2 Examples of Fokker-Planck equations with a variable diffusion coefficient and friction force
We consider a Hamiltonian system of $`N`$ particles interacting via a weak long-range binary potential $`u(|𝐫𝐫^{}|)`$. These particles can be stars in stellar clusters, electrons or ions in a plasma, point vortices in two-dimensional hydrodynamics, particles located on a ring in the HMF model … We assume that the cluster is homogeneous and in a steady state characterized by a distribution function $`f_0(𝐯)`$. In general, $`f_0`$ will be the statistical equilibrium state (thermal bath) but in certain cases it can be a slowly evolving distribution function. We introduce a “test particle” and denote by $`P(𝐯,t)`$ its velocity distribution function. Due to the interaction with the “field particles”, the velocity distribution will change and the test particle will acquire the distribution of the bath $`f_0(𝐯)`$ for $`t+\mathrm{}`$ (see for more details). The general Fokker-Planck equation describing the relaxation (“thermalization”) of the test particle in the bath of field particles can be written as :
$`{\displaystyle \frac{P}{t}}=\pi (2\pi )^dm{\displaystyle \frac{}{v^\mu }}{\displaystyle 𝑑𝐯_1𝑑𝐤k^\mu k^\nu \frac{\widehat{u}(𝐤)^2}{|ϵ(𝐤,𝐤𝐯)|^2}\delta [𝐤(𝐯𝐯_1)]\left(\frac{}{v^\nu }\frac{}{v_1^\nu }\right)f_0(𝐯_1)P(𝐯,t)},`$
where $`\widehat{u}(𝐤)`$ is the Fourier transform of $`u(𝐫)`$ and
$`ϵ(𝐤,\omega )=1+(2\pi )^d\widehat{u}(𝐤){\displaystyle \frac{𝐤\frac{f_0}{𝐯}}{\omega 𝐤𝐯}𝑑𝐯},`$ (2)
is the dielectric function. This equation can be obtained from the Lenard-Balescu equation by fixing the distribution function of the field particles $`f(𝐯_1,t)`$ to its static value $`f_0(𝐯_1)`$ . This (thermal) bath approximation transforms an integro-differential equation (Lenard-Balescu) into a differential equation (Fokker-Planck). The Lenard-Balescu equation was introduced in plasma physics for the Coulombian potential but it can apply to other systems with weak long-range interactions. Equations (2)-(2) can also be obtained from the general expression of the Fokker-Planck equation by explicitly calculating the coefficients of friction and diffusion (first and second moments of the velocity increments) using the Klimontovich approach . The fact that Eq. (2) is linear does not imply that the distribution $`P(𝐯,t)`$ is close to equilibrium. The test particle approach is different from considering a small perturbation of the Lenard-Balescu equation around equilibrium. In the first case, we describe the evolution of a single test particle (or an ensemble of test particles that do not interact among themselves) in a thermal bath while in the second case one would describe the evolution of all the particles (the system “as a whole”) close to equilibrium.
If we consider that the field particles are at statistical equilibrium with the Boltzmann distribution $`f_0(𝐯)e^{\beta mv^2/2}`$, and if we neglect collective effects taking $`|ϵ(𝐤,𝐤𝐯)|=1`$ (Landau approximation) we can rewrite the general Fokker-Planck equation (2) in the form :
$$\frac{P}{t}=\frac{1}{t_R}\frac{}{x^\mu }\left[G^{\mu \nu }(𝐱)\left(\frac{P}{x^\nu }+2Px^\nu \right)\right],$$
(3)
where we have set $`𝐱=(\beta m/2)^{1/2}𝐯`$. The diffusion coefficient is proportional to the tensor
$$G^{\mu \nu }(𝐱)=𝑑\widehat{𝐤}\widehat{k}^\mu \widehat{k}^\nu e^{(\widehat{𝐤}𝐱)^2}$$
(4)
with $`\widehat{𝐤}=𝐤/k`$, and the quantity
$$t_R^1=\left(\frac{\pi }{8}\right)^{1/2}d^{3/2}(2\pi )^d\frac{\rho m}{v_m^3}_0^+\mathrm{}k^d\widehat{u}(k)^2𝑑k$$
(5)
provides an estimate of the inverse relaxation time of the test particle toward the distribution of the bath. Here $`v_m=(d/\beta m)^{1/2}`$ denotes the r.m.s. velocity and $`\rho `$ the spatial density. A more general expression of the diffusion coefficient can be obtained by taking into account collective effects . Note, however, that for $`v\mathrm{}`$ (which is a limit that we shall be particularly interested with in the sequel), the expression (4) is asymptotically exact since $`|ϵ(𝐤,𝐤𝐯)|1`$ for $`|𝐯|+\mathrm{}`$. Note also that for weak long-range potentials of interaction for which our approach is valid, and in the Landau approximation, the precise form of the potential $`\widehat{u}(k)`$ only determines the timescale of the relaxation, through Eq. (5), not the form of the kinetic operator. Therefore, in the Landau approximation, the expression of the diffusion tensor as a function of the velocity only depends on the dimension of space $`d`$.
In dimension $`d=3`$, the diffusion tensor can be written
$`G^{\mu \nu }=(G_{}{\displaystyle \frac{1}{2}}G_{}){\displaystyle \frac{x^\mu x^\nu }{x^2}}+{\displaystyle \frac{1}{2}}G_{}\delta ^{\mu \nu },`$ (6)
where $`G_{}`$ and $`G_{}`$ are the diffusion coefficients in the directions parallel and perpendicular to the velocity of the test particle. They are explicitly given by
$`G_{}={\displaystyle \frac{2\pi ^{3/2}}{x}}G(x),G_{}={\displaystyle \frac{2\pi ^{3/2}}{x}}[\mathrm{erf}(x)G(x)],`$ (7)
with
$`G(x)={\displaystyle \frac{2}{\sqrt{\pi }}}{\displaystyle \frac{1}{x^2}}{\displaystyle _0^x}t^2e^{t^2}𝑑t={\displaystyle \frac{1}{2x^2}}\left[\mathrm{erf}(x){\displaystyle \frac{2x}{\sqrt{\pi }}}e^{x^2}\right],`$ (8)
where
$`\mathrm{erf}(x)={\displaystyle \frac{2}{\sqrt{\pi }}}{\displaystyle _0^x}e^{t^2}𝑑t,`$ (9)
is the error function. If we consider spherically symmetric distributions, noting that $`P/x^\mu =(1/x)(P/x)x^\mu `$ and $`G^{\mu \nu }x^\nu =G_{}x^\mu `$ we obtain
$$\frac{P}{t}=\frac{1}{t_R}\frac{1}{x^2}\frac{}{x}\left[x^2G_{}(x)\left(\frac{P}{x}+2Px\right)\right].$$
(10)
For the gravitational potential, this Fokker-Planck equation has been studied by Chandrasekhar in his Brownian theory of stellar dynamics . It has also been considered in plasma physics as an approximation of the Landau equation valid for sufficiently large times . We note in particular that the diffusion coefficient $`G_{}(x)`$ decreases algebraically like $`x^3`$ for $`x+\mathrm{}`$.
Alternatively, if we consider one dimensional systems ($`d=1`$), the general Fokker-Planck equation (2) simplifies into :
$$\frac{P}{t}=\frac{}{v}\left[D(v)\left(\frac{P}{v}P\frac{d}{dv}\mathrm{ln}f_0\right)\right],$$
(11)
where $`D(v)`$ is given by
$$D(v)=4\pi ^2mf_0(v)_0^+\mathrm{}𝑑k\frac{k\widehat{u}(k)^2}{|ϵ(k,kv)|^2}.$$
(12)
We note that the distribution function $`P(v,t)`$ of the test particle relaxes towards the distribution of the bath $`f_0(v)`$ on a timescale of the order $`Nt_D`$, where $`t_D`$ is the dynamical time . If we neglect collective effects, or consider the limit of large velocities, we find that the diffusion coefficient is given by
$$D(v)=4\pi ^2mf_0(v)_0^+\mathrm{}𝑑kk\widehat{u}(k)^2.$$
(13)
It is proportional to the distribution function of the bath $`f_0(v)`$. In particular, if the field particles are at statistical equilibrium with a Gaussian distribution, the diffusion coefficient decreases like $`e^{\beta mv^2/2}`$. This type of Fokker-Planck equations apply for example to the HMF model which can be viewed as the one Fourier component of a one-dimensional plasma (or self-gravitating system) . More generally, these Fokker-Planck equations (11) are valid for a wide class of one dimensional systems with long range interactions . We note that for one-dimensional systems the Lenard-Balescu collision term cancels out so that the distribution function of the field particles $`f(v,t)`$ does not evolve, i.e. $`f/t=0`$, on a timescale of order $`Nt_D`$. Since, on the other hand, the relaxation time of a test particle toward the distribution of the bath is of order $`Nt_D`$, this implies that we can assume that the distribution of the field particles is stationary $`f(v,t)=f_0(v)`$ when we study the relaxation of a test particle. This is true for any distribution function $`f_0(v)`$ that is a stable stationary solution of the Vlasov equation . This is not true in higher dimensions $`d=2`$ and $`d=3`$, except for the Maxwellian distribution $`f_e(𝐯)`$, since the distribution of the field particles $`f(𝐯,t)`$ changes on a time $`Nt_D`$ as it relaxes towards the statistical equilibrium state $`f_e(𝐯)`$. Furthermore, even if we assume that the distribution of the bath $`f_0(𝐯)`$ is approximately stationary, the distribution of the test particle $`P(𝐯,t)`$ will not relax towards $`f_0(𝐯)`$ for large times except if $`f_0(𝐯)`$ is Maxwellian .
A Fokker-Planck equation with a space dependent diffusion coefficient has been introduced by Chavanis in to describe the relaxation of a test vortex in a “sea” of field vortices with vorticity profile $`\omega _0(r)`$. For an axisymmetric distribution, the Fokker-Planck equation for $`P(r,t)`$ can be written :
$`{\displaystyle \frac{P}{t}}={\displaystyle \frac{1}{r}}{\displaystyle \frac{}{r}}\left[rD(r)\left({\displaystyle \frac{P}{r}}P{\displaystyle \frac{d}{dr}}\mathrm{ln}\omega _0\right)\right],`$ (14)
with a diffusion coefficient
$`D(r)={\displaystyle \frac{\gamma }{8}}{\displaystyle \frac{1}{|\mathrm{\Sigma }(r)|}}\mathrm{ln}N\omega _0(r),`$ (15)
where $`\mathrm{\Sigma }(r)=r\mathrm{\Omega }_0^{}(r)`$ is the local shear created by the field vortices ($`\mathrm{\Omega }_0(r)`$ represents the angular velocity related to the vorticity by $`\omega _0(r)=(1/r)(r^2\mathrm{\Omega }_0)^{}`$). For a vorticity profile $`\omega _0(r)=Ae^{\lambda r^2}`$ of the field vortices, it is easy to see that the diffusion coefficient of the test vortex decreases like $`D(r)r^2e^{\lambda r^2}`$ for $`r+\mathrm{}`$ .
Finally, Fokker-Planck equations with diffusion and friction coefficients depending on the velocity can occur in many areas of physics. For example, the motion of atoms in a one-dimensional optical lattice formed by two counter propagating laser beams with linear perpendicular polarization can be described, after spatial averaging, by a Fokker-Planck of the form :
$`{\displaystyle \frac{W}{t}}={\displaystyle \frac{}{p}}\left[D(p){\displaystyle \frac{W}{p}}WK(p)\right],`$ (16)
with
$`K(p)={\displaystyle \frac{\alpha p}{1+(p/p_c)^2}},D(p)=D_0+{\displaystyle \frac{D_1}{1+(p/p_c)^2}}.`$ (17)
For $`p+\mathrm{}`$, $`D(p)D_0`$ and $`K(p)\alpha p_c^2/p`$. This corresponds to a logarithmic potential $`U(p)(\alpha p_c^2/D_0)\mathrm{ln}p`$ defined by $`K(p)/D(p)=U^{}(p)`$. Equation (16) belongs to the general class of Fokker-Planck equations that we shall study in the sequel. The case of a logarithmic potential is treated specifically in Sec. 5.3.
## 3 General solution of the problem
The various examples discussed previously prompt us to study Fokker-Planck equations with a diffusion coefficient and friction force depending on the velocity. In particular, we can wonder how the distribution function $`f(v,t)`$ approaches the equilibrium distribution. This problem has been investigated by Potapenko et al. in the case of 3D plasmas where the diffusion coefficient is a power law. These authors found that the asymptotic behavior of the velocity distribution tail has a propagating wave appearance. The high velocity tail develops a front at which the distribution function drops to zero. This front $`v_f(t)`$ progresses with time and goes to $`v_f(t)+\mathrm{}`$ for $`t+\mathrm{}`$. The profile of the front also deforms itself as time goes on. However, if we use an appropriate system of coordinates, it can be expressed in terms of the error function. We shall here formulate the problem in a general setting, for an arbitrary form of diffusion coefficient and friction force, and we shall investigate how the evolution of the front and the profile of the high velocity tail distribution depend on the form of the diffusion coefficient.
Let us consider the Fokker-Planck equation
$`{\displaystyle \frac{f}{t}}={\displaystyle \frac{1}{v^{d1}}}{\displaystyle \frac{}{v}}\left[v^{d1}D(v)\left({\displaystyle \frac{f}{v}}+fU^{}(v)\right)\right],`$ (18)
where $`D(v)0`$ and $`U(v)`$ are arbitrary functions. If the following zero flux condition
$$v^{d1}D(v)\left(\frac{f}{v}+fU^{}(v)\right)0$$
(19)
when $`v\mathrm{}`$ is fulfilled, then the stationary solutions of this Fokker-Planck equation take the form
$`f_e(v)=Ae^{U(v)},`$ (20)
where $`A`$ is a constant (normalization). The Fokker-Planck equation (18) decreases the free energy $`F=ES`$ where $`E=fU(v)𝑑𝐯`$ is the energy and $`S=f\mathrm{ln}fd𝐯`$ is the Boltzmann entropy (the temperature has been included in the potential $`U`$). Indeed, one has
$`\dot{F}={\displaystyle \frac{D}{f}\left(\frac{f}{𝐯}+f\frac{U}{𝐯}\right)^2𝑑𝐯}0.`$ (21)
Therefore, if $`F`$ is bounded from below, the distribution will converge towards the equilibrium state (20) for $`t+\mathrm{}`$. We want to analyze the propagation of the front in the high velocity tail of the distribution function. Thus, we set
$`f(v,t)=f_e(v)u(v,t).`$ (22)
For sufficiently large times, the core of the distribution function will have reached its asymptotic value (20) so that $`u1`$ in that region. On the other hand, for sufficiently large velocities, the distribution has not relaxed yet and $`u=0`$. Therefore, we expect the formation of a front at a typical velocity value $`v_f(t)`$ where the function $`u(v,t)`$ passes from $`u=1`$ to $`u=0`$. On this phenomenological basis, $`u(v,t)`$ is the relevant function to consider in our “travelling front” analysis. Its exact evolution is governed by
$`{\displaystyle \frac{u}{t}}={\displaystyle \frac{1}{v^{d1}}}{\displaystyle \frac{}{v}}\left(v^{d1}D(v){\displaystyle \frac{u}{v}}\right)D(v)U^{}(v){\displaystyle \frac{u}{v}},`$ (23)
which is obtained from Eqs. (18)-(22). If we perform the change of variables $`dx/dv=1/\sqrt{D(v)}`$, we obtain the equivalent equation
$`{\displaystyle \frac{u}{t}}={\displaystyle \frac{^2u}{x^2}}+\left[{\displaystyle \frac{1}{2}}{\displaystyle \frac{D^{}(v)}{\sqrt{D(v)}}}+{\displaystyle \frac{d1}{v}}\sqrt{D(v)}U^{}(v)\sqrt{D(v)}\right]{\displaystyle \frac{u}{x}},`$ (24)
where $`v=v(x)`$ must be viewed as an implicit function of $`x`$. If we introduce the velocity field
$`V(v)=\sqrt{D(v)}\left[U^{}(v){\displaystyle \frac{d1}{v}}{\displaystyle \frac{1}{2}}(\mathrm{ln}D)^{}(v)\right],`$ (25)
or, equivalently,
$`V(v)=\sqrt{D(v)}{\displaystyle \frac{d}{dv}}\left\{\mathrm{ln}\left[v^{d1}e^{U(v)}D^{1/2}(v)\right]\right\},`$ (26)
Eq. (24) can be re-written as
$`{\displaystyle \frac{u}{t}}+V(v){\displaystyle \frac{u}{x}}={\displaystyle \frac{^2u}{x^2}}.`$ (27)
The structure of this equation is clear. The right hand side corresponds to a diffusion and the left hand side corresponds to an advection by a velocity field $`V(v)`$. This velocity field (in velocity space) governs the evolution of the front. To get a physical insight into the problem, let us first neglect the diffusion term. The resulting equation
$`{\displaystyle \frac{u}{t}}+V(v){\displaystyle \frac{u}{x}}=0`$ (28)
can be solved with the method of characteristics. Writing
$`{\displaystyle \frac{dx}{dt}}=V(v)={\displaystyle \frac{dx}{dv}}{\displaystyle \frac{dv}{dt}}={\displaystyle \frac{1}{\sqrt{D(v)}}}{\displaystyle \frac{dv}{dt}},`$ (29)
we get
$`{\displaystyle \frac{dv}{dt}}=\sqrt{D(v)}V(v)=D(v)\left[U^{}(v){\displaystyle \frac{d1}{v}}{\displaystyle \frac{1}{2}}(\mathrm{ln}D)^{}(v)\right].`$ (30)
This equation determines the evolution of the front $`v_f(t)`$. In the extreme approximation where the diffusion is neglected, the profile of the front is given by a step function $`u(x,t)=\eta (xx_f(t))`$ where $`\eta `$ is the Heaviside function ($`\eta (x)=1`$ for $`x<0`$ and $`\eta (x)=0`$ for $`x>0`$). As we shall see, the diffusion term will smooth out this profile.
To take into account the effect of diffusion, we return to Eq. (27) and perform the change of variables
$`z=xx_f(t),u(x,t)=\varphi (z,t),`$ (31)
where the function $`x_f(t)`$ is defined by Eq. (29). Substituting this in Eq. (27), we obtain
$`{\displaystyle \frac{\varphi }{t}}={\displaystyle \frac{^2\varphi }{z^2}}[V(v)V(v_f)]{\displaystyle \frac{\varphi }{z}}.`$ (32)
So far, no approximation has been made so that Eq. (32) is exact and bears the same information as the initial Fokker-Planck equation (18). It is just written in a more convenient form which will allow us to examine the situation in the region of the front. Indeed, far from the front the profile is stationary (for sufficiently long time) and Eq. (32) is automatically satisfied as $`\varphi =1`$ for $`vv_f(t)`$ and $`\varphi =0`$ for $`vv_f(t)`$. If we consider values of the velocity that are close to $`v_f(t)`$, we can expand the term in brackets in Taylor series (the validity of this approximation will be studied in Sec. 4). Keeping only the first term in this expansion, we get
$`{\displaystyle \frac{\varphi }{t}}={\displaystyle \frac{^2\varphi }{z^2}}g(t)z{\displaystyle \frac{\varphi }{z}},`$ (33)
where
$`g(t)=V^{}(v_f(t))\sqrt{D(v_f(t))}.`$ (34)
For future convenience, we set $`\tau =2t`$ and define $`h(\tau )=g(\tau /2)`$. Therefore, the forgoing equation becomes
$`{\displaystyle \frac{\varphi }{\tau }}={\displaystyle \frac{1}{2}}\left({\displaystyle \frac{^2\varphi }{z^2}}h(\tau )z{\displaystyle \frac{\varphi }{z}}\right).`$ (35)
The general solution of this equation, for an arbitrary initial condition, is given in Appendix A. Here, we look for particular solutions of the form $`\varphi (z,\tau )=\mathrm{\Phi }(z/\chi (\tau ))`$. In Appendix A, we derive the condition under which such solutions describe the asymptotic long time behavior of the system (independently of the initial condition) and we check that this condition is fulfilled for all the explicit examples that we shall investigate in the following. To describe more general situations, one must use the results of Appendix A. Inserting the ansatz $`\varphi (z,\tau )=\mathrm{\Phi }(z/\chi (\tau ))`$ in Eq. (35) leads to
$`\mathrm{\Phi }^{\prime \prime }+(2\chi \dot{\chi }h\chi ^2)x\mathrm{\Phi }^{}=0.`$ (36)
The variables of velocity and time separate provided that the term in parenthesis is a constant that we can arbitrarily set equal to
$`2\chi \dot{\chi }h\chi ^2=2.`$ (37)
Then, $`\mathrm{\Phi }`$ is the function
$`\mathrm{\Phi }(x)={\displaystyle \frac{1}{\sqrt{\pi }}}{\displaystyle _x^+\mathrm{}}e^{y^2}𝑑y,`$ (38)
connected to the error function by $`\mathrm{\Phi }(x)=\frac{1}{2}(1\mathrm{erf}(x))`$. It satisfies $`\mathrm{\Phi }(\mathrm{})=1`$, $`\mathrm{\Phi }(0)=1/2`$ and $`\mathrm{\Phi }(+\mathrm{})=0`$ so it reproduces the expected properties of the front (it can be seen as a smooth step function). On the other hand, solving for $`\chi ^2`$ in Eq. (37) and taking $`\chi (1)=0`$, we finally obtain
$`\chi ^2(\tau )=2{\displaystyle _1^\tau }e^{[H(\tau )H(\tau ^{})]}𝑑\tau ^{},`$ (39)
where $`H`$ is a primitive of $`h`$ with $`H(1)=0`$. The function $`\varphi (z,\tau )=\mathrm{\Phi }(z/\chi (\tau ))`$ with (38) and (39) is the solution of Eq. (35) for all times, corresponding to the initial condition $`\varphi (z,1)=\eta (z)`$ where $`\eta `$ is a step function. It also governs the long time behavior of the system for a large class of initial conditions (Appendix A).
In particular, for $`h(\tau )=\gamma /\tau `$, we have
$`{\displaystyle \frac{\varphi }{\tau }}={\displaystyle \frac{1}{2}}\left({\displaystyle \frac{^2\varphi }{z^2}}{\displaystyle \frac{\gamma }{\tau }}z{\displaystyle \frac{\varphi }{z}}\right),`$ (40)
and we recover the solution given in , namely
$`\varphi (z,\tau )=\mathrm{\Phi }\left[{\displaystyle \frac{z}{\sqrt{2}}}\left({\displaystyle \frac{1\gamma }{\tau \tau ^\gamma }}\right)^{1/2}\right].`$ (41)
For $`\gamma =1`$, a value that will frequently occur in the following examples, we have
$`\varphi (z,\tau )=\mathrm{\Phi }\left({\displaystyle \frac{z}{\tau ^{1/2}}}\right),\mathrm{for}\tau 1.`$ (42)
In terms of the function $`g`$, these analytical solutions correspond to $`g(t)=\lambda /t`$. We come back to the original variables by setting $`\tau =2t`$ and $`\gamma =2\lambda `$.
## 4 Condition of validity
We shall now investigate in greater detail the ability of our approach to describe the front structure. First, we note that, within our approximations,
$`u(v_f(t),t)=\mathrm{\Phi }(0)={\displaystyle \frac{1}{2}},`$ (43)
so that $`v_f(t)`$ gives the position of the half-height profile. On the other hand, noting that
$`V(v)V(v_f)=V^{}(v_f)(vv_f)+{\displaystyle \frac{1}{2}}V^{\prime \prime }(v_f)(vv_f)^2+\mathrm{}`$ (44)
for $`vv_f`$, our approximation (33) will be valid in the range of velocities where we can neglect the second term in the Taylor expansion with respect of the first. This corresponds to
$`|vv_f|\left|{\displaystyle \frac{2V^{}(v_f)}{V^{\prime \prime }(v_f)}}\right|.`$ (45)
In this range of velocities, our approach is accurate. Now, it will provide a good description of the whole front if this range $`|vv_f|`$ is larger than the typical front half-width $`\mathrm{\Delta }_f(t)/2`$. This condition can be written
$`\mathrm{\Delta }_f(t)\left|{\displaystyle \frac{4V^{}(v_f)}{V^{\prime \prime }(v_f)}}\right|.`$ (46)
Now, the front width can be estimated by
$`\mathrm{\Delta }_f(t)=\left|{\displaystyle \frac{u}{v}}(v_f(t),t)\right|^1.`$ (47)
Within our approximation, this can be finally rewritten as
$`\mathrm{\Delta }_f(t)=\sqrt{\pi }\chi (2t)\sqrt{D(v_f(t))}.`$ (48)
Therefore, our approach will provide a good description of the front structure if
$`ϵ(t)\left|{\displaystyle \frac{\sqrt{\pi }}{4}}{\displaystyle \frac{\chi (2t)\sqrt{D(v_f)}}{V^{}(v_f)/V^{\prime \prime }(v_f)}}\right|1.`$ (49)
Since we shall be interested by the large time limit, it is particularly important to know the asymptotic behavior of the function $`ϵ(t)`$ for $`t+\mathrm{}`$. This has to be considered case by case (see Sec. 5).
The other assumption made in our study is that the long time behavior of the system is described by $`\varphi (z,\tau )=\mathrm{\Phi }(z/\chi (\tau ))`$ with (38) and (39). In Appendix A, we show that this is the case if $`\varphi (z,1)1`$ for $`z\mathrm{}`$, $`\varphi (z,1)0`$ for $`z+\mathrm{}`$ and if
$`H_2(\tau )={\displaystyle _1^\tau }e^{H(\tau ^{})}𝑑\tau ^{}+\mathrm{}`$ (50)
for $`\tau +\mathrm{}`$. In particular, for $`h(\tau )=\gamma /\tau `$, we have
$`H(\tau )=\gamma \mathrm{ln}\tau ,H_1(\tau )={\displaystyle \frac{1}{\tau ^{\gamma /2}}},`$ (51)
and
$`\chi ^2(\tau )={\displaystyle \frac{2}{1\gamma }}(\tau \tau ^\gamma ),H_2(\tau )={\displaystyle \frac{\tau ^{1\gamma }1}{1\gamma }},`$ (52)
so that the criterion (50) is met if $`\gamma 1`$ and not met if $`\gamma >1`$.
## 5 Particular examples
We shall now discuss explicitly some particular examples of physical interest and compare our theoretical predictions with direct numerical simulations of the Fokker-Planck equation.
### 5.1 Quadratic potential
We first consider the case of a quadratic potential $`U(v)=v^2/2`$ leading to the Kramers equation
$`{\displaystyle \frac{f}{t}}={\displaystyle \frac{1}{v^{d1}}}{\displaystyle \frac{}{v}}\left[v^{d1}D(v)\left({\displaystyle \frac{f}{v}}+fv\right)\right],`$ (53)
with a variable diffusion coefficient $`D(v)`$. If we assume a zero flux condition (19), then the stationary solutions are the Maxwellian distributions
$`f_e(v)=Ae^{\frac{v^2}{2}}.`$ (54)
We now consider various forms of diffusion coefficient and use the general theory developed in Sec. 3 to characterize the front profile and its evolution. We shall only give asymptotic expressions which are valid for large velocities and large time. This will be implicit in all the following calculations.
#### 5.1.1 Power-law decay
For the diffusion coefficient
$`D(v)v^\alpha ,\alpha >0,`$ (55)
we get
$`x{\displaystyle \frac{2}{\alpha +2}}v^{\frac{\alpha +2}{2}},v_f(t)(\alpha t)^{1/\alpha },`$ (56)
and
$$\{\begin{array}{c}g(t)\frac{2\alpha }{2\alpha }\frac{1}{t},\text{if}\alpha 2\hfill \\ g(t)\frac{d2}{2}\frac{1}{t^2},\text{if}\alpha =2.\hfill \end{array}$$
(57)
For $`\alpha =2`$, we have assumed that the subdominant corrections to the $`D(v)v^2`$ behavior are of order $`v^4`$ or smaller. Note that if $`\alpha =2`$ and $`d=2`$, then $`g=0`$ and the Fokker-Planck equation can be solved exactly (see Sec. 6).
For $`\alpha 2`$, we find that $`h=\gamma /\tau `$ with $`\gamma =2/\alpha 1`$. The front profile is given by
$`u(v,t)\mathrm{\Phi }\left[{\displaystyle \frac{v^{\frac{\alpha +2}{2}}(\alpha t)^{\frac{\alpha +2}{2\alpha }}}{(\alpha +2)t^{1/2}}}\left({\displaystyle \frac{1\gamma }{1(2t)^{\gamma 1}}}\right)^{1/2}\right].`$ (58)
The criterion (50) is fulfilled only for $`\alpha 1`$, so that the above function provides the correct asymptotic behavior of the solution, for any initial condition, only in that case. For $`\alpha <1`$, it can however provide the correct asymptotic behavior if the initial condition $`\varphi (z,1)`$ is a step function (see Appendix A). Equation (58) returns the results obtained in . Of course this formula is written for $`\alpha 1`$, but it also provides the expression of the solution for $`\alpha =1`$ by passing to the limit $`\alpha 1`$ yielding:
$`u(v,t)\mathrm{\Phi }\left[{\displaystyle \frac{v^{\frac{3}{2}}t^{\frac{3}{2}}}{3(t\mathrm{ln}(2t))^{1/2}}}\right].`$ (59)
On the other hand, for $`\alpha =2`$ we find that $`h=2(d2)/\tau ^2`$ yielding $`\chi ^2(\tau )2\tau `$ and $`H_2(\tau )\tau `$ for $`\tau +\mathrm{}`$. We find that the criterion (50) is satisfied and that the front profile is given by
$`u(v,t)\mathrm{\Phi }\left[{\displaystyle \frac{v^22t}{4t^{1/2}}}\right],`$ (60)
which turns out to be consistent with Eq. (58).
Concerning the validity of this approach in relation with the criterion (49), we have
$$\{\begin{array}{c}ϵ(t)A_\alpha t^{1/\alpha },\text{if}\alpha >1\text{and}\alpha 2,\hfill \\ ϵ(t)\frac{B_\alpha }{t}\text{if}\alpha <1,\hfill \\ ϵ(t)\frac{\sqrt{\pi }}{4}\frac{(\mathrm{ln}t)^{1/2}}{t}\text{if}\alpha =1,\hfill \\ ϵ(t)\frac{3\sqrt{\pi }}{4}t^{1/2}\text{if}\alpha =2(d2),\hfill \end{array}$$
(61)
where $`A_\alpha =(\sqrt{2\pi }/8)\alpha ^{11/\alpha }(\alpha 1)^{1/2}`$ and $`B_\alpha =(\sqrt{\pi }/8)\alpha ^{11/\alpha }(1\alpha )^{1/2}2^{1/\alpha 1/2}`$. Therefore, the condition $`ϵ1`$ is always fulfilled for sufficiently large times. For sake of illustration, we show the case $`\alpha =3`$ for $`d=3`$ (plasmas and stellar systems) in Figs. 1 and 2. These results are obtained by solving numerically the Fokker-Planck equation (53) starting from a step function: $`f(v,t=0)=3/4\pi `$ if $`v<1`$ and $`f(v,t=0)=0`$ if $`v>1`$ (water-bag). In the simulation we have adopted the expression (7)-a of the diffusion coefficient which reduces to Eq. (55) with $`\alpha =3`$ for $`v+\mathrm{}`$. Other examples are given in .
#### 5.1.2 Gaussian decay
We now consider a diffusion coefficient of the form
$`D(v)e^{\gamma v^2}.`$ (62)
Considering the evolution of the high velocity tail, we have
$`x{\displaystyle _0^v}e^{\frac{\gamma }{2}y^2}𝑑y.`$ (63)
For large $`v`$, we obtain the relation (see Appendix B)
$`x{\displaystyle \frac{1}{\gamma v}}e^{\frac{\gamma }{2}v^2}.`$ (64)
The position of the front $`v_f(t)`$ is determined by
$`{\displaystyle ^{\gamma v_f^2}}{\displaystyle \frac{e^y}{y}}𝑑y2(\gamma +1)t.`$ (65)
The integral can be expressed in terms of the exponential integral $`E_i(x)`$. For large times, we get
$`{\displaystyle \frac{e^{\gamma v_f^2}}{\gamma v_f^2}}2(\gamma +1)t.`$ (66)
To leading order, we have
$`v_f(t)\left({\displaystyle \frac{\mathrm{ln}t}{\gamma }}\right)^{1/2}.`$ (67)
We note that the evolution of the front is extremely slow (logarithmic). On the other hand, we find that
$`g(t){\displaystyle \frac{1}{2t}},`$ (68)
implying that the criterion (50) is fulfilled. In order to determine the front profile, we need first to evaluate $`x_f(t)`$. An equivalent for $`t+\mathrm{}`$ is obtained by combining Eqs. (64) and (66) leading to
$`x_f(t)\sqrt{{\displaystyle \frac{2(\gamma +1)t}{\gamma }}}.`$ (69)
Therefore, the front profile is given by
$`u(v,t)\mathrm{\Phi }\left[{\displaystyle \frac{\frac{1}{\gamma v}e^{\frac{\gamma }{2}v^2}\sqrt{\frac{2(\gamma +1)t}{\gamma }}}{\sqrt{2t}}}\right].`$ (70)
Concerning the validity of this approach in relation with the criterion (49), we find that $`ϵ(t)\frac{1}{4}(\frac{\pi \gamma }{\gamma +1})^{1/2}`$ for $`t+\mathrm{}`$ so that our approximations are marginally valid: $`ϵ(t)`$ does not go to zero but it does not diverge with time neither.
We have performed numerical simulations of the Fokker-Planck equation (53) with the diffusion coefficient (62) with $`\gamma =1/2`$. As discussed in Sec. 2, this equation describes the evolution of a “test particle” in a bath of “field particles” at statistical equilibrium (with Maxwellian distribution function $`f_0e^{v^2/2}`$) for one-dimensional systems ($`d=1`$) with long-range interactions such as the HMF model. In Fig. 3, we show the “short” time evolution of the distribution function $`f(v,t)`$ starting from a step function: $`f(v,0)=1/3`$ for $`v3`$ whose width is larger than the thermal speed $`v_{th}=1`$ (dashed line). Since $`D(v)`$ rapidly decreases with the velocity for $`vv_{th}=1`$, the high velocity component of the initial condition takes time to be depleted. Indeed, the core of the distribution rapidly reaches a Maxwellian distribution while the tail keeps the memory of the initial condition and remains flat for relatively long times. For the parameters of the simulation, this can even create a non-monotonic distribution for short times as shown in Fig. 3 (see Appendix C).
In Fig. 4 we show the evolution of the normalized distribution function $`u(v,t)`$ for short and large times. In that case, we start from an initial condition: $`f(v,0)=1`$ if $`v<1`$ and $`f(v,0)=0`$ if $`v>1`$. In Fig. 5, we focus on the evolution of the front for large times and we compare the result of the numerical simulation with the theoretical prediction. In Fig. 6, we compare the evolution of the front displacement $`v_f(t)`$ with the theoretical prediction. The numerical results are in fair agreement with the theory although its domain of validity was expected to be marginal in that case according to our estimates. This fair agreement may be explained by the relatively small value of $`ϵ(+\mathrm{})=\frac{1}{4}(\pi /3)^{1/2}0.255\mathrm{}`$ for $`\gamma =1/2`$ even if this parameter does not strictly tends to zero for $`t+\mathrm{}`$. We note that for $`t+\mathrm{}`$ and $`vv_f(t)`$, we have $`u\mathrm{\Phi }(\sqrt{(\gamma +1)/\gamma })=\mathrm{\Phi }(\sqrt{\pi }/4ϵ)`$ which has not converged to $`u=1`$ precisely because $`ϵ0`$. However, for $`ϵ(+\mathrm{})0.255\mathrm{}`$, we find that $`u0.993`$ which is very close to one.
In relation with the HMF model, we would like to recall that the present approach describes the relaxation of a given test particle immersed in a bath of field particles at statistical equilibrium. The relaxation is due to finite $`N`$ effects (correlations). This does not describe the evolution of the distribution function of the system as a whole. The clearest reason is that the Fokker-Planck equation (2) resulting from a thermal bath approximation does not conserve energy contrary to the microcanonical evolution of a Hamiltonian system. In the collisionless regime, the distribution function $`f(\theta ,v,t)`$ of the HMF model is governed by the Vlasov equation coupled to the mean-field potential $`\mathrm{\Phi }(\theta ,t)=\frac{k}{2\pi }\mathrm{cos}(\theta \theta ^{})f(\theta ^{},v^{},t)𝑑\theta ^{}𝑑v^{}`$ produced by the particles. The Vlasov equation can experience a process of violent relaxation and converge toward a Quasi Stationary State (QSS) on the coarse-grained scale . Some dynamical theories of violent relaxation based on a Maximum Entropy Production Principle (MEPP) propose to model the evolution of the coarse-grained distribution function by a generalized Fokker-Planck equation of the form (for a water-bag initial condition):
$`{\displaystyle \frac{\overline{f}}{t}}+v{\displaystyle \frac{\overline{f}}{\theta }}\mathrm{\Phi }{\displaystyle \frac{\overline{f}}{v}}={\displaystyle \frac{}{v}}\left\{D(\theta ,v,t)\left[{\displaystyle \frac{\overline{f}}{v}}+\beta (t)\overline{f}(\eta _0\overline{f})v\right]\right\},`$ (71)
where $`\beta (t)`$ evolves so as to conserve energy (see for more details). An important point is that the diffusion coefficient is not constant but is related to the correlations of the fine-grained fluctuations. It can depend on position, velocity and time and can be very small in certain regions of phase-space and for large times. The vanishing or smallness of the diffusion coefficient can slow down the dynamics and lead to a confinement of the distribution function in phase space which may be qualitatively similar to what is shown in Fig. 3. However, the dynamical equation (71) and its physical interpretation are different from Eq. (11), so that their similarity is, at most, an analogy.
#### 5.1.3 Exponential decay
For the diffusion coefficient
$`D(v)e^{\gamma v},`$ (72)
we get
$`x{\displaystyle \frac{2}{\gamma }}e^{\frac{\gamma }{2}v}.`$ (73)
The evolution of the front is given by
$`{\displaystyle \frac{e^{\gamma v_f}}{\gamma v_f}}t.`$ (74)
To leading order, we have
$`v_f(t){\displaystyle \frac{1}{\gamma }}\mathrm{ln}t,`$ (75)
which shows that the progression of the front is again very slow. We also have
$`g(t){\displaystyle \frac{1}{2t}},`$ (76)
implying that the criterion (50) is fulfilled. In order to determine the front profile, we need first to evaluate $`x_f(t)`$. An equivalent for $`t+\mathrm{}`$ is obtained by combining Eqs. (73) and (74) leading to
$`x_f(t){\displaystyle \frac{2}{\gamma }}(t\mathrm{ln}t)^{1/2}.`$ (77)
However, since the evolution with time is slow, we shall work with the more precise expression $`x_f(t)\frac{2}{\gamma }(t\mathrm{ln}(t\mathrm{ln}t))^{1/2}`$ obtained by keeping the term of next order. With this expression, we find that the front profile is given by
$`u(v,t)\mathrm{\Phi }\left[\sqrt{2}{\displaystyle \frac{e^{\frac{\gamma }{2}v}(t\mathrm{ln}(t\mathrm{ln}t))^{1/2}}{\gamma t^{1/2}}}\right].`$ (78)
Concerning the validity of this approach in relation with the criterion (49), we find that $`ϵ(t)\frac{\gamma }{8}\sqrt{2\pi }(\mathrm{ln}t)^{1/2}`$ for $`t+\mathrm{}`$ so that our approximations are valid for sufficiently large times. Note that the decay is slow with time (logarithmic) but the prefactor $`\frac{\gamma }{8}\sqrt{2\pi }=0.313\mathrm{}`$ is relatively small (for $`\gamma =1`$) which can explain the very good agreement with the numerics even for moderate timescales. Figures 7 and 8 are obtained by solving the Fokker-Planck equation (53) with the diffusion coefficient (72) with $`\gamma =1`$, starting from an initial condition: $`f(v,0)=1`$ if $`v<1`$ and $`f=0`$ if $`v>1`$.
#### 5.1.4 Stretched exponential
For the diffusion coefficient
$`D(v)e^{\gamma v^\delta },\delta >0,`$ (79)
we get
$`x{\displaystyle _0^v}e^{\frac{\gamma }{2}y^\delta }𝑑y.`$ (80)
We shall briefly discuss how the results depend on the value of $`\delta `$. A more detailed analysis can be carried out as in the previous sections where the cases $`\delta =1`$ and $`\delta =2`$ are explicitly considered.
For $`\delta <2`$, the front evolution is given by
$`{\displaystyle \frac{e^{\gamma v_f^\delta }}{\gamma v_f^\delta }}\delta t.`$ (81)
To leading order, we get
$`v_f(t)\left({\displaystyle \frac{\mathrm{ln}t}{\gamma }}\right)^{1/\delta }.`$ (82)
On the other hand,
$`g(t){\displaystyle \frac{1}{2t}},`$ (83)
implying that the criterion (50) is fulfilled. Concerning the validity of this approach in relation with the criterion (49), we find that
$$ϵ(t)\frac{\sqrt{2\pi \delta }}{8}\left(\frac{\mathrm{ln}t}{\gamma }\right)^{(2\delta )/2\delta },$$
(84)
so that the criterion is satisfied for sufficiently large times (but slowly).
For $`\delta >2`$, the front evolution is given by
$`{\displaystyle \frac{e^{\gamma v_f^\delta }}{(\gamma v_f^\delta )^{2(\delta 1)/\delta }}}{\displaystyle \frac{1}{2}}\gamma ^{2/\delta }\delta ^2t.`$ (85)
To leading order, we get
$`v_f(t)\left({\displaystyle \frac{\mathrm{ln}t}{\gamma }}\right)^{1/\delta }.`$ (86)
On the other hand,
$`g(t){\displaystyle \frac{1}{2t}},`$ (87)
implying that the criterion (50) is fulfilled. Concerning the validity of this approach in relation with the criterion (49), we find that $`ϵ\frac{\sqrt{\pi }}{4}`$ so that this approach is marginally valid.
### 5.2 Linear potential
We now consider the case of a linear potential $`U(v)=\gamma v`$ leading to a Fokker-Planck equation of the form
$`{\displaystyle \frac{f}{t}}={\displaystyle \frac{1}{v^{d1}}}{\displaystyle \frac{}{v}}\left[v^{d1}D(v)\left({\displaystyle \frac{f}{v}}+\gamma f\right)\right].`$ (88)
The stationary solution is
$`f_e(v)=Ae^{\gamma v}.`$ (89)
For a diffusion coefficient decreasing algebraically with the velocity
$`D(v)v^\alpha ,\alpha >0,`$ (90)
we get
$`x{\displaystyle \frac{2}{\alpha +2}}v^{\frac{\alpha +2}{2}},v_f(t)[(\alpha +1)\gamma t]^{1/(\alpha +1)},`$ (91)
and
$`g(t){\displaystyle \frac{\alpha }{2(\alpha +1)}}{\displaystyle \frac{1}{t}},`$ (92)
implying that the criterion (50) is fulfilled. We also find that $`ϵ(t)t^{1/(2(\alpha +1))}`$ so that the validity criterion of our approach is always fulfilled for sufficiently large times.
### 5.3 Logarithmic potential
In this subsection, we consider the case of a constant diffusion coefficient $`D(v)=1`$ and a logarithmic potential $`U(v)=(\alpha /2)\mathrm{ln}(1+v^2)`$ in $`d=1`$ leading to a Fokker-Planck equation of the form
$`{\displaystyle \frac{f}{t}}={\displaystyle \frac{}{v}}\left({\displaystyle \frac{f}{v}}+\alpha f{\displaystyle \frac{v}{1+v^2}}\right).`$ (93)
This type of Fokker-Planck equations arises in the study of optical lattices . The stationary solution of this equation is of the form
$`f_e(v)={\displaystyle \frac{A}{(1+v^2)^{\alpha /2}}}.`$ (94)
This is similar to a Tsallis distribution with $`q=(\alpha 2)/\alpha `$ . However, it arises here from a linear Fokker-Planck equation (associated with the Boltzmann entropy) with a logarithmic potential, instead of a nonlinear Fokker-Planck equation (associated with the Tsallis entropy) with a quadratic potential . The distribution (94) is normalizable provided that $`\alpha >\alpha _{crit}=1`$. The case $`\alpha =2`$ corresponds to the Lorentzian. The evolution of the front $`v_f(t)`$ satisfies
$`{\displaystyle \frac{dv_f}{dt}}={\displaystyle \frac{\alpha v_f}{1+v_f^2}}.`$ (95)
The exact solution is given by $`v_f^2e^{v_f^2}=e^{2\alpha t+C}`$ where $`C`$ is a constant of integration. It can be written $`v_f^2(t)=W(e^{2\alpha t+C})`$ where $`W(x)`$ is the Lambert function which is solution of the transcendental equation $`W(x)\mathrm{exp}[W(x)]=x`$. For $`t+\mathrm{}`$, we get $`v_f(t)(2\alpha t)^{1/2}`$ and $`g(t)1/(2t)`$. Therefore, the front profile is given by
$`u(v,t)=\mathrm{\Phi }\left({\displaystyle \frac{v\sqrt{2\alpha t}}{\sqrt{2t}}}\right),`$ (96)
for large times. Concerning the validity of our approach with respect to the criterion (49), we find that $`ϵ(\frac{\pi }{4\alpha })^{1/2}`$ for $`t+\mathrm{}`$ so that our approach is marginally valid. The theoretical prediction is not very good for $`\alpha =2`$ but the agreement improves for large values of $`\alpha \alpha _{crit}=1`$ for which $`ϵ`$ is reduced. In particular, the results of numerical simulations performed with $`\alpha =10`$ are shown in Figs. 9 and 10. Concerning the evolution of the front, we find a good agreement with the Lambert function except that the measured exponent in Fig. 9 is $`20.6`$ instead of $`2\alpha =20`$. Therefore, the front increases a little bit faster than the theoretical prediction. This is confirmed in Fig. 10 which shows the evolution of the front profile. We note, however, the relatively good agreement between direct numerical simulation and theory. The slight discrepency is due to the finite value of $`ϵ0.28`$.
### 5.4 Fokker-Planck equations associated with one-dimensional systems with long-range interactions
We now consider the case of Fokker-Planck equations describing one dimensional systems with long-range interactions. In that case, the diffusion coefficient $`Df_0(v)`$ for $`v+\mathrm{}`$ and the potential $`U(v)=\mathrm{ln}f_0(v)`$ are expressed in terms of the distribution of the bath $`f_0(v)`$. The corresponding Fokker-Planck equation can be rewritten as (see Sec. 2):
$`{\displaystyle \frac{f}{t}}={\displaystyle \frac{}{v}}\left(f_0{\displaystyle \frac{f}{v}}f{\displaystyle \frac{df_0}{dv}}\right).`$ (97)
The general relations obtained in Sec. 3 take the simpler forms
$`{\displaystyle \frac{dx}{dv}}{\displaystyle \frac{1}{\sqrt{f_0(v)}}},V(v){\displaystyle \frac{3}{2}}{\displaystyle \frac{f_0^{}(v)}{\sqrt{f_0(v)}}},{\displaystyle \frac{dv_f}{dt}}{\displaystyle \frac{3}{2}}f_0^{}(v_f).`$ (98)
The case where $`f_0(v)`$ is a Gaussian distribution (thermal bath) has already been studied in Sec. 5.1.2. We shall consider other examples where $`f_0(v)`$ is not necessarily the statistical equilibrium state. This description still makes sense if $`f_0`$ is a stable stationary solution of the Vlasov equation because the relaxation of the system as a whole is slower than the relaxation of a test particle in a fixed distribution.
#### 5.4.1 Exponential distribution
For the exponential distribution
$`f_0e^{\gamma v},`$ (99)
we obtain
$`x{\displaystyle \frac{2}{\gamma }}e^{\frac{\gamma }{2}v},v_f{\displaystyle \frac{1}{\gamma }}\mathrm{ln}\left({\displaystyle \frac{3}{2}}\gamma ^2t\right),`$ (100)
and
$`g(t){\displaystyle \frac{1}{2t}},`$ (101)
implying that the criterion (50) is fulfilled. The front profile is given by
$`u(v,t)\mathrm{\Phi }\left[\sqrt{2}{\displaystyle \frac{e^{\frac{\gamma }{2}v}(\frac{3}{2}\gamma ^2t)^{1/2}}{\gamma t^{1/2}}}\right].`$ (102)
Concerning the validity of this approach in relation with the criterion (49), we find that $`ϵ\frac{1}{4}(\frac{\pi }{3})^{1/2}0.256\mathrm{}`$ for $`t+\mathrm{}`$ so that our approximations are marginally valid and are accurate since $`ϵ`$ is relatively small.
#### 5.4.2 Power-law distribution
For the power-law distribution
$`f_0v^\alpha ,`$ (103)
we obtain
$`x{\displaystyle \frac{2}{\alpha +2}}v^{\frac{\alpha +2}{2}},v_f(t)\left[{\displaystyle \frac{3}{2}}\alpha (\alpha +2)t\right]^{\frac{1}{\alpha +2}},`$ (104)
and
$`g(t){\displaystyle \frac{1}{2t}},`$ (105)
implying that the criterion (50) is fulfilled. The front profile is given by
$`u(v,t)\mathrm{\Phi }\left[\sqrt{2}{\displaystyle \frac{v^{\frac{\alpha +2}{2}}\sqrt{\frac{3}{2}\alpha (\alpha +2)t}}{(\alpha +2)t^{1/2}}}\right].`$ (106)
Concerning the validity of this approach in relation with the criterion (49), we find that $`ϵ\frac{1}{8}(\alpha +4)\sqrt{\frac{4\pi }{3\alpha (\alpha +2)}}`$ for $`t+\mathrm{}`$ so that our approximations are marginally valid.
## 6 A class of exactly solvable Fokker-Planck equations
In this section we give a class of Fokker-Planck equations (18) for which our approach turns out to be exact. More precisely, we derive a relationship between $`D(v)`$ and $`U(v)`$ under which the corresponding Fokker-Planck equation (18) can be exactly solved. First, we recall that equations (18) and (32) are equivalent and no approximation has been made in the derivation of Eq. (32) from Eq. (18). However to pass from Eq. (32) to Eq. (33) for a general velocity field $`V(v)`$, we have made an approximation and have assumed that the concerned velocities are close enough to the position of the front $`v_f(t)`$. This is in fact the only approximation in the approach developed above. We now look for a situation where such an approximation is exact. This is the case if $`V(v)=Ax+B`$ where $`A`$ and $`B`$ are arbitrary constants. Using
$$\frac{d}{dx}(V(v))=V^{}(v)\sqrt{D(v)}=A,$$
(107)
and replacing $`V`$ by its expression (25), we equivalently get
$$\sqrt{D(v)}\frac{d}{dv}\left(\sqrt{D(v)}\left[U^{}(v)\frac{d1}{v}\frac{1}{2}(\mathrm{ln}D)^{}(v)\right]\right)=A.$$
(108)
Defining $`R(v)=\sqrt{D(v)}`$, we obtain
$$R(v)\frac{d}{dv}\left[\left(U^{}(v)\frac{d1}{v}\right)RR^{}(v)\right]=A.$$
(109)
If $`U(v)`$ and $`R(v)=\sqrt{D(v)}`$ satisfy Eq. (109), as discussed in Appendix D, we can use the results of Appendix A with $`g(t)=A`$ to obtain exact solutions of Eq. (18) for all times.
Let us make these solutions more explicit. First of all, for $`h(\tau )=A`$, we have
$`H(\tau )=A(\tau 1),H_1(\tau )=e^{\frac{A}{2}(\tau 1)},`$ (110)
$`\chi ^2(\tau )={\displaystyle \frac{2}{A}}(e^{A(\tau 1)}1),H_2(\tau )={\displaystyle \frac{1}{A}}(1e^{A(\tau 1)}).`$ (111)
We note that the criterion (50) is met only for $`A0`$. On the other hand, integrating
$`{\displaystyle \frac{dx}{dt}}=V(v)=Ax+B,`$ (112)
with $`x_f(1/2)=0`$, we get
$`x_f(t)={\displaystyle \frac{B}{A}}(e^{\frac{A}{2}(2t1)}1).`$ (113)
Then, substituting in Eq. (152), we obtain the general solution
$`u(v,t)={\displaystyle \frac{1}{\sqrt{\frac{2\pi }{A}(1e^{2At})}}}{\displaystyle _{\mathrm{}}^+\mathrm{}}e^{\frac{\{e^{At}[x\frac{B}{A}(e^{At}1)]y\}^2}{\frac{2}{A}(1e^{2At})}}u(y,0)𝑑y,`$ (114)
where we have taken the origin of times at $`t=0`$. If we start from a step function $`u(x,0)=\eta (x)`$, the above expression reduces to
$`u(v,t)=\mathrm{\Phi }\left[{\displaystyle \frac{x\frac{B}{A}(e^{At}1)}{\sqrt{\frac{2}{A}}(e^{2At}1)^{1/2}}}\right].`$ (115)
On the other hand, for $`A=0`$, we get
$`h(\tau )=0,H(\tau )=0,H_1(\tau )=1,H_2(\tau )=\tau 1,`$ (116)
$`\chi ^2(\tau )=2(\tau 1),x_f(t)={\displaystyle \frac{B}{2}}(2t1).`$ (117)
The general solution is
$`u(v,t)={\displaystyle \frac{1}{\sqrt{4\pi t}}}{\displaystyle _{\mathrm{}}^+\mathrm{}}e^{\frac{(xBty)^2}{4t}}u(y,0)𝑑y,`$ (118)
and if we start from a step function $`u(x,0)=\eta (x)`$, we get
$`u(v,t)=\mathrm{\Phi }\left({\displaystyle \frac{xBt}{2t^{1/2}}}\right).`$ (119)
We recall that the above profiles are the exact solutions of the Fokker-Planck equation (18) when the functions $`U(v)`$ and $`D(v)`$ satisfy the differential equation (109). Let us examine some particular cases. If $`A=0`$, Eq. (109) can be integrated at once leading to the relation
$$D(v)=B^2\left(_v^+\mathrm{}w^{d1}e^{U(w)}𝑑w\right)^2\frac{e^{2U(v)}}{v^{2(d1)}}.$$
(120)
For the wide class of potentials $`U(v)`$ satisfying the following condition:
$$(d1)\frac{U^{}(v)}{v}U^{\prime \prime }(v)U^{}(v)^2$$
(121)
for large $`v`$, the behavior of the diffusion coefficient at infinity is given by
$$D(v)\frac{C}{U^{}(v)^2}.$$
(122)
In particular, when $`U(v)=v^2/2`$ the diffusion coefficient $`D(v)`$ behaves like $`C/v^2`$ for $`v+\mathrm{}`$ and when $`U(v)=\gamma v`$ it tends to a constant. Note that condition (121) is not only satisfied by any non constant polynomial, but also by potentials $`U(v)`$ that dominate $`\mathrm{ln}(v)`$ for large velocities, as for instance $`U(v)=(\mathrm{ln}(v))^\alpha `$, with $`\alpha >1`$. Note finally that for $`v0`$, the diffusion coefficient behaves like $`1/v^{2(d1)}`$ which is divergent except for $`d=1`$.
In the case of a quadratic potential $`U(v)=v^2/2`$, the integral in Eq. (120) can be expressed in terms of the error function. There is a simplification in $`d=2`$ leading to $`D(v)=1/v^2`$. This is precisely the case that we have found in Sec. 5.1.1. In that case, we have $`x=v^2/2`$ so that Eq. (118) can be explicitly written in terms of $`v`$. Inversely, for $`D(v)=1/v^2`$ and $`g=A`$, we find from Eq. (109) that $`U(v)=(A/8)v^4+(B/2)v^2+(d2)\mathrm{ln}v`$ so for this type of potential the solution of (114) applies with $`x=v^2/2`$. In the case of a linear potential $`U(v)=\gamma v`$, the integral in Eq. (120) can be expressed in terms of the incomplete Gamma function. There is a simplification in $`d=1`$ leading to $`D(v)=1`$. In that case, we have $`x=v`$ so that Eq. (118) can be explicitly written in terms of $`v`$. Inversely, for $`D(v)=1`$ and $`g=A`$, we find from Eq. (109) that $`U(v)=(A/2)v^2+Bv+(d1)\mathrm{ln}v`$ so for this type of potential the solution of (114) applies with $`x=v`$.
For $`d=1`$, $`D(v)=1`$ and $`U(v)=(A/2)v^2+Bv`$, we have $`V(v)=Ax+B`$ with $`x=v`$ and we are in the situation mentioned above. If we take $`f(v,0)=\delta (vv_0)`$ and recall that $`u(v,t)=f(v,t)/f_e(v)`$, Eq. (114) gives after simplification
$`f(v,t)={\displaystyle \frac{1}{\sqrt{\frac{2\pi }{A}(1e^{2At})}}}e^{\frac{A}{2(1e^{2At})}(vv_0e^{At}+\frac{B}{A}(1e^{At}))^2}.`$ (123)
In particular, for $`B=0`$ and $`A=1`$, we recover the well-known solution :
$`f(v,t)={\displaystyle \frac{1}{\sqrt{2\pi (1e^{2t})}}}e^{\frac{(vv_0e^t)^2}{2(1e^{2t})}}.`$ (124)
In that case, $`v(t)=v_0e^t`$. Alternatively, for $`A=0`$ and $`B=1`$, we get
$`f(v,t)={\displaystyle \frac{1}{\sqrt{4\pi t}}}e^{\frac{(vv_0+t)^2}{4t}}.`$ (125)
In that case, $`v(t)=v_0t`$. Note that there is no normalizable stationary state for the Fokker-Planck equation (18) when $`U(v)=Bv`$ so that $`f(v,t)`$ tends to zero for large times and spreads so as to conserve the total mass. By contrast, when $`A0`$, the distribution function relaxes towards $`f_e(v)`$ for $`t+\mathrm{}`$. From Eq. (123), we can obtain the time evolution of the distribution function for any initial condition $`f_0(v)=f(v,0)`$ by multiplying Eq. (123) by $`f_0(v_0)`$ and integrating over $`v_0`$.
Note finally that for $`d>1`$, the velocity $`v=|𝐯|`$ is restricted to positive values so that the preceding results must be slightly revised. The idea is to extend $`f(v,t)`$ by parity to negative values of $`v`$. We replace $`D(v)`$, $`U(v)`$ and $`v^{d1}`$ in Eq. (18) by $`D(|v|)`$, $`U(|v|)`$ and $`|v|^{d1}`$. With this extension, Eq. (18) is invariant by the transformation $`vv`$. Therefore, if $`f(v,t_0)`$ is even, $`f(v,t)`$ will remain even for all times and its values for $`v0`$ correspond to those of the distribution function that is solution of the original Eq. (18). Thus, we can apply the preceding results with almost no modification. We note however that since $`V(v)`$ and $`x(v)=_0^v𝑑w/\sqrt{D(w)}`$ are odd, the case $`V(v)=Ax+B`$ is only possible if $`B=0`$. By contrast, in $`d=1`$, we can consider cases where $`f(v,t)`$, $`D(v)`$ and $`U(v)`$ have no special parity.
## 7 Conclusion
In this paper, we have developed a general formalism to characterize the evolution of the distribution function tail for systems described by a Fokker-Planck equation with a diffusion coefficient and a friction force depending on the velocity. Our analytical results give good agreement with the numerics even in cases where the validity of our approach is marginal. When the diffusion coefficient decreases algebraically with the velocity, the progression of the front is also algebraic. When the diffusion coefficient decreases like an exponential or like a Gaussian, the progression of the front is logarithmic. The high velocity component of the distribution function keeps the memory of the initial condition for a long time and is slowly depleted. There are several applications of our formalism to, e.g., stellar dynamics, plasma physics, vortex dynamics, the HMF model, optical lattices etc. In future works, we shall study more specifically the dynamics of point vortices and investigate the evolution of the front profile and the time evolution of the correlation functions .
## Appendix A General solution of Eq. (35)
In this Appendix, we provide the general solution of the PDE:
$`2{\displaystyle \frac{\varphi }{\tau }}={\displaystyle \frac{^2\varphi }{z^2}}h(\tau )z{\displaystyle \frac{\varphi }{z}},`$ (126)
for an arbitrary initial condition $`\varphi _1(z)=\varphi (z,1)`$. Taking the Fourier transform of Eq. (126) with the conventions
$`\varphi (z)={\displaystyle \widehat{\varphi }(\xi )e^{i\xi z}𝑑\xi },\widehat{\varphi }(\xi )={\displaystyle \varphi (z)e^{i\xi z}\frac{dz}{2\pi }},`$ (127)
and using the relation
$`z{\displaystyle \frac{\varphi }{z}}={\displaystyle zi\xi \widehat{\varphi }(\xi )e^{i\xi z}𝑑\xi }`$
$`={\displaystyle \xi \widehat{\varphi }(\xi )\frac{}{\xi }(e^{i\xi z})𝑑\xi }={\displaystyle \frac{}{\xi }(\xi \widehat{\varphi }(\xi ))e^{i\xi z}𝑑\xi },`$ (128)
we get
$`2{\displaystyle \frac{\widehat{\varphi }}{\tau }}=(h(\tau )\xi ^2)\widehat{\varphi }+h(\tau )\xi {\displaystyle \frac{\widehat{\varphi }}{\xi }}.`$ (129)
We introduce the change of variables
$`f(y,\tau )=\widehat{\varphi }(H_1(\tau )y,\tau ),\xi =H_1(\tau )y`$ (130)
and choose the function $`H_1(\tau )`$ such that
$`{\displaystyle \frac{H_1^{}(\tau )}{H_1(\tau )}}={\displaystyle \frac{h(\tau )}{2}}.`$ (131)
Substituting Eq. (130) in Eq. (129), we find that $`f(y,\tau )`$ satisfies
$`{\displaystyle \frac{f}{\tau }}+{\displaystyle \frac{1}{2}}[H_1^2(\tau )y^2h(\tau )]f=0.`$ (132)
Let $`H(\tau )`$ be the primitive of $`h(\tau )`$ such that
$`H(\tau )={\displaystyle _1^\tau }h(\tau ^{})𝑑\tau ^{}.`$ (133)
Then, we choose $`H_1`$, solution of Eq. (131), such that
$`H_1(\tau )=e^{\frac{H(\tau )}{2}}.`$ (134)
By convention, $`H(1)=0`$ and $`H_1(1)=1`$. Equation (132) can be integrated leading to
$`f(y,\tau )=f(y,1)e^{\frac{H(\tau )}{2}}e^{\frac{1}{2}H_2(\tau )y^2},`$ (135)
where we have defined
$`H_2(\tau )={\displaystyle _1^\tau }H_1(\tau ^{})^2𝑑\tau ^{}={\displaystyle _1^\tau }e^{H(\tau ^{})}𝑑\tau ^{}.`$ (136)
Returning to original variables, we obtain
$`\widehat{\varphi }(\xi ,\tau )=\widehat{\varphi }_1\left({\displaystyle \frac{\xi }{H_1(\tau )}}\right)e^{\frac{H(\tau )}{2}}e^{\frac{1}{2}\frac{H_2(\tau )}{H_1^2(\tau )}\xi ^2}.`$ (137)
We now observe that
$`{\displaystyle \frac{H_2(\tau )}{H_1^2(\tau )}}={\displaystyle _1^\tau }e^{H(\tau )H(\tau ^{})}𝑑\tau ^{}{\displaystyle \frac{1}{2}}\chi ^2(\tau ),`$ (138)
where $`\chi (\tau )`$ is the function introduced in Eq. (39). Therefore the general solution of Eq. (126) in Fourier space can be written
$`\widehat{\varphi }(\xi ,\tau )={\displaystyle \frac{1}{H_1(\tau )}}\widehat{\varphi }_1\left({\displaystyle \frac{\xi }{H_1(\tau )}}\right)e^{\frac{1}{4}\chi ^2(\tau )\xi ^2}.`$ (139)
Defining
$`q(z)=\varphi _1(H_1(\tau )z)\widehat{q}(\xi )={\displaystyle \frac{1}{H_1(\tau )}}\widehat{\varphi }_1\left({\displaystyle \frac{\xi }{H_1(\tau )}}\right)`$ (140)
$`g(z)=G(z/\chi (\tau ))\widehat{g}(\xi )=\chi (\tau )\widehat{G}(\chi (\tau )\xi )`$ (141)
where
$`G(z)=e^{z^2}\widehat{G}(\xi )={\displaystyle \frac{1}{2\sqrt{\pi }}}e^{\xi ^2/4}`$ (142)
we can rewrite Eq. (139) in the form
$`\widehat{\varphi }(\xi ,\tau )={\displaystyle \frac{2\sqrt{\pi }}{\chi (\tau )}}\widehat{q}(\xi )\widehat{g}(\xi ).`$ (143)
Taking the inverse Fourier transform, we can express the solution of Eq. (126) as a convolution
$`\varphi (z,\tau )={\displaystyle \frac{2\sqrt{\pi }}{\chi (\tau )}}{\displaystyle q(zz^{})g(z^{})\frac{dz^{}}{2\pi }},`$ (144)
or, equivalently
$`\varphi (z,\tau )={\displaystyle \frac{1}{\sqrt{\pi }}}{\displaystyle _{\mathrm{}}^+\mathrm{}}e^{x^2}\varphi _1(H_1(\tau )(z\chi (\tau )x))𝑑x.`$ (145)
By direct substitution, we can check that Eq. (145) is indeed solution of Eq. (126).
If $`\varphi _1(z)=\eta (z)`$ is a step function with $`\eta (z)=1`$ for $`z<0`$ and $`\eta (z)=0`$ for $`z>0`$, we immediately find that
$`\varphi (z,\tau )={\displaystyle \frac{1}{\sqrt{\pi }}}{\displaystyle _{z/\chi (\tau )}^+\mathrm{}}e^{x^2}𝑑x=\mathrm{\Phi }\left({\displaystyle \frac{z}{\chi (\tau )}}\right),`$ (146)
and we recover the result of Sec. 3. Then, the general solution of Eq. (126) can be put in the form
$`\varphi (z,\tau )=\mathrm{\Phi }\left({\displaystyle \frac{z}{\chi (\tau )}}\right)+{\displaystyle \frac{1}{\sqrt{\pi }}}{\displaystyle _{\mathrm{}}^+\mathrm{}}e^{x^2}(\varphi _1\eta )(H_1(\tau )(z\chi (\tau )x))𝑑x.`$ (147)
We introduce the new variable $`z^{}=z/\chi (\tau )`$ and consider the limit $`t+\mathrm{}`$. The integral depending on the initial condition is given by
$`I={\displaystyle \frac{1}{\sqrt{\pi }}}{\displaystyle _{\mathrm{}}^+\mathrm{}}e^{x^2}(\varphi _1\eta )(\sqrt{2H_2(\tau )}(z^{}x))𝑑x.`$ (148)
We shall assume that $`\varphi _1(z)1`$ for $`z\mathrm{}`$ and $`\varphi _1(z)0`$ for $`z+\mathrm{}`$. Then, if $`H_2(\tau )+\mathrm{}`$ for $`\tau +\mathrm{}`$, the function $`(\varphi _1\eta )(\sqrt{2H_2(\tau )}(z^{}x))`$ will be very peaked around $`x=z^{}`$ and we can approximate the integral by
$`I{\displaystyle \frac{e^{z^2}}{\sqrt{2\pi H_2(\tau )}}}{\displaystyle _{\mathrm{}}^+\mathrm{}}(\varphi _1\eta )(x)𝑑x0.`$ (149)
Therefore, the condition that the solution of (126) tends asymptotically to the function (146) for $`\tau +\mathrm{}`$ is that $`(\varphi _1\eta )(\pm \mathrm{})=0`$ and
$`H_2(\tau )={\displaystyle _1^\tau }e^{H(\tau ^{})}𝑑\tau ^{}+\mathrm{}`$ (150)
for $`\tau +\mathrm{}`$. We can also write the general solution (145) in the form
$`\varphi (z,\tau )={\displaystyle \frac{1}{\sqrt{2\pi H_2(\tau )}}}{\displaystyle _{\mathrm{}}^+\mathrm{}}e^{\frac{(H_1(\tau )zx)^2}{2H_2(\tau )}}\varphi _1(x)𝑑x.`$ (151)
Returning to original variables, we get
$`u(v,t)={\displaystyle \frac{1}{\sqrt{2\pi H_2(2t)}}}{\displaystyle _{\mathrm{}}^+\mathrm{}}e^{\frac{(H_1(2t)(xx_f(t))y)^2}{2H_2(2t)}}u(y,1/2)𝑑y.`$ (152)
where $`x=x(v)`$ and we have taken by convention $`x_f(t=1/2)=0`$.
## Appendix B Asymptotic behavior of Eq. (63)
Setting $`z=(\gamma /2)y^2`$, we can rewrite Eq. (63) in the equivalent form
$`x{\displaystyle \frac{1}{\sqrt{2\gamma }}}e^{\frac{\gamma }{2}v^2}{\displaystyle _0^{\frac{\gamma }{2}v^2}}{\displaystyle \frac{e^{z\frac{\gamma }{2}v^2}}{\sqrt{z}}}𝑑z.`$ (153)
With the change of variables $`y=z+\frac{\gamma }{2}v^2`$, we get
$`x{\displaystyle \frac{1}{\gamma v}}e^{\frac{\gamma }{2}v^2}{\displaystyle _0^{\frac{\gamma }{2}v^2}}{\displaystyle \frac{e^y}{\sqrt{1\frac{2y}{\gamma v^2}}}}𝑑y.`$ (154)
Then taking the limit $`v+\mathrm{}`$, we find that
$`x{\displaystyle \frac{1}{\gamma v}}e^{\frac{\gamma }{2}v^2}{\displaystyle _0^+\mathrm{}}e^y𝑑y{\displaystyle \frac{1}{\gamma v}}e^{\frac{\gamma }{2}v^2}.`$ (155)
## Appendix C Criterion for the monotonicity of $`f(v,t)`$
In this Appendix, we establish a criterion which guarantees the monotonicity of $`f(v,t)`$ for all times, if the distribution function is initially monotonic (decreasing). Let us first rewrite the Fokker-Planck equation (18) in the form
$`{\displaystyle \frac{f}{t}}=D{\displaystyle \frac{^2f}{v^2}}+\left(D^{}+DU^{}+{\displaystyle \frac{d1}{v}}D\right){\displaystyle \frac{f}{v}}+\left[(DU^{})^{}+{\displaystyle \frac{d1}{v}}DU^{}\right]f.`$ (156)
Because of the positivity of the diffusion coefficient, a classical comparison principle states that if $`f(v,t)`$ is initially positive, it will remain positive for all times . The idea is now to apply the same argument to $`g=f/v`$. Taking the derivative of Eq. (156) with respect to $`v`$, we get
$`{\displaystyle \frac{g}{t}}=D{\displaystyle \frac{^2g}{v^2}}+\left(2D^{}+DU^{}+{\displaystyle \frac{d1}{v}}D\right){\displaystyle \frac{g}{v}}+\left[D^{\prime \prime }+2(DU^{})^{}+{\displaystyle \frac{d1}{v}}DU^{}+(d1)\left({\displaystyle \frac{D}{v}}\right)^{}\right]g`$
$`+\left[(DU^{})^{\prime \prime }+(d1)\left({\displaystyle \frac{DU^{}}{v}}\right)^{}\right]f.`$ (157)
If
$`(DU^{})^{\prime \prime }+(d1)\left({\displaystyle \frac{DU^{}}{v}}\right)^{}0,`$ (158)
then
$`{\displaystyle \frac{g}{t}}D{\displaystyle \frac{^2g}{v^2}}+\left(2D^{}+DU^{}+{\displaystyle \frac{d1}{v}}D\right){\displaystyle \frac{g}{v}}+\left[D^{\prime \prime }+2(DU^{})^{}+{\displaystyle \frac{d1}{v}}DU^{}+(d1)\left({\displaystyle \frac{D}{v}}\right)^{}\right]g.`$
(159)
Again we use a comparison principle and deduce from this inequality that if $`g(v,t)`$ is initially negative, it will remain negative for all times. Therefore, if $`f(v,t)`$ is initially decreasing, it will remain decreasing for all times if the criterion (158) is satisfied. This criterion is just a sufficient condition of monotonicity.
Let us give some particular examples of application. In the case of a quadratic potential $`U(v)=v^2/2`$, the criterion (158) becomes
$`(Dv)^{\prime \prime }+(d1)D^{}0.`$ (160)
For a diffusion coefficient of the form $`D=v^\alpha `$, we find the condition $`\alpha d`$. In particular, for the case $`\alpha =d=3`$, the criterion is satisfied. We can check that the criterion (160) is also satisfied for all $`v`$ for the diffusion coefficient (7-a) corresponding to the Coulombian (or Newtonian) potential of interaction. Alternatively, for a diffusion coefficient of the form $`D=e^{\gamma v^2}`$, we find that the criterion (160) is not satisfied for large values of $`v`$ (more precisely $`v>(d+\sqrt{d^2+8\gamma })/4\gamma `$). Therefore, the profile $`f(v,t)`$ can be non-monotonic even if the initial condition is monotonic, as in the case of Fig. 3.
## Appendix D General solution of Eq. (109)
In this Appendix, we show that Eq. (109) can be solved explicitly. First of all, we note that it can be written
$`\left(U^{}(v){\displaystyle \frac{d1}{v}}\right)RR^{}(v)=Ax+B.`$ (161)
Since $`R=dv/dx`$, we obtain
$`\left(U^{}(v){\displaystyle \frac{d1}{v}}{\displaystyle \frac{R^{}(v)}{R(v)}}\right)dv=(Ax+B)dx`$ (162)
which leads to
$`e^{U(v)}v^{d1}R(v)=Ke^{(\frac{A}{2}x^2+Bx)},`$ (163)
where $`K`$ is a constant. Equation (163) can again be integrated into
$`{\displaystyle _v^+\mathrm{}}e^{U(w)}w^{d1}𝑑w=K{\displaystyle _{x(v)}^+\mathrm{}}e^{(\frac{A}{2}y^2+By)}𝑑y.`$ (164)
For $`A=0`$, we obtain
$`{\displaystyle _v^+\mathrm{}}e^{U(w)}w^{d1}𝑑w={\displaystyle \frac{K}{B}}e^{Bx}.`$ (165)
Substituting this relation in Eq. (163), we recover the result of Eq. (120). On the other hand, for $`A0`$, we get
$`{\displaystyle _v^+\mathrm{}}e^{U(w)}w^{d1}𝑑w=K\sqrt{{\displaystyle \frac{2\pi }{A}}}e^{\frac{B^2}{2A}}\mathrm{\Phi }\left[\sqrt{{\displaystyle \frac{A}{2}}}\left({\displaystyle \frac{B}{A}}+x\right)\right],`$ (166)
where $`\mathrm{\Phi }(x)`$ is defined in Eq. (38). Substituting the foregoing relation in Eq. (163), we find that
$`D(v)={\displaystyle \frac{e^{2U(v)}}{v^{2(d1)}}}F\left[{\displaystyle _v^+\mathrm{}}Ce^{U(w)}w^{d1}𝑑w\right],`$ (167)
where
$`F(x)=\mathrm{exp}\left\{2(\mathrm{\Phi }^1(x))^2\right\},`$ (168)
and $`C`$ is a constant.
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# 1 Introduction
## 1 Introduction
Investigation of various aspects of Non-Anti-Commutative (NAC) superspace and related deformations of supersymmetric field theories are the subject of intensive studies during the last two years (see, e.g., refs. directly related to our title, and the references therein for the earlier work in the NAC-deformed N=1 superspace). Our general motivation is to enhance understanding of the role of spacetime in supersymmetry, while keeping globally supersymmetric field theory under control. Yet additional motivation is provided by string theory, where the non-anti-commutativity can be related to a constant (self-dual) gravi-photon background (see e.g., ref. and references therein).
This paper is devoted to the NAC-deformed supersymmetric Non-Linear Sigma-Models (NLSM) in four dimensions. The bosonic NLSM we consider are the most general scalar field theories, without any gauge fields or higher derivatives. Their (undeformed) N=1 supersymmetric extension is well known to require Kähler geometry of the NLSM metric and a holomorphic scalar superpotential (see e.g., the pioneering paper and ref. for a review).
We work in four-dimensional Euclidean <sup>6</sup><sup>6</sup>6The use of Atiyah-Ward spacetime of the signature $`(+,+,,)`$ is another possibility . N=1 superspace $`(x^\mu ,\theta ^\alpha ,\overline{\theta }^\stackrel{_{{}_{}{}^{}}}{\alpha })`$, and use the standard notation . The NAC deformation is given by
$$\{\theta ^\alpha ,\theta ^\beta \}_{}=C^{\alpha \beta },$$
$`(1.1)`$
where $`C^{\alpha \beta }`$ are some constants. The remaining superspace coordinates in the chiral basis ($`y^\mu =x^\mu +i\theta \sigma ^\mu \overline{\theta }`$, $`\mu ,\nu =1,2,3,4`$ and $`\alpha ,\beta ,\mathrm{}=1,2`$) still (anti)commute,
$$y^\mu ,y^\nu =\{\overline{\theta }^\stackrel{_{{}_{}{}^{}}}{\alpha },\overline{\theta }^\stackrel{_{{}_{}{}^{}}}{\beta }\}=\{\theta ^\alpha ,\overline{\theta }^\stackrel{_{{}_{}{}^{}}}{\beta }\}=y^\mu ,\theta ^\alpha =y^\mu ,\overline{\theta }^\stackrel{_{{}_{}{}^{}}}{\alpha }=0.$$
$`(1.2)`$
The $`C^{\alpha \beta }0`$ explicitly break the four-dimensional ‘Lorentz’ invariance at the fundamental level. The NAC nature of $`\theta `$’s can be fully taken into account by using the Moyal-Weyl-type (star) product of superfields ,
$$f(\theta )g(\theta )=f(\theta )\mathrm{exp}\left(\frac{C^{\alpha \beta }}{2}\frac{\stackrel{}{}}{\theta ^\alpha }\frac{\stackrel{}{}}{\theta ^\beta }\right)g(\theta ),$$
$`(1.3)`$
which respects the N=1 superspace chirality. The star product (1.3) is polynomial in the deformation parameter ,
$$f(\theta )g(\theta )=fg+(1)^{\mathrm{deg}f}\frac{C^{\alpha \beta }}{2}\frac{f}{\theta ^\alpha }\frac{g}{\theta ^\beta }detC\frac{^2f}{\theta ^2}\frac{^2g}{\theta ^2},$$
$`(1.4)`$
where we have used the identity
$$detC=\frac{1}{2}\epsilon _{\alpha \gamma }\epsilon _{\beta \delta }C^{\alpha \beta }C^{\gamma \delta },$$
$`(1.5)`$
and the notation
$$\frac{^2}{\theta ^2}=\frac{1}{4}\epsilon ^{\alpha \beta }\frac{}{\theta ^\alpha }\frac{}{\theta ^\beta }.$$
$`(1.6)`$
We also use the following book-keeping notation for 2-component spinors:
$$\theta \chi =\theta ^\alpha \chi _\alpha ,\overline{\theta }\overline{\chi }=\overline{\theta }_\stackrel{_{{}_{}{}^{}}}{\alpha }\overline{\chi }^\stackrel{_{{}_{}{}^{}}}{\alpha },\theta ^2=\theta ^\alpha \theta _\alpha ,\overline{\theta }^2=\overline{\theta }_\stackrel{_{{}_{}{}^{}}}{\alpha }\overline{\theta }^\stackrel{_{{}_{}{}^{}}}{\alpha }.$$
$`(1.7)`$
The spinorial indices are raised and lowered by the use of two-dimensional Levi-Civita symbols . Grassmann integration amounts to Grassmann differentiation. The anti-chiral covariant derivative in the chiral superspace basis is $`\overline{D}_\stackrel{_{{}_{}{}^{}}}{\alpha }=/\overline{\theta }^\stackrel{_{{}_{}{}^{}}}{\alpha }`$. The field component expansion of a chiral superfield $`\mathrm{\Phi }`$ reads
$$\mathrm{\Phi }(y,\theta )=\varphi (y)+\sqrt{2}\theta \chi (y)+\theta ^2M(y).$$
$`(1.8)`$
An anti-chiral superfield $`\overline{\mathrm{\Phi }}`$ in the chiral basis is given by
| $`\overline{\mathrm{\Phi }}(y^\mu 2i\theta \sigma ^\mu \overline{\theta },\overline{\theta })=`$ | $`\overline{\varphi }(y)+\sqrt{2}\overline{\theta }\overline{\chi }(y)+\overline{\theta }^2\overline{M}(y)`$ |
| --- | --- |
| | $`+\sqrt{2}\theta \left(i\sigma ^\mu _\mu \overline{\chi }(y)\overline{\theta }^2i\sqrt{2}\sigma ^\mu \overline{\theta }_\mu \overline{\varphi }(y)\right)+\theta ^2\overline{\theta }^2\mathrm{}\overline{\varphi }(y),`$ |
$`(1.9)`$
where $`\mathrm{}=_\mu _\mu `$. The bars over fields serve to distinguish between the ‘left’ and ‘right’ components that are truly independent in Euclidean spacetime.
The non-anticommutativity $`C_{\alpha \beta }0`$ also explicitly breaks half of the original N=1 supersymmetry . Only the chiral subalgebra generated by the chiral supercharges (in the chiral basis) $`Q_\alpha =/\theta ^\alpha `$ is preserved, with $`\{Q_\alpha ,Q_\beta \}_{}=0`$, thus defining what is now called N=1/2 supersymmetry. The use of the NAC-deformed superspace allows one to keep N=1/2 supersymmetry manifest. The N=1/2 supersymmetry transformation laws of the chiral and anti-chiral superfield components in eqs. (1.8) and (1.9) are as follows:
$$\delta \varphi =\sqrt{2}\epsilon ^\alpha \chi _\alpha ,\delta \chi _\alpha =\sqrt{2}\epsilon _\alpha M,\delta M=0,$$
$`(1.10)`$
and
$$\delta \overline{\varphi }=0,\delta \overline{\chi }^\stackrel{_{{}_{}{}^{}}}{\alpha }=i\sqrt{2}(\stackrel{~}{\sigma }_\mu )^{\stackrel{_{{}_{}{}^{}}}{\alpha }\beta }\epsilon _\beta _\mu \overline{\varphi },\delta \overline{M}=i\sqrt{2}_\mu \overline{\chi }_\stackrel{_{{}_{}{}^{}}}{\alpha }(\stackrel{~}{\sigma }_\mu )^{\stackrel{_{{}_{}{}^{}}}{\alpha }\beta }\epsilon _\beta ,$$
$`(1.11)`$
respectively, where we have introduced the N=1/2 supersymmetry (chiral) parameter $`\epsilon ^\alpha `$.
The most general four-dimensional N=1 supersymmetric NLSM action (without any gauge fields) is given by
$$S[\mathrm{\Phi },\overline{\mathrm{\Phi }}]=d^4y[d^2\theta d^2\overline{\theta }K(\mathrm{\Phi }^i,\overline{\mathrm{\Phi }}{}_{}{}^{\overline{j}})+d^2\theta W(\mathrm{\Phi }^i)+d^2\overline{\theta }\overline{W}(\overline{\mathrm{\Phi }}{}_{}{}^{\overline{j}})].$$
$`(1.12)`$
This action is completely specified by the Kähler superpotential $`K(\mathrm{\Phi },\overline{\mathrm{\Phi }})`$, the scalar superpotential $`W(\mathrm{\Phi })`$, and the anti-chiral superpotential $`\overline{W}(\overline{\mathrm{\Phi }})`$, in terms of some number $`n`$ of chiral and anti-chiral superfields, $`i,\overline{j}=1,2,\mathrm{},n`$. In Euclidean superspace the functions $`W(\mathrm{\Phi })`$ and $`\overline{W}(\overline{\mathrm{\Phi }})`$ are independent upon each other.
The problem of computing the NAC-deformed extension of eq. (1.12) in four dimensions, after a ‘Seiberg-Witten map’ (i.e. after evaluating all the star products), was solved in ref. in the case of a single (anti)chiral superfield. That solution was extended to the case of several (anti)chiral superfields in ref. . In both cases the perturbative solutions were found, i.e. in terms of the infinite sums with respect to the deformation parameter and the auxiliary fields. A similar problem in two dimensions was perturbatively solved in ref. . A non-perturbative solution (i.e. in a closed form, in terms of finite functions) was found in ref. in two dimensions, and in ref. in four dimensions. The calculations behind each result appear to be quite extensive, so it is non-trivial to make a comparison. In our opinion, an explicit summation is necessary anyway, both for comparison and elimination of the auxiliary fields. Getting the results in a closed form is also crucial for any physical applications of the NAC-deformation and its geometrical interpretation, as well as in investigating concrete examples (see below). In this paper we use our four-dimensional results that are in precise agreement with the basis formulae of ref. , after dimensional reduction to two dimensions. The non-perturbative results of refs. and presumably amount to the full summation of the perturbative results in refs. and , respectively.
We use the chiral basis that is most suitable for computing NAC-deformation, by reducing the most non-trivial problem of calculation of the NAC-deformed Kähler superpotential to that for the NAC-deformed scalar superpotential . The simple general results about the NAC-deformed scalar superpotentials were found in refs. . Our major concern in this paper is getting solutions to the auxiliary fields that control splitting or fuzziness in the NLSM target space.
Our paper is organized as follows. In sect. 2 we summarize the earlier results that represent our setup here. In sect. 3 we eliminate the auxiliary fields in a generic NAC-deformed NLSM, by using implicit functions. The case of a single (anti)chiral superfield is also considered separately. In sect. 4 we specify our results to the case of the $`CP^n`$ NLSM. The $`CP^1`$ case is fully investigated. Sect. 5 is our conclusion.
## 2 The NAC Kähler potential and superpotentials
We use the following notation valid for any function $`F(\varphi ,\overline{\varphi })`$:
$$F_{,i_1i_2\mathrm{}i_s\overline{p}_1\overline{p}_2\mathrm{}\overline{p}_t}=\frac{^{s+t}F}{\varphi ^{i_1}\varphi ^{i_2}\mathrm{}\varphi ^{i_s}\overline{\varphi }^{\overline{p}_1}\overline{\varphi }^{\overline{p}_2}\mathrm{}\overline{\varphi }^{\overline{p}_t}},$$
$`(2.1)`$
and the Grassmann integral normalisation $`d^2\theta \theta ^2=1`$. The actual deformation parameter, in the case of the NAC-deformed field theory (1.12), appears to be
$$c=\sqrt{detC},$$
$`(2.2)`$
where we have used the definition
$$detC=\frac{1}{2}\epsilon _{\alpha \gamma }\epsilon _{\beta \delta }C^{\alpha \beta }C^{\gamma \delta }.$$
$`(2.3)`$
As a result, unlike the case of the NAC-deformed supersymmetric gauge theories , the NAC-deformation of the NLSM field theory (1.12) appears to be ‘Lorentz’-invariant .
A simple non-perturbative formula, describing an arbitrary NAC-deformed scalar superpotential $`V`$ depending upon a single chiral superfield $`\mathrm{\Phi }`$, was found in ref. ,
$$d^2\theta V_{}(\mathrm{\Phi })=\frac{1}{2c}\left[V(\varphi +cM)V(\varphi cM)\right]\frac{\chi ^2}{4cM}\left[V_{,\varphi }(\varphi +cM)V_{,\varphi }(\varphi cM)\right].$$
$`(2.4)`$
The NAC-deformation in the single superfield case thus gives rise to the split of the scalar potential, which is controlled by the auxiliary field $`M`$. When using an elementary identity
$$f(x+a)f(xa)=a\frac{}{x}_1^{+1}𝑑\xi f(x+\xi a),$$
$`(2.5)`$
valid for any function $`f`$, we can rewrite eq. (2.4) to the equivalent form, as in ref. ,
$$d^2\theta V_{}(\mathrm{\Phi })=\frac{1}{2}M\frac{}{\varphi }_1^{+1}𝑑\xi V(\varphi +\xi cM)\frac{1}{4}\chi ^2\frac{^2}{\varphi ^2}_1^{+1}𝑑\xi V(\varphi +\xi cM),$$
$`(2.6)`$
Similarly, in the case of several chiral superfields, one finds
$$d^2\theta V_{}(\mathrm{\Phi }^I)=\frac{1}{2}M^I\frac{}{\varphi ^I}\stackrel{~}{V}(\varphi ,M)\frac{1}{4}(\chi ^I\chi ^J)\frac{^2}{\varphi ^I\varphi ^J}\stackrel{~}{V}(\varphi ,M),$$
$`(2.7)`$
in terms of the auxiliary pre-potential $`\stackrel{~}{V}`$ ,
$$\stackrel{~}{V}(\varphi ,M)=_1^{+1}𝑑\xi V(\varphi ^I+\xi cM^I).$$
$`(2.8)`$
Hence the NAC-deformation of a generic scalar superpotential $`V`$ results in its smearing or fuzziness controlled by the auxiliary fields $`M^I`$ .
A calculation of the NAC deformed Kähler potential
$$d^4yL_{\mathrm{kin}.}d^4yd^2\theta d^2\overline{\theta }K(\mathrm{\Phi }^i,\overline{\mathrm{\Phi }}{}_{}{}^{\overline{j}})_{}$$
$`(2.9)`$
can be reduced to eqs. (2.4) or (2.7), when using a chiral reduction in superspace, with the following result :
| $`L_{\mathrm{kin}.}=`$ | $`\frac{1}{2}M^iY_{,i}+\frac{1}{2}^\mu \overline{\varphi }{}_{}{}^{\overline{p}}_{\mu }^{}\overline{\varphi }{}_{}{}^{\overline{q}}Z_{,\overline{p}\overline{q}}^{}+\frac{1}{2}\mathrm{}\overline{\varphi }{}_{}{}^{\overline{p}}Z_{,\overline{p}}^{}\frac{1}{4}(\chi ^i\chi ^j)Y_{,ij}`$ |
| --- | --- |
| | $`\frac{1}{2}i(\chi ^i\sigma ^\mu \overline{\chi }{}_{}{}^{\overline{p}})_\mu \overline{\varphi }{}_{}{}^{\overline{q}}Z_{,i\overline{p}\overline{q}}^{}\frac{1}{2}i(\chi ^i\sigma ^\mu _\mu \overline{\chi }{}_{}{}^{\overline{p}})Z_{,i\overline{p}},`$ |
$`(2.10)`$
where we have introduced the (component) smeared Kähler pre-potential
$$Z(\varphi ,\overline{\varphi },M)=_1^{+1}d\xi K^\xi \mathrm{with}K^\xi K(\varphi ^i+\xi cM^i,\overline{\varphi }{}_{}{}^{\overline{j}}),$$
$`(2.11)`$
as well as the extra (auxiliary) pre-potential
$$Y(\varphi ,\overline{\varphi },M,\overline{M})=\overline{M}{}_{}{}^{\overline{p}}Z_{,\overline{p}}^{}\frac{1}{2}(\overline{\chi }{}_{}{}^{\overline{p}}\overline{\chi }{}_{}{}^{\overline{q}})Z_{,\overline{p}\overline{q}}+c_1^{+1}d\xi \xi [^\mu \overline{\varphi }{}_{}{}^{\overline{p}}_{\mu }^{}\overline{\varphi }{}_{}{}^{\overline{q}}K_{,\overline{p}\overline{q}}^{\xi }+\mathrm{}\overline{\varphi }{}_{}{}^{\overline{p}}K_{,\overline{p}}^{\xi }],$$
$`(2.12)`$
as in ref. . It is not difficult to check that eq. (2.10) does reduce to the standard (Kähler) N=1 supersymmetric NLSM in the limit $`c0`$. Also, in the case of a free (bilinear) Kähler potential $`K=\delta _{i\overline{j}}\mathrm{\Phi }^i\overline{\mathrm{\Phi }}^{\overline{j}}`$, there is no deformation at all.
The NAC-deformed scalar superpotentials $`W(\mathrm{\Phi })_{}`$ and $`\overline{W}(\overline{\mathrm{\Phi }})_{}`$ imply, via eqs. (2.7) and (2.8), that the following component terms are to be added to eq. (2.10):
$$L_{\mathrm{pot}.}=\frac{1}{2}M^i\stackrel{~}{W}_{,i}\frac{1}{4}(\chi ^i\chi ^j)\stackrel{~}{W}_{,ij}+\overline{M}^{\overline{p}}\overline{W}_{,\overline{p}}\frac{1}{2}(\overline{\chi }^{\overline{p}}\overline{\chi }^{\overline{q}})\overline{W}_{,\overline{p}\overline{q}}.$$
$`(2.13)`$
where we have introduced the smeared scalar pre-potential
$$\stackrel{~}{W}(\varphi ,M)=_1^{+1}𝑑\xi W(\varphi ^i+\xi cM^i).$$
$`(2.14)`$
The anti-chiral scalar superpotential terms are inert under the NAC-deformation .
The $`\xi `$-integrations in the equations above represent the smearing effects. However, the smearing is merely apparent in the case of a single chiral superfield, which gives rise to the splitting (2.4) only. This can also be directly demonstrated from eq. (2.10) when using the identity (2.5) together with the related identity
$$f(x+a)+f(xa)=_1^{+1}𝑑\xi f(x+\xi a)+a\frac{}{x}_1^{+1}𝑑\xi \xi f(x+\xi a).$$
$`(2.15)`$
The single superfield case thus appears to be special, so that a sum of eq. (2.10) and (2.13) can be rewritten to the bosonic contribution
| $`L_{\mathrm{bos}.}=`$ | $`+\frac{1}{2}^\mu \overline{\varphi }_\mu \overline{\varphi }\left[K_{,\overline{\varphi }\overline{\varphi }}(\varphi +cM,\overline{\varphi })+K_{,\overline{\varphi }\overline{\varphi }}(\varphi cM,\overline{\varphi })\right]`$ |
| --- | --- |
| | $`+\frac{1}{2}\mathrm{}\overline{\varphi }\left[K_{,\overline{\varphi }}(\varphi +cM,\overline{\varphi })+K_{,\overline{\varphi }}(\varphi cM,\overline{\varphi })\right]`$ |
| | $`+{\displaystyle \frac{\overline{M}}{2c}}\left[K_{,\overline{\varphi }}(\varphi +cM,\overline{\varphi })K_{,\overline{\varphi }}(\varphi cM,\overline{\varphi })\right]`$ |
| | $`+{\displaystyle \frac{1}{2c}}\left[W(\varphi +cM)W(\varphi cM)\right]+\overline{M}{\displaystyle \frac{\overline{W}}{\overline{\varphi }}},`$ |
$`(2.16)`$
supplemented by the following fermionic terms :
| $`L_{\mathrm{ferm}.}=`$ | $`{\displaystyle \frac{1}{4c}}\overline{\chi }^2\left[K_{,\overline{\varphi }\overline{\varphi }}(\varphi +cM,\overline{\varphi })K_{,\overline{\varphi }\overline{\varphi }}(\varphi cM,\overline{\varphi })\right]`$ |
| --- | --- |
| | $`{\displaystyle \frac{i}{2cM}}(\chi \sigma ^\mu \overline{\chi })_\mu \overline{\varphi }\left[K_{,\overline{\varphi }\overline{\varphi }}(\varphi +cM,\overline{\varphi })K_{,\overline{\varphi }\overline{\varphi }}(\varphi cM,\overline{\varphi })\right]`$ |
| | $`{\displaystyle \frac{i}{2cM}}(\chi \sigma ^\mu _\mu \overline{\chi })\left[K_{,\overline{\varphi }}(\varphi +cM,\overline{\varphi })K_{,\overline{\varphi }}(\varphi cM,\overline{\varphi })\right]`$ |
| | $`{\displaystyle \frac{\overline{M}}{4cM}}\chi ^2\left[K_{,\varphi \overline{\varphi }}(\varphi +cM,\overline{\varphi })K_{,\varphi \overline{\varphi }}(\varphi cM,\overline{\varphi })\right]`$ |
| | $`{\displaystyle \frac{1}{4M}}\chi ^2^\mu \overline{\varphi }_\mu \overline{\varphi }\left[K_{,\varphi \overline{\varphi }\overline{\varphi }}(\varphi +cM,\overline{\varphi })+K_{,\varphi \overline{\varphi }\overline{\varphi }}(\varphi cM,\overline{\varphi })\right]`$ |
| | $`+{\displaystyle \frac{1}{4cM^2}}\chi ^2^\mu \overline{\varphi }_\mu \overline{\varphi }\left[K_{,\overline{\varphi }\overline{\varphi }}(\varphi +cM,\overline{\varphi })K_{,\overline{\varphi }\overline{\varphi }}(\varphi cM,\overline{\varphi })\right]`$ |
| | $`{\displaystyle \frac{1}{4M}}\chi ^2\mathrm{}\overline{\varphi }\left[K_{,\varphi \overline{\varphi }}(\varphi +cM,\overline{\varphi })+K_{,\varphi \overline{\varphi }}(\varphi cM,\overline{\varphi })\right]`$ |
| | $`+{\displaystyle \frac{1}{4cM^2}}\chi ^2\mathrm{}\overline{\varphi }\left[K_{,\overline{\varphi }}(\varphi +cM,\overline{\varphi })K_{,\overline{\varphi }}(\varphi cM,\overline{\varphi })\right]`$ |
| | $`+{\displaystyle \frac{1}{8cM}}\chi ^2\overline{\chi }^2\left[K_{,\varphi \overline{\varphi }\overline{\varphi }}(\varphi +cM,\overline{\varphi })K_{,\varphi \overline{\varphi }\overline{\varphi }}(\varphi cM,\overline{\varphi })\right]`$ |
| | $`{\displaystyle \frac{1}{4cM}}\chi ^2\left[W_{,\varphi }(\varphi +cM)W_{,\varphi }(\varphi cM)\right]\frac{1}{2}\overline{\chi }^2\overline{W}_{,\overline{\varphi }\overline{\varphi }}.`$ |
$`(2.17)`$
The anti-chiral auxiliary fields $`\overline{M}^{\overline{p}}`$ enter the action (2.10) linearly (as Lagrange multipliers), while their algebraic equations of motion,
$$\frac{1}{2}M^iZ_{,i\overline{p}}\frac{1}{4}(\chi ^i\chi ^j)Z_{,ij\overline{p}}+\overline{W}_{,\overline{p}}=0,$$
$`(2.18)`$
represent the non-linear set of equations on the auxiliary fields $`M^i=M^i(\varphi ,\overline{\varphi })`$<sup>7</sup><sup>7</sup>7Equation (2.18) is not a linear system because the function $`Z`$ is $`M`$-dependent — see eq. (2.11). As a result, the bosonic scalar potential in components is given by
$$V_{\mathrm{scalar}}(\varphi ,\overline{\varphi })=\frac{1}{2}M^i\stackrel{~}{W}_{,i}|_{M=M(\varphi ,\overline{\varphi })}.$$
$`(2.19)`$
## 3 Solving for the auxiliary fields in a generic case
The NAC-deformed NLSM in sect. 2 is completely specified by a Kähler function $`K(\mathrm{\Phi },\overline{\mathrm{\Phi }})`$, a chiral function $`W(\mathrm{\Phi })`$, an anti-chiral function $`\overline{W}(\overline{\mathrm{\Phi }})`$ and a constant deformation parameter $`c`$. As a matter of fact, we didn’t really used the constancy of $`c`$, so our results in sect. 2 are still valid even for a coordinate-dependent NAC deformation $`c(y)`$. The very possibility of such extension preserving N=1/2 supersymmetry was conjectured in ref. and further investigated in ref. . In this section we assume, however, that $`c=const`$ for simplicity.
A NAC-deformation is only possible in Euclidean spacetime. However, we may try to press further, by analytically continue our results into Minkowski spacetime (this would require $`detC<0`$). Then the natural place for possible physical applications of NAC-deformation could be given by confining supersymmetric gauge field theories whose low-energy (IR) effective action takes the form of eq. (1.12) in terms some colorless composite scalar superfields $`\mathrm{\Phi }`$ and $`\overline{\mathrm{\Phi }}`$ known as the glueball superfields. The use of Konishi anomaly equations allows one to construct exact effective superpotentials of the glueball superfields in various N=1 supersymmetric quantum gauge theories with some number of colors and flavors (see e.g., ref. for details), generalizing the standard Veneziano-Yankielowicz effective potential . This may provide some natural input for the NAC deformation presumably describing some supergravitational corrections to the low-energy effective field theory, as well as for a possible dynamical supersymmetry breaking (see ref. for some attempts in this direction). <sup>8</sup><sup>8</sup>8The Konishi anomaly in N=1/2 supersymmetry was recently studied in ref. . From that physical point of view, the fact that the NAC-deformation gives rise to a ‘Lorentz’-invariant action is clearly a positive feature, whereas the apparent non-Hermiticity of the NAC-deformed action is clearly a negative feature. In this paper we are going to keep the chiral and anti-chiral functions to be arbitrary and concentrate on an investigation of the NAC-deformed kinetic terms.
Solving for the auxiliary fields in eq. (2.10) represents not only a technical but also a conceptual problem because of the smearing effects described by the $`\xi `$-integrations. To bring the kinetic terms in eqs. (2.10) and (2.16) to the standard NLSM form (i.e. without the second order spacial derivatives), one has to integrate by parts that leads to the appearance of the spacial derivatives of the auxiliary fields. This implies that one has to solve eq. (2.18) before integration by parts. Let $`M^i=M^i(\varphi ,\overline{\varphi })`$ be a solution to eq. (2.18), and let’s ignore fermions for simplicity $`(\chi _\alpha ^i=\overline{\chi }_\stackrel{_{{}_{}{}^{}}}{\alpha }^{\overline{p}}=0)`$. Substituting the solution back to the Lagrangian (2.10) and integrating by parts yields
| $`L_{\mathrm{kin}.}(\varphi ,\overline{\varphi })=`$ | $`\frac{1}{2}(_\mu \overline{\varphi }^{\overline{p}}_\mu \varphi ^q){\displaystyle _1^{+1}}𝑑\xi \left[K_{,\overline{p}q}^\xi +2c\xi M_{,q}^iK_{,\overline{p}i}^\xi +c\xi M^iK_{,\overline{p}iq}^\xi +c^2\xi ^2M^iK_{,\overline{p}ij}^\xi M_{,q}^j\right]`$ |
| --- | --- |
| | $`\frac{1}{2}(_\mu \overline{\varphi }^{\overline{p}}_\mu \overline{\varphi }^{\overline{q}}){\displaystyle _1^{+1}}𝑑\xi \left[2c\xi K_{,\overline{p}i}^\xi M_{,\overline{q}}^i+c^2\xi ^2M^iK_{,\overline{p}ij}^\xi M_{,\overline{q}}^j\right].`$ |
$`(3.1)`$
It is now apparent that the NAC-deformation does not preserve the original Kähler geometry in eq. (1.12), though the absence of $`(_\mu \varphi )^2`$ terms and the particular structure of various contributions to eq. (3.1) are remarkable. The action (3.1) takes the form of a generic NLSM, being merely dependent upon mixed derivatives of the Kähler function, so that the original Kähler gauge invariance of eq. (1.12),
$$K(\varphi ,\overline{\varphi })K(\varphi ,\overline{\varphi })+f(\varphi )+\overline{f}(\overline{\varphi }),$$
$`(3.2)`$
with arbitrary gauge functions $`f(\varphi )`$ and $`\overline{f}(\overline{\varphi })`$ is still preserved.
However, a problem arises with the $`\xi `$-integrations in eq. (3.1), because the fields $`M^i`$ are no longer independent and, hence, eq. (2.5) cannot be applied. Instead, eq. (2.5) may have to be replaced by a more general identity
| $`{\displaystyle \frac{d}{dx}}{\displaystyle _1^{+1}}𝑑\xi f(x+\xi a(x))=`$ | $`{\displaystyle \frac{1}{a(x)}}\left[f(x+a(x))f(xa(x))\right]`$ |
| --- | --- |
| | $`+{\displaystyle \frac{1}{a(x)}}\left[f(x+a(x))+f(xa(x))\right]{\displaystyle \frac{da}{dx}}`$ |
| | $`{\displaystyle \frac{1}{a^2(x)}}\left[F(x+a(x))F(xa(x))\right]{\displaystyle \frac{da}{dx}},`$ |
$`(3.3)`$
where we have defined $`F(x)=^xf(\xi )𝑑\xi `$. We believe, however, that this way of doing is not quite correct from the viewpoint of the original star product that has nothing to do with dynamics. In other words, the elimination of the auxiliary fields and the $`\xi `$-integrations do not commute, while, in our opinion, all the $`\xi `$-integrations have to be done before solving for the auxiliary fields. This procedure will be adopted in our explicit calculations for the $`CP^n`$ case in the next sect. 4. It is worth mentioning that the ordering problem does not arise in the case of a sinlge chiral superfield, because there is no need for $`\xi `$-integrations there (i.e. when smearing is replaced by splitting). In the last case, when using eq. (2.16), we find
| $`L_{\mathrm{kin}.}(\varphi ,\overline{\varphi })=`$ | $`\frac{1}{2}_\mu \overline{\varphi }_\mu \varphi \left[K_{,\varphi \overline{\varphi }}(\varphi +cM,\overline{\varphi })+K_{,\varphi \overline{\varphi }}(\varphi cM,\overline{\varphi })\right]`$ |
| --- | --- |
| | $`\frac{1}{2}_\mu \overline{\varphi }_\mu \varphi \left[cK_{,\varphi \overline{\varphi }}(\varphi +cM,\overline{\varphi })cK_{,\varphi \overline{\varphi }}(\varphi cM,\overline{\varphi })\right]{\displaystyle \frac{M}{\varphi }}`$ |
| | $`\frac{1}{2}_\mu \overline{\varphi }_\mu \overline{\varphi }\left[cK_{,\varphi \overline{\varphi }}(\varphi +cM,\overline{\varphi })cK_{,\varphi \overline{\varphi }}(\varphi cM,\overline{\varphi })\right]{\displaystyle \frac{M}{\overline{\varphi }}},`$ |
$`(3.4a)`$
where $`M(\varphi ,\overline{\varphi })`$ is supposed to be a solution to eq. (2.18). Equation (3.4a) can be rewritten to some other equivalent forms, viz.
$$L_{\mathrm{kin}.}(\varphi ,\overline{\varphi })=\frac{1}{2}(_\mu \overline{\varphi }_\mu \varphi )\frac{}{\varphi }\left[K_{,\overline{\varphi }}(\varphi +cM(\varphi ,\overline{\varphi }),\overline{\varphi })+K_{,\overline{\varphi }}(\varphi cM(\varphi ,\overline{\varphi }),\overline{\varphi })\right]$$
$$\frac{1}{2}(_\mu \overline{\varphi }_\mu \overline{\varphi })\left[cK_{,\varphi \overline{\varphi }}(\varphi +cM(\varphi ,\overline{\varphi }),\overline{\varphi })cK_{,\varphi \overline{\varphi }}(\varphi cM(\varphi ,\overline{\varphi }),\overline{\varphi })\right]\frac{M(\varphi ,\overline{\varphi })}{\overline{\varphi }},$$
$`(3.4b)`$
or
$$L_{\mathrm{kin}.}(\varphi ,\overline{\varphi })=\frac{1}{2}(_\mu \overline{\varphi }_\mu \varphi )\frac{}{\varphi }\frac{}{\overline{\varphi }}\left[K(\varphi +cM(\varphi ,\overline{\varphi }),\overline{\varphi })+K(\varphi cM(\varphi ,\overline{\varphi }),\overline{\varphi })\right]$$
$$+\frac{1}{2}(_\mu \overline{\varphi }_\mu \varphi )\frac{}{\varphi }\left[cK_{,\varphi }(\varphi +cM,\overline{\varphi })cK_{,\varphi }(\varphi cM,\overline{\varphi })\right]\frac{M(\varphi ,\overline{\varphi })}{\overline{\varphi }}$$
$$\frac{1}{2}(_\mu \overline{\varphi }_\mu \overline{\varphi })\left[cK_{,\varphi \overline{\varphi }}(\varphi +cM(\varphi ,\overline{\varphi }),\overline{\varphi })cK_{,\varphi \overline{\varphi }}(\varphi cM(\varphi ,\overline{\varphi }),\overline{\varphi })\right]\frac{M(\varphi ,\overline{\varphi })}{\overline{\varphi }}.$$
$`(3.4c)`$
For instance, the first line of eq. (3.4c) gives some Kähler-like terms with the NAC-deformed (splitted) Kähler potential, whereas the second and third lines of eq. (3.4c) apparently violate both Kählerian and Hermitian structures of the original (undeformed) NLSM when $`c0`$ and $`M(\varphi ,\overline{\varphi })/\overline{\varphi }0`$ (cf. the last ref. ).
Therefore, the NAC deformation of the Kähler kinetic terms in eq. (1.12) amounts to a non-Kählerian and non-Hermitian deformation of the original Kählerian and Hermitian NLSM, which is controlled by the auxiliary field solution to eq. (2.18). In the case of a single chiral superfield, the deformed NLSM metric can be read off from eq. (3.4). In the case of several superfields, the deformed NLSM can be read off from eq. (3.1), when assuming all the $`\xi `$-integrations to be performed with the auxiliary fields considered as the parameters or spectators.
Being unable to explicitly solve eq. (2.18) in a generic case, in the next sect. 4 we consider the simplest non-trivial example provided by the $`CP^n`$ NLSM with an (undeformed) Kähler, Hermitian and symmetric target space.
## 4 NAC-deformed $`CP^n`$ NLSM
The $`CP^n`$ (principal) NLSM is characterized by a Kähler potential
$$K(\varphi ,\overline{\varphi })=\alpha \mathrm{ln}(1+\kappa ^2\varphi \overline{\varphi }),$$
$`(4.1)`$
where we have used the notation $`\varphi \overline{\varphi }\left|\varphi \right|{}_{}{}^{2}\varphi ^i\overline{\varphi }^i`$. Due to a Grassmann nature of the fermionic fields, getting explicit results with fermions is possible. However, those results appear to be cumbersome and not very illuminating. Therefore, we ignore fermions for simplicity in this section, except of the simplest $`CP^1`$ case.
Equations (2.18) in the $`CP^n`$ case read as follows:
$$\left(\varphi ^i\frac{(1+\kappa ^2|\varphi |{}_{}{}^{2})M^i}{\kappa ^2P}\right)\left(\frac{1}{1+\kappa ^2\left|\varphi \right|{}_{}{}^{2}+c\kappa ^2P}\frac{1}{1+\kappa ^2\left|\varphi \right|{}_{}{}^{2}c\kappa ^2P}\right)$$
$$+\frac{2c\overline{W}_{,\overline{\varphi }^i}}{\alpha \kappa ^2}=0,$$
$`(4.2)`$
where we have used the notation $`P=M^i\overline{\varphi }^i`$.
Multiplying eq. (4.2) with $`\overline{\varphi }^i`$ gives rise to a quadratic equation on P,
$$\frac{\alpha \kappa ^2P}{(\kappa ^2+|\varphi |{}_{}{}^{2})^2c^2P^2}+U=0,$$
$`(4.3)`$
where we have introduced yet another notation $`U=\overline{\varphi }^i\overline{W}_{,\overline{\varphi }^i}`$.
A solution to eq. (4.3) is given by
$$P=\frac{\alpha \kappa ^2\sqrt{\alpha ^2\kappa ^4+4U^2c^2(\kappa ^2+|\varphi |{}_{}{}^{2})^2}}{2c^2U},$$
$`(4.4)`$
where we have chosen the minus sign in front of the square root, in order to assure the existence of the $`c0`$ limit. In the case of the $`CP^1`$ model, eq. (4.4) already gives a solution for $`M=\overline{\varphi }^1P`$ .
In the $`CP^n`$ case with $`n>1`$, substituting the solution (4.4) back into eq. (4.2) yields the following result:
$$M^i=\frac{P\overline{\varphi }^i}{\kappa ^2+\left|\varphi \right|^2}\frac{(\kappa ^2+|\varphi |{}_{}{}^{2})^2c^2P^2}{\alpha (\kappa ^2+|\varphi |{}_{}{}^{2})}U.$$
$`(4.5)`$
In the anti-commutative limit $`c0`$, eq. (4.4) yields
$$\underset{c0}{lim}P=\frac{(\kappa ^2+|\varphi |{}_{}{}^{2})^2}{\alpha \kappa ^2}\left(\overline{\varphi }^i\overline{W}_{,\overline{\varphi }^i}\right).$$
$`(4.6)`$
In the case of vanishing anti-chiral superpotential, $`\overline{W}=0`$, we have $`U=P=0`$, so that the bosonic terms in the auxiliary field solution also vanish and, hence, there is no deformation of the bosonic $`CP^n`$ kinetic terms at all. Given an arbitrary anti-chiral superpotential $`\overline{W}(\overline{\mathrm{\Phi }})0`$, the deformed $`CP^n`$ metric is rather complicated (see e.g., the $`CP^1`$ result at the end of this section).
To give an explicit example of the NAC-deformed structure of the NLSM fermionic terms, let’s consider the $`CP^1`$ model in the case of vanishing anti-chiral superpotential. Equation (2.18) now reads
$$M_1^{+1}𝑑\xi K_{,\varphi \overline{\varphi }}(\varphi +c\xi M,\overline{\varphi })\frac{1}{2}\chi ^2_1^{+1}𝑑\xi K_{,\varphi \varphi \overline{\varphi }}(\varphi +c\xi M,\overline{\varphi })=0,$$
$`(4.7)`$
whose Kähler function is given by eq. (4.1). The integrations over $`\xi `$ in eq. (4.7) yield
$$M\left[\frac{\alpha \kappa ^2}{cM\overline{\varphi }(\kappa ^2+\overline{\varphi }(\varphi +cM))}+\frac{\alpha \kappa ^2}{cM\overline{\varphi }(\kappa ^2+\overline{\varphi }(\varphi cM))}\right]$$
$`(4.8)`$
$$\frac{1}{2}\chi ^2\left[\frac{\alpha \kappa ^2}{cM(\kappa ^2+\overline{\varphi }(\varphi +cM))^2}\frac{\alpha \kappa ^2}{cM(\kappa ^2+\overline{\varphi }(\varphi cM))^2}\right]=0,$$
which gives rise to a simple cubic equation on $`M`$<sup>9</sup><sup>9</sup>9In the case of a non-vanishing anti-chiral superpotential, one gets a quartic equation.
$$M^3(c\overline{\varphi })^2(\kappa ^2+|\varphi |{}_{}{}^{2})^2M\frac{\chi ^2}{c^2\overline{\varphi }}(\kappa ^2+|\varphi |{}_{}{}^{2})=0.$$
$`(4.9)`$
The use of Cardano formula for the roots of a cubic equation gives us three solutions as follows:
$$M=\omega ^m\left[\frac{A_0}{2}+\sqrt{\frac{A_0^2}{4}+\frac{A_1^3}{27}}\right]^{1/3}+\omega ^{3m}\left[\frac{A_0}{2}\sqrt{\frac{A_0^2}{4}+\frac{A_1^3}{27}}\right]^{1/3},$$
$`(4.10)`$
where $`\omega =(1+\sqrt{3})/2`$, $`m=0,1,2`$ and
| $`A_0`$ | $`={\displaystyle \frac{\chi ^2}{c^2\overline{\varphi }}}(\kappa ^2+|\varphi |{}_{}{}^{2}),`$ |
| --- | --- |
| $`A_1`$ | $`=(c\overline{\varphi })^2(\kappa ^2+|\varphi |{}_{}{}^{2})^2.`$ |
$`(4.11)`$
However, due to the nilpotency property of fermions, $`(\chi ^2)^2=0`$ in the $`CP^1`$ case, <sup>10</sup><sup>10</sup>10In the $`CP^n`$ case, we merely have $`(\chi ^i\chi ^j)^{n+1}=0`$. eq. (4.10) is greatly simplified to the followong three solutions:
$$M_0=\frac{\overline{\varphi }\chi ^2}{\kappa ^2+\left|\varphi \right|^2}$$
$`(4.12a)`$
and
$$M_\pm =\frac{\pm (\kappa ^2+|\varphi |{}_{}{}^{2})}{c\overline{\varphi }}+\frac{\chi ^2\overline{\varphi }}{2(\kappa ^2+|\varphi |{}_{}{}^{2})}.$$
$`(4.12b)`$
The solution (4.12a) is clearly the same as that in the undeformed case, whereas the other two solutions are singular in the undeformed limit $`c0`$. We are unaware of any possible physical significance of the singular solutions.
Inserting the solution (4.12a) into eqs. (2.16) and (2.17), and using the nilpotency condition, $`(\chi ^2)^2=0`$, give rise to the standard (undeformed) four-dimensional N=1 supersymmetric $`CP^1`$ NLSM in components. It appears to coincide with the result in the $`CP^1`$ case, where also no NAC-deformation of the $`CP^1`$ supersymmetric NLSM was discovered, though our NAC-deformation is different from the one considered there (see sect. 5).
It is of interest to get an explicit NAC-deformed NLSM metric in the case of a non-vanishing anti-chiral potential, $`\overline{W}0`$. Even, in the $`CP^1`$ case, it gives rise to a rather complicated metric, when using eqs. (3.4), (4.1) and the auxiliary field solution
$$M=\frac{\alpha \sqrt{\alpha ^2+\left(2c\overline{\varphi }(1+\kappa ^2\varphi \overline{\varphi })\overline{W}_{,\overline{\varphi }}\right)^2}}{2c^2\kappa ^2\overline{\varphi }^2\overline{W}_{,\overline{\varphi }}}$$
$`(4.13)`$
that follows from eq. (4.4), where we have used the notation $`\overline{W}_{,\overline{\varphi }}=\overline{W}/\overline{\varphi }`$. A straightforward calculation (both ‘by hand’ and by the use of a Mathematica computer program) in the $`CP^1`$ case yields the following deformed NLSM kinetic terms:
$$L_{\mathrm{kin}.}=g_{\varphi \varphi }_\mu \varphi _\mu \varphi 2g_{\varphi \overline{\varphi }}_\mu \varphi _\mu \overline{\varphi }g_{\overline{\varphi }\overline{\varphi }}_\mu \overline{\varphi }_\mu \overline{\varphi },$$
$`(4.14)`$
where
| $`g_{\varphi \overline{\varphi }}=`$ | $`{\displaystyle \frac{\alpha \kappa ^2c^2\overline{\varphi }^2(\overline{W}_{,\overline{\varphi }})^2}{\left(\alpha +\sqrt{\alpha ^2+\left(2c\overline{\varphi }(1+\kappa ^2\varphi \overline{\varphi })\overline{W}_{,\overline{\varphi }}\right)^2}\right)\sqrt{\alpha ^2+\left(2c\overline{\varphi }(1+\kappa ^2\varphi \overline{\varphi })\overline{W}_{,\overline{\varphi }}\right)^2}}},`$ |
| --- | --- |
| $`g_{\varphi \varphi }=`$ | $`0,`$ |
| $`g_{\overline{\varphi }\overline{\varphi }}=`$ | $`{\displaystyle \frac{2\alpha ^1c^2(1+\kappa ^2\varphi \overline{\varphi })\overline{W}_{,\overline{\varphi }}}{\left(\alpha \sqrt{\alpha ^2+\left(2c\overline{\varphi }(1+\kappa ^2\varphi \overline{\varphi })\overline{W}_{,\overline{\varphi }}\right)^2}\right)\sqrt{\alpha ^2+\left(2c\overline{\varphi }(1+\kappa ^2\varphi \overline{\varphi })\overline{W}_{,\overline{\varphi }}\right)^2}}}\times `$ |
| | $`\times [4c^2\overline{\varphi }^2(\overline{W}_{,\overline{\varphi }})^3(1+\kappa ^2\varphi \overline{\varphi })`$ |
| | $`+\alpha (\alpha \sqrt{\alpha ^2+\left(2c\overline{\varphi }(1+\kappa ^2\varphi \overline{\varphi })\overline{W}_{,\overline{\varphi }}\right)^2})(2\overline{W}_{,\overline{\varphi }}+\overline{\varphi }\overline{W}_{,\overline{\varphi }\overline{\varphi }})].`$ |
$`(4.15)`$
It is worth noticing that $`detg=(g_{\varphi \overline{\varphi }})^2`$. The most apparent feature $`g_{\varphi \varphi }=0`$ is also valid in the case of a generic NAC-deformed NLSM (in a given parametrization).
## 5 Conclusion
As is clear from our discussion and explicit examples, the problem of solving for the auxiliary fields is highly non-trivial in the context of NAC-deformed N=1/2 supersymmetric NLSM under consideration. Special care should be exercised when considering this problem with the smearing effects that should be calculated first.
One should also distinguish between a NAC-deformation and N=1/2 supersymmetry. Though the NAC-deformation we considered is N=1/2 supersymmetric, the former is stronger than the latter. Requiring merely N=1/2 supersymmetry of a four-dimensional NLSM would give rise to much weaker restrictions on the NLSM target space, as is pretty obvious from eqs. (1.10) and (1.11).
Some comments are in order about the relation between our results in this paper and those obtained in ref. for the N=1/2 supersymmetric NAC-deformed $`CP^n`$ NLSM, in the absence of chiral and anti-chiral scalar superpotentials.
As is well known in the theory of NLSM (see e.g., ref. ), the so-called quotient construction (or gauging isometries of the ‘flat’ NLSM target space) can be used to represent some NLSM with homogeneous target spaces as the gauge theories. It was used in ref. to construct the NAC-deformed N=1/2 supersymmetric NLSM in four dimensions with the $`CP^n`$ target space, by combining the quotient construction and the results of ref. about the NAC-deformed supersymmetric gauge theories. Unlike the NAC-deformation of the non-gauge supersymmetric field theory (1.12) governed by the scalar deformation parameter $`c`$, the NAC-deformed N=1/2 supersymmetric field theory has $`C_{\alpha \beta }`$ as the deformation parameters, while its action is not ‘Lorentz’-invariant. In addition, there are more auxiliary fields in the supersymmetric quotient construction, while it does not lead to any splitting or smearing of the Kähler potential. As was demonstrated in ref. , the quotient construction of the NAC-deformed N=1/2 supersymmetric $`CP^n`$ model does not modify the Kähler geometry at all, while the only new term in the deformed Lagrangian takes the form
$$g_{p\overline{q}}g_{r\overline{s}}C^{\alpha \beta }(\sigma ^{\mu \nu })_\beta {}_{}{}^{\gamma }\chi _{\alpha }^{p}\chi _\gamma ^r_\mu \overline{\varphi }^{\overline{q}}_\nu \overline{\varphi }^{\overline{s}},$$
$`(5.1)`$
where $`g_{p\overline{q}}(\varphi ,\overline{\varphi })`$ is the Fubini-Study metric associated with the Kähler potential (4.1). As is clear from eq. (5.1), this deformation is linear in $`C^{\alpha \beta }`$, while it is not ‘Lorentz’-invariant. The quotient construction itself is also limited to the homogeneous NLSM, and it does not allow a scalar superpotential.
Our approach to the NAC-deformed NLSM is very general, while the NAC deformation of a Kähler potential is controlled by the auxiliary fields $`M^i`$ entering the deformed Kähler potential in the highly non-linear way. In the absence of an anti-chiral superpotential, when requiring the smooth undeformed limit, the bosonic part of the auxiliary field solution vanishes, so that we are left with the undeformed bosonic Kähler potential as well. However, when the auxiliary fields are eliminated, the structure of the fermionic terms in our approach is apparently not the same as that of ref. , even modulo possible field redefinitions.
At the same time, the N=1/2 supersymmetry transformation laws are not modified by a NAC-deformation, both in our approach and in ref. . Our conclusion is that the two methods give rise to the inequivalent N=1/2 supersymmetric extensions of the Kähler NLSM. It does not seem to be very surprising because there are many N=1/2 supersymmetric extensions of the NLSM kinetic terms.
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# 1 Introduction
## 1 Introduction
### 1.1 Notes about unconventional superconductivity
The phenomenon of unconventional Cooper pairing of fermions, i.e., the formation of Cooper pairs with nonzero angular momentum was theoretically predicted in 1959 as a mechanism of superfluidity in Fermi liquids. In 1972 the same phenomenon - unconventional superfluidity due to a $`p`$-wave (spin triplet) Cooper pairing of <sup>3</sup>He atoms, was experimentally discovered in the mK range of temperatures; for details and theoretical description see Refs. . Note that, in contrast to the standard $`s`$-wave pairing in usual (conventional) superconductors where the electron pairs are formed by an attractive electron-electron interaction due to a virtual phonon exchange, the widely accepted mechanism of the Cooper pairing in superfluid <sup>3</sup>He is based on an attractive interaction between the fermions (<sup>3</sup>He atoms) as a result of a virtual exchange of spin fluctuations. Certain spin fluctuation mechanisms of unconventional Cooper pairing of electrons are assumed also for the discovered in 1979 heavy fermion superconductors (see, e.g., Refs. ) as well as for some classes of high-temperature superconductors (see, e.g., Refs. ).
The possible superconducting phases in unconventional superconductors are described in the framework of the general Ginzburg-Landau (GL) effective free energy functional with the help of the symmetry groups theory. Thus a variety of possible superconducting orderings were predicted for different crystal structures . A detailed thermodynamic analysis of the homogeneous (Meissner) phases and a renormalization group investigation of the superconducting phase transition up to the two-loop approximation have been also performed (for a three-loop renormalization group analysis, see Ref. ; for effects of magnetic fluctuations and disorder, see ). We shall essentially use these results in our present consideration.
### 1.2 Experimental predictions
In 2000, experiments at low temperatures ($`T1`$ K) and high pressure ($`P1`$ GPa) demonstrated the existence of spin triplet superconducting states in the metallic compound UGe<sub>2</sub>. The superconductivity is triggered by the spontaneous magnetization of the ferromagnetic phase that occurs at much higher temperatures. It coexists with the superconducting phase in the whole domain of its existence below $`T1`$ K; see also experiments from Refs. , and the discussion in Ref. . The same phenomenon of existence of superconductivity at low temperatures and high pressure in the domain of the $`(T,P)`$ phase diagram where the ferromagnetic order is present was observed in other ferromagnetic metallic compounds (ZrZn<sub>2</sub> and URhGe ) soon after the discovery of superconductivity in UGe<sub>2</sub>.
In superconducting ternary and Chevrel compounds the influence of magnetic order on superconductivity is also substantial (see, e.g., ) but in the newly found ferromagnetic substances the phase transition temperature ($`T_f`$) to the ferromagnetic state is much higher than the phase transition temperature ($`T_{FS})`$ from ferromagnetic to a mixed state of coexistence of ferromagnetism and superconductivity. For example, in UGe<sub>2</sub>, $`T_{FS}=0.8`$ K while the critical temperature of the phase transition from paramagnetic to ferromagnetic state in the same material is $`T_f=35`$. It can be assumed that in these substances the material parameter $`T_s`$ defined as the usual critical temperature of the second order phase transition from normal to uniform (Meissner) supercondicting state in a zero external magnetic field is much lower than the phase transition temperature $`T_{FS}`$. The above mentioned experiments on the compounds UGe<sub>2</sub>, URhGe, and ZrZn<sub>2</sub> do not give any evidence for the existence of a standard normal-to-superconducting phase transition in a zero external magnetic field.
It seems that the superconductivity in the metallic compounds mentioned above always coexists with the ferromagnetic order and is enhanced by it. In these systems, as claimed in Ref. , the superconductivity probably arises from the same electrons that create the band magnetism and can be most naturally understood rather as a triplet than spin-singlet pairing phenomenon. Metallic compounds UGe<sub>2</sub>, URhGe, and ZrZn<sub>2</sub>, are itinerant ferromagnets. An unconventional superconductivity is also suggested as a possible outcome of recent experiments in Fe , in which a superconducting phase has been discovered at temperatures below $`2`$ K and pressures between 15 and 30 GPa. There both vortex and Meissner superconductivity phases are found in the high-pressure crystal modification of Fe with a hexagonal close-packed lattice for which the strong ferromagnetism of the usual bcc iron crystal probably disappears . It can be hardly claimed that in hexagonal Fe the ferromagnetism and superconductivity coexist but the clear evidence for a superconductivity is also a remarkable achievement.
The reasonable question whether these examples of superconductivity and coexistence of superconductivity and ferromagnetism are bulk or surface effects can be stated. The earlier experiments performed before 2004 do not answer this question. Recent experiments show that surface superconductivity appears in ZrZn<sub>2</sub> and its presence depends essentially on the way of preparation of the sample. But in our study it is important that bulk superconductivity can be considered well established in this substance.
### 1.3 Ferromagnetism versus superconductivity
The important point in all discussions of the interplay of superconductivity and ferromagnetism is that a small amount of magnetic impurities can destroy superconductivity in conventional ($`s`$-wave) superconductors by breaking up the ($`s`$-wave) electron pairs with opposite spins (paramagnetic impurity effect ). In this aspect the phenomenological arguments and the conclusions on the basis of the microscopic theory of magnetic impurities in $`s`$-wave superconductors are in a complete agreement with each other; see, e.g., Refs. . In fact, a total suppression of conventional ($`s`$-wave) superconductivity should occur in the presence of an uniform spontaneous magnetization $`𝑴`$, i.e., in a standard ferromagnetic phase . The physical reason for this suppression is the same as in the case of magnetic impurities, namely, the opposite electron spins in the $`s`$-wave Cooper pair turn over along the vector $`𝑴`$ in order to lower their Zeeman energy and, hence, the pairs break down. Therefore, the ferromagnetic order can hardly coexist with conventional superconducting states. Especially, this is valid for the coexistence of uniform superconducting and ferromagnetic states where the superconducting order parameter $`\psi (𝒙)`$ and the magnetization $`𝑴`$ do not depend on the spatial vector $`𝒙`$.
But yet a coexistence of $`s`$-wave superconductivity and ferromagnetism may appear in uncommon materials and under quite special circumstances. Furthermore, let us emphasize that the conditions for the coexistence of nonuniform (“vertex”, “spiral”, “spin-sinosoidal” or “helical” ) superconducting and ferromagnetic states are less restrictive than those for the coexistence of uniform superconducting and ferromagnetic orders. Coexistence of nonuniform phases has been discussed in details, both experiment and theory, in ternary and Chevrel-phase compounds where such a coexistence seems quite likely; for a comprehensive review, see, for example, Refs. .
In fact the only two superconducting systems for which the experimental data allow assumptions in a favor of a coexistence of superconductivity and ferromagnetism are the rare earth ternary boride compound ErRh<sub>4</sub>B<sub>4</sub> and the Chervel phase compound HoMo<sub>6</sub>S<sub>8</sub>; for a more extended review, see Refs. . In these compounds the phase of coexistence appears in a very narrow temperature region just below the Curie temperature $`T_f`$ of the ferromagnetic phase transition. At lower temperatures the magnetic moments of the rare earth 4$`f`$ electrons become better aligned, the magnetization increases and the $`s`$-wave superconductivity pairs formed by the conduction electrons disintegrate.
### 1.4 Unconventional superconductivity triggered by ferromagnetic order
We shall not consider all important aspects of the long standing problem of coexistence of superconductivity and ferromagnetism rather we shall concentrate our attention on the description of the newly discovered coexistence of ferromagnetism and unconventional (spin-triplet) superconductivity in the itinerant ferromagnets UGe<sub>2</sub>, ZrZn<sub>2</sub>, and URhGe. Here we wish to emphasize that the main object of our discussion is the superconductivity of these compounds and at a second place in the rate of importance we put the problem of coexistence. The reason is that the existence of superconductivity in itinerant ferromagnets is a highly nontrivial phenomenon. As noted in Ref. the superconductivity in these materials appears to be difficult to explain in terms of previous theories and requires new concepts to interpret the experimental data.
We have already mentioned that in ternary compounds the ferromagtetism comes from the localized 4$`f`$ electrons while the s-wave Cooper pairs are formed by conduction electrons. In UGe<sub>2</sub> and URhGe the 5$`f`$ electrons of U atoms form both superconductivity and ferromagnetic order . In ZrZn<sub>2</sub> the same double role is played by the 4$`d`$ electrons of Zr. Therefore, the task is to describe this behavior of the band electrons at a microscopic level. One may speculate about a spin-fluctuation mediated unconventional Cooper pairing as is in case of <sup>3</sup>He and heavy fermion superconductors. These important issues have not yet a reliable answer and for this reason we shall confine our consideration to a phenomenological level.
In fact, a number of reliable experimental data as the coherence length and the superconducting gap measurements are in favor of the conclusion about a spin-triplet Cooper pairing in these metallic compounds, although the mechanism of pairing remains unclear. We shall essentially use this reliable conclusion. This point of view is consistent with the experimental observation of coexistence of superconductivity only in a low temperature part of the ferromagnetic domain of the phase diagram ($`T,P`$), which means that a pure (non-ferromagnetic) superconducting phase is not observed, a circumstance, that is also in favor of the assumption of a spin-triplet superconductivity. Our investigation leads to results which confirm this general picture.
On the basis of the experimental data and conclusions presented for the first time in Refs. and shortly afterwards confirmed in Refs. one may reliably accept that the superconductivity in these magnetic compounds is considerably enhanced by the ferromagnetic order parameter $`𝑴`$ and, perhaps, it could not exist without this “mechanism of ferromagnetic trigger,” or, in short, “$`𝑴`$-trigger”; see Refs. where this concept has been introduced for the first time. The trigger phenomenon is possible for spin-triplet Cooper pairs where the electron spins point parallel to each other and their turn along the vector of the spontaneous magnetization $`𝑴`$ does not produce a breakdown of the spin-triplet Cooper pairs but rather stabilizes them and, perhaps, stimulates their creation. We shall describe this phenomenon at a phenomenological level.
### 1.5 Phenomenological studies
A phenomenological theory that explains the coexistence of ferromagnetism and unconventional spin-triplet superconductivity of GL type has been developed recently in where possible low-order couplings between the superconducting and ferromagnetic order parameters are derived with the help of general symmetry group arguments. On this basis several important features of the superconducting vortex state of unconventional ferromagnetic superconductors were demonstrated .
In our review we shall follow the approach from Refs. to investigate the conditions for the occurrence of the Meissner phase and to demonstrate that the presence of ferromagnetic order enhances the $`p`$-wave superconductivity. We also establish the phase diagram of ferromagnetic superconductors in a zero external magnetic field and show that the phase transition to the superconducting state can be either of first or second order depending on the particular substance. We confirm the predictions made in Refs. about the symmetry of the ordered phases.
In our study we use the mean-field approximation and known results about the possible phases in nonmagnetic superconductors with triplet ($`p`$-wave) pairing . Our results show that taking into account the anisotropy of the spin-triplet Cooper pairs modifies but does not drastically change the thermodynamic properties of the coexistence phase, especially in the temperature domain above the superconducting critical temperature $`T_s`$. The effect of crystal anisotropy is similar but we shall not make an overall thermodynamic analysis of this problem because we have to consider concrete systems and crystal structures for which there is no enough information from experiment to make conclusions about the parameters of the theory. Our results confirm the general concept that the anisotropy reduces the degree of ground state degeneration, and depending on the symmetry of the crystal, picks up a crystal direction for the ordering.
There exists a formal similarity between the phase diagram we obtain and the phase diagram of certain improper ferroelectrics . We shall make use of the concept in the theory of improper ferroelectrics, where the trigger of the primary order parameter by a secondary order parameter (the electric polarization) has been initially introduced and exploited; see Ref. . The mechanism of the M-triggered superconductivity in itinerant ferromagnets is formally identical to the mechanism of appearance of structural order triggered by the electric polarization in improper ferroelectrics (see, e.g., Refs. ).
Our investigation is based on the GL free energy functional of unconventional ferromagnetic superconductors presented in Sec. 2.1 and we shall establish the uniform phases which are described by it. More information about the justification of this investigation is presented in Sec. 2.2. We work with a quite general GL free energy and the problem is that there is no enough information about the values of the parameters of the model for concrete compounds (UGe<sub>2</sub>, URhGe, ZrZn<sub>2</sub>) where the ferromagnetic superconductivity has been discovered. On the one hand the lack of information makes impossible a detailed comparison of the theory to the available experimental data but on the other hand our results are not bound to one or more concrete substances and can be applied to any unconventional ferromagnetic superconductor. In Sec. 3 we discuss the phases in nonmagnetic unconventional superconductors. In Sec. 4 the M-trigger effect will be described when only a linear coupling of the magnetization $`𝑴`$ to the superconducting order parameter $`\psi `$ is considered; here the spatial dependence of order parameters and all anisotropy effects are ignored. In Sec. 5 we analyze the influence of quadratic coupling of magnetization to the superconducting order parameter on the thermodynamics of the ferromagnetic superconductors. The application of our results to experimental $`(T,P)`$ phase diagrams is discussed in Sec. 5.3. In Sec. 6 the anisotropy effects are outlined. Our main attention is focussed on the Cooper-pair anisotropy. Note, that certain types of crystal anisotropy may produce more than one ferromagnetic phase but here we shall not dwell on this interesting topic. In Sec. 7 we summarize our conclusions.
## 2 Ginzburg-Landau free energy
Following Refs. in this Chapter we discuss the phenomenological theory of spin-triplet ferromagnetic superconductors and justify our consideration in Sections 3–6.
### 2.1 Model
The general GL free energy functional, we shall use in our analysis, is
$$F[\psi ,𝑴]=d^3xf(\psi ,𝑴),$$
(1)
where the free energy density $`f(\psi ,𝑴)`$ ( hereafter called “free energy”) of a spin-triplet ferromagnetic superconductor is a sum of five terms , namely,
$$f(\psi ,𝑴)=f_\text{S}(\psi )+f_\text{F}^{}(𝑴)+f_\text{I}(\psi ,𝑴)+\frac{𝑩^2}{8\pi }𝑩\mathbf{.}𝑴.$$
(2)
In Eq. (2) the three dimensional complex vector $`\psi =\{\psi _j;j=1,2,3\}`$ represents the superconducting order parameter, $`𝑩=(𝑯+4\pi 𝑴)=\times 𝑨`$ is the magnetic induction; $`𝑯`$ is the external magnetic field, $`𝑨=\{A_j;j=1,2,3\}`$ is the magnetic vector potential. The last two terms on r.h.s. of Eq. (2) are related with the magnetic energy which includes both diamagnetic and paramagnetic effects in the superconductor; see, e.g., .
The energy part $`f_\text{S}(\psi )`$ in Eq. (2) describes the superconductivity for $`𝑯=𝑴0`$. It can be written in the form
$$f_\text{S}(\psi )=f_{grad}(\psi )+a_s|\psi |^2+\frac{b_s}{2}|\psi |^4+\frac{u_s}{2}|\psi ^2|^2+\frac{v_s}{2}\underset{j=1}{\overset{3}{}}|\psi _j|^4.$$
(3)
Here
$`f_{grad}(\psi )`$ $`=`$ $`K_1(D_i\psi _j)^{}(D_i\psi _j)+K_2[(D_i\psi _i)^{}(D_j\psi _j)`$
$`+(D_i\psi _j)^{}(D_j\psi _i)]+K_3(D_i\psi _i)^{}(D_i\psi _i),`$
where a summation over the indices $`i,j=1,2,3`$ is assumed and the symbol
$$D_j=i\mathrm{}\frac{}{x_i}+\frac{2|e|}{c}A_j$$
(5)
denotes a covariant differentiation. In Eq. (3), $`b_s>0`$ and $`a_s=\alpha _s(TT_s)`$, where $`\alpha _s`$ is a positive material parameter and $`T_s`$ is the critical temperature of the standard second order phase transition which may occur at $`H==0`$; $`H=|𝑯|`$, and $`=|𝑴|`$. The quantities $`u_s`$ and $`v_s`$ describe the anisotropy of the spin-triplet Cooper pair and the crystal anisotropy, respectively, . Parameters $`K_j`$, $`(j=1,2,3)`$ in Eq. (4) are related with the effective mass tensor of anisotropic Cooper pairs .
The superconducting part (3) of the free energy $`f(\psi ,M)`$ is derived from symmetry group arguments and is independent of particular microscopic models; see, e.g., Refs. . According to classifications the $`p`$-wave superconductivity in the cubic point group $`O_h`$ can be realized through one-, two-, and three-dimensional representations of the order parameter. The expressions (3) and (5) incorporate all three possible cases. The coefficients $`b_s`$, $`u_s`$, and $`v_s`$ in Eq. (3) are different for weak and strong spin-orbit couplings but in our investigation they are considered as undetermined material parameters which depend on the particular substance.
The free energy of a standard isotropic ferromagnet is given by the term $`f_\text{F}^{}(𝑴)`$ in Eq. (2),
$$f_\text{F}^{}(𝑴)=c_f\underset{j=1}{\overset{3}{}}|_j𝑴_j|^2+a_f(T_f^{})𝑴^2+\frac{b_f}{2}𝑴^4,$$
(6)
where $`_j=/x_j`$ and $`b_f>0`$. The quantity $`a_f(T_f^{})=\alpha _f(TT_f^{})`$ is expressed by the material parameter $`\alpha _f>0`$ and the temperature $`T_f^{}`$ which is different from the critical temperature $`T_f`$ of the ferromagnet and this point will be discussed below. We have already added a negative term ($`2\pi ^2`$) to the total free energy $`f(\psi ,𝑴)`$ and that is obvious by setting $`H=0`$ in Eq. (2). The negative energy ($`2\pi ^2`$) should be added to $`f_\text{F}^{}(𝑴)`$. In this way one obtains the total free energy $`f_\text{F}(𝑴)`$ of the ferromagnet in a zero external magnetic field that is given by a modification of Eq. (6) according to the rule
$$f_\text{F}(a_f)=f_\text{F}^{}\left[a_f(T_f^{})a_f(T_f)\right],$$
(7)
where $`a_f=\alpha _f(TT_f)`$ and
$$T_f=T_f^{}+\frac{2\pi }{\alpha _f}$$
(8)
is the critical temperature of a standard ferromagnetic phase transition of second order. This scheme was used in studies of rare earth ternary compounds . Alternatively , one may use from the beginning the total ferromagnetic free energy $`f_\text{F}(a_f,𝑴)`$ as given by Eqs. (6) - (8) but in this case the magnetic energy included in the last two terms on r.h.s. of Eq. (2) should be replaced with $`H^2/8\pi `$. Both approaches are equivalent.
The interaction between the ferromagnetic order parameter $`𝑴`$ and the superconducting order parameter $`\psi `$ is given by
$$f_\text{I}(\psi ,𝑴)=i\gamma _0𝑴.(\psi \times \psi ^{})+\delta 𝑴^2|\psi |^2.$$
(9)
The $`\gamma _0`$-term in the above expression is the most substantial for the description of experimentally found ferromagnetic superconductors and the $`\delta 𝑴^2|\psi |^2`$–term makes the model more realistic in the strong coupling limit as it gives the opportunity to enlarge the phase diagram including both positive and negative values of the parameter $`a_s`$. In this way the domain of the stable ferromagnetic order is extended down to zero temperatures for a wide range of values of material parameters and the pressure $`P`$, a situation that corresponds to the experiments in ferromagnetic superconductors.
In Eq. (9) the coupling constant $`\gamma _0>0`$ can be represented in the form $`\gamma _0=4\pi J`$, where $`J>0`$ is the ferromagnetic exchange parameter . In general, the parameter $`\delta `$ for ferromagnetic superconductors may take both positive and negative values. The values of the material parameters ($`T_s`$, $`T_f`$, $`\alpha _s`$, $`\alpha _f`$, $`b_s`$, $`u_s`$, $`v_s`$, $`b_f`$, $`K_j`$, $`\gamma _0`$ and $`\delta `$) depend on the choice of the concrete substance and on thermodynamic parameters as temperature $`T`$ and pressure $`P`$.
### 2.2 Way of treatment
The total free energy (2) is a quite complex object of theoretical investigation. The possible vortex and uniform phases of the model cannot be investigated within a single calculation, rather one should focus on concrete problems. In Ref. the vortex phase was discussed with the help of the criterion for a stability of this state near the phase transition line $`T_{c2}(H)`$; see also, Ref. . In case of $`H=0`$ one should apply the same criterion with respect to the magnetization $``$ for small values of $`|\psi |`$ near the phase transition line $`T_{c2}()`$ as performed in Ref. .
Here we shall be interested in the uniform phases when the order parameters $`\psi `$ and $`𝑴`$ do not depend on the spatial vector $`𝒙V`$, where $`V`$ is the volume of the superconductor. We shall restrict our analysis to the consideration of the coexistence of uniform (Meissner) phases and ferromagnetic order. We shall make a detailed investigation in order to show that the main properties of the uniform phases can be well determined within an approximation when the crystal anisotropy is neglected. Even, some of the main features of the uniform phases in unconventional ferromagnetic superconductors can be reliably outlined when the Cooper pair anisotropy is neglected, too.
The assumption of a uniform magnetization $`𝑴`$ is always reliable outside a quite close vicinity of the magnetic phase transition and under the condition that the superconducting order parameter $`\psi `$ is also uniform, i.e. that vortex phases are not present at the respective temperature domain. This conditions are directly satisfied in type I superconductors but in type II superconductors the temperature should be sufficiently low and the external magnetic field should be zero. Nevertheless, the mentioned conditions for type II superconductors may turn insufficient for the appearance of uniform superconducting states in materials with quite high values of the spontaneous magnetization. In such situation the uniform (Meissner) superconductivity and, hence, the coexistence of this superconductivity with uniform ferromagnetic order may not appear even at zero temperature. Up to now type I unconventional ferromagnetic superconductors have not been found whereas the experimental data for the recently discovered compounds UGe<sub>2</sub>, URhGe, and ZrZn<sub>2</sub> are not enough to conclude definitely either about the lack or the existence of uniform superconducting states at low and ultra-low temperatures.
If real materials can be modelled by the general GL free energy (1) - (9), the ground state properties will be described by uniform states, which we shall investigate. The problem about the availability of such states in real materials at finite temperatures is quite subtle at the present stage of research when the experimental data are not enough. Recently in Ref. an experimental evidence was given for the coexistence of uniform superconductivity and ferromagnetism in UGe<sub>2</sub>. Thus we shall assume that uniform phases may exist in some unconventional ferromagnetic superconductors. Moreover, we have to emphasize that these phases appear as solutions of the GL equations corresponding to the free energy (1) - (9). These arguments completely justify our study.
In case of a strong easy axis type of magnetic anisotropy, as is in UGe<sub>2</sub> , the overall complexity of mean-field analysis of the free energy $`f(\psi ,𝑴)`$ can be avoided by performing an “Ising-like” description: $`𝑴=(0,0,)`$, where $`=\pm |𝑴|`$ is the magnetization along the “$`z`$-axis.” Because of the equivalence of the “up” and “down” physical states $`(\pm 𝑴)`$ the thermodynamic analysis can be performed within the “gauge” $`0`$. When the magnetic order has a continuous symmetry we can take advantage of the symmetry of the total free energy $`f(\psi ,𝑴)`$ and avoid the consideration of equivalent thermodynamic states that occur as a result of the respective symmetry breaking at the phase transition point but have no effect on thermodynamics of the system. In the isotropic system one may again choose a gauge, in which the magnetization vector has the same direction as $`z`$-axis ($`|𝑴|=M_z=`$) and this will not influence the generality of thermodynamic analysis. Here we shall prefer the alternative description within which the ferromagnetic state may appear through two equivalent up and down domains with magnetizations $``$ and ($``$), respectively.
We shall make the mean-field analysis of the uniform phases and the possible phase transitions between such phases in a zero external magnetic field ($`𝑯=0)`$, when the crystal anisotropy is neglected ($`v_s0`$). The only exception will be the consideration in Sec. 3, where we briefly discuss the nonmagnetic superconductors ($`0`$).
We shall use notations in which the number of parameters is reduced. For this reason we introduce
$$b=(b_s+u_s+v_s)$$
(10)
and redefine the order parameters and the other quantities in the following way:
$`\phi _j=b^{1/4}\psi _j=\varphi _je^{i\theta _j},M=b_f^{1/4},`$ (11)
$`r={\displaystyle \frac{a_s}{\sqrt{b}}},t={\displaystyle \frac{a_f}{\sqrt{b_f}}},w={\displaystyle \frac{u_s}{b}},v={\displaystyle \frac{v_s}{b}},`$
$`\gamma ={\displaystyle \frac{\gamma _0}{b^{1/2}b_f^{1/4}}},\gamma _1={\displaystyle \frac{\delta }{(bb_f)^{1/2}}}.`$
Having in mind that the order parameters $`\psi `$ and $`𝑴`$ are considered uniform and using Eqs. (10) and (11), we can write the free energy density $`f(\psi ,M)=F(\psi ,M)/V`$ in the form
$`f(\psi ,M)`$ $`=`$ $`r\varphi ^2+{\displaystyle \frac{1}{2}}\varphi ^4+2\gamma \varphi _1\varphi _2M\text{sin}(\theta _2\theta _1)`$
$`+\gamma _1\varphi ^2M^2+tM^2+{\displaystyle \frac{1}{2}}M^4`$
$`2w[\varphi _1^2\varphi _2^2\text{sin}^2(\theta _2\theta _1)+\varphi _1^2\varphi _3^2\text{sin}^2(\theta _1\theta _3)`$
$`+\varphi _2^2\varphi _3^2\text{sin}^2(\theta _2\theta _3)]`$
$`v[\varphi _1^2\varphi _2^2+\varphi _1^2\varphi _3^2+\varphi _2^2\varphi _3^2].`$
Note, that in the above expression the order parameters $`\psi `$ and $`𝑴`$ are defined per unit volume.
The equilibrium phases are obtained from the equations of state
$$\frac{f(\mu _0)}{\mu _\alpha }=0,$$
(13)
where the series of symbols $`\mu `$ can be defined as, for example, $`\mu =\left\{\mu _\alpha \right\}=(M,\varphi _1,\mathrm{},\varphi _3,`$ $`\theta _1,\mathrm{},\theta _3)`$; $`\mu _0`$ denotes an equilibrium phase. The stability matrix $`\stackrel{~}{F}`$ of the phases $`\mu _0`$ is defined by
$$\widehat{F}(\mu _0)=\left\{F_{\alpha \beta }(\mu _0)\right\}=\frac{^2f(\mu _0)}{\mu _\alpha \mu _\beta }.$$
(14)
An alternative treatment can be done in terms of real ($`\psi _j^{}`$) and imaginary ($`\psi _j^{\prime \prime }`$) parts of the complex numbers $`\psi _j=\psi _j^{}+i\psi _j^{\prime \prime }`$. The calculation with moduli $`\varphi _j`$ and phase angles $`\theta _j`$ of $`\psi _j`$ is more simple but in cases of strongly degenerate phases some of the angles $`\theta _j`$ remain unspecified. Then an alternative analysis with the help of the components $`\psi _j^{}`$ and $`\psi _j^{\prime \prime }`$ should be done.
The thermodynamic stability of the phases that are solutions of Eqs. (13) is checked with the help of the matrix (14). An additional stability analysis is done by the comparison of free energies of phases that satisfy (13) and render the stability matrix (14) positive in one and the same domain of parameters $`\{r,t,\gamma ,\gamma _1,w,v\}`$. This step is important because the complicated form of the free energy generates a great number of solutions of Eqs. (13) and we have to sift out the stable from metastable phases that correspond either to global or local minima of the free energy, respectively .
Some solutions of Eqs. (13) have a marginal stability, i.e., their stability matrix (14) is neither positively nor negatively definite. This is often a result of the degeneration of phases with broken continuous symmetry. If the reason for the lack of a clear positive definiteness of the stability matrix is precisely the mentioned degeneration of the ground state, one may reliably conclude that the respective phase is stable. If there is another reason, the analysis of the matrix (14) will be insufficient to determine the respective stability property. These cases are quite rare and occur for particular values of the parameters $`\{r,t,\gamma ,\mathrm{}\}`$.
## 3 Pure superconductivity
Let us set $`M0`$ in Eq. (12) and briefly summarize the known results for the “pure superconducting case” when the magnetic order cannot appear and magnetic effects do not affect the stability of the uniform (Meissner) superconducting phases. The possible phases can be classified by the structure of the complex vector order parameter $`\psi =(\psi _1,\psi _2,\psi _2)`$. We shall often use the moduli vector $`(\varphi _1,\varphi _2,\varphi _3)`$ with magnitude $`\varphi =(\varphi _1^2+\varphi _2^2+\varphi _3^2)^{1/2}`$ and the phase angles $`\theta _j`$.
The normal phase (0,0,0) is always a solution of the Eqs. (13). It is stable for $`r0`$, and corresponds to a free energy $`f=0`$. Under certain conditions, six ordered phases occur for $`r<0`$. Here we shall not repeat the detailed description of these phases but we shall briefly mention their structure.
The simplest ordered phase is of type $`(\psi _1,0,0)`$ with equivalent domains: $`(0,\psi _2,0)`$ and $`(0,0,\psi _3)`$. Multi- domain phases of more complex structure also occur, but we shall not always enumerate the possible domains. For example, the “two-dimensional” phases can be fully represented by domains of type $`(\psi _1,\psi _2,0)`$ but there are also other two types of domains: $`(\psi _1,0,\psi _3)`$ and $`(0,\psi _2,\psi _3)`$. As we consider the general case when the crystal anisotropy is present $`(v0)`$, this type of phases possesses the property $`|\psi _i|=|\psi _j|`$.
The two-dimensional phases are two and have different free energies. To clarify this point let us consider, for example, the phase $`(\psi _1,\psi _2,0)`$. The two complex numbers, $`\psi _1`$ and $`\psi _2`$ can be represented either as two-component real vectors, or, equivalently, as rotating vectors in the complex plane. One can easily show that Eq. (12) yields two phases: a collinear phase, when $`(\theta _2\theta _1)=\pi k(k=0,\pm 1,\mathrm{})`$, i.e. when the vectors $`\psi _1`$ and $`\psi _2`$ are collinear, and another (non-collinear) phase when the same vectors are perpendicular to each other: $`(\theta _2\theta _1)=\pi (k+1/2)`$. Having in mind that $`|\varphi _1|=|\varphi _2|=\varphi /\sqrt{2}`$, the domain $`(\psi _1,\psi _2,0)`$ of the collinear phase is given by $`(\pm 1,1,0)\varphi /\sqrt{2}`$ whereas the same domain for the non-collinear phase will be $`(\pm i,1,0)\varphi /\sqrt{2}`$. The two domains of these phases have similar representations.
In addition to the mentioned three ordered phases, three other ordered phases exist. For these phases all three components $`\psi _j`$ have nonzero equilibrium values. Two of them have equal to one another moduli $`\varphi _j`$, i.e., $`\varphi _1=\varphi _2=\varphi _3`$. The third phase is of the type $`\varphi _1=\varphi _2\varphi _3`$ and is unstable so it cannot occur in real systems. The two three-dimensional phases with equal moduli of the order parameter components have different phase angles and, hence, different structure. The difference between any couple of angles $`\theta _j`$ is given by $`\pm \pi /3`$ or $`\pm 2\pi /3`$. The characteristic vectors of this phase can be of the form $`(e^{i\pi /3},e^{i\pi /3},1)\varphi /\sqrt{3}`$ and $`(e^{2i\pi /3},e^{i2\pi /3},1)\varphi /\sqrt{3}`$. The second stable three dimensional phase is “real”, i.e. the components $`\psi _j`$ lie on the real axis; $`(\theta _j\theta _j)=\pi k`$ for any couple of angles $`\theta _j`$ and the characteristic vectors are $`(\pm 1,\pm 1,1)\varphi /\sqrt{3}`$. The stability properties of these five stable ordered phases were presented in details in Refs. .
When the crystal anisotropy is not present ($`v=0`$) the picture changes. The increase of the level of degeneracy of the ordered states leads to an instability of some phases and to a lack of some noncollinear phases. Both two- and three-dimensional real phases, where $`(\theta _j\theta _j)=\pi k`$, are no more constrained by the condition $`\varphi _i=\varphi _j`$ but rather have the freedom of a variation of the moduli $`\varphi _j`$ under the condition $`\varphi ^2=r>0`$. The two-dimensional noncollinear phase exists but has a marginal stability . All other noncollinear phases even in the presence of a crystal anisotropy $`(v0)`$ either vanish or are unstable; for details, see Ref. . This discussion demonstrates that the crystal anisotropy stabilizes the ordering along the main crystallographic directions, lowers the level of degeneracy of the ordered state related with the spontaneous breaking of the continuous symmetry and favors the appearance of noncollinear phases.
The crystal field effects related to the unconventional superconducting order were established for the first time in Ref. . In our consideration of unconventional ferromagnetic superconductors in Sec. 4–7 we shall take advantage of these effects of the crystal anisotropy. In both cases $`v=0`$ and $`v0`$ the matrix (14) indicates an instability of three-dimensional phases (all $`\varphi _j0)`$ with an arbitrary ratios $`\varphi _i/\varphi _j`$. As already mentioned, for $`v0`$ the phases of type $`\varphi _1=\varphi _2\varphi _3`$ are also unstable whereas for $`v=0`$, even the phase $`\varphi _1=\varphi _2=\varphi _3>0`$ is unstable.
## 4 M-triggered superconductivity
We shall consider the Walker-Samokhin model when only the $`M\varphi _1\varphi _2`$coupling between the order parameters $`\psi `$ and $`M`$ is taken into account ($`\gamma >0`$, $`\gamma _1=0`$) and the anisotropies $`(w=v=0)`$ are ignored. The uniform phases and the phase diagram in this case were investigated in Refs. .
Here we summarize the main results in order to make a clear comparison with the results presented in Sec. 5 and Sec. 6. Our main aim is the description of a trigger effect which consists of the appearance of a “compelled superconductivity” caused by the presence of ferromagnetic order (here, this is a standard uniform ferromagnetic order); see also Refs. where this effect has been already established and briefly discussed. As mentioned in the Introduction, a similar trigger effect is known in the physics of improper ferroelectrics. We shall set $`\theta _30`$ and use the notation $`\theta \mathrm{\Delta }\theta =(\theta _2\theta _1)`$.
### 4.1 Phases
The possible (stable, metastable and unstable) phases are given in Table 1 together with the respective existence and stability conditions. The normal or disordered phase, denoted in Table 1 by $`N`$, always exists (for all temperatures $`T0)`$ and is stable for $`t>0`$, $`r>0`$. The superconducting phase denoted in Table 1 by SC1 is unstable. The same is valid for the phase of coexistence of ferromagnetism and superconductivity denoted in Table 1 by CO2. The N–phase, the ferromagnetic phase (FM), the superconducting phases (SC1–3) and two of the phases of coexistence (CO1–3) are generic phases because they appear also in the decoupled case $`\left(\gamma 0\right)`$. When the $`M\varphi _1\varphi _2`$–coupling is not present, the phases SC1–3 are identical and represented by the order parameter $`\varphi `$ with components $`\varphi _j`$ that participate on equal footing. The asterisk attached to the stability condition of the second superconductivity phase (SC2) indicates that our analysis is insufficient to determine whether this phase corresponds to a minimum of the free energy.
It will be shown that the phase SC2, two other purely superconducting phases and the coexistence phase CO1, have no chance to become stable for $`\gamma 0`$. This is so, because the phase of coexistence of superconductivity and ferromagnetism (FS in Table 1) that does not occur for $`\gamma =0`$ is stable and has a lower free energy in their domain of stability. A second domain $`\left(M<0\right)`$ of the FS phase is denoted in Table 1 by FS. Here we shall describe only the first domain FS. The domain FS is considered in the same way.
The cubic equation for magnetization of FS-phase (see Table 1) is shown in Fig. 1 for $`\gamma =1.2`$ and $`t=0.2`$. For any $`\gamma >0`$ and $`t`$, the stable FS thermodynamic states are given by $`r\left(M\right)<r_m=r\left(M_m\right)`$ for $`M>M_m>0`$, where $`M_m`$ corresponds to the maximum of the function $`r\left(M\right)`$. The dependence of $`M_m\left(t\right)`$ and $`M_0\left(t\right)=\left(t+\gamma ^2/2\right)^{1/2}=\sqrt{3}M_m\left(t\right)`$ on $`t`$ is drawn in Fig. 2 for $`\gamma =1.2`$. Functions $`r_m\left(t\right)=4M_m^3\left(t\right)/\gamma `$ for $`t<\gamma ^2/2`$ (depicted by the line of circles in Fig. 3) and
$$r_e\left(t\right)=\gamma \left|t\right|^{1/2},$$
(15)
for $`t<0`$ define the borderlines of stability and existence of FS.
### 4.2 Phase diagram
We have outlined the domain in the ($`t`$, $`r`$) plane where the FS phase exists and is a minimum of the free energy. For $`r<0`$ the cubic equation for $`M`$ (see Table 1) and the existence and stability conditions are satisfied for any $`M0`$ provided $`t\gamma ^2`$. For $`t<\gamma ^2`$ the condition $`MM_0`$ have to be fulfilled, here the value $`M_0=\left(t+\gamma ^2/2\right)^{1/2}`$ of $`M`$ is obtained from $`r\left(M_0\right)=0`$. Thus for $`r=0`$ the N-phase is stable for $`t\gamma ^2/2`$, and FS is stable for $`t\gamma ^2/2`$. For $`r>0`$, the requirement for the stability of FS leads to the inequalities
$$max(\frac{r}{\gamma },M_m)<M<M_0,$$
(16)
where $`M_m=\left(M_0/\sqrt{3}\right)`$ and $`M_0`$ should be the positive solution of the cubic equation of state from Table 1; $`M_m>0`$ gives a maximum of the function $`r\left(M\right)`$; see also Figs. 1 and 2.
The further analysis defines the existence and stability domain of FS below the line AB denoted by circles (see Fig. 3). In Fig. 3 the curve of circles starts from the point A with coordinates ($`\gamma ^2/2`$, $`0`$) and touches two other (solid and dotted) curves at the point B with coordinates ($`t_B=\gamma ^2/4`$, $`r_B=\gamma ^2/2`$). Line of circles represents the function $`r\left(M_m\right)r_m\left(t\right)`$ where
$$r_m\left(t\right)=\frac{4}{3\sqrt{3}\gamma }\left(\frac{\gamma ^2}{2}t\right)^{3/2}.$$
(17)
Dotted line represents $`r_e\left(t\right)`$ defined by Eq. (15). The inequality $`r<r_m\left(t\right)`$ is a condition for the stability of FS, whereas the inequality $`rr_e\left(t\right)`$ for $`\left(t\right)\gamma ^2/4`$ is a condition for the existence of FS as a solution of the respective equation of state. This existence condition for FS is obtained from $`\gamma M>r`$ (see Table 1).
In the region on the left of the point B in Fig. 3, the FS phase satisfies the existence condition $`\gamma M>r`$ only below the dotted line. In the domain confined between the lines of circles and the dotted line on the left of the point B the stability condition for FS is satisfied but the existence condition is broken. The inequality $`rr_e\left(t\right)`$ is the stability condition of FM for $`0\left(t\right)\gamma ^2/4`$. For $`\left(t\right)>\gamma ^2/4`$ the FM phase is stable for all $`rr_e\left(t\right)`$.
In the region confined by the line of circles AB, the dotted line for $`0<\left(t\right)<\gamma ^2/4`$, and the $`t`$axis, the phases N, FS and FM have an overlap of stability domains. The same is valid for FS, the SC phases and CO1 in the third quadrant of the plane ($`t`$, $`r`$). The comparison of the respective free energies for $`r<0`$ shows that the stable phase is FS whereas the other phases are metastable within their domains of stability.
The part of the $`t`$-axis given by $`r=0`$ and $`t>\gamma ^2/2`$ is a phase transition line of second order which describes the N-FS transition. The same transition for $`0<t<\gamma ^2/2`$ is represented by the solid line AC which is the equilibrium transition line of a first order phase transition. The equilibrium transition curve is given by the function
$$r_{eq}\left(t\right)=\frac{1}{4}\left[3\gamma \left(\gamma ^2+16t\right)^{1/2}\right]M_{eq}\left(t\right).$$
(18)
Here
$$M_{eq}\left(t\right)=\frac{1}{2\sqrt{2}}\left[\gamma ^28t+\gamma \left(\gamma ^2+16t\right)^{1/2}\right]^{1/2}$$
(19)
is the equilibrium jump of the magnetization. The order of the N-FS transition changes at the tricritical point A.
The domain above the solid line AC and below the line of circles for $`t>0`$ is the region of a possible overheating of FS. The domain of overcooling of the N-phase is confined by the solid line AC and the axes ($`t>0`$, $`r>0`$). At the triple point C with coordinates \[0, $`r_{eq}\left(0\right)=\gamma ^2/4`$\] the phases N, FM, and FS coexist. For $`t<0`$ the straight line
$$r_{eq}^{}\left(t\right)=\frac{\gamma ^2}{4}+\left|t\right|,t_B<t<0,$$
(20)
describes the extension of the equilibrium phase transition line of the N-FS first order transition to negative values of $`t`$. For $`t<t_B`$ the equilibrium phase transition FM-FS is of second order and is given by the dotted line on the left of the point B which is the second tricritical point in this phase diagram. Along the first order transition line $`r_{eq}^{}\left(t\right)`$ given by Eq. (20) the equilibrium value of $`M`$ is $`M_{eq}=\gamma /2`$, which implies an equilibrium order parameter jump at the FM-FS transition equal to ($`\gamma /2\sqrt{\left|t\right|}`$). On the dotted line of the second order FM-FS transition the equilibrium value of $`M`$ is equal to that of the FM phase ($`M_{eq}=\sqrt{\left|t\right|}`$). The FM phase does not exist below $`T_s`$ and this is a shortcoming of the model (2.2) with $`\gamma _1=0`$.
The equilibrium FM-FS and N-FS phase transition lines in Fig. 3 can be expressed by the respective equilibrium phase transition temperatures $`T_{eq}`$ defined by the equations $`r_e=r\left(T_{eq}\right)`$, $`r_{eq}=r\left(T_{eq}\right)`$, $`r_{eq}^{}=r\left(T_{eq}\right)`$, and with the help of the relation $`M_{eq}=M\left(T_{eq}\right)`$. This limits the possible variations of parameters of the theory. For example, the critical temperature ($`T_{eq}T_c`$) of the FM-FS second order transition ($`\gamma ^2/4<t`$) is obtained in the form $`T_c=\left(T_s+4\pi J/\alpha _s\right)`$, or, using $`=\left(a_f/b_f\right)^{1/2}`$,
$$T_c=T_s\frac{T^{}}{2}+\left[\left(\frac{T^{}}{2}\right)^2+T^{}\left(T_fT_s\right)\right]^{1/2}.$$
(21)
Here $`T_f>T_s`$, and $`T^{}=\left(4\pi J\right)^2\alpha _f/\alpha _s^2b_f`$ is a characteristic temperature of the model (2.2) with $`\gamma _1=w=v=0`$. A discussion of Eq. (21) is given in Sec. 5.3.
The investigation of the conditions for the validity of Eq. (21) leads to the conclusion that the FM-FS continuous phase transition (at $`\gamma ^2<t)`$ will be possible only if the following condition is satisfied:
$$T_fT_s>=\left(\varsigma +\sqrt{\varsigma }\right)T^{},$$
(22)
where $`\varsigma =b_f\alpha _s^2/4b_s\alpha _f^2`$. Therefore, the second order FM-FS transition should disappear for a sufficiently large $`\gamma `$–coupling. Such a condition does not exist for the first order transitions FM-FS and N-FS.
The inclusion of the gradient term (4) in the free energy (2) should lead to a depression of the equilibrium transition temperature. As the magnetization increases with the decrease of the temperature, the vortex state should occur at temperatures which are lower than the equilibrium temperature $`T_{eq}`$ of the Meissner state. For example, the critical temperature ($`\stackrel{~}{T}_c`$) corresponding to the vortex phase of FS-type has been evaluated to be lower than the critical temperature (21): $`\left(T_c\stackrel{~}{T}_c\right)=4\pi \mu _B/\alpha _s`$, where $`\mu _B=\left|e\right|\mathrm{}/2mc`$ is the Bohr magneton. For $`J\mu _B`$, we have $`T_c\stackrel{~}{T}_c`$.
For $`r>0`$, namely, for temperatures $`T>T_s`$ the superconductivity is triggered by the magnetic order through the $`\gamma `$-coupling. The superconducting phase for $`T>T_s`$ is entirely in the $`(t,r)`$ domain of the ferromagnetic phase. Therefore, the uniform supeconducting phase can occur for $`T>T_s`$ only through a coexistence with the ferromagnetic order.
The properties of the magnetic susceptibility and the specific heat near the phase transition lines shown in Fig. 3 have been investigated in Refs. and here we shall not dwell on these topics. Note that the results for the thermodynamic quantities should be extended to include the physical effects considered in the next parts of this review. Such a consideration requires a numerical analysis.
In the next Sections we shall focus on the temperature range $`T>T_s`$ which seems to be of main practical interest. We shall not dwell on the superconductivity in the fourth quadrant $`\left(t>0,r<0\right)`$ of the $`(t,r)`$ diagram where pure superconducting phases can occur for systems with $`T_s>T_f`$, but this is not the case for UGe<sub>2</sub>, URhGe and ZrZn<sub>2</sub>. Also we shall not discuss the possible metastable phases in the third quadrant $`\left(t<0,r<0\right)`$ of the $`(t,r)`$ diagram.
### 4.3 Note about a simplified theory
The analysis in this Section can be done following an approximate scheme known from the theory of improper ferroelectrics; see, e.g., Ref. . In this approximation the order parameter $`M`$ is considered small enough which makes possible to ignore $`M^4`$-term in the free energy. Then one easily obtains from the data for FS presented in Table 1 or by a direct calculation of the respective reduced free energy that the order parameters $`\varphi `$ and $`M`$ of FS–phase are described by the simple equalities $`r=\left(\gamma M\varphi ^2\right)`$ and $`M=\left(\gamma /2t\right)\varphi ^2`$. For ferroelectrics working with oversimplified free energy gives a substantial departure of theory from experiment . The same approximation has been recently applied to ferromagnetic Bose-Einstein condensates .
For ferromagnetic superconductors the domain of reliability of this approximation could be the close vicinity of the ferromagnetic phase transition, i.e., for temperatures near the critical temperature $`T_f`$. This discussion can be worthwhile if only the primary order parameter also exists in the same narrow temperature domain ($`\varphi >0`$). Therefore, the application of the simplified scheme can be useful in systems, where $`T_sT_f`$.
For $`T_s<T_f`$, the analysis can be simplified if we suppose a relatively small value of the modulus $`\varphi `$ of the superconducting order parameter. This approximation should be valid in some narrow temperature domain near the line of second order phase transition from FM to FS.
## 5 Effect of symmetry conserving coupling
Here we shall include in our consideration both linear and quadratic couplings of magnetization to the superconducting order parameter which means that both parameters $`\gamma `$ and $`\gamma _1`$ in free energy (12) are different from zero. In this way we shall investigate the effect of the symmetry conserving $`\gamma _1`$-term in the free energy on the thermodynamics of the system. When $`\gamma `$ is equal to zero but $`\gamma _10`$ the analysis is easy and the results are known from the theory of bicritical and tetracritical points . For the problem of coexistence of conventional superconductivity and ferromagnetic order the analysis $`\left(\gamma =0,\gamma _10\right)`$ was made in Ref. .
At this stage we shall not take into account any anisotropy effects because we do not want to obscure the influence of quadratic interaction by considering too many parameters. For $`\gamma ,\gamma _10`$ and $`w=0`$, $`v=0`$ the results again can be presented in an analytical form, only a small part of phase diagram should be calculated numerically.
### 5.1 Phases
The calculations show that for temperatures $`T>T_s`$, i.e., for $`r>0`$, we have again three stable phases. Two of them are quite simple: the normal ($`N`$-) phase with existence and stability domains shown in Table 1, and the FM phase with the existence condition $`t<0`$ as shown in Table 1, and a stability domain defined by the inequality $`r_e^{\left(1\right)}r`$. Here
$$r_e^{\left(1\right)}=\gamma _1t+\gamma \sqrt{t},$$
(23)
and one can compare it with the respective expression (15) for $`\gamma _1=0`$. In this paragraph we shall retain the same notations as in Sec. 4, but with a superscript $`\left(1\right)`$ in order to distinguish them from the case $`\gamma _1=0`$ The third stable phase for $`r>0`$ is a more complex variant of the mixed phase FS and its domain FS, discussed in Sec. 4. The symmetry of the FS phase coincides with that found in .
We have to mention that for $`r<0`$ there are five pure superconducting ($`M=0`$, $`\varphi >0`$) phases. Two of them, $`\left(\varphi _1>0,\varphi _2=\varphi _3=0\right)`$ and $`\left(\varphi _1=0,\varphi _2>0,\varphi _3>0\right)`$ are unstable. Two other phases, $`\left(\varphi _1>0,\varphi _2>0,\varphi _3=0,\theta _2=\theta _1+\pi k\right)`$ and $`(\varphi _1>0,\varphi _2>0,\varphi _3>0,\theta _2=\theta _1+\pi k,\theta _3`$ – arbitrary; $`k=0,\pm 1,\mathrm{})`$ show a marginal stability for $`t>\gamma _1r`$.
Only one of the five pure superconducting phases, the phase SC3, given in Table 1, is stable. In case of $`\gamma _10`$ the values of $`\varphi _j`$ and the existence domain of SC3 are the same as shown in Table 1 for $`\gamma _1=0`$ but the stability domain is different and is given by $`t>\gamma _1r`$. When the anisotropy effects are taken into account the phases exhibiting marginal stability within the present approximation may become stable. Besides, three other mixed phases $`\left(M0,\varphi >0\right)`$ exist for $`r<0`$ but one of them is metastable (for $`\gamma _1^2>1,t<\gamma _1r`$, and $`r<\gamma _1t`$) and the other two are absolutely unstable.
Here the thermodynamic behavior for $`r<0`$ is much more abundant in phases than for improper ferroelectrics with two component primary order parameter . However, at this stage of experimental needs about the properties of unconventional ferromagnetic superconductors the investigation of the phases for temperatures $`T<T_s`$ is not of primary interest and for this reason we shall focus our attention on the temperature domain $`r>0`$.
The FS phase for $`\gamma _10`$ is described by the following equations:
$$\varphi _1=\varphi _2=\frac{\varphi }{\sqrt{2}},\varphi _3=0,$$
(24)
$$\varphi ^2=\left(\pm \gamma Mr\gamma _1M^2\right),$$
(25)
$$\left(1\gamma _1^2\right)M^3\pm \frac{3}{2}\gamma \gamma _1M^2+\left(t\frac{\gamma ^2}{2}\gamma _1r\right)M\pm \frac{\gamma r}{2}=0,$$
(26)
and
$$\left(\theta _2\theta _1\right)=\frac{\pi }{2}+2\pi k,$$
(27)
($`k=0,\pm 1,\mathrm{}`$). The upper sign in Eqs. (24) - (27) corresponds to the FS domain where $`\text{sin}\left(\theta _2\theta _1\right)=1`$ and the lower sign corresponds to the FS domain with $`\text{sin}\left(\theta _2\theta _1\right)=1`$. This is a generalization of the two-domain FS phase discussed in Sec. 3. The analysis of the stability matrix (14) for these phase domains shows that FS is stable for $`M>0`$ and FS is stable for $`M<0`$, just like our result in Sec. 4. As these domains belong to the same phase, namely, have the same free energy and are thermodynamically equivalent, we shall consider one of them, for example, FS.
### 5.2 Phase stability and phase diagram
In order to outline the ($`t,r`$) phase diagram we shall use the information given above for the other two phases which have their own domains of stability in the $`(t,r)`$ plane: N and FM. The FS stability conditions when $`\gamma _10`$ become
$$2\gamma Mr\gamma _1M^20,$$
(28)
$$\gamma M0,$$
(29)
$$3\left(1\gamma _1^2\right)M^2+3\gamma \gamma _1M+t\gamma _1r\gamma ^2/20.$$
(30)
and we prefer to treat Eqs. (28) - (30) together with the existence condition $`\varphi ^20`$, with $`\varphi `$ given by Eq. (25), with the help of the picture shown in Fig. 4.
The most direct approach to analyze the existence and stability of FS phase is to express $`r`$ as of function of $`(M,t)`$ from the equation of state (26),
$`r_{eq}^{\left(1\right)}\left(t\right)=`$
$`{\displaystyle \frac{M_{eq}}{\left(\gamma _1M_{eq}\gamma /2\right)}}\left[\left(1\gamma _1^2\right)M_{eq}^2+{\displaystyle \frac{3}{2}}\gamma \gamma _1M_{eq}+\left(t{\displaystyle \frac{\gamma ^2}{2}}\right)\right],`$
and to substitute the above expression in the existence and stability conditions of FS-phase. It is obvious that there is a special value of $`M`$
$$M_{S1}=\frac{\gamma }{2\gamma _1}$$
(32)
that is a solution of Eq. (26) for any value of $`r`$ and
$$t_{S1}=\frac{\gamma ^2}{4\gamma _1^2},$$
(33)
for which this procedure cannot be applied and should be considered separately. Note, that $`M_{S1}`$ is given by the respective horizontal dashed line in Fig. 4. The analysis shows that in the interval $`t_B^{\left(1\right)}<t<\gamma ^2/2`$ the phase transition is again of first order; here
$$t_B^{\left(1\right)}=\frac{\gamma ^2}{4\left(1+\gamma _1\right)^2}.$$
(34)
To find the equilibrium magnetization of first order phase transition, depicted by the thick line $`ACB`$ in Fig. 4 we need the expression for equilibrium free energy of FS-phase. It is obtained from Eq. (12) by setting $`\left(w=0,v=0\right)`$ and substituting $`r,\varphi _i`$ as given by Eqs. (24), (25) and (31). The result is
$`f_{FS}^{\left(1\right)}`$ $`=`$ $`{\displaystyle \frac{M^2}{2\left(M\gamma _1\gamma /2\right)^2}}\times \{(1\gamma _1^2)M^4+\gamma \gamma _1M^3`$
$`+2[t(1\gamma _1^2){\displaystyle \frac{\gamma ^2}{8}}]M^22\gamma \gamma _1tM+t(t{\displaystyle \frac{\gamma ^2}{2}})\},`$
where $`MM_{eq}`$.
For the phase transition from N to FS phase ($`0<t<\gamma ^2/2`$), $`M_{eq}`$ is found by setting the FS free energy from the above expression equal to zero, as we have by convention that the free energy of the normal phase is zero. The value of $`M_{eq}^{\left(1\right)}`$ for positive $`t`$ is obtained numerically and is illustrated by thick black curve $`AC`$ in Fig. 4. When $`t_B^{\left(1\right)}t<0`$ the transition is between FM and FS phases and we obtain $`M_{eq}^{\left(1\right)}`$ from the equation $`f_{FS}=f_{FM}=\left(t^2/2\right)`$, where $`f_{FM}`$ is the free energy of FM phase. The equilibrium magnetization in the above $`t`$-interval is given by the formula
$$M_{eq}^{\left(1\right)}=\frac{\gamma }{2\left(1+\gamma _1\right)},$$
(36)
and is drawn by thick line $`CB`$ in Fig. 4. The existence and stability analysis shows that for $`r>0`$ the equilibrium magnetization of the first order phase transition should satisfy the condition $`M_m^{\left(1\right)}<M_{eq}^{\left(1\right)}<M_0^{\left(1\right)}`$.
By $`M_0^{\left(1\right)}`$ we denote the positive solution of $`r^{\left(1\right)}\left(M_{eq}\right)=0`$ and its $`t`$-dependence is drawn in Fig. 4 by the curve with circles. $`M_m^{\left(1\right)}`$ is the smaller positive root of stability condition (30) and also gives the maximum of the function $`r_{eq}^{\left(1\right)}\left(M\right)`$; see Eq. (31). The function $`M_m^{\left(1\right)}`$ is depicted by the dotted curve $`AB`$ in Fig. 4. When $`t_{S1}<t<t_B^{\left(1\right)}`$ the existence and stability conditions are fulfilled if $`\sqrt{t}<M<M_{S1}`$, where $`\sqrt{t}`$ is the magnetization of ferromagnetic phase and is drawn by a thin black line on the left of point B in Fig. (4). Here we have two possibilities: $`r>0`$ for $`\sqrt{t}<M<M_0^{\left(1\right)}`$ and $`r<0`$ for $`M_0^{\left(1\right)}<M<M_{S1}`$. To the left of $`t_{S1}`$ and $`t>t_{S2}`$, where
$$t_{S2}=\left(\frac{\gamma }{\gamma _1}\right)^2.$$
(37)
the FS phase is stable and exists for $`M_{S1}<M<\sqrt{t}`$. Here $`r`$ will be positive when $`M_0^{\left(1\right)}<M<\sqrt{t}`$ and $`r<0`$ for $`M_0^{\left(1\right)}>M>M_{S1}`$. When $`t<t_{S2}`$, $`M<\sqrt{t}`$ and $`r`$ is always negative.
On the basis of the existence and stability analysis we draw in Fig. 5 the $`(t,r)`$-phase diagram for concrete values of $`\gamma `$ and $`\gamma _1`$. As we have mentioned above the order of phase transitions is the same as for $`\gamma _1=0`$, see Fig. 3, Sec. 4. The phase transition between the normal and FS phases is of first order and goes along the equilibrium line $`AC`$ in the interval ($`t_A=\gamma ^2/2`$ and $`t_C=0`$). The function $`r_{eq}^{\left(1\right)}\left(t\right)`$ is given by Eq. (31) with $`M_{eq}^{\left(1\right)}`$ from Fig. 4.
N, FM, and FS phases coexist at the triple point $`C`$ with coordinates $`t=0`$, and $`r_C^{\left(1\right)}=\gamma ^2/4\left(\gamma _1+1\right)`$. On the left of $`C`$ for $`t_B^{\left(1\right)}<t<0`$ the phase transition line of first order $`r_{eq}^{\left(1\right)}\left(t\right)`$ is found by substituting in Eq. (31) the respective equilibrium magnetization, given by Eq. (35). In result we obtain
$$r_{eq}^{\left(1\right)}\left(t\right)=\frac{\gamma ^2}{4\left(1+\gamma _1\right)}t.$$
(38)
This function is illustrated by the line $`BC`$ in Fig. 5 that terminates at the tricritical point $`B`$ with coordinates $`t_B^{\left(1\right)}`$ from Eq. (34), and
$$r_B^{\left(1\right)}=\frac{\gamma ^2\left(2+\gamma _1\right)}{4\left(1+\gamma _1\right)^2}.$$
(39)
To the left of the tricritical point $`B`$ the second order phase transition curve is given by the relation (23). Here the magnetization is $`M=\sqrt{t}`$ and the superconducting order parameter is equal to zero ($`\varphi =0`$). This line intersects t-axis at $`t_{S2}`$ and is well defined also for $`r<0`$. The function $`r_e^{\left(1\right)}\left(t\right)`$ has a maximum at the point $`(t_{S1},\gamma ^2/4\gamma _1)`$; here $`M=M_{S1}`$. When this point is approached the second derivative of the free energy with respect to $`M`$ tends to infinity. The result for the curves $`r_{eq}^{\left(1\right)}\left(t\right)`$ of equilibrium phase transitions (N-FS and FM-FS) can be used to define the respective equilibrium phase transition temperatures $`T_{FS}`$.
We shall not discuss the region, $`t>0`$, $`r<0`$, because we have supposed from the very beginning that the transition temperature for the ferromagnetic ordering T<sub>f</sub> is higher then the superconducting transition temperature T<sub>s</sub>, as is for the known unconventional ferromagnetic superconductors. But this case may become of substantial interest when, as one may expect, materials with T$`{}_{f}{}^{}<`$T<sub>s</sub> may be discovered experimentally.
### 5.3 Discussion
The shape of the equilibrium phase transition lines corresponding to the phase transitions N-SC, N-FS, and FM-FS is similar to that of the more simple case $`\gamma _1=0`$ and we shall not dwell on the variation of the size of the phase domains with the variations of the parameter $`\gamma _1`$ from zero to values constrained by the condition $`\gamma _1^2<1`$. We shall draw the attention to the important qualitative difference between the equilibrium phase transition lines shown in Figs. 3 and 5. The second order phase transition line $`r_e\left(t\right)`$, shown by the dotted line on the left of point $`B`$ in Fig. 3, tends to large positive values of $`r`$ for large negative values of $`t`$ and remains in the second quadrant ($`t<0,r>0)`$ of the plane ($`t,r`$) while the respective second order phase transition line $`r_e^{\left(1\right)}\left(t\right)`$ in Fig. 5 crosses the $`t`$-axis at the point $`t_{S2}`$ and is located in the third quadrant ($`t<0,r<0`$) for all possible values $`t<t_{S2}`$. This means that the ground state (at 0 K) of systems with $`\gamma _1=0`$ will be always the FS phase while two types of ground states, FM and FS, can exist for systems with $`0<\gamma _1^2<1`$. The latter seems more realistic when we compare theory and experiment, especially, in ferromagnetic compounds like UGe<sub>2</sub>, URhGe, and ZrZn<sub>2</sub> where the presence of FM phase is observed at very low temperatures and relatively low pressure $`P`$.
The final aim of the phase diagram investigation is the outline of the ($`T,P`$) diagram. Important conclusions about the shape of the $`(T,P)`$ diagram can be made from the form of the $`(t,r)`$ diagram without an additional information about the values of the relevant material parameters $`(a_s`$, $`a_f,\mathrm{}`$) and their dependence on the pressure $`P`$. One should know also the characteristic temperature $`T_s`$, which has a lower value than the experimentally observed phase transition temperature $`\left(T_{FS}1K\right)`$ to the coexistence FS–phase. A supposition about the dependence of the parameters $`a_s`$ and $`a_f`$ on the pressure $`P`$ was made in Ref. . Our results for $`T_fT_s`$ show that the phase transition temperature $`T_{FS}`$ varies with the variation of the system parameters $`(\alpha _s,\alpha _f,\mathrm{})`$ from values which are higher than the characteristic temperature $`T_s`$ down to zero temperature. This is seen from Fig. 5.
In systems where a pure superconducting phase is not observed for temperatures $`TT_f`$ or $`TT_{FS}`$, we can set $`T_s0`$ in Eq. (21). Neglecting $`T_s`$ in Eq. (21) and assuming that $`\left(T^{}/T_f\right)1`$ we obtain that $`T_cT_{FS}\left(T^{}T_f\right)^{1/2}`$. Note that the first $`\left(T^{}/T_f\right)^{1/2}`$-correction to this result has a negative sign which means that a suitable dependence of the characteristic temperature $`T^{}`$ on the pressure P may be used in attempts to describe the experimental shape of the FM-FS phase transition line in the $`(T,P)`$ diagrams of UGe<sub>2</sub> and ZrZn<sub>2</sub>; see, for example, Fig. 2 in Ref. , Fig. 3 in Ref. , Fig. 4 in Ref. . The experimental phase diagrams indicate that $`T_f\left(P\right)`$ is a smooth monotonically decreasing function of the pressure $`P`$ and $`T_f\left(P\right)`$ tends to zero when the pressure $`P`$ exceeds some critical value $`P_c1`$ GPa. Postulating the respective experimental shape of the function $`T_f\left(P\right)`$ one may try to give a theoretical prediction for the shape of the curve $`T_{FS}`$.
The lack of experimental data about important parameters of the theory forces us to make some suppositions about the behavior of the function $`T^{}\left(P\right)`$. The phase transition temperature $`T_{FS}`$ will qualitatively follow the shape of $`T_f\left(P\right)`$ provided the dependence $`T^{}\left(P\right)`$ is very smooth. This is in accord with the experimental shapes of these curves near the critical pressure $`P_c`$ where both $`T_f`$ and $`T_{FS}`$ are very small. The substantial difference between $`T_f`$ and $`T_{FS}`$ at lower pressure ($`P<P_c`$) can be explained with the negative sign of the correction term to the leading dependence $`T_{FS}\left(P\right)\left[T^{}\left(P\right)T_f\left(P\right)\right]^{1/2}`$ mentioned above and a convenient supposition for the form of the function $`T^{}\left(P\right)`$.
Eq. (21) presents a rather simplified theoretical result for $`T_CT_{FS}`$ because the effect of $`M^2\left|\psi \right|^2`$ coupling is not taken into account. But following the same ideas, used in our discussion of Eq. (21), a more reliable theoretical prediction of the shape of FM-FS phase transition line can be given on the basis of Eq. (23). With the help of the experimentally found shape of $`T_f\left(P\right)`$ and the definition of the parameters $`r`$ and $`t`$ by Eq. (11) we can substitute $`T=T_{FS}\left(P\right)`$ in Eq. (23). In doing this we have applied the following approximations, namely, that $`T_s0`$ for any pressure $`P`$, $`T_{FS}\left(P_c\right)T_f\left(P_c\right)0`$ and for substantially lower pressure ($`P<P_c`$), $`T_f\left(P\right)T_{FS}\left(P\right)`$. Then near the critical pressure $`P_c`$, we easily obtain the transition temperature $`T_{FS}0`$, as should be. For substantially lower values of the pressure there exists an experimental requirement $`\left(T_{FS}T_s\right)\left(T_fT_{FS}\right)`$. Using the latter we establish the approximate formula
$$\left(T_fT_{FS}\right)=\gamma ^2b_f^{1/2}/\gamma _1^2\alpha _f.$$
(40)
The same formula for $`\left(T_fT_{FS}\right)`$ can be obtained from the parameter $`t_{S2}\left(T_{FS}\right)`$ given by Eq. (37). The pressure dependence of the parameters included in this formula defines two qualitatively different types of behavior of $`T_{FS}\left(P\right)`$ at relatively low pressures ($`PP_c`$): (a) $`T_{FS}\left(P\right)0`$ below some (second) critical value of the pressure ($`P_c^{}<P_c`$), and (b) finite $`T_{FS}\left(P\right)`$ up to $`P0`$. Therefore, we can estimate the value of the pressure $`P_c^{}<P_c`$ in UGe<sub>2</sub>, where $`T_{FS}\left(P_c^{}\right)0`$. It can be obtained from the equation
$$T_f\left(P_c^{}\right)=\left(\gamma ^2b_f^{1/2}/\gamma _1^2\alpha _f\right)$$
(41)
provided the pressure dependence of the respective material parameters is known. So, the above consideration is consistent with the theoretical prediction that the dashed line in Fig. 5 crosses the axis $`r=0`$ and for this reason we have the opportunity to describe two ordered phases at low temperatures and broad variations of the pressure. Our theory allows also a description of the shape of the transition line $`T_{FS}\left(P\right)`$ in ZrZn<sub>2</sub> and URhGe, where the transition temperature $`T_{FS}`$ is finite at ambient pressure. To avoid a misunderstanding, let us note that the diagram in Fig. 5 is quite general and the domain containing the point $`r=0`$ of the phase transition line for negative $`t`$ may not be permitted in some ferromagnetic compounds.
Up to now we have discussed experimental curves of second order phase transitions. Our analysis gives the opportunity to describe also first order phase transition lines. Our investigation of the free energy (12) leads to the prediction of triple ($`C`$) and tricritical points ($`A`$ and $`B`$); see Figs. 3 and 5. We shall not consider the possible application of these results to the phase diagrams of real substances, for which first order phase transitions and multicritical phenomena occur; see, e.g., Refs. , where first order phase transitions and tricritical points have been observed. The explanation of the phase transition lines in Refs. requires further theoretical studies that can be done on the basis of a convenient extension of the free energy (12). For example, the investigation of vortex phases in Ref. requires to take into account the gradient terms (4). Another generalization should be done in order to explain the observation of two FM phases . Note, that the experimentalists are not completely certain whether the FS phase is a uniform or a vortex phase, and this is a crucial point for the further investigations. But we find quite encouraging that our studies naturally lead to the prediction of the same variety of phase transition lines and multicritial points that has been observed in recent experiments .
## 6 Anisotropy
Our analysis demonstrates that when the anisotropy of Cooper pairs is taken into account, there will be no drastic changes in the shape the phase diagram for $`r>0`$ and the order of the respective phase transitions. Of course, there will be some changes in the size of the phase domains and the formulae for the thermodynamic quantities. It is readily seen from Figs. 6 and 7 that the temperature domain of first order phase transitions and the temperature domain of stability of FS above $`T_s`$ essentially vary with the variations of the anisotropy parameter $`w`$. The parameter $`w`$ will also insert changes in the values of the thermodynamic quantities like the magnetic susceptibility and the entropy and specific heat jumps at the phase transition points.
Besides, and this seems to be the main anisotropy effect, the $`w`$\- and $`v`$-terms in the free energy lead to a stabilization of the order along the main crystal directions which, in other words, means that the degeneration of the possible ground states (FM, SC, and FS) is considerably reduced. This means also a smaller number of marginally stable states.
The dimensionless anisotropy parameter $`w=u_s/\left(b_s+u_s\right)`$ can be either positive or negative depending on the sign of $`u_s`$. Obviously when $`u_s>0`$, the parameter $`w`$ will be positive too and will be in the interval $`0<`$ w$`<1`$ to ensure the positiveness of parameter $`b`$ from Eq. (10). When $`w<0`$, the latter condition is obeyed if the original parameters of free energy (3) satisfy the inequality $`b_s<u_s<0`$.
We should mention here that a new phase of coexistence of superconductivity and ferromagnetism occurs as a solution of Eqs. (13). It is defined in the following way:
$$\varphi _1^2+\varphi _2^2=\frac{1}{1\gamma _1^2}\left[\gamma _1\left(t+\frac{\gamma ^2}{2w}\right)r\right],$$
(42)
$$M^2=\frac{1}{1\gamma _1^2}\left[\gamma _1r\left(t+\frac{\gamma ^2}{2w}\right)\right],$$
(43)
and
$$2w\mathrm{sin}\left(\theta _2\theta _1\right)=\gamma M,\mathrm{cos}\left(\theta _2\theta _1\right)0.$$
(44)
In the present approximation the phase (42) - (44) is unstable, but this may be changed when the crystal anisotropy is taken into account.
We shall write the equations for order parameters $`M`$ and $`\varphi _j`$ of FS phase in order to illustrate the changes when $`w0`$
$$\varphi ^2=\frac{\pm \gamma Mr\gamma _1M^2}{\left(1w\right)}0,$$
(45)
and
$`\left(1w\gamma _1^2\right)M^3`$
$`\pm {\displaystyle \frac{3}{2}}\gamma \gamma _1M^2+\left[t\left(1w\right){\displaystyle \frac{\gamma ^2}{2}}\gamma _1r\right]M\pm {\displaystyle \frac{\gamma r}{2}}=0,`$
where the meaning of the upper and lower sign is the same as explained just below Eq. (27). The difference in the stability conditions is more pronounced and gives new effects that will be explained further,
$$\frac{\left(2w\right)\gamma Mr\gamma _1M^2}{1w}0,$$
(47)
$$\gamma Mwrw\gamma _1M^20,$$
(48)
and
$$\frac{3\left(1w\gamma _1^2\right)M^2+3\gamma \gamma _1M+t\left(1w\right)\gamma ^2/2\gamma _1r}{1w}0.$$
(49)
The calculations of the phase diagram in ($`t,r`$) parameter space are done in the same way as in case of $`w=0`$ and show that for $`w>0`$ there is no qualitative change of the phase diagram. Quantitatively, the region of first order phase transition widens both with respect to $`t`$ and $`r`$ as illustrated in Fig. 6. On the contrary, when $`w<0`$ the first order phase transition region becomes more narrow but the condition (47) limits the stability of FS for $`r<0`$. This is seen from Fig. 7 where FS is stable above the straight dotted line for $`r<0`$ and $`t<0`$. So, purely superconducting (Meissner) phases occur also as ground states together with FS and FM phases.
## 7 Conclusion
We investigated the M-trigger effect in unconventional ferromagnetic superconductors. This effect arises from the $`M\psi _1\psi _2`$-coupling term in the GL free energy and provokes the appearance of superconductivity in a domain of the system’s phase diagram that is entirely occupied by the ferromagnetic phase. The coexistence of unconventional superconductivity and ferromagnetic order is possible for temperatures above and below the critical temperature $`T_s`$, that corresponds to the standard second-order phase transition from normal to Meissner phase – usual uniform superconductivity in a zero external magnetic field which occurs outside the domain of existence of ferromagnetic order. Our investigation is mainly intended to clarify the thermodynamic behavior at temperatures $`T_s<T<T_f`$ where the superconductivity cannot appear without the mechanism of M-triggering. We describe the possible ordered phases (FM and FS) in this most interesting temperature interval.
The Cooper pair and crystal anisotropies are investigated and their main effects on the thermodynamics of the triggered phase of coexistence is established. Of course, in discussions of concrete real materials the respective crystal symmetry should be considered. But the low symmetry and low order (in both $`M`$ and $`\psi `$) $`\gamma `$-term in the free energy determines the leading features of the coexistence phase and the dependence of essential thermodynamic properties on the type of crystal symmetry is not so considerable.
Below the superconducting critical temperature $`T_s`$ a variety of pure superconducting and mixed phases of coexistence of superconductivity and ferromagnetism exists and the thermodynamic behavior at these relatively low temperatures is more complex than in improper ferroelectrics. The case $`T_f<T_s`$ also needs a special investigation.
Our results are referred to the possible uniform superconducting and ferromagnetic states. Vortex and other nonuniform phases need a separate study.
The relation of the present investigation to properties of real ferromagnetic compounds, such as UGe<sub>2</sub>, URhGe, and ZrZn<sub>2</sub>, has been discussed throughout the text. In these compounds the ferromagnetic critical temperature is much larger than the superconducting critical temperature $`\left(T_fT_s\right)`$ and that is why the M-triggering of the spin-triplet superconductivity is very strong. Moreover, the $`\gamma _1`$-term is important to stabilize the FM order up to the absolute zero (0 K), as is in the known spin-triplet ferromagnetic superconductors. Ignoring the symmetry conserving $`\gamma _1`$-term does not allow a proper description of the real substances of this type. More experimental information about the values of the material parameters ($`a_s,a_f,\mathrm{}`$) included in the free energy (12) is required in order to outline the thermodynamic behavior and the phase diagram in terms of thermodynamic parameters $`T`$ and $`P`$. In particular, a reliable knowledge about the dependence of the parameters $`a_s`$ and $`a_f`$ on the pressure $`P`$, the value of the characteristic temperature $`T_s`$ and the ratio $`a_s/a_f`$ at zero temperature are of primary interest.
## Acknowledgments
D.I.U. thanks the hospitality of MPI-PKS (Dresden), where a part of this work has been made. Financial support by SCENET (Parma) is also acknowledged.
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# Forward Physics and BRAHMS results
## 1 Introduction
The first collisions of gold ions at nucleon-nucleon center of mass energies $`\sqrt{s_{NN}}=130GeV`$ at RHIC showed a dramatic drop in the production of pions at intermediate $`p_T`$ compared to an incoherent sum of pions produced in p+p collisions at the same energy (approximately a factor of 5 for the most central collisions) . The measured suppression could be the result of energy degradation of jets traversing a newly formed dense medium or could also be related to modifications to the wave function of the ions brought in by the high energy of the collisions together with the high atomic number A of the nuclei . Collisions between deuterium and gold ions at $`\sqrt{s_{NN}}=200GeV`$ were produced during the third RHIC run to resolve the apparent conflict between the above mentioned explanations of the suppression. Particle production from d+Au collisions around mid-rapidity do not show the suppression seen in Au+Au collisions . What is seen instead is an enhancement that has been associated with the so called “Cronin effect” , where partons undergo multiple incoherent scatterings that increase their transverse momenta as they traverse through the target. These results constitute evidence that a new dense and highly opaque medium has been formed in Au+Au collisions at RHIC and the suppression of intermediate to high $`p_T`$ leading particles is directly related to their interaction with that medium be it collisional or by induced gluon radiation.
But the unexpected low overall multiplicity seen in Au+Au collisions points to a strong degree of coherence compatible with the onset of saturation in the wave function of the ions. A saturation that appears below a transverse momentum scale whose value increases with the energy of the collisions and the atomic number (A=197) of the colliding nuclei.
Asymmetric reactions like d+Au run in collider mode are fertile ground for QCD studies because the projectile and target fragmentation regions are well separated, and the detection of particles in the fragmentation regions skews even more the kinematics at the partonic level. (See appendix A). These asymmetric systems are thus ideal to study the small-x components of the Au target wave function. BRAHMS, one of the RHIC experiments specially designed to study particle production at high rapidity, was thus particularly well suited to produce measurements that are considered a first indication of saturation in the gluon density of the Au ion.
## 2 Lower energy p+A measurements
Before proceeding with the description of BRAHMS studies of d+Au collisions, a brief review of previous measurements is necessary to emphasize the novelty of the results that we are going to present. Early studies of p+A collisions were conducted with the primary aim of extracting the proton energy loss in nuclear matter. All those measurements were done in fixed target mode. I base this review in two well know papers . Compared with the RHIC measurements, all these experiments suffer from the fact that in fixed target mode the projectile fragmentation region (close to the proton rapidity) is boosted to very small angles and a detailed transverse momentum dependent studies were just not possible.
Many of these measurements have concentrated in particle production as function of polar angle or pseudo-rapidity $`\eta =log(tan\frac{\theta }{2})`$.
Figure 1 is a compilation of the NA35 pA program at CERN in the form of multiplicity distributions in laboratory rapidity space. These distributions have two main features: the first being the fact that most of the yield (in this particular case, negative particles or mostly pions) from p+A collisions appears close to the target fragmentation regions (y or $`\eta 0`$ in the lab. reference frame) and the second is the fact that within $`1`$ unit of rapidity close to the beam rapidity, the projectile proton has no more “memory” of the target. Similar behavior has been shown to be independent of the beam energy (one such compilation in the projectile reference frame can be found in ref. ). If the yields from p+A collisions are compared to those from p+p at the same energy, one obtains a characteristic wedge like distribution that has been explained in the context of multiple parton interactions; the projectile appears as if it had only one interaction but that interaction can involve several partons from the target nucleons. The ratio starts at $`\eta =0`$ with a value equal to the numbers of partonic interactions and ends at the rapidity of the projectile with a value equal to one. .
The measurements of hadrons at large transverse momentum done at Fermilab have shown what is now called the Cronin effect; an enhancement that is widely considered as incoherent multiple elastic interactions as the projectile moves through the target, each one of these interactions modifies the transverse momentum distributions by shifting counts from low $`p_T`$ values up to intermediate values ($`45GeV/c`$). The nuclear modification factor defined as a ratio of differential cross sections normalized by the atomic number A of the targets is used to compare to an incoherent sum of p+p collisions at the same energy. One such comparison is shown in the left panel of Fig. 2. The ratio for pions has a clear enhancement that starts above 2 and extends to 7 GeV/c. Anti-protons show an strikinly different behavior when compared to the above described pions. The difference between baryon and meson present at this energy ($`\sqrt{s_{NN}}=27.4GeV`$) has also been seen at RHIC. The right panel of the figure shows a comparison of abundances of baryons (protons and anti-protons) and mesons (pions). The ratio of anti-proton to negative pions is small and consistent with hadronization in the vacuum, the ratio of protons to positive pions, is greater and approaches one for the heaviest target. Because the energy of these collisions is not that high, it may be, that this ratio is affected by beam protons, as some degree of stopping is expected to transfer protons to mid-rapidity where this ratio was constructed.
Figure 3 shows results from collisions at an even lower energy. This time, the nuclear modification factor is defined as a ratio of invariant differential cross sections for positive pion production, scaled by the atomic number of the targets. The comparison is done between a heavy target (Au) and a light one (Be). The data was collected in fixed target mode at the AGS with the E802 spectrometer. The large acceptance coverage in rapidity provides some access to both target and projectile fragmentations regions. The rapidity of the 14.6 GeV/c proton beam is equal to 3.4 and the data points in the figure are labeled with the rapidities in the nucleon-nucleon center of mass ($`y_{CM}=1.7`$).
The most striking feature of this figure is the fact that the curves are arranged in the same descending order as the above mentioned “triangular distribution”. Only the ratio calculated close to target rapidities ($`1.1<y<0.9`$) crosses the value of 1 at $`p_T1GeV/c`$. The ratios corresponding to mid-rapidity $`y=0`$ and $`y1`$ have values smaller than one. According to the scattering models used to explain the Cronin effect the low momentum (ratio below 1) is depleted because each rescattering shifts the event to higher $`p_T`$ bins. The depletion is stronger as the rapidity increases because the reach into lower values of x in the target wave function is greater; more scattering centers are thus available. Naive scattering models imply that at higher rapidities, the Cronin peak appears at higher values of $`p_T`$, however, data extending to high $`p_T`$ values is not available. It should also be said that available phase space limits the applicability of such arguments.
## 3 Intermediate $`p_T`$ studies and the nuclear modification factors
BRAHMS is one of the four RHIC experiments with the unique capability to measure identified hadrons with transverse momenta that can reach moderately high values ($`5GeV/c`$) and can access rapidities close to the beam rapidity (y=5.4 for the 100 GeV/c per nucleon beam). The data that is described in this presentation is thus a first detailed study of particle production in the beam fragmentation region in d+Au collisions at the highest energy in the center of mass. As such, the data may be a window to new phenomena and in particular, it has been listed as a first indication that the small-x components of the target wave function have entered a non-linear mode. The data collected from d+Au collisions is compared to p+p using the so called nuclear modification factor defined as: $`R_{dAu}=\frac{1}{N_{coll}}\frac{\frac{dN^{dAu}}{dp_Td\eta }}{\frac{dN^{pp}}{dp_Td\eta }}`$. If the target is already a saturated system of gluons, the ratio is expected to show a decrease in value as the rapidity of the detected particles increases. If the target is a dilute system of gluons, the ratio should grow with rapidity because, at higher rapidities, the detected particle is related to a parton that has interacted with a greater number of small-x gluons, each contributing a finite amount of transverse momentum, such that the ratio beyond some value of $`p_T`$ grows greater than one and then tends to one from above.
Figure 4 shows the nuclear modification factor with the number of binary collisions set to $`N_{coll}=7.2\pm 0.6`$ for minimum biased d+Au collisions. This particular study was done without identifiying the particles. Each panel shows the ratio calculated at a different pseudo-rapidity $`\eta `$ values. At mid-rapidity ($`\eta =0`$), the nuclear modification factor exceeds 1 for transverse momenta greater than 2 GeV/c in similar way as the pions in Fig. 2.
One unit of rapidity towards the deuteron rapidity is enough to make the enhancement disappear, and then become consistently smaller than 1 for the next two values of pseudo-rapidity ($`\eta `$ = 2.2 and 3.2) indicating a suppression in d+Au collisions compared to scaled p+p systems at the same energy.
The novelty of this result stands mainly on the fact that the measurement extends to moderate transverse momenta ($`3.5GeV/c`$) and the factor appears consistently suppressed for all value of $`p_T`$, specially at $`\eta =3.2`$.
The four panels of Fig. 5 show the central $`R_{CP}^{central}`$ (filled symbols) and semi-central $`R_{CP}^{semicentral}`$ (open symbols) ratios for the four $`\eta `$ settings. Central events have a higher number of target nucleons participating in the interaction with the deuteron projectile. This higher number of nucleons translates into an increased number of gluons present in the system and, if the conditions are set for saturation, central collisions would have stronger suppression as function of rapidity. If the target is still a linear dilute system the $`R_{CP}^{central}`$ would be enhanced as fuction of rapidity because of the higher number of scattering centers that become available in the target. In the left panel of Fig. 5 corresponding to $`\eta =0`$, the yield from the central sample of events (filled symbols) is systematically higher than those of the semi-central events, but at the highest pseudo-rapidity $`\eta =3.2`$, the trend is reversed; the yields from central events are $`60\%`$ lower than the semi-central events at all values of $`p_T`$. More details on these results can be found in . These results have been described within the context of the Color Glass Condensate ; the evolution of the nuclear modification factor with rapidity and centrality is consistent with a description of the Au target where the rate of gluon fusion becomes comparable with that of gluon emmission as the rapidity increases and it slows down the overall growth of the gluon density. The measured nuclear modification factor compares the slowed down growth of the numerator to a sum of incoherent p+p collisions, considered as dilute systems, whose gluon densities grow faster with rapidity because of the abscence of gluon fusion in dilute systems . Other explanations for the measured suppression have been proposed and they also reproduce the data .
The nuclear modification factor of baryons is different from the one calculated with mesons, whenever the factor shows the so called Cronin enhancement, baryons show a stronger enhancement. Such difference has been seen ot lower energies and is shown for pions and anti-protons in Fig. 2, it has also been found at RHIC energies at all rapidities, in particular, Fig. 6 presents the minimum bias nuclear modification $`R_{dAu}`$ for anti-protons and negative pions at $`\eta =3.2`$. These ratios were obtained making use of ratios of raw counts of identified particles compared to those of charged particles in each $`p_T`$ bin:
$$R_{dAu}^{\overline{p}}=R_{dAu}^h^{}\frac{(\frac{\overline{p}}{h^{}})^{dAu}}{(\frac{\overline{p}}{h^{}})^{pp}}=\frac{1}{N_{coll}}\frac{\frac{dn^{dAu}}{dp_Td\eta })^h^{}}{\frac{dn^{pp}}{dp_Td\eta })^h^{}}\frac{\frac{\frac{dn^{dAu}}{dp_Td\eta })^{\overline{p}}}{\frac{dn^{dAu}}{dp_Td\eta })^h^{}}}{\frac{\frac{dn^{pp}}{dp_Td\eta })^{\overline{p}}}{\frac{dn^{pp}}{dp_Td\eta })^h^{}}}=\frac{1}{N_{coll}}\frac{\frac{dn^{dAu}}{dp_Td\eta })^{\overline{p}}}{\frac{dn^{pp}}{dp_Td\eta })^{\overline{p}}}$$
No attempt was made to estimate the contributions from anti-lambda feed down to the anti-proton result. The remarkable difference between baryons and mesons has been related to parton recombination .
The ratios shown in Fig. 7 show a new aspect of particle production at forward rapidities. The left panel shows that in d+Au collisions at $`\eta =3.2`$ the yield of protons is comparable to the pions yield ($`80\%`$) while the yield of positive kaons hovers around $`40\%`$. These results indicate that eventhough pion production is well described at all rapidities in p+p collisions , the presence of so many baryons at that rapidity brings additional complications to NLO pQCD calculations, which cannot be reconciled with the data if standard fragmentation functions are used . The abundance of baryons at this high energy and rapidity doesn’t support the idea of baryon suppression in the fragmentation region where, because of their high energy, the quarks of the beam would fragment independently mostly into mesons. But if that suppression was actually present much closer to beam rapidity, baryon number conservation would force the transfer of beam protons to lower rapidities, what remains a mystery is the mechanism that gives these protons the high transverse momentum ($`2GeV/c`$) that is measured. Panel b of Fig. 7 shows a similar baryon excess in p+p collisions at the same high rapidity. This time the comparison is made to Pythia simulations.
In summary, particle production from d+Au and p+p collisions at $`\sqrt{s_{NN}}=200GeV`$ and at different rapidities with the BRAHMS setup offers a window to the small-x components of the Au wave function. The suppression found in the particle production at high rapidities from d+Au collisions may be the first indication of the onset of saturation in the gluon distribution function of the Au target.
## 4 Acknowledgments
This work was supported by the Office of Nuclear Physics of the U.S. Department of Energy, the Danish Natural Science Research Council, the Research Council of Norway, the Polish State Committee for Scientific Research (KBN) and the Romanian Ministry of Research.
## References
## Appendix A Kinematics of pA
The Parton model describes hadrons moving at very high momentum as combinations of systems of massles partons in numbers that grow as the energy of the probe increases. Each parton carries a fraction x of the hadron momemtum P. (Because of the very high momentum of the hadron any transverse motion can be neglected.)
What follows is the derivation of the connection between the longitudinal fractions $`x_1`$ and $`x_2`$ of the two partons in the initial state and the rapidity and transverse mass of the measured particle as well as the overall energy of the collision. This derivation is done for the simpler 2-¿1 case but similar results are obtained for the more general 2-¿2 interactions. These derivations were done with much help from Chelis Chasman.
Let S be the total center-of-mass energy squared: $`S=(𝐏_𝐀+𝐏_𝐁)^\mathrm{𝟐}`$ where $`𝐏_𝐀`$ and $`𝐏_𝐁`$ are four vectors. (Naturally the beam momentum defines one prefered direction, and from now on we state that the four-momenta $`𝐏_𝐀`$ and $`𝐏_𝐁`$ have only one component along that direction. If the momenta of the colliding nucleons is high compared to their masses, one can neglect the masses and write:
$$S=(E_A+E_B)^2(\stackrel{}{P_A}+\stackrel{}{P_B})^2=m_A^2+m_B^2+2E_AE_B2\stackrel{}{P_A}\stackrel{}{P_B}2𝐏_𝐀𝐏_𝐁$$
The masses of the nucleons can be neglected (E = P for both nuclei and $`P_A=P_B`$) and one writes in the center of mass of the collision:
$$S=4P_AP_B$$
The fraction x of the hadron’s longitudinal momentum carried by each parton is defined as the ratio of the appropriate light-cone momenta: $`x\frac{e_{parton}+p_{parton}^{}}{E_{hadron}+P_{hadron}}`$ for a left-to-right moving hadron and $`x\frac{e_{parton}p_{parton}^{}}{E_{hadron}P_{hadron}}`$ for the right-to-left hadron. Using the same labels shown in Fig. A1 we call $`x_1`$ the fraction of longitudinal momentum of a beam (proton or deuteron) parton:
$$x_1=\frac{e_b+p_b^{}}{E_B+P_B}=\frac{e_b+p_b^{}}{\sqrt{S}}$$
where we neglected the mass of the hadrons and wrote: $`E=P=\frac{\sqrt{S}}{2}`$.
For the parton in the heavy nuclei (Au) we write the fraction $`x_2`$ emphasizing the fact that the hadron now moves from right-to-left i.e. $`\stackrel{}{P_A}<0`$ :
$$x_2=\frac{e_ap_a^{}}{E_AP_A}=\frac{e_a+p_a^{}}{\sqrt{S}}$$
$$x_1\sqrt{S}=e_b+p_b^{}\text{ and }x_2\sqrt{S}=e_a+p_a^{}$$
$$x_1\sqrt{S}(e_bp_b^{})=m_b^2+k_b^2=(m_{}^b)^2\text{ and }x_2\sqrt{S}(e_ap_a^{})=m_a^2+k_a^2=(m_{}^a)^2$$
where $`k_a`$ and $`k_b`$ are the transverse momentum components of the interacting partons and $`m_{}^a`$ and $`m_{}^b`$ are called transverse masses of the a and b partons respectively. If we multiply the squares of those transverse masses we get:
$$(m_{}^a)^2(m_{}^b)^2=Sx_1x_2(2e_ae_b+2p_a^{}p_b^{}(e_a+p_a^{})(e_b+p_b^{}))$$
where we used the fact that the partons move in opposite directions along the longitudinal axis.
If we define $`\widehat{s}`$ as total energy in the center-of-mass of the interaction partons a and b:
$$\widehat{s}=(e_a+e_b)^2(p_a^{}+p_b^{})^2(k_a+k_b)^2=M^2$$
where M is the mass of the new system formed by the ineraction of the partons a and b. That system has transverse motion given by: $`\stackrel{}{p_T}=\stackrel{}{k_a}+\stackrel{}{k_b}`$ and we can write $`\widehat{s}`$ in the following way:
$$\widehat{s}=(m_{}^a)^2+(m_{}^b)^2+2e_ae_b+2p_a^{}p_b^{}p_T^2$$
$$2e_ae_b+2p_a^{}p_b^{}=\widehat{s}(m_{}^a)^2(m_{}^b)^2+p_T^2$$
when we replace the left side of this equation in the product of transverse masses derived a few lines above, we have:
$$\widehat{s}(m_{}^a)^2(m_{}^b)^2+p_T^2=\frac{(m_{}^a)^2(m_{}^b)^2}{Sx_1x_2}$$
$$\widehat{s}+p_T^2=M_T^2=x_1x_2S+\frac{(m_{}^a)^2(m_{}^b)^2}{Sx_1x_2}+(m_{}^a)^2+(m_{}^b)^2$$
If one neglects the masses of the partons and their transverse momentum motion:
$$x_1x_2S=M_T^2$$
On the other hand, if the system formed at the interaction of partons a and b has rapidity y and longitudinal momentum $`p_L`$, its longitudinal momentum fraction is defined in the center-of-mass frame as: $`x_L=\frac{2p_L}{\sqrt{S}}`$ longitudinal momentum conservation for the 2 to 1 process is written as:
$$x_1x_2=x_L=\frac{2M_T}{\sqrt{s}}sinhy$$
and together with the relation between energy in the nucleon-nucleon centrer of mass S and the one in the parton-parton system, we have a system of two equations with two unknowns $`x_1`$ and $`x_2`$.
$$x_1=\frac{M_T}{\sqrt{S}}e^y$$
$$x_2=\frac{M_T}{\sqrt{S}}e^y$$
It is now clear how the work at high rapidities opens a window into the low values of $`x_2`$ in the target wave function.
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# Equilibrium cluster formation and gelation
## 1 Introduction
The van der Waals theory is the foundation for our current understanding of fluid structure and phase equilibria . It predicts, for an interaction potential consisting of a short-range repulsion and a long-range attraction, a first-order transition (below $`T_c`$) between a dilute vapour and a dense fluid or liquid. Here we consider the unusual situation which pertains when the range of the repulsive and attractive forces are interchanged.
Experiments , theory and simulation suggest that a combination of a very short-ranged attraction and a much longer range repulsion leads to a highly nontrivial phase behaviour. At low densities, there is compelling evidence for the existence of a fluid phase of finite-sized equilibrium aggregates or clusters . The size and shape of these clusters are expected to be very sensitive to the relative range of the attractive and repulsive components ($`V_R(r)`$) of the total interparticle potential. Estimates of the ground state energies of clusters of different sizes suggest that the equilibrium shape changes from spherical to linear as the cluster diameter exceeds the decay length of $`V_R(r)`$ . In addition recent experiments have revealed that, at high densities, the competition inherent in this system can also stabilize new particle architectures. Three-dimensional confocal microscopy in a system of weakly charged colloidal spheres show, at high densities, the formation of a macroscopically percolating gel composed of linear aggregates of face-sharing tetrahedra of particles. The colloidal particles are arranged in the motif of a helical spiral in which each particle is connected to six neighbours, analogous to the structure first proposed by Bernal in his classic study of hard-sphere glasses .
In this paper, we present an experimental study of a carefully-characterized model system of colloidal spheres. The particles interact through an effective potential which consists of a short-range attraction, generated by the exclusion of a non-adsorbing polymer, and a long-range screened electrostatic repulsion. We focus on what happens to the size and shape of the self-assembled clusters as the gel boundary is approached.
## 2 Experimental details
We use a well-studied model system: random-coil polystyrene (radius of gyration $`r_g=92`$ nm) mixed with PMMA spheres (diameter $`\sigma =1700`$ nm) dispersed in a near-density and near-index matched mixture of 78:22 % by weight cycloheptyl bromide and cis-decalin. The PMMA spheres were labelled with the dye DiIC<sub>18</sub> to make them visible by fluorescence microscopy. Exclusion of the non-adsorbing polymer from a region between two neighbouring particles generates a short-range attraction. The range and strength of which is controlled by $`r_g`$ and the free polymer concentration $`\varphi _p^{}`$. For the cases studied here the width of the attractive range is roughly $`0.1\sigma `$. Measurements of the particle electrophoretic mobility using phase analysis light scattering revealed that the PMMA particles develop a small positive charge, $`Q+145e`$, in the density-matched solvent used in this study. The Debye screening length was estimated from measurements of the solvent conductivity ($`170\pm 30`$ pS cm<sup>-1</sup>) as of the order of $`\kappa \sigma 0.5`$.
Mixtures with different colloid ($`\varphi _c`$) and polymer concentrations ($`\varphi _p^{}`$) were prepared and homogenized by extensive tumbling. To minimize decomposition and to ensure reproducibility all samples, once prepared, were stored in light-tight containers at 4C, prior to measurement. The size and internal structure of clusters were studied using laser scanning confocal microscopy. The suspensions were sealed in a cylindrical cell of $`50`$ $`\mu `$L volume. Three-dimensional confocal microscopy was used to distinguish between a phase of isolated clusters and the formation of a macroscopic percolating network of particles. Our observations are summarized in figure 1. To follow the transition between isolated clusters and a gel phase in detail we focus here on the sample sequence A–G, figure 1, at $`\varphi _p^{}=0.95`$. A stack of typically 345 images, spaced by 0.16 $`\mu `$m vertically, was collected from a representative volume of 73 $`\mu `$m x 73 $`\mu `$m x 55 $`\mu `$m within each suspension. Each stack contained, depending on $`\varphi _c`$, the images of between 1500 and 45000 spheres. The images were analyzed quantitatively using algorithms similar to those described in to extract the three-dimensional coordinates of each sphere. The position of each sphere was measured with a precision of approximately 70 nm.
## 3 Data analysis
The number, size and shape of the clusters were identified from the coordinate list. We define particles which are ‘bonded’ to each other by their separation $`r`$. If $`rr_0`$, where the cut-off distance $`r_0`$ was identified with the position of the first minimum in the pair distribution function $`g(r)`$, then the particles were considered neighbours. Standard algorithms were used to partition particles into clusters of size $`s`$ and to evaluate the cluster size distribution $`n_s`$. The second moment of the cluster distribution $`s_2`$ was used to characterize the average cluster size ,
$$s_2=\frac{_ss^2n_s}{_ssn_s}.$$
(1)
The geometric properties of the clusters were evaluated from the eigenvalues $`\lambda _1\lambda _2\lambda _3`$ of the radius of gyration tensor $`\mathrm{S}`$, defined as
$$S_{\alpha \beta }=\frac{1}{2N^2}\underset{i,j=1}{\overset{N}{}}[r_{i,\alpha }r_{j,\alpha }][r_{i,\beta }r_{j,\beta }]$$
(2)
with $`\mathrm{r}_\mathrm{i}`$ the position of the $`i`$th sphere in the cluster of size $`N`$ and $`(\alpha ,\beta )=`$ 1, 2, 3 denoting the corresponding Cartesian components. We characterize the size and shape of the cluster in terms of the rotational invariants of the tensor $`\mathrm{S}`$. The trace $`_i\lambda _i`$ gives the usual isotropic squared radius of gyration $`r_g^2`$ while the anisotropy of the cluster is measured by the asphericity $`A_2`$ ,
$$A_2=\frac{(\lambda _1\lambda _2)^2+(\lambda _2\lambda _3)^2+(\lambda _3\lambda _1)^2}{2(\lambda _1+\lambda _2+\lambda _3)^2}.$$
(3)
The asphericity is normalized so that $`A_2`$ ranges from zero for a spherically symmetric cluster to one for a rod-like cluster.
## 4 Results
The left panel of figure 2 shows the evolution in the proportion of single unbonded particles $`f`$ as a function of $`\varphi _c`$. The boundary between a non-percolating cluster phase and the percolating gel occurs at $`\varphi _g0.1`$. The most striking feature is the abrupt change in the number of single particles approaching the gel boundary. At low colloid concentrations $`f`$ is approximately constant at $`f0.22`$, independent of $`\varphi _c`$ \- the remaining 78% of particles being incorporated into clusters. The proportion of single particles however drops rapidly in the gel phase, although it remains finite even for dense gels ($`f10^2`$ at $`\varphi _c=0.3`$, for instance). The coexistence evident in these samples between free and bonded particles highlights the dynamic nature of the assembly process. Visual observation revealed particles continuously joining and leaving clusters. To provide further evidence that equilibrium had been reached we studied the time-dependence of $`f`$. We found no qualitative change on the period of a week, during which time particle sedimentation was negligible.
Figure 2(b) shows the growth in the average cluster size $`s_2`$, defined in (1), with increasing $`\varphi _c`$ for all non-percolating states. Clearly the the average cluster size increases rapidly approaching the gel boundary. This rapid growth is mirrored in the number of particles $`s_{max}`$ contained in the largest observed cluster, which increases from 37 ($`\varphi _c=0.011`$), to 155 ($`\varphi _c=0.076`$), to reach $`s_{max}=2340`$ at $`\varphi _c=0.086`$. To quantify these changes in greater detail we determined the size dependence of the ensemble-averaged radius of gyration $`r_g`$. For fractal aggregates $`r_g`$ displays a power law dependence on the number of particles $`s`$, $`r_gs^{1/d_f}`$, where $`d_f`$ is the fractal dimension. Figure 3 shows plots of $`r_g`$ versus $`s`$, at three values of $`\varphi _c`$. Interestingly we find that at all densities the clusters show fractal scaling but that the fractal dimension $`d_f`$ shows a marked dependence on $`\varphi _c`$. At low $`\varphi _c`$, the fractal dimension found of $`d_f=1.73\pm 0.07`$ is consistent with the diffusion-limited cluster-cluster aggregation (DLCA) value in three dimensions ($`d_f=1.78`$ ). While at high $`\varphi _c`$, the value of $`d_f=2.14\pm 0.03`$ is consistent with the reaction-limited cluster-cluster aggregation (RLCA) value in three dimensions ($`d_f=2.1`$ ). This apparent crossover between DLCA and RLCA kinetics with $`\varphi _c`$ is unexpected. However it is unclear if fractal models for non-equilibrium growth are applicable to the equilibrium case of assembly studied here for which, as we indicate below, the cluster shape is not truly self-similar.
A simple indicator of the anisotropy of the clusters is provided by the asphericity $`A_2`$, defined in (3). In contrast to the $`r_g`$ data we find little $`\varphi _c`$ dependence of the average cluster asphericity, plotted in figure 4(a). Both low and high density samples show a similar monotonic decrease in $`A_2`$ with $`s`$, consistent with a progressive change in shape with increasing cluster size. The clusters formed are therefore not self-similar. Clusters smaller than about 30–50 monomers show a definite tendency to be elongated in one dimension as evidenced by the increasing value of $`A_2`$ found at small $`s`$. This preference for a linear or rod-like shape seems to be more pronounced the smaller the cluster. These observations are supported by images of individual clusters which reveal characteristic one-dimensional bundles of particles. By contrast, while large clusters remain anisotropic the levels of asymmetries are significantly lower and typical of clusters formed from either random percolation ($`A_2=0.25`$ ) or DLCA models $`A_2=0.32`$ ). This crossover suggests that in large clusters the short one-dimensional bundles have branched many times generating clusters with geometries controlled by random percolation.
The tendency seen for one-dimensional aggregation at short scales might indicate that the cluster is constructed locally from segments of Bernal spiral structures. To check for this possibility, we calculated the average number of nearest neighbours $`n`$. As shown in the right hand panel of figure 4, $`n`$ grows abruptly with increasing $`\varphi _c`$ approaching the gel boundary. Above $`\varphi _g`$, the mean coordination number is insensitive to $`\varphi _c`$ and is approximately six, consistent with the characteristic six-fold coordination of a Bernal spiral. Moreover figure 4 reveals that $`n`$ is less than six in the non-percolating cluster samples (A–C). This is particularly evident in sample B ($`\varphi _c=0.076`$) which while it contains clusters of upto 100 monomers has a mean coordination number of only $`n=2.7`$. The lack of spiral ordering, shown by the low values of $`n`$, is also supported by direct inspection which reveals little evidence for Bernal ordering in any of the non-percolating cluster samples.
## 5 Conclusions
In this work we have explored the change in the size and shape of the clusters formed in a colloidal system of short range attractions and long range repulsions, approaching the gelation transition. In contrast to a purely attractive system the clusters are equilibrium entities and their properties time-independent. Our most noteworthy finding is that the average size $`s_2`$ of the cluster distribution, the proportion of unattached particles, and the mean coordination number all change abruptly as $`\varphi \varphi _g`$. Although we have studied only a limited number of concentrations (at one specific choice of potential parameters) these observation suggest that, for this system, gelation has many of the hallmarks of an equilibrium transition. Qualitatively our observations appear consistent with the suggestion by Groenewold and Kegel that under the right conditions a fluid of isolated clusters will undergo a spinodal decomposition, driven by counter-ion condensation, to form a macroscopic gel. This raises the enticing possibility of an ideal thermally-reversible gel – one in which the gel-line might be approached infinitely closely and the dynamics of arrest studied precisely and reproducibly, free from the dynamical complications of purely attractive systems. Further experiments are planned to test these ideas.
It is a pleasure to thank Drs A I Cambell, J S van Duijneveldt and M Faers for assistance. This work was supported financially by EPSRC and MCRTN-CT-2003-504712.
## References
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# On weak maps between 2-groups
## 1. Introduction
There are several incarnations of 2-groups in mathematics. To mention a few:<sup>1</sup><sup>1</sup>1To see some more explicit example see the last section of \[BaLa\]. 1) They appear abstractly as special classes of monoidal categories, e.g., in the context of bitorsors \[Bre1\] or Picard categories \[SGA4\]; 2) They classify connected pointed homotopy 2-types (via the fundamental 2-group), as discovered by MacLane and Whitehead; 3) They appear as symmetries of objects in a 2-category. Let us dissect (3) a bit further by giving some examples.
The auto-equivalences of a stack from a 2-group. Thus, 2-groups play an important rule in the study of (2-) group actions on stacks \[BeNo\]. Along similar lines, 2-groups appear in the representation theory of 2-vector spaces \[El\], \[GaKa\], and also in the theory of gerbes and higher local systems \[Al\], \[Deb\], \[Bre1-4\], \[BrMe\], \[PoWa\], \[BDR\]. The latter is indeed deeply connected with physics through higher gauge theory: various kind of gerbes that come up in gauge theory are to be thought of as 2-principal-bundles for a certain 2-group. For example: Giraud’s $`G`$-gerbes are 2-principal-bundles for the 2-group $`𝔄𝔲𝔱(G)`$ \[Bre1\]; $`S^1`$-gerbes \[Bry\] are 2-principal-bundles for the 2-group whose 2-morphisms are $`S^1`$ and has no nontrivial 1-morphisms; string bundles are essentially the same as 2-principal-bundles for the string 2-groups \[BCSS\], \[BaSch\]. Loop groups are also closely related to 2-groups \[BCSS\].
Most 2-groups that arise in nature are weak, in the sense that, either the multiplication is only associative up to higher coherences, or inverses exist only in a weak sense (or both). Furthermore, interesting morphisms between 2-groups are also often weak, in the sense that they preserve products only up to higher coherences. It is a standard fact that one can strictify weak 2-groups, but not the weak functors. Put differently, the 2-category of strict 2-groups and weak functors between them is the “homotopically correct” habitat for 2-groups.
The (strict) 2-groups by themselves can be codified conveniently using crossed-modules. Weak morphisms between strict 2-groups, however, are more complicated and to write them down results in somewhat inconvenient cocycles. It is therefore desirable to find a clean and way to deal with weak morphisms between 2-groups so as to make them more tractable in geometric situations.
The aim of this paper is to do exactly this. Namely, we give a concrete and manageable cocycle-free model for the space of weak morphisms between two crossed-modules. (It is easy to see that this space is a 1-type, so its homotopy type is described by a groupoid.) Theorem 8.4 (also see Theorem 13.1) gives us a functorial model for this groupoid in terms of what we call butterflies. Butterflies indeed furnish a neat bicategory structure on crossed-modules (Theorem 10.1), therefore giving rise to a model for the homotopy category of pointed connected 2-types. We also discuss the braided and abelian versions of this bicategory.
In a future paper we generalize these result to the case where everything is relative to a Grothendieck site. As a consequence, the main results of the present paper will apply to the case of Lie 2-groups, 2-group schemes, and so on. Bear in mind that in these geometric settings the cocycle approach to weak morphisms is very inconvenient and sometimes hopeless (see below).
Some applications.
For an application of butterflies to the HRS-tilting theory \[HRS\] we refer the reader to \[No2\]. In another application (joint work with E. Aldrovandi), we investigate butterflies from the point of view of gr-stacks. We also employ butterflies (over a Grothendieck site) to study the “change of the structure 2-group along a weak morphism of 2-groups” for 2-principal-bundles and its effect on their geometry. (This has been previously looked at only for strict 2-group morphism \[Bre1\].)
2-group actions on stacks. Let us spell out in more detail an application of our approach to weak morphisms – this will also serve to explain why a cocycle-free approach could be advantageous sometimes.
We propose a systematic way to study 2-group actions on stacks. Given a stack $`𝒳`$ (say, topological, differentiable, analytic, algebraic, etc.), the set $`\mathrm{Aut}𝒳`$ of self-equivalences of $`𝒳`$ is naturally a weak 2-group. (It is weak because equivalences are not strictly invertible; the associativity, however, remains strict.) To have a 2-group $``$ act on $`𝒳`$ is the same thing as to have a weak map from $``$ to $`\mathrm{Aut}𝒳`$. If two such maps are related by a (pointed) transformation, they should be regarded as giving the “same” action of $``$ on $`𝒳`$. So the question is to classify such equivalence classes of weak maps $`f:\mathrm{Aut}𝒳`$.
The weakness of $`\mathrm{Aut}𝒳`$, and of the map $`f`$, are, however, disturbing and we would like to make things as strict as possible. Using a bit of homotopy theory, and some standard strictification procedures, it can be shown that what we are looking for is $`[,𝔾]_{\mathrm{𝟐}𝐆𝐩}`$, where $`𝔾`$ is a strict model for $`\mathrm{Aut}𝒳`$. Here, $`[,𝔾]_{\mathrm{𝟐}𝐆𝐩}`$ stands for the set of morphisms from $``$ to $`𝔾`$ in the homotopy category $`\mathrm{Ho}(\mathrm{𝟐}𝐆𝐩)`$ of the category of strict 2-groups and strict maps. Equivalence of 2-groups and crossed-modules ($`\mathrm{\S }`$ 3.3) now enables us to translate the problem to the language of crossed-modules, in which case we have an explicit description of $`[,𝔾]_{\mathrm{𝟐}𝐆𝐩}`$ in simple group theoretic terms thanks to Theorem 8.4 (also see Theorem 13.1 and Corollary 13.2).
All we need is to run this method is to find a crossed-module model for $`\mathrm{Aut}𝒳`$
Thanks to the very explicit nature of the above procedure, we are able to give solid constructions with stacks, circumventing a lot of “weaknesses” and coherence conditions that arise in studying group actions on stacks. An application of this strictification method is given in \[BeNo\], where the covering theory of stacks is used to classify smooth Deligne-Mumford analytic curves, and also to give an explicit description of them as quotient stacks. For instance, using these 2-group theoretic techniques we obtain the following completely geometric result: every smooth analytic (respectively, algebraic) Deligne-Mumford stack of dimension one is the quotient stack for the action of either a finite group or a central finite extension of $`^{}`$ on a complex manifold (respectively, complex variety).
Organization of the paper
Sections 3 to 6 are devoted to recalling some standard facts about 2-groups and crossed-modules and fixing the notation. Essential for more easily reading the paper is the fact that the category of 2-groups is equivalent to the category of crossed-modules ($`\mathrm{\S }`$3.3). The reader will find it beneficial to keep in mind how this equivalence works, as we will freely switch back and forth between 2-groups and crossed-modules throughout the paper (sometimes even using the two terms synonymously). For us, 2-groups are the conceptual side of the story, whereas crossed-modules provide the computational framework.
Viewed as 2-groupoids with one object, 2-groups (hence, also crossed-modules) can be treated via the Moerdijk-Svensson model structure \[MoSe\]. This is briefly recalled in Section 6. We point out that, all we need from closed model categories is the notion of fibrant/cofibrant resolution and the way it can be used to compute hom-sets in the homotopy category. Taking this for granted, the reader unfamiliar with closed model categories can proceed with no difficulty.
Section 7 concerns some elementary constructions from group theory. We introduce a pushout construction for crossed-modules and work out its basic properties. This section is perhaps is not so interesting by itself, but it provides the technical tools required in the proof of Theorem 8.4.
Section 8 is the core of the paper. In it we state and prove our main result (Theorem 8.4). This is based on the notion of butterfly (Definition 8.1). We investigate butterflies some more in Section 9. In $`\mathrm{\S }`$10 we give an explicit model $`𝒞`$ for the 2-category of crossed-modules and weak morphisms in terms of butterflies.
In Section 11 we indicate how the notions of kernel, cokernel, exact sequence, and so on of weak morphisms of 2-groups find a natural and simple form in the world of butterflies.
Braided and abelian butterflies are briefly discussed in Section 12. We use the latter to give a simple description of the derived category of complexes of length 2 in an abelian category $`𝖠`$.
In Section 13 we consider a special case of Theorem 8.4 in which the source 2-group is an honest group (Theorem 13.1) and discuss its connection with results of Dedecker and Blanco-Bullejos-Faro. In Section 14 we discuss a cohomological version of this (Theorem 14.6). This cohomological classification is not a new result (with some diligence, the reader can verify that it is a special case of the work of \[AzCe\]), and our emphasis is only to make precise the way it relates to Theorem 13.1, as it was used in \[BeNo\].
In Sections 15 and 16 we explain the homotopical meaning of Theorem 13.1 in terms the Postnikov decomposition of the classifying space of a crossed-module; again, this is folklore. We make this precise using the notion of difference fibration which is an obstruction theoretic construction introduced in \[Ba\] used in studying the liftings of a map into a fibration (from the base to the total space).
In the appendix we review basic general facts about 2-categories and 2-groupoids.
Acknowledgement. I am grateful to E. Aldrovandi, H-J. Baues, L. Breen, J. Elgueta, M. Jibladze, M. Kamgarpour, and F. Muro for many useful discussions. I would like to thank Bertrand Toën for pointing me to Elgueta’s work \[El\]. The idea of Theorem 13.1 was conceived in conversations with Kai Behrend in our joint work \[BeNo\]. The relation between butterflies and cocycles for weak morphisms was pointed out to me by Roman Mikhailov. Finally, I would like to thank Max-Planck-Institut für Mathematik, Bonn, where the research for this work was done, for providing pleasant working conditions.
## 2. Notation and terminology
We list some of the notations and conventions used throughout the paper.
The structure map of a crossed-module $`𝔾=[G_2G_1]`$ is usually denoted by $`\alpha \underset{¯}{\alpha }`$. The components of a morphism $`P:𝔾`$ of crossed-modules are denoted by $`p_2:H_2G_2`$ and $`p_1:H_1G_1`$. The action of $`G_1`$ on $`G_2`$ is denoted by $`^a`$, and so is the conjugation action of $`G_1`$ on itself.
In a semi-direct product $`AB`$, respectively $`BA`$, the group $`B`$ acts on $`A`$ on the left, respectively right.
For objects $`A`$ and $`B`$ in a category $`𝐂`$ with a notion of weak equivalence, we denote the the set of morphisms in the homotopy category from $`B`$ to $`A`$, that is $`\mathrm{Hom}_{\mathrm{Ho}(𝐂)}(B,A)`$, by $`[B,A]_𝐂`$.
We usually denote a short exact sequence
$$1NE\mathrm{\Gamma }1$$
simply by $`E`$. Quotient maps, such as $`E\mathrm{\Gamma }`$ in the above sequence, or $`G_1\pi _1𝔾`$, are usually denoted by $`x\overline{x}`$.
Some abuse of terminology.
We tend to use the term “map” where it is perhaps more appropriate to use the term “morphism”. Also, we use the term “equivalence” (of 2-groups, crossed-modules, etc.) for what should really be called a “weak equivalence” (or “quasi-isomorphism”).
For 2-groups $`𝔾`$ and $``$, when we say the homotopy class of a weak map from $``$ to $`𝔾`$, we mean a map in $`\mathrm{Ho}(\mathrm{𝟐}𝐆𝐩)`$ from $``$ to $`𝔾`$; this terminology is justified by Theorem 6.7 and Proposition 6.8.
## 3. Quick review of 2-groups and crossed-modules
### 3.1. Quick review of 2-groups
We recall some basic facts about 2-groups and crossed-modules. Our main references are \[Bro, Wh, McWh, MoSe, Lo, BaLa\].
A 2-group $`𝔊`$ is a group object in the category of groupoids. Alternatively, we can define a 2-group to be a groupoid object in the category of groups, or also, as a (strict) 2-category with one object in which all 1-morphisms and 2-morphisms are invertible (in the strict sense). We will try to stick with the ‘group in groupoids’ point of view throughout the paper, but occasionally switching back and forth between different points of view is inevitable. Therefore, the reader will find it rewarding to master how the equivalence of these three point of views works.
A (strict) morphism $`f:𝔊`$ of 2-groups is a map of groupoids that respects the group operation. If we view $`𝔊`$ and $``$ as 2-categories with one object, such $`f`$ is nothing but a strict 2-functor. The category of 2-groups is denoted by $`\mathrm{𝟐}𝐆𝐩`$.
To a 2-group $`𝔊`$ we associate the groups $`\pi _1𝔊`$ and $`\pi _2𝔊`$ as follows. The group $`\pi _1𝔊`$ is the set of isomorphism classes of object of the groupoid $`𝔊`$. The group structure on $`\pi _1𝔊`$ is induced from the group operation of $`𝔊`$. The group $`\pi _2𝔊`$ is the group of automorphisms of the identity object $`e𝔊`$. This is an abelian group. A morphism $`f:𝔊`$ of 2-groups is called an equivalence if it induces isomorphisms on $`\pi _1`$ and $`\pi _2`$. The homotopy category of 2-groups is the category obtained by inverting all the equivalences in $`\mathrm{𝟐}𝐆𝐩`$. We denote it by $`\mathrm{Ho}(\mathrm{𝟐}𝐆𝐩)`$.
Caveat: an equivalence between 2-groups need not have an inverse. Also, two equivalent 2-groups may not be related by an equivalence, but only a zig-zag of equivalences.
### 3.2. Quick review of Crossed-modules
A crossed-module $`𝔊=[G_2G_1]`$ is a pair of groups $`G_1,G_2`$, a group homomorphism $`_𝔾:G_2G_1`$, and a (right) action of $`G_1`$ on $`G_2`$, denoted $`^a`$, which lifts to $`G_2`$ the conjugation action of $`G_1`$ on the image of $`_𝔾`$ and descends to $`G_1`$ the conjugation action of $`G_2`$ on itself. The kernel of $``$ is a central (in particular abelian) subgroup of $`G_2`$ and is denoted by $`\pi _2𝔊`$. The image of $``$ is a normal subgroup of $`G_1`$ whose cokernel is denoted by $`\pi _1𝔊`$. A (strict) morphism of crossed-modules is a pair of group homomorphisms which commute with the $``$ maps and respect the actions. A morphism is called an equivalence if it induces isomorphisms on $`\pi _1`$ and $`\pi _2`$.
Notation. Elements of $`G_2`$ are usually denoted by Greek letters and those of $`G_1`$ by lower case Roman letters. The components of a map $`P:𝔾`$ of crossed-modules are denoted by $`p_2:H_2G_2`$ and $`p_1:H_1G_1`$. We usually use $``$ instead of $`_𝔾`$ if the 2-group $`𝔾`$ is clear from the context. We sometimes suppress $``$ from the notation and denote $`(\alpha )`$ by $`\underset{¯}{\alpha }`$. For elements $`g`$ and $`a`$ in a group $`G`$ we sometimes denote $`a^1ga`$ by $`g^a`$. The compatibility assumptions built in the definition of a crossed-module make this unambiguous. With this notation, the two compatibility axioms of a crossed-modules can be written in the following way:
* $`\alpha ,\beta G_2,\beta ^{\underset{¯}{\alpha }}=\beta ^\alpha `$;
* $`\beta G_2,aG_1,\underset{¯}{\beta }^a=\underset{¯}{\beta ^a}`$.
### 3.3. Equivalence of 2-groups and crossed-modules
There is a natural pair of inverse equivalences between the category $`\mathrm{𝟐}𝐆𝐩`$ of 2-groups and the category $`\mathrm{𝐗𝐌𝐨𝐝}`$ of crossed-modules. Furthermore, these functors preserve $`\pi _1`$ and $`\pi _2`$. They are constructed as follows.
Functor from 2-groups to crossed-modules. Let $`𝔊`$ be a 2-group. Let $`G_1`$ be the group of objects of $`𝔊`$, and $`G_2`$ the set of arrows emanating from the identity object $`e`$; the latter is also a group (namely, it is a subgroup of the group of arrows of $`𝔊`$).
Define the map $`:G_2G_1`$ by sending $`\alpha G_2`$ to $`t(\alpha )`$.
The action of $`G_1`$ on $`G_2`$ is given by conjugation. That is, given $`\alpha G_2`$ and $`gG_1`$, the action is given by $`g^1\alpha g`$. Here were are thinking of $`g`$ as an identity arrow and multiplication takes place in the group of arrows of of $`𝔾`$. It is readily checked that $`[:G_2G_1]`$ is a crossed-module.
Functor from crossed-modules to 2-groups. Let $`[:G_2G_1]`$ be a crossed-module. Consider the groupoid $`𝔊`$ whose underlying set of objects is $`G_1`$ and whose set of arrows is $`G_1G_2`$. The source and target maps are given by $`s(g,\alpha )=g`$, $`t(g,\alpha )=g(\alpha )`$. Two arrows $`(g,\alpha )`$ and $`(h,\beta )`$ such that $`g(\alpha )=h`$ are composed to $`(g,\alpha \beta )`$. Now, taking into account the group structure on $`G_1`$ and the semi-direct product group structure on $`G_1G_2`$, we see that $`𝔊`$ is indeed a groupoid object in the category of groups, hence a 2-group.
The above discussion shows that there is a pair of inverse functors inducing an equivalence between $`\mathrm{𝐗𝐌𝐨𝐝}`$ and $`\mathrm{𝟐}𝐆𝐩`$. These functors respect $`\pi _1`$ and $`\pi _2`$. Therefore, we have an equivalence
## 4. Transformations between morphisms of crossed-modules
We go over the notions of transformation and pointed transformation between maps of crossed-modules. These are adaptations of the usual 2-categorical notions, translated to the crossed-module language via the equivalence $`\mathrm{𝐗𝐌𝐨𝐝}\mathrm{𝟐}𝐆𝐩`$ (see Appendix, $`\mathrm{\S }`$17.1). The idea is to think of a crossed-module as a 2-group ($`\mathrm{\S }`$3.3), which is itself thought of as a 2-groupoid with one object.
###### Definition 4.1.
Let $`𝔾=[G_2G_1]`$ and $`=[H_2H_1]`$ be crossed-modules, and let $`P,Q:𝔾`$ be morphisms between them. A transformation $`T:QP`$ consists of a pair $`(a,\theta )`$ where $`aG_1`$ and $`\theta :H_1G_2`$ is a crossed homomorphism for the induced action, via $`p_1`$, of $`H_1`$ on $`G_2`$ (that is, $`\theta (hh^{})=\theta (h)^{p_1(h^{})}\theta (h^{})`$). We require the following:
* $`q_1(h)^a=p_1(h)\underset{¯}{\theta (h)}`$, for every $`hH_1`$;
* $`q_2(\beta )^a=p_2(\beta )\theta (\underset{¯}{\beta })`$, for every $`\beta H_2`$.
We say $`T`$ is pointed if $`a=1`$; in this case, we denote $`T`$ simply by $`\theta `$. When $`\theta `$ is the trivial map, the transformation $`T`$ is called conjugation by $`a`$; in this case, we use the notation $`P=Q^a`$ or $`P=a^1Qa`$.
###### Remark 4.2.
Given a 2-group $`𝔊`$ and an element $`a`$ in $`G_1`$ (the group of objects) we define the morphism $`c_a:𝔾𝔾`$, called conjugation by $`a`$, to be the map that sends an object $`g`$ (respectively, an arrow $`\alpha `$) to $`a^1ga`$ (respectively, $`a^1\alpha a`$). If we consider the corresponding crossed-module $`[G_2G_1]`$, the conjugation morphism $`c_a`$ sends $`gG_1`$ to $`a^1ga`$ and $`\alpha G_2`$ to $`\alpha ^a`$. In the notation of Definition 4.1, it is easy to see that $`Q^a=c_aQ`$.
###### Lemma 4.3.
Let $`P,Q:𝔾`$ be maps of 2-groups and $`(a,\theta )`$ a transformation from $`Q`$ to $`P`$. Then $`\pi _iP=(\pi _iQ)^a`$, $`i=1,2`$. In particular, if $`P`$ and $`Q`$ are related by a pointed transformation, then they induce the same map on homotopy groups.
###### Proof.
Obvious. ∎
Let $`P,Q,R:𝔾`$ be maps of crossed-modules. Given homotopies $`(b,\sigma ):RQ`$ and $`(a,\theta ):QP`$, consider the pointwise product $`\theta \sigma :H_1G_2`$. It is easily checked that $`(ba,\theta \sigma )`$ is a transformation from $`R`$ to $`P`$. This construction, of course, corresponds to the usual composition of weak 2-transformation between 2-functors.
A transformation $`(a,\theta ):QP`$ has an inverse $`(a^1,\theta ^1):PQ`$, where $`\theta ^1:H_1G_2`$ is defined by $`\theta ^1(h):=\theta (h)^1`$.
###### Definition 4.4.
Let $`𝔾`$ and $``$ be crossed-modules. We define the mapping groupoid $`\underset{¯}{om}_{}(,𝔾)`$ to be the groupoid whose objects are crossed-module maps $`𝔾`$ and whose morphisms are pointed transformations.
###### Remark 4.5.
Observe that in the definition above of $`\underset{¯}{om}_{}(,𝔾)`$ we have not used ‘modifications’ and, in particular, the outcome is a groupoid and not a 2-groupoid. This is because between two pointed transformations there is no non-trivial pointed modification.
With hom-groupoids being $`\underset{¯}{om}_{}(,𝔾)`$, the category $`\mathrm{𝐗𝐌𝐨𝐝}`$ is enriched over groupoids. We denote the resulting 2-category by $`\underset{¯}{\mathrm{𝐗𝐌𝐨𝐝}}`$. Similarly, $`\mathrm{𝟐}𝐆𝐩`$ can be enriched over groupoids by taking 2-morphisms to be pointed transformations (Definition 4.1). We denote the resulting 2-category by $`\underset{¯}{\mathrm{𝟐}𝐆𝐩}`$. The equivalence
$$\mathrm{𝐗𝐌𝐨𝐝}\mathrm{𝟐}𝐆𝐩$$
of $`\mathrm{\S }`$3.3 now becomes a biequivalence of 2-categories
###### Proposition 4.6.
The construction of $`\mathrm{\S }`$3.3 gives rise to a biequivalence
$$\underset{¯}{\mathrm{𝐗𝐌𝐨𝐝}}\underset{¯}{\mathrm{𝟐}𝐆𝐩}.$$
of 2-categories.
The biequivalence of Proposition 4.6 is a biequivalence in a strong sense: it induces isomorphisms on hom-groupoids.
The mapping space $`\underset{¯}{om}_{}(,𝔾)`$ in not the homotopically “correct” mapping space, as it lacks the expected homotopy invariance property. That is, an equivalence $`^{}`$ of crossed-modules does not necessarily induce an equivalence $`\underset{¯}{om}_{}(,𝔾)\underset{¯}{om}_{}(^{},𝔾)`$ of groupoids. We explain in $`\mathrm{\S }`$6 how this failure can be fixed by making use of cofibrant replacements in the category of crossed-modules (especially, see Definition 6.6).
## 5. Weak morphisms between 2-groups
A 2-group $`𝔾`$ can equivalently be defined to be a strict4 monoidal groupoid in which multiplication (on the left, and on the right) by any object is an isomorphism of categories. Given two 2-groups $`𝔾`$ and $``$, the mapping groupoid $`\underset{¯}{om}_{}(,𝔾)`$ is then naturally isomorphic to the groupoid whose object are strict monoidal functors and whose morphisms are natural monoidal transformations.
We define a weak morphism $`f:𝔾`$ to be a monoidal functor which respects the unit objects strictly (also see \[No1\], $`\mathrm{\S }`$7 and $`\mathrm{\S }`$8). With monoidal transformations between them, these are objects of a groupoid which we denote by $`\underset{¯}{𝒪}_{}(,𝔾)`$.<sup>2</sup><sup>2</sup>2For the interpretation of weak morphisms in the crossed-module language see \[No1\], $`\mathrm{\S }`$8.
The goal of the paper is to give an explicit model for $`\underset{¯}{𝒪}_{}(,𝔾)`$. Our strategy is to use the fact that $`\underset{¯}{𝒪}_{}(,𝔾)`$ is equivalent to the derived mapping groupoid $`\underset{¯}{om}_{}(,𝔾)`$; see Theorem 6.7. We then give an explicit model for the derived mapping space $`\underset{¯}{om}_{}(,𝔾)`$ in terms of butterflies (Definition 8.1 and Theorem 8.4).
## 6. Moerdijk-Svensson closed model structure and crossed-modules
It has been known since \[Wh\] that crossed-modules model pointed connected homotopy 2-types. That is, the pointed homotopy type of a connected pointed CW-complex with $`\pi _iX=0`$, $`i3`$, is determined by (the equivalence class of) a crossed-module. In particular, the homotopical invariants of such a CW-complex can be read off from the corresponding crossed-module.
The approach in \[Wh\] and \[McWh\] to the classification of 2-types is, however, not functorial. To have a functorial classification of homotopy 2-types (i.e. one that also accounts for maps between such objects), it is best to incorporate closed model categories. To do so, recall that a crossed-module can be regarded as a 2-group, and a 2-group is in turn a 2-groupoid with one object.
In \[MoSe\], Moerdijk and Svensson introduce a closed model structure on the (strict) category of (strict) 2-groupoids, and show that there is a Quillen pair between the closed model category of 2-groupoids and the closed model category of CW-complexes, which induce an equivalence between the homotopy category of 2-groupoids and the homotopy category of CW-complexes with vanishing $`\pi _i`$, $`i3`$. We use this model structure to deduce some results about crossed-modules.
We emphasize that, in working with crossed-modules, what we are using is the pointed homotopy category. So we need to adopt a pointed version of the Moerdijk-Svensson structure. But this does not cause any additional difficulty as everything in \[MoSe\] carries over to the pointed case. For a quick review of the Moerdijk-Svensson structure see Appendix.
It is easy to see that a weak equivalence between 2-groups in the sense of Moerdijk-Svensson is the same as a weak equivalence between crossed-modules in the sense of $`\mathrm{\S }`$3. Let us see what the fibrations look like.
###### Definition 6.1.
A map $`(f_2,f_1):[H_2H_1][G_2G_1]`$ of crossed-modules is called a fibration if $`f_2`$ and $`f_1`$ are both surjective. It is called a trivial fibration if, furthermore, the map $`H_2H_1\times _{G_1}G_2`$ is an isomorphism.
We leave it to the reader to translate these to the language of 2-groups and verify that they coincide with Moerdijk-Svensson definition of (trivial) fibration.
Let us now look at cofibrations. In fact, we will only describe what the cofibrant objects are, because that is all we need in this paper.
###### Definition 6.2.
A crossed-module $`[G_2G_1]`$ is cofibrant if $`G_1`$ is a free group.
Observe that this is much weaker than Whitehead’s notion of a free crossed-module. However, this is the one that corresponds to Moerdijk-Svensson’s definition.
###### Proposition 6.3.
A crossed-module $`𝔊=[G_2G_1]`$ is cofibrant in the sense of Definition 6.2 if and only if its corresponding 2-group is cofibrant in the Moerdijk-Svensson structure.
###### Proof.
This follows immediately from the Remark on page 194 of \[MoSe\], but we give a direct proof. A 2-group $`𝔾`$ is cofibrant in Moerdijk-Svensson structure, if and only if every trivial fibration $`𝔾`$, where $``$ is a 2-groupoid, admits a section. But, we can obviously restrict ourselves to 2-groups $``$. So, we can work entirely within crossed-modules, and use the notion of trivial fibration as in Definition 6.1.
Assume $`G_1`$ is free. Let $`(f_2,f_1):[H_2H_1][G_2G_1]`$ be a trivial fibration. Since $`G_1`$ is free and $`f_1`$ is surjective, there is a section $`s_1:G_1H_1`$. Using the fact that $`H_2H_1\times _{G_1}G_2`$, we also get a natural section $`s_2:G_2H_2`$ for the projection $`H_1\times _{G_1}G_2G_2`$, namely, $`s_2(\alpha )=(s_1(\underset{¯}{\alpha }),\alpha )`$. It is easy to see that $`(s_2,s_1):[G_2G_1][H_2H_1]`$ is a map of crossed-modules.
To prove the converse, choose a free group $`F_1`$ and a surjection $`f_1:F_1G_1`$. Form the pull back crossed-module $`[F_2F_1]`$ by setting $`F_2=F_1\times _{G_1}G_2`$. Then, we have a trivial fibration $`[F_2F_1][G_2G_1]`$. By assumption, this has a section, so in particular we get a section $`s_1:G_1F_1`$ which embeds $`G_1`$ as a subgroup of $`F_1`$. It follows from Nielsen’s theorem that $`G_1`$ is free. ∎
###### Remark 6.4.
It is easy to see that a 2-group $`𝔾`$ is cofibrant in the Moerdijk-Svensson structure if an only if the inclusion $`𝔾`$ is a cofibration. So the definition of cofibrant is the same in the pointed category. Also, in the pointed category, all 2-groupoids are fibrant.
###### Example 6.5.
* Let $`𝔾=[G_2G_1]`$ be an arbitrary crossed-module. Let $`F_1G_1`$ be a surjective map from a free group $`F_1`$, and set $`F_2:=F_1\times _{G_1}G_2`$. Consider the crossed-module $`𝔽=[F_2F_1]`$. Then $`𝔽`$ is cofibrant, and the natural map $`𝔽𝔾`$ is a trivial fibration (Definition 6.1). In other words, $`𝔽𝔾`$ is a cofibrant replacement for $`𝔾`$.
* Let $`\mathrm{\Gamma }`$ be a group, and $`F/R\mathrm{\Gamma }`$ be a presentation of $`\mathrm{\Gamma }`$ as a quotient of a free group $`F`$. Then the map of crossed-modules $`[RF][1\mathrm{\Gamma }]`$ is a cofibrant replacement for $`\mathrm{\Gamma }`$.
###### Definition 6.6.
Let $``$ and $`𝔾`$ be 2-groups (or crossed-modules). Choose a cofibrant replacement $`𝔽`$ for $``$, as in Example 6.5.$`\mathrm{𝟏}`$. The derived mapping groupoid $`\underset{¯}{om}_{}(,𝔾)`$ is defined to be $`\underset{¯}{om}_{}(𝔽,𝔾)`$, where $`\underset{¯}{om}_{}`$ is as in Definition 4.4.
Observe that in the Moerdijk-Svensson structure all 2-groups are automatically fibrant, so in the above definition we do not need a fibrant replacement for $`𝔾`$.
The derived mapping groupoid $`\underset{¯}{om}_{}(,𝔾)`$ depends on the choice of the cofibrant replacement $`𝔽`$, but it is unique up to an equivalence of groupoids (which is itself unique up to transformation). Another way of thinking about the derived mapping groupoid $`\underset{¯}{om}_{}(,𝔾)`$ is that it gives a model for the groupoid of weak morphisms from $``$ to $`𝔾`$ and pointed weak transformations between them.<sup>3</sup><sup>3</sup>3Since we are working in the pointed category, the modification are trivial. This is due to he fact that, whenever $`X`$ and $`Y`$ are pointed connected homotopy 2-types, the pointed mapping space $`\mathrm{𝐇𝐨𝐦}_{}(X,Y)`$ is a 1-type.
###### Theorem 6.7 (\[No1\], Proposition 8.1).
Let $`𝔾`$ and $``$ be 2-groups. Then, there is a natural (up to homotopy) equivalence of groupoids
$$\underset{¯}{om}_{}(,𝔾)\underset{¯}{𝒪}_{}(,𝔾).$$
The fact that derived mapping groupoids are the correct models for homotopy invariant mapping spaces is justified by the following.
###### Proposition 6.8.
Let $``$ and $`𝔾`$ be 2-groups (or crossed-modules). We have a natural bijection
$$\pi _0\underset{¯}{om}_{}(,𝔾)[,𝔾]_{\mathrm{𝟐}𝐆𝐩}.$$
###### Proof.
First let us remark that the proposition is not totally obvious because the category of 2-groupoids with the Moerdijk-Svensson model structure is not a monoidal model category.
By Proposition 17.12, we have
$$\pi _0\underset{¯}{om}_{}(,𝔾)\pi _0\mathrm{𝐇𝐨𝐦}_{}(N,N𝔾)[N,N𝔾]_{\mathrm{𝐒𝐒𝐞𝐭}_{}}.$$
By Proposition 17.9, $`[N,N𝔾]_{\mathrm{𝐒𝐒𝐞𝐭}_{}}[,𝔾]_{\mathrm{𝟐}𝐆𝐩}`$. ∎
In $`\mathrm{\S }`$8 we give a canonical explicit model for $`\underset{¯}{om}_{}(,𝔾)`$, which is what we want.
## 7. Some group theory
In this section we introduce some basic group theoretic lemmas which will be used in the proof of Theorem 8.4. The results in this section are mostly of technical nature.
### 7.1. Generalized semi-direct products
We define a generalized notion of semi-direct product of groups, and use that to introduce a pushout construction for crossed-modules.
Let $`H`$, $`G`$ and $`K`$ be groups, each equipped with a right action of $`K`$, the one on $`K`$ itself being conjugation. We denote all the actions by $`^k`$ (even the conjugation one). Assume we are given a $`K`$-equivariant diagram
in which we require the compatibility condition $`g^{d(h)}=g^{p(h)}\left(:=p(h)^1gp(h)\right)`$ is satisfied for every $`hH`$ and $`gG`$.
###### Definition 7.1.
The semi-direct product $`K^HG`$ of $`K`$ and $`G`$ along $`H`$ is defined to be $`KG/N`$, where
$$N=\left\{(d(h)^1,p(h)),hH\right\}.$$
For this definition to make sense, we have to verify that $`N`$ is a normal subgroup of $`KG`$. This is left to the reader. Hint: show that $`G`$ centralizes $`N`$, and an element $`kK`$ acts by
$$(d(h)^1,p(h))(d(h^k)^1,p(h^k)).$$
There are natural group homomorphisms $`p^{}:KK^HG`$ and $`d^{}:GK^HG`$, making the following diagram commute
There is also an action of $`K^HG`$ on $`G`$ which makes the above diagram equivariant. An element $`(k,g)K^HG`$ acts on $`xG`$ by sending it to $`g^1x^kg`$. Indeed, $`[d^{}:GK^HG]`$ is a crossed-module.
###### Lemma 7.2.
In the above square, the induced map $`\mathrm{Coker}(d)\mathrm{Coker}(d^{})`$ is an isomorphism and the induced map $`\mathrm{Ker}(d)\mathrm{Ker}(d^{})`$ is surjective. The kernel of the latter is equal to $`\mathrm{Ker}(p)\mathrm{Ker}(d)`$.
###### Proof.
Straightforward. ∎
The relative semi-direct product construction satisfies the obvious universal property. Namely, to give a homomorphism $`K^HGT`$ to an arbitrary group $`T`$ is equivalent to giving a pair of homomorphisms $`\delta :GT`$ and $`\varpi :KT`$ such that
*
* $`\delta (g^k)=\delta (g)^{\varpi (k)}`$, for every $`gG`$ and $`kK`$.
###### Remark 7.3.
Consider the the subgroup $`I=p\left(\mathrm{Ker}(d)\right)=\mathrm{Ker}d^{}G`$. It is a $`K`$-invariant central subgroup of $`G`$. If in the relative semi-direct product construction we replace $`G`$ by $`G/I`$ the outcome will be the same.
In the following lemma we slightly modify the notation and denote $`K^HG`$ by $`K^{H,p}G`$.
###### Lemma 7.4.
Notation being as above, let $`\theta :KG`$ be a crossed homomorphism (i.e. $`\theta (kk^{})=\theta (k)^k^{}\theta (k)`$). Consider the group homomorphism $`q:HG`$ defined by $`q(h)=p(h)\theta (d(h))`$, and use it to form $`K^{H,q}G`$. Also consider the new action of $`K`$ on $`G`$ given by $`g^k:=\theta (k)^1g^k\theta (k)`$. (The $``$ is used just to differentiate the new action from the old one.) Then, there is a natural isomorphism
$$\begin{array}{cccc}\hfill \theta _{}:& K^{H,q}G& \stackrel{}{}& K^{H,p}G\\ & (k,g)& & (k,\theta (k)g).\end{array}$$
The map $`\theta _{}`$ makes the following triangle commute:
Furthermore, we have the following commutative triangle of isomorphisms
where the top row is induced by $`\theta _{}`$.
###### Proof.
We use the universal property of the relative semi-direct product. To give a map from $`K^{H,p}G`$ to a group $`T`$ is equivalent to giving a pair $`(\delta ,\varpi )`$ of maps $`\delta :GT`$ and $`\varpi :KT`$ satisfying the two conditions described in the paragraph just before the lemma. To such a pair, we can associate a new pair $`(\delta ^{},\varpi ^{})`$, with $`\delta ^{}:=\delta `$ and $`\varpi ^{}(k):=\varpi (k)\delta (\theta (k))`$. It is easy to see that the pair $`(\delta ^{},\varpi ^{})`$ satisfies the two conditions required by the universal property of $`K^{H,q}G`$. Similarly, we can go backwards from a pair $`(\delta ^{},\varpi ^{})`$ for $`K^{H,q}G`$ to a pair $`(\delta ,\varpi )`$ for $`K^{H,p}G`$. It is easy to see that this correspondence is realized by $`\theta _{}`$. This proves that $`\theta _{}`$ is an isomorphism.
Commutativity of the triangles is obvious. (Also see Lemma 7.2.) ∎
We have the following converse for Lemma 7.4
###### Lemma 7.5.
Consider the semi-direct product diagrams
where the $`K`$-actions on $`H`$ are the same. Assume we are given an isomorphism of groups $`\vartheta :K^{H,q}G\stackrel{}{}K^{H,p}G`$ such that the triangles of Lemma 7.4 commute. Denote $`p(\mathrm{Ker}d)G`$ by $`I`$.
* If $`\mathrm{Ker}d\mathrm{Ker}p`$, then there is a unique crossed homomorphism $`\theta :KG`$ such that $`\vartheta =\theta _{}`$ (see Lemma 7.4, and Remark 7.3).
* If $`K`$ is a free group, then there exists a (not necessarily unique) crossed homomorphism $`\theta :KG`$ such that $`\vartheta =\theta _{}`$.
* If for crossed-homomorphisms $`\theta `$ and $`\theta ^{}`$ we have $`\theta _{}=\theta _{}^{}`$, then the difference of $`\theta `$ and $`\theta ^{}`$ factors through $`IG`$. That is, for every $`kK`$, $`\theta ^1(k)\theta ^{}(k)`$ lies in $`I)`$.
###### Proof of.
($`𝐢`$). Pick an element $`(k,g)K^{H,q}G`$, and let $`\theta (k,1)=(k^{},g^{})`$. (Note that $`(k,1)`$ and $`(k^{},g^{})`$ are just representatives for actual elements in the corresponding relative semi-direct product groups). By the commutativity of the above triangle, the images of $`k`$ and $`k^{}`$ are the same in $`\mathrm{Coker}(d)`$; that is, there exists $`hH`$ such that $`kd(h)=k^{}`$. So, after adjusting $`(k^{},g^{})`$ by the $`(d(h)^1,p(h))N`$ (see Definition 7.1), we may assume $`k^{}=k`$; that is $`\vartheta (k,1)=(k,g^{})`$. Define $`\theta (k)`$ to be $`g^{}`$. It is easily verified that $`\theta `$ is a crossed homomorphism, and that $`\theta _{}=\vartheta `$.
Proof of ($`\mathrm{𝐢𝐢}`$). Replace $`G`$ by $`G/I`$ and apply ($`𝐢`$) to obtain $`\theta :KG/I`$. Then use freeness of $`K`$ to lift $`\theta `$ to $`G`$.
Proof of ($`\mathrm{𝐢𝐢𝐢}`$). Easy. ∎
### 7.2. A pushout construction for crossed-modules
Continuing with the set-up of the previous section, we now bring crossed-modules into the picture. Namely, we assume that $`[d:HK]`$ is a crossed-module. (Note that the condition $`\mathrm{𝐂𝐌𝟐}`$ of crossed-modules ($`\mathrm{\S }`$3.2) is already part of the hypothesis.) To be compatible with our crossed-module notation, let us denote $`H`$, $`K`$, $`G`$ and $`d`$ by $`H_2`$, $`H_1`$, $`G_2`$ and <sub>-</sub> respectively. Recall that $`[G_2H_1^{H2}G_2]`$ is again a crossed-module.
###### Definition 7.6.
Let $`=[H_2H_1]`$ be a crossed-module. Let $`G_2`$ be a group with an action of $`H_1`$, and $`p:H_2G_2`$ an $`H_1`$-equivariant group homomorphism such that for every $`\beta H_2`$ and $`\alpha G_2`$ we have $`\alpha ^{\underset{¯}{\beta }}=\alpha ^{p(\beta )}`$. We call the crossed-module $`[G_2H_1^{H2}G_2]`$ the pushout of $``$ along $`p`$ and denote it by $`p_{}`$.
###### Lemma 7.7.
There is a natural induced map of crossed-modules $`p_{}:p_{}`$. Furthermore, $`\pi _1(p_{})`$ is an isomorphism and $`\pi _2(p_{})`$ is surjective. The kernel of $`\pi _2(p_{})`$ is equal to
$$\{\beta H_2|\underset{¯}{\beta }=1,p(\beta )=1\}.$$
In particular, if $`p:H_2G_2`$ is injective, then $`p_{}:[H_2H_1][G_2H_1^{H2}G_2]`$ is an equivalence of crossed-modules.
###### Proof.
Straightforward. ∎
Assume now that we are given two crossed-modules $`𝔾=[G_2G_1]`$, $`=[H_2H_1]`$ and a morphism $`P:𝔾`$ between them. This gives us a diagram
like the one in the beginning of this section. We also have an action of $`H_1`$ on $`G_2`$ with respect to which $`p_2`$ is $`H_1`$-equivariant. Namely, for $`\alpha G_2`$ and $`hH_1`$, we define $`\alpha ^h`$ to be $`\alpha ^{p_1(h)}`$, the latter being the action in $`𝔾`$. So, we can form the crossed-module $`[G_2H_1^{H2}G_2]`$.
Define the map $`\rho :H_1^{H2}G_2G_1`$ by $`\rho (h,\alpha ):=p_1(h)\underset{¯}{\alpha }`$. It is easily seen to be well-defined. We obtain the following commutative diagram of crossed-modules:
Note that the front-left square is almost an equivalence of crossed-modules (Lemma 7.7); it is an actual equivalence if and only if $`\pi _2P:\pi _2\pi _2𝔾`$ is injective. If this is the case, the above diagram means that, up to equivalence, we have managed to replace our crossed-module map $`P:𝔾`$ with one, i.e. $`P_{}𝔾`$, in which $`p_2`$ is the identity map (the front-right square).
Notation. If $`P:𝔾`$ is a map of crossed-modules, we use the notation $`P_{}`$ instead of $`p_{2,}`$.
The next thing we consider is, how the pushout construction for crossed-modules behaves with respect to pointed transformations between maps.
###### Lemma 7.8.
Let $`P,Q:𝔾`$ be maps of crossed-modules and $`\theta :QP`$ a pointed transformation between them (Definition 4.1). Then, we have the following commutative diagram of maps of crossed-modules:
in which the front faces compose to $`P`$ and the back faces compose to $`Q`$. Here $`\theta _{}`$ is obtained by the construction of Lemma 7.4 applied to $`\theta :H_1G_2`$. Furthermore, the following triangle commutes:
###### Proof.
This is basically a restatement of Lemma 7.4. Only proof of the equality $`\rho _Q=\rho _P\theta _{}`$ is missing. To prove this, pick $`(h,\alpha )H_1^{H_2,Q}G_2`$. Since $`\theta _{}(h,\alpha )=(h,\theta (h)\alpha )`$, we have
$`\rho _P(\theta _{}(h,\alpha ))=p_1(h)\underset{¯}{\theta (h)}\underset{¯}{\alpha }=q_1(h)\underset{¯}{\alpha }=\rho _Q(h,\alpha ).`$
###### Lemma 7.9.
Consider the commutative diagrams of Lemma 7.8, but with $`\vartheta `$ instead of $`\theta _{}`$. Assume $`\pi _2P:\pi _2\pi _2𝔾`$ is the zero homomorphism. Then, there is a unique transformation $`\theta :QP`$ such that $`\vartheta =\theta _{}`$.
###### Proof.
This is more or less a restatement of Lemma 7.5.$`𝐢`$, with $`H_1`$, $`H_2`$ and $`G_2`$ playing the roles of $`K`$, $`H`$ and $`G`$, respectively. More explicitly, construct $`\theta :H_1G_2`$ as in Lemma 7.5. It automatically satisfies condition $`\mathrm{𝐓𝟐}`$ of Definition 4.1. The fact that it satisfies $`\mathrm{𝐓𝟏}`$ follows from the definition of $`\theta _{}`$; see Lemma 7.4. ∎
## 8. Butterflies as weak morphisms
In this section, we give a description for the groupoid of weak maps between two crossed-modules (Theorem 8.4). The key is the following definition.
###### Definition 8.1.
Let $`𝔾=[G_2G_1]`$ and $`=[H_2H_1]`$ be crossed-modules. By a butterfly from $``$ to $`𝔾`$ we mean a commutative diagram of groups
in which both diagonal sequences are complexes, and the NE-SW sequence, that is, $`G_2EH_1`$, is short exact. We require $`\rho `$ and $`\sigma `$ satisfy the following compatibility with actions. For every $`xE`$, $`\alpha G_2`$, and $`\beta H_2`$,
$$\iota (\alpha ^{\rho (x)})=x^1\iota (\alpha )x,\kappa (\beta ^{\sigma (x)})=x^1\kappa (\beta )x.$$
We denote the above butterfly by the tuple $`(E,\rho ,\sigma ,\iota ,\kappa )`$. A morphism between two butterflies $`(E,\rho ,\sigma ,\iota ,\kappa )`$ and $`(E^{},\rho ^{},\sigma ^{},\iota ^{},\kappa ^{})`$ is an isomorphism $`f:EE^{}`$ commuting with all four maps. We define $`(,𝔾)`$ to be the groupoid of butterflies from $``$ to $`𝔾`$.
###### Remark 8.2.
Our notion of butterfly should not be confused with J. Pradines’. The latter correspond to morphisms in the localized category of topological groupoids obtained by inverting Morita equivalences.
###### Lemma 8.3.
In a butterfly $`(E,\rho ,\sigma ,\iota ,\kappa )`$, every element in the image of $`\iota `$ commutes with every element in $`\mathrm{Ker}\rho `$. Similarly, every element in the image of $`\kappa `$ commutes with every element in $`\mathrm{Ker}\sigma `$. In particular, the elements in the images of $`\iota `$ commute with the elements in the image of $`\kappa `$.
###### Proof.
Easy. ∎
The following theorem explains why we butterflies are interesting objects: they correspond to weak morphisms between crossed-modules ($`\mathrm{\S }`$5).
###### Theorem 8.4.
There is an equivalence of groupoids, natural up to a natural homotopy,
$$\mathrm{\Omega }:\underset{¯}{om}_{}(,𝔾)(,𝔾).$$
###### Proof.
We begin with the following observation. Consider the product crossed-module $`[H_2\times G_2H_1\times G_1]`$. Then, to give a butterfly from $``$ to $`𝔾`$ is equivalent to giving a triangle
which has the property that $`\rho `$ intertwines the action on $`H_2\times G_2`$ of $`E`$ (by conjugation) with the action of $`H_1\times G_1`$ (from the crossed-module structure) and that the projection map $`[H_2\times G_2E][H_2H_1]`$ is an equivalence of crossed-modules. (For this we use Lemma 8.3.)
Let us now construct the functor $`\mathrm{\Omega }:\underset{¯}{om}_{}(,𝔾)(,𝔾)`$. Choose a cofibrant replacement $`𝔽=[F_2F_1]`$ for $``$. Then, $`\underset{¯}{om}_{}(,𝔾)\underset{¯}{om}_{}(𝔽,𝔾)`$. So, we have to construct $`\underset{¯}{om}_{}(𝔽,𝔾)(,𝔾)`$. We use the pushout construction of $`\mathrm{\S }`$7.1. Namely, given a morphisms $`P:𝔽𝔾`$, we pushout $`𝔽`$ along the map $`F_2H_2\times G_2`$ to obtain the crossed module $`H_2\times G_2E`$. This is defined to be $`\mathrm{\Omega }(P)`$. This way we obtain a triangle
Lemma 7.7 implies that the induced map $`𝔽[H_2\times G_2E]`$ is a weak equivalence of crossed modules. Therefore, $`[H_2\times G_2E][H_2H_1]`$ is also an equivalence of crossed-modules. Thus, we see that $`\mathrm{\Omega }(P):=[H_2\times G_2E]`$ is really a butterfly from $``$ to $`𝔾`$. This defines the effect of $`\mathrm{\Omega }`$ on objects. The effect on morphisms is defined in the obvious way (see Lemma 7.8).
We need to show that $`\mathrm{\Omega }`$ is essentially surjective, full and faithful. Essential surjectivity follows from the fact that $`𝔽`$ is cofibrant, and fullness follows from Lemma 7.5.$`\mathrm{𝐢𝐢}`$. Let us prove the faithfulness. Let $`\theta ,\theta ^{}:F_1G_2`$ be transformations between $`P,Q:𝔽𝔾`$ such that $`\mathrm{\Omega }(\theta )=\mathrm{\Omega }(\theta ^{})`$. Define $`\widehat{\theta }:F_1H_2\times G_2`$ by $`\widehat{\theta }=(1,\theta )`$. Define $`\widehat{\theta ^{}}:F_1H_2\times G_2`$ in the similar way. By hypothesis, we have $`\widehat{\theta }_{}=\widehat{\theta ^{}}_{}:E_QE_P`$. By Lemma 7.5.$`\mathrm{𝐢𝐢𝐢}`$, for every $`xF_1`$, the element $`\widehat{\theta }^1(x)\widehat{\theta ^{}}(x)`$ lies in the image of $`\pi _2𝔽`$ under the map $`F_2H_2\times G_2`$. Since $`𝔽`$ induces an isomorphism on $`\pi _2`$, the only element in this image which is of the form $`(1,\alpha )`$ is $`(1,1)`$. On the other hand, every element $`\widehat{\theta }^1(x)\widehat{\theta ^{}}(x)`$ is of the form $`(1,\alpha )`$. We conclude that $`\widehat{\theta }^1(x)\widehat{\theta ^{}}(x)=(1,1)`$. Therefore, $`\theta (x)=\theta ^{}(x)`$. This completes the proof. ∎
###### Remark 8.5.
We can think of the category of crossed-modules and weak maps as the localization of $`\mathrm{𝟐}𝐆𝐩`$ with respect to equivalences of crossed-modules. Theorem 8.4 says that every weak map $`f:𝔾`$ can be canonically written as a fraction
where
$$𝔎:=\sigma ^{}=[E_{H_1}H_2\stackrel{pr_1}{}E][H_2\times G_2\stackrel{(\kappa ,\iota )}{}E].$$
###### Example 8.6.
Let $`X`$ be a topological space, and $`A`$ an abelian group. We can use the above theorem to give a description of the second cohomology $`H^2(X,A)`$. Recall that this cohomology group is in bijection with the set of homotopy classes of maps $`XK(A,2)`$, where $`K(A,2)`$ stands for the Eilenberg-MacLane space. Since $`K(A,2)`$ is simply connected, we can, equivalently, work with pointed homotopy classes. Assume the 2-type of $`X`$ is represented by the crossed-module $`=[H_2H_1]`$. Then, there is a natural bijection $`H^2(X,A)[,A]_{\mathrm{𝟐}𝐆𝐩}`$, where we think of $`A`$ as the crossed module $`[A1]`$. From Theorem 8.4 we conclude that $`H^2(X,A)`$ is in natural bijection with the set of isomorphism classes of pairs $`(E,\kappa )`$, where $`E`$ is a central extension of $`H_1`$ by $`A`$, and $`\kappa :H_2E`$ is an $`E`$-equivariant homomorphism. Here, $`E`$ acts on itself by conjugation and on $`H_2`$ via the projection $`EH_1`$.
Of course, we can take any crossed-module $`𝔾`$ instead of $`K(A,2)`$ and this way we find a description of the non-abelian cohomology $`H^1(X,𝔾)`$.
### 8.1. Cohomological point of view
We quote the following cohomological description of the set of pointed homotopy classes of weak maps from $``$ to $`𝔾`$, that is, $`[,𝔾]_{\mathrm{𝐗𝐌𝐨𝐝}}`$ from \[El\] (which apparently goes back to Joyal and Street; see \[ibid.\], $`\mathrm{\S }`$3.5).
###### Theorem 8.7 (\[El\], Theorem 2.7).
Let $`𝔾`$ and $``$ be crossed-modules, and assume they are represented by triples $`(\pi _1𝔾,\pi _2𝔾,[\alpha ])`$ and $`(\pi _1,\pi _2,[\beta ])`$, respectively, where $`\alpha `$ and $`\beta `$ are 3-cocycles representing the corresponding Postnikov invariants (see $`\mathrm{\S }`$15). Then, there is a bijection between the set of pointed homotopy classes of weak maps from $``$ to $`𝔾`$, that is, $`[,𝔾]_{\mathrm{𝐗𝐌𝐨𝐝}}`$, and the set of triples $`(\chi ,\lambda ,[c])`$, where $`\chi :\pi _1\pi _1𝔾`$ is a group homomorphism, $`\lambda :\pi _2\pi _2𝔾`$ is a $`\chi `$-equivariant homomorphism such that $`[\chi ^{}(\alpha )]=[\lambda _{}(\beta )]`$ in $`H^3(\pi _1,\pi _2𝔾^\chi )`$, and $`[c]`$ is the class, modulo coboundary, of a 2-cochain on $`\pi _1`$ with values in $`\pi _2𝔾^\chi `$, such that $`c=\chi ^{}(\alpha )\lambda _{}(\beta )`$. Here, $`\pi _2𝔾^\chi `$ stands for $`\pi _2𝔾`$ made into a $`\pi _1`$-module via $`\chi `$.
We will say more about this in the case where $`=\mathrm{\Gamma }`$ is an ordinary group in sections 14-16.
## 9. Basic facts about butterflies
In this section we investigate butterflies in more detail.
### 9.1. Butterflies vs. spans
Theorem 8.4 can be interpreted as saying that once we invert strict equivalences of crossed-modules in $`\mathrm{𝐗𝐌𝐨𝐝}`$, the morphisms of the resulting localized category can be presented, in a canonical way, by butterflies. In fact, every butterfly $`(E,\rho ,\sigma ,\iota ,\kappa )`$ gives rise to a span of strict crossed-module morphisms
The map $`\mu :H_2\times G_2E`$ is defined by $`(\alpha ,\beta )\kappa (\alpha )\iota (\beta )`$. Note that, by Lemma 8.3, the images of $`\iota `$ and $`\kappa `$ in $`E`$ commute. The componentwise action of $`E`$ on $`H_2`$ and $`G_2`$ makes $`𝔼:=[\mu :H_2\times G_2E]`$ into a crossed-module. It is easy to verify that $`𝔼`$ is an equivalence of crossed-modules.
### 9.2. Butterflies vs. cocycles
In (\[No1\], Definition 8.4), we give a definition of a weak morphism between crossed-modules $`[H_2H_1]`$ and $`[G_2G_1]`$ in terms of certain cocycles. By definition, such a cocycle consists of a triple $`(p_1,p_2,\epsilon )`$, where $`p_1:H_1G_1`$, $`p_2:H_2G_2`$, and $`\epsilon :H_1\times H_1G_2`$ are pointed set maps satisfying certain axioms \[ibid.\].<sup>4</sup><sup>4</sup>4There is an error in the set of axioms in the published version of \[ibid.\]. This is corrected in an erratum which is available on my webpage. Given such a triple, one can recover the corresponding butterfly as follows.
Let $`E`$ be the group that has $`H_1\times G_2`$ as the underlying set and whose product is defined by
$$(h,g)(h^{},g^{}):=(hh^{},\epsilon _{h,h^{}}^1g^{p_1(h^{})}g^{}).$$
Define the group homomorphisms $`\rho :EG_1`$ and $`\kappa :H_2E`$ by
$$\rho (h,g)=p_1(h)\underset{¯}{g},\kappa (h)=(\underset{¯}{h},p_2(h)^1).$$
The homomorphisms $`\iota :G_2E`$ and $`\sigma :EH_1`$ are the inclusion and the projection maps on the corresponding components.
It is easy to verify that this gives rise to a butterfly, and that this construction takes a (pointed) transformation of weak maps of crossed-modules (see loc. cit.) to a morphism of butterflies.
Conversely, given a butterfly $`(E,\rho ,\sigma ,\iota ,\kappa )`$, any choice of a set theoretic section $`s:H_1E`$ for $`\sigma :EH_1`$ gives rise to a cocycle $`(p_1,p_2,\epsilon )`$ defined by:
$$\begin{array}{c}p_1:H_1G_1,h\rho s(h),\hfill \\ p_2:H_2G_2,\alpha \kappa (\alpha )^1s(\underset{¯}{\alpha }),\hfill \\ \epsilon _{h,h^{}}=s(h^{})^1s(h)^1s(hh^{}).\hfill \end{array}$$
If we choose a different section $`s^{}`$ for $`\sigma `$, we find a pointed transformation $`\theta `$ between $`(p_1,p_2,\epsilon )`$ and $`(p_1^{},p_2^{},\epsilon ^{})`$ given by
$$\theta :H_1G_2,hs(h)^1s^{}(h).$$
We can summarize this discussion by saying that, to give a triple $`(p_1,p_2,\epsilon )`$ as in (\[No1\], Definition 8.4) is the same thing as to give a butterfly $`(E,\rho ,\sigma ,\iota ,\kappa )`$ together with a set theoretic splitting $`s:H_1E`$ of $`\sigma `$. In particular, we have an equivalence of groupoids
$$(,𝔾)\mathrm{Hom}_{Weak}(,𝔾)$$
where the right hand side is the groupoid whose morphisms of triples $`(p_1,p_2,\epsilon )`$ and whose morphisms of pointed transformations $`\theta `$ as in loc. cit.
### 9.3. The induced map on homotopy groups
By Theorem 8.4, a butterfly $`𝒫=(E,\rho ,\sigma ,\iota ,\kappa )`$ from $``$ to $`𝔾`$ induces a well-defined map $`NN𝔾`$ in the homotopy category of simplicial sets, which should be thought of as the “nerve of the weak map $`𝒫`$”. Indeed, any choice of a set-theoretic section $`s`$ for the map $`\sigma :EH_1`$ gives rise to a natural simplicial map $`N_s𝒫:NN𝔾`$. Furthermore, if $`s^{}`$ is another choice of a section, there is a natural simplicial homotopy between $`N_s𝒫`$ and $`N_s^{}𝒫`$.
In particular, a butterfly $`𝒫`$ induces natural homomorphisms on $`\pi _1`$ and $`\pi _2`$. We can in fact describe these maps quite explicitly.
To define $`\pi _1𝒫:\pi _1\pi _1𝔾`$, let $`x`$ be an element in $`\pi _1`$ and choose $`yE`$ such that $`\sigma (y)=x`$. Define $`\pi _1𝒫(x)`$ to be the class of $`\rho (y)`$ in $`\pi _1𝔾`$. This is easily seen to be well-defined. To define $`\pi _2𝒫:\pi _2\pi _2𝔾`$, let $`\beta `$ be an element in $`\pi _2`$ and consider $`y=\kappa (\beta )E`$. Since $`\sigma (y)=(\beta )=1`$, there exists a unique $`\alpha G_2`$ such that $`\iota (\alpha )=y`$. We define $`\pi _2𝒫(\beta )`$ to be $`\alpha `$. Note that $`(\alpha )=\rho (y)=1`$, so $`\alpha `$ is indeed in $`\pi _2𝔾`$.
### 9.4. Homotopy fiber of a butterfly
We saw in $`\mathrm{\S }`$9.3 that a butterfly $`(E,\rho ,\sigma ,\iota ,\kappa )`$ from $``$ to $`𝔾`$ induces a map $`NN𝔾`$ of simplicial sets which is well-defined up to simplicial homotopy. It is natural to ask what is the homotopy fiber of this map. In this subsection we see that this homotopy fiber can be recovered from the NW-SE sequence of the butterfly. We describe two ways of doing this.
Let $`𝒫:(E,\rho ,\sigma ,\iota ,\kappa )`$ be a butterfly from $``$ to $`𝔾`$. We define its homotopy fiber $`𝔉=𝔉_𝒫`$ to be the following 2-groupoid. The set of objects of $`𝔉`$ is $`G_1`$. Given $`g,g^{}G_1`$, the set of 1-morphisms in $`𝔉`$ from $`g`$ to $`g^{}`$ is the set of all $`xE`$ such that $`g\rho (x)=g^{}`$. For every two such $`x,yE`$, the set of 2-morphisms from $`x`$ to $`y`$ is the set of all $`\gamma H_2`$ such that $`x\kappa (\gamma )=y`$. We depict this 2-cell by
It is clear how to define compositions rules in $`𝔉`$, except perhaps for horizontal composition of 2-morphisms. Consider two 2-morphisms
We define their composition to be
There is a natural (strict) morphism of 2-groupoids $`\mathrm{\Phi }:𝔉`$. To describe this map, we need to think of $``$ as a 2-groupoid with one object, as in $`\mathrm{\S }`$3.3. We recall how this works. The unique object of $``$ is denoted $``$. The set of 1-morphisms of $``$ is $`H_1`$. Given two 1-morphisms $`h,h^{}H_1`$, the set of 2-morphisms from $`h`$ to $`h^{}`$ is the set of all $`\gamma H_2`$ such that $`h(\gamma )=h^{}`$.
The natural map $`\mathrm{\Phi }:𝔉`$ is described as follows:
###### Theorem 9.1.
The sequence $`N𝔉\stackrel{N\mathrm{\Phi }}{}N\stackrel{N𝒫}{}N𝔾`$, which is well-defined in the homotopy category of simplicial sets, is a homotopy fiber sequence.
In order to prove Theorem 9.1, we recall a few facts about homotopy fibers in $`\mathrm{𝟐}𝐆𝐩𝐝`$. Given a strict morphism $`P:𝔾`$ of 2-groupoids, and a base point $``$ in $`𝔾`$, there is a standard model for the homotopy fiber of $`P`$ which is given by the following strict fiber product:
$$\mathrm{Fib}_P:=\times _{,𝔾,s}𝔾^I\times _{t,𝔾,P}.$$
Here $`𝔾^I:=\underset{¯}{om}(01,𝔾)`$ is the 2-groupoid of 1-morphisms of $`𝔾`$, and $`s,t:𝔾^I𝔾`$ are the source and target functors. (For a more precise definition of $`𝔾^I`$ see Definition 17.1.)
In the case where $`𝔾`$ and $``$ are 2-groups associated to crossed-modules $`[G_2G_1]`$ and $`[H_2H_1]`$, the 2-groupoid $`\mathrm{Fib}_P`$ is described more explicitly as follows:
* Objects are elements of $`G_1`$.
* 1-morphisms from $`g`$ to $`g^{}`$ are pairs $`(x,\alpha )H_1\times G_2`$, as in the 2-cell
This means, $`gp_1(x)\underset{¯}{\alpha }=g^{}`$.
* 2-morphisms from $`(x,\alpha )`$ to $`(y,\beta )`$ are elements $`\gamma H_2`$ making the following diagram commutative:
More precisely, we want $`x\underset{¯}{\gamma }=y`$ and $`p_2(\gamma )\beta =\alpha `$.
There is a natural projection functor $`pr:\mathrm{Fib}_P=\times _{,𝔾,s}𝔾^I\times _{t,𝔾,f}`$ which fits in the following fiber homotopy sequence:
$$\mathrm{Fib}_P\stackrel{pr}{}\stackrel{P}{}𝔾.$$
In the case where the map $`p_2:H_2G_2`$ is surjective, there is a smaller model for the homotopy fiber of $`P`$ which we now describe. Let $`\mathrm{Fib}_P^{}`$ be the full sub-2-category of $`\mathrm{Fib}_P`$ in which for 1-morphisms we only take the ones for which $`\alpha `$ is the identity. It is easily checked that, in this case, the inclusion $`\mathrm{Fib}_P^{}\mathrm{Fib}_P`$ is an equivalence of 2-groupoids.
Let us record this as a lemma.
###### Lemma 9.2.
Let $`𝔾=[G_2G_1]`$ and $`=[H_2H_1]`$ be crossed-modules, viewed as 2-groups. Let $`P:𝔾`$ be a strict morphism of 2-groups, and let $`\mathrm{Fib}_P^{}`$ be the 2-groupoid defined in the previous paragraph. If $`p_2:H_2G_2`$ is surjective, then the sequence
$$\mathrm{Fib}_P^{}\stackrel{pr}{}\stackrel{P}{}𝔾$$
is a homotopy fiber sequence in the category of 2-groupoids.
We can now prove Theorem 9.1.
###### Proof of Theorem 9.1.
We apply Lemma 9.2 to the morphism of crossed-modules $`P:[H_2\times G_2E][G_2G_1]`$. Observe that the crossed-module on the left is a model for $``$ (via the natural weak equivalence $`(pr_1,\sigma ):[H_2\times G_2E][H_2H_1]`$), so $`P`$ is a strict model for $`𝒫`$.
It is easily verified that the 2-groupoid $`𝔉`$ of the theorem is canonically isomorphic to the 2-groupoid $`\mathrm{Fib}_P^{}`$ of Lemma 9.2. So, the homotopy fiber sequence of Lemma 9.2 is naturally isomorphic, in the homotopy category of 2-groupoids, to the sequence
$$𝔉\stackrel{\mathrm{\Phi }}{}\stackrel{𝒫}{}𝔾.$$
Taking the nerves gives the homotopy fiber sequence of Theorem 9.1. ∎
We now derive some corollaries of Theorem 9.1.
###### Proposition 9.3.
Let $`𝒫:(E,\rho ,\sigma ,\iota ,\kappa )`$ be a butterfly from $``$ to $`𝔾`$. Consider the nerve $`N𝒫:NN𝔾`$, which is well-defined in the homotopy category of simplicial sets, and let $`F`$ be its homotopy fiber (note that $`F`$ is naturally pointed). Let
$$C:H_2\stackrel{\kappa }{}E\stackrel{\rho }{}G_1$$
be the NW-SE sequence of $`𝒫`$. Then, there are natural isomorphisms $`H_i(C)\pi _iF`$, $`i=0,1,2`$. (Of course, for $`i=0`$ this means an isomorphism of pointed sets.)
###### Proof.
By Theorem 9.1, $`F`$ is naturally homotopy equivalent to $`N𝔉`$. It is easily verified that $`H_i(C)\pi _i𝔉`$, $`i=0,1,2`$. The result follows. ∎
###### Remark 9.4.
Note that the truncated sequence $`H_2E`$ is indeed naturally a crossed-module (take the action of $`E`$ on $`H_2`$ induced via $`\sigma `$). The 2-groupoid $`𝔉`$ of Theorem 9.1 can be recovered from the sequence $`C`$ together with this extra structure.
###### Corollary 9.5.
Notation being as in Proposition 9.3, there is a long exact sequence
###### Proof.
Immediate. ∎
###### Proposition 9.6.
A butterfly $`(E,\rho ,\sigma ,\iota ,\kappa )`$ from $``$ to $`𝔾`$ is a weak equivalence (i.e. induces a homotopy equivalence on the nerves), if and only if the NW-SE sequence
$$H_2\stackrel{\kappa }{}E\stackrel{\rho }{}G_1$$
is short exact. A weak inverse for this butterfly in obtained by simply flipping the butterfly, as in the following diagram
###### Proof.
Immediate. ∎
There is another way of describing the homotopy fiber of a butterfly. Define $`𝔽`$ to be the following crossed-module in groupoids:
$$𝔽:=[\underset{gG_1}{}H_2\stackrel{f}{}G_1\times E\stackrel{s,t}{}G_1].$$
The groupoid $`[G_1\times E\stackrel{}{}G_1]`$ is the translation groupoid of the right multiplication action of $`E`$ on $`G_1`$ via $`\rho `$. Its source and target maps are $`s:(g,x)g`$ and $`t:(g,x)g\rho (x)`$. The restriction of $`f`$ to the copy of $`H_2`$ corresponding to the index $`gG_1`$ sends the element $`\beta H_2`$ to $`(g,\kappa (\beta ))`$.
It can be shown that the nerve of $`𝔽`$ is naturally equivalent to $`𝔉`$.
### 9.5. Weak morphisms vs. strict morphisms
We saw in the Theorem 8.4 that $`(,𝔾)`$ is a model for the groupoid of weak morphism from $``$ to $`𝔾`$. Strict morphisms $`𝔾`$ form a subgroupoid of $`(,𝔾)`$. The objects of this subgroupoid are precisely the butterflies for which the NE-SW short exact sequence is split.
More precisely, given a strict morphisms $`P:𝔾`$, the corresponding butterfly looks as follows:
where $`\iota =(1,\mathrm{id})`$, $`\sigma =pr_1`$, $`\kappa (\beta )=(\underset{¯}{\beta },p_2(\beta ^1))`$, and $`\rho (x,\alpha )=p_1(x)\underset{¯}{\alpha }`$. Here, the action of $`H_1`$ on $`G_2`$ is obtained via $`p_1`$ from that of $`G_1`$ on $`G_2`$.
Observe that, any butterfly coming from a strict morphism has a canonical splitting. A morphism of butterflies does not necessarily respect this splitting. In fact, the difference between the resulting splittings determines a (pointed) transformation (Definition 4.1) between the corresponding strict functors, and vice versa. Let us state this in the following.
###### Proposition 9.7.
Let $`𝒮(,𝔾)`$ be the groupoid whose objects are butterflies with a splitting and whose morphisms are the morphisms of the underlying butterflies. Then, we have a natural equivalence of groupoids
$$\underset{¯}{om}_{}(,𝔾)𝒮(,𝔾).$$
### 9.6. Butterflies in the differentiable or algebraic contexts
Our discussion of butterflies can be generalized to a global setting in which everything happens over a Grothendieck site $`𝖲`$. A group is now replaced by a sheaf of groups, and a crossed-module $`[G_2G_1]`$ will have $`G_1`$ and $`G_2`$ sheaves of groups. Also, in the definition of a butterfly we require that the NW-SE sequence is a short exact sequence of sheaves of groups.
It can be shown that the discussion of the previous, and the subsequent, sections on butterflies goes through more or less verbatim in the relative case. I particular, it can be shown that butterflies model morphisms in the homotopy category of crossed-modules over $`𝖲`$.
As a consequence of this global approach, we obtain theories of butterflies for crossed-modules in Lie groups, topological groups, group schemes and so on. For example, a Lie butterfly between two Lie crossed-modules is one in which $`E`$ is a Lie group, all the maps $`\iota `$, $`\kappa `$, $`\rho `$, and $`\sigma `$ are differentiable homomorphisms, and the sequence
$$1G_2\stackrel{\iota }{}E\stackrel{\sigma }{}H_11$$
is a short exact sequence of Lie groups.
The general theory of butterflies over a Grothendieck site will appear in a joint paper with E. Aldrovandi.
###### Remark 9.8.
The savvy reader may complain that in the example of Lie butterflies mentioned above, one can only expect $`E`$ to be a sheaf of groups over the site $`𝖲`$ of differentiable manifolds. However, it is not hard to see that when $`E`$ is a sheaf of groups which sits in a short exact sequence whose both ends are Lie groups, then $`E`$ itself is necessarily representable by a Lie groups.
## 10. Bicategory of crossed-modules and weak maps
In this section we construct a bicategory $`𝒞`$ whose objects are crossed-modules and whose 1-morphisms are butterflies. The bicategory $`𝒞`$ is a model for the homotopy category of pointed connected 2-types.
### 10.1. The bicategory $`𝒞`$
Theorem 8.4 enables us to give a concrete model for the 2-category of crossed-modules, weak morphisms and weak transformation. More precisely, define $`𝒞`$ to be the bicategory whose objects are crossed-modules and whose morphism groupoids are $`(,𝔾)`$. The composition functors $`(𝕂,)\times (,𝔾)(𝕂,𝔾)`$ are defined as follows. Given butterflies
we define their composite to be the butterfly
where $`F\underset{H_1}{\overset{H_2}{\times }}E`$ is a Baer type product. More precisely, it is the quotient of the group
$$L:=\{(y,x)F\times E|\rho ^{}(y)=\sigma (x)H_1\}$$
modulo the subgroup
$$I=\{(\iota ^{}(\beta ),\kappa (\beta ))F\times E|\beta H_2\}.$$
We have the following theorem (see $`\mathrm{\S }`$4 for notation).
###### Theorem 10.1.
With morphism groupoids being $`(,𝔾)`$, crossed-modules form a bicategory $`𝒞`$. There is a natural weak functor $`\underset{¯}{\mathrm{𝐗𝐌𝐨𝐝}}𝒞`$ which is fully faithful on morphism groupoids.
###### Proof.
It is straightforward to verify that $`𝒞`$ is a bicategory. The functor $`\underset{¯}{\mathrm{𝐗𝐌𝐨𝐝}}𝒞`$ is the one constructed in $`\mathrm{\S }`$9.5. Fully faithfulness on morphism groupoids is a restatement of Proposition 9.7. ∎
###### Remark 10.2.
The bicategory $`𝒞`$ is a model for (i.e., is naturally biequivalent to) the 2-category of 2-groups and weak morphisms between them; see \[No1\], especially Propositions 8.1 and 7.8. More precisely, the latter is the 2-category whose objects are 2-groups, whose 1-morphisms are weak morphisms between 2-groups, and whose 2-morphisms are pointed transformations between them.
The advantage of working with $`𝒞`$ is that 1-morphisms and 2-morphisms in $`𝒞`$ are quite explicit, and in order to describe weak morphism one does not need to deal with complicated cocycles and coherence conditions (like the ones coming from the hexagon axiom) which in practice make explicit computations intractable. This is a great advantage in geometric contexts, say, when dealing with weak morphisms of Lie or algebraic 2-groups (see 9.6).
### 10.2. A special case of composition of butterflies
There is a simpler, and quite useful, description for the composition of two butterflies in the case where one of them is strict ($`\mathrm{\S }`$9.5).
When the first morphisms is strict, say
then the composition is
Here, $`q_1^{}(F)`$ stands for the pull back of the extension $`F`$ along $`q_1:K_1H_1`$. More precisely, $`q_1^{}(F)=K_1\times _{H_1}F`$ is the fiber product.
When the second morphisms is strict, say
then the composition is
Here, $`p_{2,}(E)`$ stands for the push forward of the extension $`E`$ along $`p_2:H_2G_2`$. More precisely, $`p_{2,}(E)=E\times _{H_2}G_2`$ is the pushout.
## 11. Kernels and cokernels of butterflies
In this section we give an overview of how certain basic notions of group theory, such as kernels, cokernels, complexes, extensions, and so on, can be generalized to the setting of 2-groups and weak morphisms through the language of butterflies. These notions have, of course, been studied by various authors already, but our emphasis is on showing how our approach via butterflies makes things remarkably simple and reveals structures that are not easy to see using the cocycle approach.
This section is meant as a general overview and we do not give many proofs. The proofs are not at all hard, but we believe this topic deserves a systematic treatment which is out of the scope of this paper. This is being worked out in \[No3\].
### 11.1. Kernels and cokernels of butterflies
For a weak morphism $`𝔾`$ of 2-groups, one can define a kernel and a cokernel using a certain universal property. The kernel is always expected to be a 2-group again. The cokernel, however, is only a pointed groupoid in general. In this section, we give an explicit description of the kernel and the cokernel of a weak morphism of 2-groups using the formalism of butterflies.
Consider the butterfly $`𝒫:𝔾`$ given by
We define the kernel of $`𝒫`$ by
$$\mathrm{Ker}𝒫:=[H_2\stackrel{\kappa }{}\mathrm{Ker}\rho ],$$
and the cokernel of $`𝒫`$ by
$$\mathrm{Coker}𝒫:=[\mathrm{Coker}\kappa \stackrel{\rho }{}G_1].$$
Note that $`\mathrm{Coker}𝒫`$ is just a group homomorphism and may not be a crossed-module in general (also see $`\mathrm{\S }`$12.1).
###### Remark 11.1.
The 2-group associated to the crossed-module $`\mathrm{Ker}𝒫`$ is naturally equivalent to the 2-group of automorphisms of the base point of the homotopy fiber $`𝔉_𝒫`$ of $`𝒫`$ defined in $`\mathrm{\S }`$9.4. The pointed set $`\pi _1(\mathrm{Coker}𝒫)`$ is in natural bijection with the set of connected components of $`𝔉_𝒫`$. We have a natural isomorphism
$$\pi _1\mathrm{Ker}𝒫\pi _2\mathrm{Coker}𝒫.$$
(Note: using the notation $`\pi _1`$ and $`\pi _2`$ for the cokernel and the kernel of a group homomorphism $`[KL]`$ which is not a crossed-module (e.g., for $`[KL]=\mathrm{Coker}𝒫`$) is not quite appropriate. We have used it here for the sake of notational consistency. This will not appear elsewhere in the paper except for this subsection.)
There is a (strict) morphism $`I_𝒫:=(\mathrm{id}_{H_2},\sigma ):\mathrm{Ker}𝒫`$. The morphism $`I_𝒫`$ comes with a natural 2-morphism $`\theta _𝒫:𝒫I_𝒫1`$ in $`\mathrm{Hom}_𝒞(\mathrm{Ker}𝒫,𝔾)`$, where $`1`$ stands for the trivial morphism. The following proposition follows from Proposition 11.6 below. We leave it to the reader to supply the proof and also to formulate the corresponding statement for the cokernel of a butterfly.
###### Proposition 11.2.
The pair $`(I_𝒫,\theta _𝒫)`$ satisfies the universal property of a kernel for $`𝒫:𝔾`$ in the bicategory $`𝒞`$. That is, for every crossed-module $`𝔏`$, $`(I_𝒫,\theta _𝒫)`$ induces an equivalence of groupoids between $`\mathrm{Hom}_𝒞(𝔏,\mathrm{Ker}𝒫)`$ and the homotopy fiber of $`𝒫_{}:\mathrm{Hom}_𝒞(𝔏,)\mathrm{Hom}_𝒞(𝔏,𝔾)`$ over the trivial map $`1\mathrm{Hom}_𝒞(𝔏,𝔾)`$.
###### Proposition 11.3.
Let $`𝒫:𝔊`$ be a butterfly. The map $`I_𝒫:\mathrm{Ker}𝒫`$ gives rise to a long exact sequence:
###### Proof.
Exercise. ∎
###### Corollary 11.4.
A butterfly $`𝒫:𝔊`$ is an equivalence if and only if $`\mathrm{Ker}𝒫`$ and $`\mathrm{Coker}𝒫`$ are trivial (i.e., equivalent to a point).
###### Remark 11.5.
One can also define the image and the coimage of a butterfly. Kernel, image, cokernel, and the coimage of a butterfly correspond to the top-left, top-right, bottom-right, and the bottom-left arrows in the butterfly, respectively.
In fact, careful study of butterflies shows that each of the notions kernel, image, cokernel, and the coimage can be defined in two ways, one of which we have given above. So, we have eight different crossed-modules associated to a butterfly, and these eight crossed-modules interact in an interesting way. This seems to suggest that one ought to consider all eight at once. This is discussed in detail in \[No3\].
### 11.2. Exact sequences of crossed-modules in $`𝒞`$
The following proposition is the key in studying complexes of 2-groups.
###### Proposition 11.6.
Let $`𝒬`$ and $`𝒫`$ be butterflies as in the following diagram:
Then, the composition $`𝒫𝒬`$ is the trivial butterfly if and only if there exists a group homomorphism $`\delta :FE`$ making the above diagram commutative such that the sequence
$$1\stackrel{}{}K_2\stackrel{\kappa ^{}}{}F\stackrel{\delta }{}E\stackrel{\rho }{}G_1\stackrel{}{}1$$
is a complex (i.e., the composition of consecutive maps is the trivial map).
Let $`𝒬:𝕂`$ and $`𝒫:𝔾`$ be butterflies. Proposition 11.6 gives a necessary and sufficient condition for the composition $`𝒫𝒬`$ to be zero. Observe that when this condition is satisfied, we have a natural morphism $`𝕂\mathrm{Ker}𝒫`$ given by the butterfly
The next question to ask is when the sequence
$$𝕂\stackrel{𝒬}{}\stackrel{𝒫}{}𝔾$$
is exact at $``$. The correct definition of exactness seems to be the following.
###### Definition 11.7.
We say that the sequence
$$𝕂\stackrel{𝒬}{}\stackrel{𝒫}{}𝔾$$
is exact at $``$ if $`𝒫𝒬`$ is zero and $`\mathrm{Coker}(𝕂\mathrm{Ker}𝒫)`$ is trivial.
###### Remark 11.8.
If we pretend for a moment that $`\mathrm{Coker}𝒬=[\mathrm{Coker}\kappa ^{}\stackrel{\rho ^{}}{}H_1]`$ is a crossed-module, then we obtain a “butterfly” $`\mathrm{Coker}𝒬𝔾`$ given by
Definition 11.7 is now equivalent to $`\mathrm{Ker}(\mathrm{Coker}𝒬𝔾)`$ being trivial. This argument is actually a valid argument if we assume $``$ is a braided crossed-module because in this case $`\mathrm{Coker}𝒬`$ is, indeed, a crossed-module and the above diagram is an honest butterfly; see $`\mathrm{\S }`$12.
The following proposition is immediate from the definition.
###### Proposition 11.9.
The sequence $`𝕂\stackrel{𝒬}{}\stackrel{𝒫}{}𝔾`$ is exact at $``$ (respectively, short exact) if and only if the sequence
$$1K_2\stackrel{\kappa ^{}}{}F\stackrel{\delta }{}E\stackrel{\rho }{}G_11$$
of Proposition 11.6 is exact at $`F`$ and $`E`$ (respectively, exact everywhere).
Using the above proposition the reader can work out what a complex (or an exact sequence) of crossed-modules in $`𝒞`$ looks like.
Proposition 11.9 also provides a new perspective on the problem of studying extensions of a 2-group $`𝔾`$ by a 2-group $`𝕂`$ (see \[Bre2\] and \[Rou\]). This will be studied in more detail in \[No3\].
## 12. Braided and abelian butterflies
In this section we will discuss butterflies in an abelian category $`𝖠`$. This is intended to be an overview of the main features of braided butterflies and we will not give proofs.
We obtain an explicit description of the derived category of complexes of length two in $`𝖠`$; compare \[SGA4\]. We begin by discussing braided butterflies.
### 12.1. Braided butterflies
Assume $``$ and $`𝔾`$ are braided crossed-modules \[Co\], and let $`\{,\}_{}`$ and $`\{,\}_𝔾`$ denote the corresponding braidings. (Our braiding convention is that $`\{x,y\}=xyx^1y^1`$.)
###### Definition 12.1.
A butterfly $`𝒫=(E,\rho ,\sigma ,\iota ,\kappa ):𝔾`$ is called braided if the following conditions are satisfied:
$$x,yE,\kappa \{\sigma (x),\sigma (y)\}_{}=xyx^1\iota \{\rho (x),\rho (y^1)\}_𝔾y^1.$$
In the case where $`𝒫`$ comes from a strict morphism of 2-groups ($`\mathrm{\S }`$9.5) these correspond to the usual braided morphisms of crossed-modules (which, in turn, corresponds to the braided morphisms of 2-groups under the equivalence of $`\mathrm{\S }`$3.3).
It is easy to verify that two braided butterflies compose to a braided butterfly. Braided crossed-modules and braided butterflies form a bicategory $`𝒞`$. There is a forgetful bifunctor $`𝒞𝒞`$. The braided version of Theorem 10.1 is also true.
### 12.2. Kernels and cokernels of braided butterflies
The most interesting aspect of braided butterflies is that the NE-SW of such a butterfly has a natural structure of a 2-crossed-module.
Also, the kernel, the cokernel, and the image of a braided butterfly $`𝒫:𝔾`$ (see $`\mathrm{\S }`$11.1) between braided crossed-modules are naturally braided crossed-modules, and the natural maps $`\mathrm{Ker}𝒫`$ and $`𝔾\mathrm{Coker}𝒫`$ are braided morphisms of crossed-modules. The action of $`G_2`$ on $`E`$ given by $`x^g:=x\iota (\{\rho (x),g\}_𝔾)`$.
### 12.3. Butterflies in an abelian category
Let $`𝖠`$ be an abelian category. The results of the previous sections are also valid for the category of complexes of length 2 in $`𝖠`$. More precisely, let $`𝒞h^{[1,0]}(𝖠)`$ be the bicategory whose objects are complexes $`X^1X^0`$, whose morphisms are abelian<sup>5</sup><sup>5</sup>5Abelian means that, since there are no actions, we are dropping the requirement for compatibility of actions. butterflies
and whose 2-morphisms are isomorphisms of butterflies. The composition of butterflies is defined as in the case of usual butterflies ($`\mathrm{\S }`$10). The criterion for strictness of a butterfly is also valid ($`\mathrm{\S }`$9.5).
Let us describe the additive structure on $`𝒞h^{[1,0]}(𝖠)`$. Given two morphisms $`P`$, $`P^{}`$ in $`𝒞h^{[1,0]}(𝖠)`$
with the same source and target, we define $`P+P^{}`$ to be the butterfly
We define $`P`$ to be the butterfly
Using the above addition, the morphism groupoids in $`𝒞h^{[1,0]}(𝖠)`$ become symmetric monoidal categories, and the composition rule becomes a symmetric monoidal functor.
The category obtained by identifying 2-isomorphic morphisms in $`𝒞h^{[1,0]}(𝖠)`$ is naturally equivalent to the full subcategory of the derived category $`𝒟(𝖠)`$ consisting of complexes sitting in degrees $`[1,0]`$. Under this equivalence, the diagonal sequence
$$X^1\stackrel{\kappa }{}E\stackrel{\rho }{}Y^0,$$
viewed as a complex sitting in degrees $`[2,0]`$, corresponds to the mapping cone.
###### Remark 12.2.
As in the non-abelian case, an abelian butterfly as above comes from a chain map if and only if the NE-SW sequence is split. This is automatically the case if either $`X_0`$ is projective or $`Y^1`$ is injective.
Let us define $`𝒞h_{st}^{[1,0]}(𝖠)`$ to be the category whose objects are complexes sitting in degrees $`[1,0]`$, whose morphisms are morphisms of complexes, and whose 2-morphisms are chain homotopies. Note that morphism groupoids in $`𝒞h_{st}^{[1,0]}(𝖠)`$ are (strict) symmetric monoidal. There is natural (weak) functor $`𝒞h_{st}^{[1,0]}(𝖠)𝒞h^{[1,0]}(𝖠)`$ defined as in $`\mathrm{\S }`$9.5. (This is the 2-categorified version of the quotient functor from the homotopy category to the derived category.) This functor is a bijection on objects and faithful (and symmetric monoidal) on morphism groupoids; see Theorem 10.1.
Let $`𝕏`$ and $`𝕐`$ be objects in $`𝒞h^{[1,0]}(𝖠)`$. The next proposition gives us some information about the groupoid $`\mathrm{Hom}_{𝒞h^{[1,0]}(𝖠)}(𝕏,𝕐)`$.
###### Proposition 12.3.
Let $`𝕏=[X^1X^0]`$ and $`𝕐=[Y^1Y^0]`$ be objects in $`𝒞h^{[1,0]}(𝖠)`$. Let $`xt(X^0,Y^1)`$ be the groupoid whose objects are extensions of $`X^0`$ by $`Y^1`$ and whose morphisms are isomorphisms of extensions. Then, the forgetful map $`\mathrm{Hom}_{𝒞h^{[1,0]}(𝖠)}(𝕏,𝕐)xt(X^0,Y^1)`$ which sends a butterfly to its NE-SW sequence is a fibration of groupoids whose fiber is equivalent to $`\mathrm{Hom}_{𝒞h_{st}^{[1,0]}(𝖠)}(𝕏,𝕐)`$. Indeed, the sequence
$$0\mathrm{Hom}_{𝒞h_{st}^{[1,0]}(𝖠)}(𝕏,𝕐)\mathrm{Hom}_{𝒞h^{[1,0]}(𝖠)}(𝕏,𝕐)xt(X^0,Y^1)$$
is an exact sequence of symmetric monoidal categories.
## 13. A special case of Theorem 8.4; Dedecker’s theorem
In this section we look at the special case of Theorem 8.4 with $`=[1\mathrm{\Gamma }]`$, where $`\mathrm{\Gamma }`$ is a group. The groupoid $`(\mathrm{\Gamma },𝔾)`$ looks as follows. The object of $`(\mathrm{\Gamma },𝔾)`$ are diagrams of the form
Here $`E`$ is an extension of $`\mathrm{\Gamma }`$ by $`G_2`$, and for every $`xE`$ and $`\alpha G_2`$ we require that $`\alpha ^{\rho (x)}=x^1\alpha x`$. In fact, the short exact sequence
$$1G_2E\mathrm{\Gamma }1$$
is also part of the data, but we suppress it from the notation and denote such a diagram simply by $`(E,\rho )`$.
A morphism in $`(\mathrm{\Gamma },𝔾)`$ from $`(E,\rho )`$ to $`(E^{},\rho ^{})`$ is an isomorphism $`f:EE^{}`$ of extensions (so it induces identity on $`G_2`$ and $`\mathrm{\Gamma }`$) such that $`\rho =\rho ^{}f`$.
###### Theorem 13.1.
The functor
$$\mathrm{\Omega }:\underset{¯}{om}_{}(\mathrm{\Gamma },𝔾)(\mathrm{\Gamma },𝔾)$$
is an equivalence of groupoids. That is, $`(\mathrm{\Gamma },𝔾)`$ is naturally equivalent to the groupoid of weak functors from $`\mathrm{\Gamma }`$ to $`𝔾`$.
###### Proof.
Follows immediately from Theorem 8.4. ∎
###### Corollary 13.2 (Dedecker, \[Ded\]).
There is a natural bijection
$$\pi _0:\mathrm{Hom}_{\mathrm{Ho}(\mathrm{𝟐}𝐆𝐩)}(\mathrm{\Gamma },𝔾)\stackrel{}{}\pi _0(\mathrm{\Gamma },𝔾).$$
In other words, the homotopy classes of weak maps from $`\mathrm{\Gamma }`$ to $`𝔾`$ are in a natural bijection with isomorphism classes of diagrams of the form
where the $`E`$ is an extension of $`\mathrm{\Gamma }`$ by $`G_2`$, and for every $`xE`$ and $`\alpha G_2`$ the equality $`\alpha ^{\rho (x)}=x^1\alpha x`$ is satisfied.
###### Example 13.3.
Let $`A`$ be an abelian group, and consider the crossed-module $`𝔄=[A1]`$. The groupoid of central extensions of $`\mathrm{\Gamma }`$ by $`A`$ is naturally equivalent to the groupoid $`(\mathrm{\Gamma },𝔄)`$, which is itself naturally equivalent to the derived mapping groupoid $`\underset{¯}{om}_{}(\mathrm{\Gamma },𝔄)`$, by Theorem 13.1. In particular, by Corollary 13.2, the set of homotopy classes of weak maps from $`\mathrm{\Gamma }`$ to $`𝔄`$ is in natural bijection with isomorphism classes of central extensions of $`\mathrm{\Gamma }`$ by $`A`$, which is itself in natural bijection with $`H^2(\mathrm{\Gamma },A)`$.<sup>6</sup><sup>6</sup>6This is of course not surprising, since $`𝒜`$ is an algebraic model for the Eilenberg-MacLane space $`K(A,2)`$.
###### Example 13.4.
Let $`K`$ be a group. Let $`𝔄𝔲𝔱(K)`$ be the crossed-module $`[K\mathrm{Aut}(K)]`$.<sup>7</sup><sup>7</sup>7The corresponding 2-group is the 2-group of self-equivalences of $`K`$, where $`K`$ is viewed as a category with one object. The groupoid of extensions of $`\mathrm{\Gamma }`$ by $`K`$ is naturally equivalent to the groupoid $`(\mathrm{\Gamma },𝔄𝔲𝔱(K))`$, which is itself naturally equivalent to the derived mapping groupoid $`\underset{¯}{om}_{}(\mathrm{\Gamma },𝔄𝔲𝔱(K))`$, by Theorem 13.1. In particular, by Corollary 13.2, the set of homotopy classes of weak maps from $`\mathrm{\Gamma }`$ to $`𝔄𝔲𝔱(K)`$ is in natural bijection with isomorphism classes of extensions of $`\mathrm{\Gamma }`$ by $`K`$.
###### Remark 13.5.
The above example is identical to Theorem 2 (and Proposition 3) of \[BBF\] in the case where $`𝒢`$ (notation as in loc. cit.) has only one object. It seems that, indeed, the general case of Theorem 2 (and also Proposition 3) of \[BBF\], with $`𝔊`$ an arbitrary groupoid, is equivalent to this special case. To see this, note that we may assume $`𝒢`$ is connected. On the other hand, since both sides of the equality in Theorem 2 (and also Proposition 3) of \[BBF\] are functorial in $`𝒢`$, we may assume $`𝒢`$ has only one object.
###### Remark 13.6.
Corollary 13.2 was first discovered by Dedecker \[Ded\]. It is discussed in detail in \[Bre1\] in the more general case where everything takes place in a Grothendieck site. In analogy with Example 13.4, an element in $`\pi _0(\mathrm{\Gamma },𝔾)`$ is called a $`𝔾`$-extension of $`\mathrm{\Gamma }`$. Such extensions can be thought of as “$`𝔾`$-torsors” over the classifying space $`B\mathrm{\Gamma }`$, the same way that group extensions of $`\mathrm{\Gamma }`$ by $`K`$ can be regarded as $`𝔄𝔲𝔱(K)`$-torsors (i.e., $`K`$-gerbes) over $`B\mathrm{\Gamma }`$. Therefore, one can alternatively describe elements of $`\pi _0(\mathrm{\Gamma },𝔾)`$ as certain first cohomology classes of $`\mathrm{\Gamma }`$ with coefficients in the crossed-module $`𝔾`$. For more on this, the reader is encouraged to consult \[Bre1\], especially, $`\mathrm{\S }`$8.
In the rest of this section we present three isolated facts for which we have no use, but we found them interesting nevertheless!
### 13.1. Side note 1: split crossed-modules
We present a cute (and presumably well-known) application of Corollary 13.2.
A 2-group $`𝔊`$ is called split if it is completely determined by $`\pi _1𝔊`$, $`\pi _2𝔊`$, and the action of the former on the latter. More precisely, $`𝔊`$ is split if it is isomorphic in $`\mathrm{Ho}(\mathrm{𝐗𝐌𝐨𝐝})`$ to the crossed-module $`[:\pi _2𝔊\pi _1𝔊]`$, where $``$ is the trivial homomorphism. From the homotopical point of view, the following proposition is straightforward (see the beginning of $`\mathrm{\S }`$15). However, to give a purely algebraic proof seems to be tricky. We give a proof that makes use of Corollary 13.2.
###### Proposition 13.7.
Let $`𝔊=[G_2G_1]`$ be a crossed-module, and assume that the map $`𝔾\pi _1𝔾`$, viewed in $`\mathrm{Ho}(\mathrm{𝐗𝐌𝐨𝐝})`$, admits a section. Then $`𝔊`$ is split.
###### Proof.
By Corollary 13.2, there exists an extension
$$1\stackrel{}{}G_2\stackrel{}{}E\stackrel{f}{}\pi _1𝔾\stackrel{}{}1$$
and a map $`\rho :EG_1`$ satisfying the conditions stated therein. Consider the semi-direct product $`G_2\pi _2𝔾`$ where $`G_2`$ acts on $`\pi _2𝔾`$ by conjugation. It fits in a crossed-module $`𝔾^{}=[G_2\pi _2𝔾E]`$ where the map $`G_2\times \pi _2𝔾E`$ is obtained from the first projection map, and the action of $`E`$ on both factors is by conjugation. We have a homomorphism $`\sigma :G_2\pi _2𝔾G_2`$ which on the first factor is just the identity and on the second factor is the inclusion map. It is easy to see that $`(\sigma ,\rho ):𝔾^{}𝔾`$ is a crossed-module map; in fact, it is an equivalence of crossed-modules. On the other hand, $`(pr_2,f):[G_2\pi _2𝔾E][\pi _2𝔾\pi _1𝔾]`$ is also an equivalence of crossed-modules. So we have constructed a zigzag of equivalences that connect $`𝔾`$ to the trivial crossed-module $`[\pi _2𝔾\pi _1𝔾]`$. So $`𝔾`$ is split. ∎
### 13.2. Side note 2: Hopf’s formula
It is interesting to compute $`\pi _0\underset{¯}{om}_{}(\mathrm{\Gamma },𝔄)`$ straight from definition of the derived mapping groupoid (Definition 6.6). This leads to Hopf’s formula for $`H^2(\mathrm{\Gamma },A)`$.
Choose a presentation $`F/R\mathrm{\Gamma }`$ of $`\mathrm{\Gamma }`$, where $`F`$ is free. The crossed-module maps $`[RF][A1]`$ are precisely the group homomorphism $`g:RA`$ that are constant on the conjugacy classes (under the $`F`$-action) of elements in $`R`$. In other words, $`g(x)=1`$ for every $`x[F,R]`$. Two such homomorphisms $`g`$ and $`g^{}`$ are homotopic, if there is a group homomorphism $`h:FA`$ such that $`h|_R=g^{}g^1`$. So, $`\pi _0\underset{¯}{om}_{}(\mathrm{\Gamma },𝔄)`$ is in natural bijection with
$$\mathrm{Coker}\left\{\mathrm{Hom}(F,A)\mathrm{Hom}(R/[F,R],A)\right\}.$$
This is Hopf’s famous formula for $`H^2(\mathrm{\Gamma },A)`$.
## 14. Cohomological point of view
In this section we give a cohomological characterization of $`(\mathrm{\Gamma },𝔾)`$, Theorem 14.6, and compare it to Theorem 13.1. Theorem 14.6 is well-known<sup>8</sup><sup>8</sup>8The diligent reader can dig it out of \[AzCe\], $`\mathrm{\S }`$5., but the explicit way in which it relates to Theorem 13.1 is what we are interested in. Since this construction has been quoted in \[BeNo\], we feel obliged to include the precise account. We begin by recalling some standard facts about group extensions.
### 14.1. Groups extensions and $`H^2`$; review
We recall some basic facts about classification of group extensions via cohomological invariants. Our main reference is \[Bro\].<sup>9</sup><sup>9</sup>9Also see Example 13.4.
Let $`N`$ and $`\mathrm{\Gamma }`$ be groups (not necessarily abelian). We would like to classify extensions
$$1NE\mathrm{\Gamma }1,$$
up to isomorphism. First of all, notice that such an extension gives rise to a group homomorphism $`\psi :\mathrm{\Gamma }\mathrm{Out}(N)`$. So we might as well fix $`\psi `$ as part of the data, and classify extension which induce $`\psi `$. Denote the set of such extensions by $`(\mathrm{\Gamma },N,\psi )`$. Let $`C`$ denote the center of $`N`$, made into a $`\mathrm{\Gamma }`$-module through $`\psi `$. We have the following theorem.
###### Theorem 14.1 (\[Bro\], Theorem 6.6).
The set $`(\mathrm{\Gamma },N,\psi )`$ admits a natural free, transitive action by the abelian group $`H^2(\mathrm{\Gamma },C)`$. Hence, either $`(\mathrm{\Gamma },N,\psi )=\mathrm{}`$, or else there is a bijection $`(\mathrm{\Gamma },N,\psi )H^2(\mathrm{\Gamma },C)`$. This bijection depends on the choice of a particular element of $`(\mathrm{\Gamma },N,\psi )`$.
###### Remark 14.2.
If we are given a lift $`\stackrel{~}{\psi }:\mathrm{\Gamma }\mathrm{Aut}(N)`$ of $`\psi `$, we obtain a distinguished element in $`(\mathrm{\Gamma },N,\psi )`$, namely, the semi-direct product $`N\mathrm{\Gamma }`$. This is automatically the case if, for instance, $`N`$ is abelian. Therefore, when $`N`$ is abelian, we have a canonical bijection $`(\mathrm{\Gamma },N,\psi )H^2(\mathrm{\Gamma },C)`$.
The meaning of this theorem is that, given two elements $`E_0`$, $`E`$ in $`(\mathrm{\Gamma },N,\psi )`$, one can produce their difference as an element in $`H^2(\mathrm{\Gamma },C)`$. Notice that $`C`$ is now abelian, so, by Remark 14.2, every element in $`H^2(\mathrm{\Gamma },C)`$ gives rise to a canonical extension of $`\mathrm{\Gamma }`$ by $`C`$. Below we will explain how this extension can be explicitly constructed from $`E_0`$ and $`E`$.
Let us call a complex $`ME\mathrm{\Gamma }`$ semi-exact if the left map is injective, the right map is surjective, and the kernel $`K`$ of $`E\mathrm{\Gamma }`$ is generated by $`M`$ and $`C_K(M)`$ (the the centralizer of $`M`$ in $`K`$). This last condition guarantees that there is a well-defined homomorphism $`\psi :\mathrm{\Gamma }\mathrm{Out}(M)`$.
Assume we are given two semi-exact sequences
$$ME_0\mathrm{\Gamma },ME\mathrm{\Gamma }$$
such that $`\psi _0=\psi `$. Define $`L`$ to be the group of pairs $`(x,y)E_0\times E`$ such that $`\overline{x}=\overline{y}\mathrm{\Gamma }`$, and that conjugation by $`x`$ and $`y`$ induce the same automorphism of $`M`$. Observe that $`I=\{(a,a)E_0\times E|aM\}`$ is a normal subgroup of $`L`$.
###### Definition 14.3 (Baer product).
Notation being as above, define $`E_0\underset{\Gamma }{\overset{𝑀}{\times }}E:=L/I`$.
There is an obvious surjective homomorphism $`E_0\underset{\Gamma }{\overset{𝑀}{\times }}E\mathrm{\Gamma }`$. It fits in a natural exact sequence
$$1CC_{K_0}(M)\times C_K(M)E_0\underset{\Gamma }{\overset{𝑀}{\times }}E\mathrm{\Gamma }1,$$
where $`C`$ is the center of $`M`$ mapping diagonally to $`C_K(M)\times C_K^{}(M)`$. (This gives us two semi-exact sequences:
$$C_{K_0}(M)E_0\underset{\Gamma }{\overset{𝑀}{\times }}E\mathrm{\Gamma },C_K(M)E_0\underset{\Gamma }{\overset{𝑀}{\times }}E\mathrm{\Gamma }$$
which are somehow mirror to each other.)
If the two semi-exact sequences that we started with were actually exact, we would have $`C_K(M)=C_K^{}(M)=C`$. So, by identifying $`C`$ as the cokernel of the diagonal map $`CC\times C`$, we obtain the following exact sequence
$$1CE_0\underset{\Gamma }{\overset{𝑀}{\times }}E\mathrm{\Gamma }1.$$
More explicitly, we define the map $`CE_0\underset{\Gamma }{\overset{𝑀}{\times }}E`$ by sending $`a`$ to $`(a,1)`$.
###### Definition 14.4.
Let $`E_0,E(\mathrm{\Gamma },N,\psi )`$. Define the difference $`D(E_0,E)`$ to be the the sequence $`1CE_0\underset{\Gamma }{\overset{𝑁}{\times }}E\mathrm{\Gamma }1`$ defined above.
###### Remark 14.5.
Observe that $`E_0\underset{\Gamma }{\overset{𝑀}{\times }}E`$ is symmetric with respect to $`E_0`$ and $`E`$, but $`D(E_0,E)`$ is not. What determines the sign in the above construction is the map $`CE_0\underset{\Gamma }{\overset{𝑁}{\times }}E`$. So, if instead of $`a(a,1)`$ we used $`a(1,a)`$ we would obtain $`D(E,E_0)`$, because in $`CE_0\underset{\Gamma }{\overset{𝑀}{\times }}E`$ the elements $`(1,a)`$ and $`(a,1)`$ are inverse to each other.
Conversely, given $`E_0(\mathrm{\Gamma },N,\psi )`$ and an extension $`1CH\mathrm{\Gamma }1`$ (recall that $`C`$ is the center of $`N`$), we can recover $`E`$ as the difference $`E:=D(E_0,H)`$. In other words, consider the group $`E_0\underset{\Gamma }{\overset{𝐶}{\times }}H`$. This contains $`N=N\underset{1}{\overset{𝐶}{\times }}C`$ as a normal subgroup. The sequence
$$1NE_0\underset{\Gamma }{\overset{𝐶}{\times }}H\mathrm{\Gamma }1$$
is the desired extension.
### 14.2. Cohomological classification of maps into a 2-group
In this subsection we prove the following cohomological classification of the homotopy classes of weak maps from a group $`\mathrm{\Gamma }`$ to a 2-group $`𝔾`$. The result itself is well-known, but the way in which it relates to the classification theorem (Theorem 13.1) is what we are interested in.
###### Theorem 14.6.
Let $`𝔊`$ be a 2-group, and let $`\mathrm{\Gamma }`$ be a discrete group. Fix a homomorphism $`\chi :\mathrm{\Gamma }\pi _1𝔊`$, and let $`[\mathrm{\Gamma },𝔊]_{\mathrm{𝟐}𝐆𝐩}^\chi `$ be the set of homotopy classes of weak maps $`\mathrm{\Gamma }𝔾`$ inducing $`\chi `$ on $`\pi _1`$. Then, either $`[\mathrm{\Gamma },𝔊]_{\mathrm{𝟐}𝐆𝐩}^\chi `$ is empty, or it is naturally a transitive $`H^2(\mathrm{\Gamma },\pi _2𝔊)`$-set. (Here, $`\pi _2𝔊`$ is made into a $`\mathrm{\Gamma }`$ module via $`\chi `$.)
In $`\mathrm{\S }`$15 we see exactly when $`[\mathrm{\Gamma },𝔊]_{\mathrm{𝟐}𝐆𝐩}^\chi `$ is non-empty (see Remark 15.1) .
###### Remark 14.7.
If in the above proposition we take $`𝔾=𝔄𝔲𝔱(N)`$ we recover Theorem 14.1; see Example 13.4.
Before proving Theorem 14.6 we need some preliminaries.
Conventions for this subsection. Throughout this subsection, we fix $`\chi `$ (and consequently an action of $`\mathrm{\Gamma }`$ on $`\pi _2𝔾`$). We will think of $`H^2(\mathrm{\Gamma },\pi _2𝔊)`$ as the group of isomorphism classes of extensions of $`\mathrm{\Gamma }`$ by $`\pi _2𝔊`$ for which the induced action of $`\mathrm{\Gamma }`$ on $`\pi _2𝔾`$ is the one we have fixed. Whenever we talk about an extension of $`\mathrm{\Gamma }`$ by $`\pi _2𝔊`$ we assume that this condition is satisfied.
Construction of the action of $`H^2(\mathrm{\Gamma },\pi _2𝔊)`$ on $`[\mathrm{\Gamma },𝔊]_{\mathrm{𝟐}𝐆𝐩}^\chi `$.
Suppose we are given $`(E_0,\rho _0)`$ in $`[\mathrm{\Gamma },𝔊]_{\mathrm{𝟐}𝐆𝐩}^\chi `$ and an extension
$$1\pi _2𝔾K\mathrm{\Gamma }1$$
in $`H^2(\mathrm{\Gamma },\pi _2𝔊)`$. Set $`E:=E_0\underset{\Gamma }{\overset{\pi _2𝔊}{\times }}K`$ (see Definition 14.3). The inclusion $`G_2E_0`$ induces a natural homomorphism $`G_2E`$ which identifies $`G_2`$ with a normal subgroup of $`E`$. The quotient is $`\mathrm{\Gamma }`$. We have a natural map $`\rho :EG_1`$ defined by $`\rho (x,a)=\rho _0(x)`$. This is easily seen to be well-defined. Finally, it is easy to check that the action of $`E`$ on $`G_2`$ induced via $`\rho `$ is equal to the conjugation action of $`E`$ on $`G_2`$. This gives us the desired diagram:
Transitivity of the action. To prove the transitivity of the action of $`H^2(\mathrm{\Gamma },\pi _2𝔊)`$ on $`[\mathrm{\Gamma },𝔊]_{\mathrm{𝟐}𝐆𝐩}^\chi `$ we employ a ‘difference construction’ similar to the ones of the previous subsection. It takes two elements in $`[\mathrm{\Gamma },𝔊]_{\mathrm{𝟐}𝐆𝐩}^\chi `$ and produces their difference, which is an element in $`H^2(\mathrm{\Gamma },\pi _2𝔊)`$.
Assume $`[\mathrm{\Gamma },𝔊]_{\mathrm{𝟐}𝐆𝐩}^\chi `$ is non-empty, and fix an element $`(E_0,\rho _0)`$ in it as in the following diagram (see Theorem 13.1):
Let $`(E,\rho )`$ be another such diagram. Define the group $`L`$ by
$$L=\{(x,y)E_0\times E|\overline{x}=\overline{y},\rho _0(x)=\rho (y)\}.$$
There is a natural surjective group homomorphism $`L\mathrm{\Gamma }`$ sending $`(x,y)`$ to $`\overline{x}=\overline{y}`$. The kernel of this map is the following group:
$$\{(\alpha ,\beta )G_2\times G_2|\alpha \beta ^1\pi _2𝔊\}=I\times \pi _2𝔊$$
where
$$I:=\{(\beta ,\beta )E_0\times E|\beta G_2\},$$
and $`\pi _2𝔾`$ is identified with the subgroup of elements of the form $`(\alpha ,1)`$, $`\alpha \pi _2𝔾`$. It is easy to check that $`I`$ is normal in $`L`$.
###### Definition 14.8.
Define $`E_0\times _𝔾E=L/I`$. The map $`\alpha (\alpha ,1)`$ identifies $`\pi _2𝔾`$ with a normal subgroup of $`L`$ with cokernel $`\mathrm{\Gamma }`$. The extension
$$1\pi _2𝔾E_0\times _𝔾E\mathrm{\Gamma }1,$$
or its class in $`H^2(\mathrm{\Gamma },\pi _2𝔾)`$, is called the difference of $`(E_0,\rho _0)`$ and $`(E,\rho )`$ and is denoted by $`D((E_0,\rho _0),(E,\rho ))`$.<sup>10</sup><sup>10</sup>10 Definition 14.4 is a special case of Definition 14.8 with $`𝔾=𝔄𝔲𝔱(N)`$; see Example 13.4.
###### Proof of Theorem 14.6.
We leave it to the reader to verify that the construction of the action of $`H^2(\mathrm{\Gamma },\pi _2𝔊)`$ on $`[\mathrm{\Gamma },𝔊]_{\mathrm{𝟐}𝐆𝐩}^\chi `$ and the difference construction of Definition 14.8 are inverse to each other. ∎
An interesting special case is when $`\chi :\mathrm{\Gamma }\pi _1𝔾`$ can be lifted to $`\stackrel{~}{\chi }:\mathrm{\Gamma }G_1`$. In the following corollary we fix such a lift.
###### Corollary 14.9.
With the hypothesis of the preceding paragraph, every class $`f[\mathrm{\Gamma },𝔊]_{\mathrm{𝟐}𝐆𝐩}^\chi `$ is uniquely characterized by (the isomorphism class of) an extension
$$1\pi _2𝔊K\mathrm{\Gamma }1.$$
More explicitly, given such an extension we obtain $`(E,\rho )`$, where $`E:=K^{\pi _2𝔊}G_2`$ (Definition 7.1), and $`\rho (k,\alpha ):=\stackrel{~}{\chi }(\overline{k})\underset{¯}{\alpha }`$, for $`(k,\alpha )E`$. Here the action of $`K`$ on $`G_2`$ is obtained via $`\stackrel{~}{\chi }`$ from that of $`G_1`$ on $`G_2`$.
###### Proof.
Set $`E_0=\mathrm{\Gamma }G_2`$, the action being obtained through $`\stackrel{~}{\chi }`$, and define $`\rho _0:E_0G_1`$ by $`\rho (x,\alpha ):=\chi (x)\underset{¯}{\alpha }`$. (Keep in kind that $`K\underset{\Gamma }{\overset{\pi _2𝔊}{\times }}(\mathrm{\Gamma }G_2)=K^{\pi _2𝔊}G_2`$.) ∎
An important special case of the above corollary is when there exists a section for the map $`G_1\pi _1𝔊`$. In this case, after fixing such a section, we have an explicit description of the elements in $`[\mathrm{\Gamma },𝔊]_{\mathrm{𝟐}𝐆𝐩}`$, for any group $`\mathrm{\Gamma }`$.
## 15. Homotopical point of view
The results of the previous sections are best understood if viewed from the point of homotopy theory of 2-types. Recall that 2-groups (respectively, weak functors between 2-groups, weak natural transformations) are algebraic models for pointed and connected homotopy 2-types (respectively, pointed continuous maps, pointed homotopies). Using this dictionary, the problem of classification of weak maps between 2-groups translates to the problem of classification of (pointed) homotopy classes of continuous maps between (pointed) homotopy 2-types. The latter can be solved using standard technique from obstruction theory.
In this section we explain how the homotopical approach works and give a homotopical proof of Theorem 14.6. In the next section we compare the homotopical approach with the approaches of the previous section. These all are presumably folklore and well-known.
### 15.1. The classifying space functor and homotopical proof of Theorem 14.6
Consider the classifying space functor $`B:\mathrm{Ho}(\mathrm{𝟐}𝐆𝐩)\mathrm{Ho}(\mathrm{𝐓𝐨𝐩}_{})`$ defined by $`B(𝔾)=|N(𝔾)|`$. By Corollary 17.10 of Appendix, we know that this gives rise to an equivalence between the homotopy category of 2-groups and the homotopy category of pointed connected CW-complexes with vanishing $`\pi _i`$, $`i3`$ (the so called connected homotopy 2-types), a result essentially due to Whitehead \[Wh\]. So it would be most natural to view the algebraic results of the previous sections from this homotopic perspective. To complete the picture, we recall MacLane-Whitehead characterization of pointed connected homotopy 2-types from \[McWh\].
Mac Lane and Whitehead show that, to give a connected pointed homotopy 2-types is equivalent to giving a triple $`(\pi _1,\pi _2,\kappa )`$ where $`\pi _1`$ is an arbitrary group, $`\pi _2`$ is an abelian group endowed with a $`\pi _1`$ action, and $`\kappa H^3(\pi _1,\pi _2)`$. To see where such a triple comes from, let $`X`$ be a pointed connected CW-complex such that $`\pi _iX=0`$, $`i3`$. Consider the Postnikov decomposition of $`X`$:
The triple corresponding to $`X`$ is $`(\pi _1X,\pi _2X,\kappa )`$, where $`\kappa `$ is the Postnikov invariant corresponding to the above picture. Recall that this Postnikov invariant is the obstruction to existence of a section for the fibration $`p:X_2X_1`$. We know from obstruction theory that, if this obstruction vanishes, then for any choice of base points $`x_1X_1`$ and $`x_2p^1(x_1)`$, there exist a pointed section $`X_1X_2`$. Furthermore, after fixing such a section, the pointed homotopy classes<sup>11</sup><sup>11</sup>11We should actually be considering fiberwise homotopy classes, but since in our case the fiber is simply connected we get the same thing. of such section are in bijection with $`H^2(\pi _1X,\pi _2X)`$.
###### Homotopical proof of Theorem 14.6.
Passing to classifying space induces a bijection between $`[\mathrm{\Gamma },𝔊]_{\mathrm{𝟐}𝐆𝐩}`$ and pointed homotopy class of maps $`B\mathrm{\Gamma }B𝔊`$. Let
be the Postnikov tower of $`B𝔊`$. To give a group homomorphism $`\chi :\mathrm{\Gamma }\pi _1𝔾`$ is equivalent to giving a pointed homotopy class from $`B\mathrm{\Gamma }`$ to $`X_1`$. Fix such a class $`F:B\mathrm{\Gamma }X_1`$. The question is now to classify the pointed (equivalently, fiberwise – because the fiber is simply connected) homotopy classes of lifts $`f:B\mathrm{\Gamma }X_2`$ of $`F`$. We know from obstruction theory that the obstruction to existence of such a lift is precisely $`F^{}(\kappa )=\chi ^{}(\kappa )H^3(\mathrm{\Gamma },\pi _2X)`$, where $`\pi _2X`$ is made into a $`\mathrm{\Gamma }`$-module via $`\chi `$. By obstruction theory, whenever this obstruction vanishes, the set of pointed (equivalently, fiberwise) homotopy classes of lifts of such lifts $`f`$ is a transitive $`H^2(\mathrm{\Gamma },\pi _2X)`$-set. ∎
###### Remark 15.1.
As we saw in the above proof, for a fixed $`\chi :\mathrm{\Gamma }\pi _1𝔾`$, the set $`[\mathrm{\Gamma },𝔊]_{\mathrm{𝟐}𝐆𝐩}^\chi `$ is non-empty if and only if $`\chi ^{}(\kappa )H^3(\mathrm{\Gamma },\pi _2𝔾)`$ is zero, where $`\kappa `$ is the Postnikov invariant of $`𝔾`$.
###### Example 15.2.
Let $`\mathrm{\Gamma }`$ and $`N`$ be discrete groups. For a given $`\chi :\mathrm{\Gamma }\mathrm{Out}(N)`$, the obstruction to lifting $`\chi `$ to a map $`\mathrm{\Gamma }𝔄𝔲𝔱(N)`$ (equivalently, to finding an extension of $`\mathrm{\Gamma }`$ by $`N`$ giving giving rise to $`\chi `$ – see Example 13.4), is the element $`\chi ^{}(\kappa )H^3(\mathrm{\Gamma },C(N))`$, where $`C(N)`$ is the center of $`N`$ and $`\kappa H^3(\mathrm{Out}(N),C(N))`$ is the Postnikov invariant of $`𝔄𝔲𝔱(N)`$. (In other words, the Postnikov invariant $`\kappa `$ of $`𝔄𝔲𝔱(N)`$ is the universal obstruction class for the existence of group extensions.) When $`\chi ^{}(\kappa )=0`$, the set of lifts of $`\chi `$ to $`𝔄𝔲𝔱(N)`$ (equivalently, extensions of $`\mathrm{\Gamma }`$ by $`N`$ giving rise to $`\chi `$) admits a natural transitive action of $`H^2(\mathrm{\Gamma },C)`$, where $`C=\pi _2𝔄𝔲𝔱(N)`$ is the center of $`N`$.
For the sake of amusement, we also include a homotopical proof for Proposition 13.7.
###### Homotopical proof of Proposition 13.7.
Let $`𝔾`$ be as in Proposition 13.7. We have to show that the 2-type $`B𝔾`$ is split. That is, it is homotopy equivalent to the product $`K(\pi _1𝔊,1)\times K(\pi _2𝔊,2)`$ of Eilenberg-MacLane spaces. We know that $`𝔊`$ is classified by the triple $`(\pi _1𝔊,\pi _2𝔊,\kappa )`$. By assumption, the action of $`\pi _1𝔊`$ on $`\pi _2𝔊`$ is trivial. Also, by assumption, the corresponding Postnikov tower $`X_2X_1`$ has a section, so $`\kappa `$ vanishes. Since the 2-type $`K(\pi _1𝔊,1)\times K(\pi _2𝔊,2)`$ also gives rise to the same triple, it must be (pointed) homotopy equivalent to $`B𝔾`$, which is what we wanted to prove. ∎
## 16. Compatibility of different approaches
Let $`p:YX`$ be a fibration of CW-complexes (or simplicial sets). In this section we recall the notion of difference fibration for two liftings of a map $`F:AX`$, and use it to clarify Definition 14.8, as well as to explain why the cohomological classification of maps into a 2-group ($`\mathrm{\S }`$14.2) is compatible with the homotopical approach ($`\mathrm{\S }`$15).
A more detailed discussion of the difference fibration construction, and its application to obstruction theory, can be found in \[Ba\].
### 16.1. Difference fibration construction
The difference constructions of $`\mathrm{\S }`$14.1 and $`\mathrm{\S }`$14.2 are special cases of (and were originally motivated by) a general difference construction for maps of simplicial sets (or topological spaces). In this section we review this general construction and explain how it relates to the algebraic versions of it that we have already encountered in previous sections.
Let $`X`$ and $`Y`$ be simplicial sets and $`p:YX`$ a simplicial map. Let $`A`$ be another simplicial set and $`F:AX`$ a map of simplicial sets. We are interested in the classification of lifts of $`f`$ to $`Y`$, if such lifts exist. In our case of interest, $`A`$ and $`X`$ are going to be 1-types and $`Y`$ a 2-type, so everything is explicit and easy.
A useful tool in the study of such lifting problems is the difference construction, as in (\[Ba\], page 293).<sup>12</sup><sup>12</sup>12We will use the simplicial version of Baues’s definition though.
Let $`f_0,f:AY`$ be liftings of $`F`$. To measure the difference between $`f_0`$ and $`f`$, we construct the simplicial set $`D_p(f_0,f)`$ as in the following cartesian diagram:
Here, $`c:AX^{\mathrm{\Delta }^1}`$ is the map that sends $`aA`$ to the constant path at $`F(a)`$. We are actually interested in the case where $`p:YX`$ is a fibration.
The following lemma justifies the terminology, but since we will not need it here we will not give the proof (except in a special case – see Proposition 16.7). In the case of topological spaces, a proof is can be found in \[Ba\].
###### Lemma 16.1.
If $`p:YX`$ is a fibration, then $`D_p(f_0,f)A`$ is also a fibration.
The usefulness of the difference construction is justified by the following proposition.
###### Proposition 16.2.
* The simplicial set of sections to the map $`D_p(f_0,f)A`$ is naturally isomorphic to the simplicial set of fiberwise homotopies between $`f_0`$ and $`f`$.
* The primary difference of the sections $`f_0`$ and $`f`$ is precisely the primary obstruction to the existence of a section to $`D_p(f_0,f)A`$.
###### Proof.
Part ($`𝐢`$) follows from the definition. Part ($`\mathrm{𝐢𝐢}`$) is proved in \[Ba\] (see page 295 loc. cit.) in the case of topological spaces. The proof can be adopted to the simplicial situation. ∎
When all spaces and maps are pointed, $`D_p(f_0,f)`$ is also naturally pointed. In this case we have:
###### Proposition 16.3.
* The simplicial set of pointed sections to the map $`D_p(f_0,f)A`$ is naturally isomorphic to the simplicial set of fiberwise homotopies between $`f_0`$ and $`f`$ which fix the base point.
* The primary difference of the sections $`f_0`$ and $`f`$ is precisely the primary obstruction to the existence of a section to $`D_p(f_0,f)A`$ (everything pointed).
Difference fibration construction can indeed be performed in any category with fiber products in which there is an interval, and we have a notion of internal hom.<sup>13</sup><sup>13</sup>13We are not asking for a monoidal structure here. For instance, $`\mathrm{𝐂𝐚𝐭}`$, $`\mathrm{𝟐}𝐂𝐚𝐭`$, and $`\mathrm{𝟐}𝐆𝐩𝐝`$ are examples of such a category, where for the interval we take $`𝐈_1=\{01\}`$, and the internal homs are given by $`\underset{¯}{om}`$, as in Definition 17.1. So, given a diagram
of 2-categories, we can talk about the difference $`D_p(f_0,f)`$ of $`f_0`$ and $`f`$. This is a 2-category with a natural functor $`D_p(f_0,f)A`$. Furthermore, if $`𝔄`$, $``$ and $`𝔇`$ are pointed (i.e. have a chosen object), then so is $`D_p(f_0,f)`$. The following 2-categorical version of Proposition 16.2 is also valid.
###### Proposition 16.4.
Suppose $`𝔄`$, $``$ and $`𝔇`$ are 2-categories (respectively, pointed 2-categories) as above. Then the 2-category of sections (respectively, pointed sections) to the functor $`D_p(f_0,f)A`$ is naturally isomorphic to the 2-category of fiberwise transformations (respectively, pointed fiberwise transformations) between $`f_0`$ and $`f`$.
Clearly if $`𝔄`$, $``$ and $`𝔇`$ are 2-groupoids, then so is $`D_p(f_0,f)`$. So, we can talk about difference construction for 2-groupoids. It is also true that, if $`p`$ is a fibration (Appendix, Definition 17.2), then so is $`D_p(f_0,f)A`$. The latter statement, whose simplicial counterpart we did not prove (Lemma 16.1), can be proved easily by verifying the conditions of Definition 17.2.
###### Definition 16.5.
In the above situation, assume $`𝔄`$, $``$ and $`𝔇`$ are 2-groupoids with one object (i.e. 2-groups), and let $`\mathrm{Ob}D_p(f_0,f)`$ be the canonical base point of $`D_p(f_0,f)`$. We define $`D_p(f_0,f)_{}`$ to be the 2-group of automorphisms of the object $``$.
###### Remark 16.6.
Proposition 16.4 remains valid when $`𝔄`$, $``$ and $`𝔇`$ are 2-groups, and $`D_p(f_0,f)`$ is replaced by $`D_p(f_0,f)_{}`$. However, the natural functor $`D_p(f_0,f)_{}𝔄`$ is not in general a fibration anymore.
Finally, observe that the nerve functor $`N:\mathrm{𝟐}𝐂𝐚𝐭\mathrm{𝐒𝐒𝐞𝐭}`$ respects the difference construction. That is, $`ND_p(f_0,f)`$ is naturally homotopy equivalent to $`D_{Np}(Nf_0,Nf)`$. This is because $`N`$ preserves fiber products (Appendix, Proposition 17.5) and path spaces (Appendix, Proposition 17.11).
### 16.2. Difference fibrations for crossed-modules
We saw in the previous subsection that we can perform difference construction for 2-groups. We will make this more explicit using the language of crossed-modules. So assume $`𝔽`$, $`𝔾`$ and $``$ are crossed-modules, as in the following picture:
To avoid notational complications, we have used $`f^{}`$ instead of $`f_0`$. In what follows, it would be helpful to think of elements of $`G_1`$ (respectively, $`G_2`$) as 1-cells (respectively, 2-cells) of the nerve $`N𝔾`$ (see Appendix).
Let $`[D_2D_1]`$ be the crossed-module presentation of $`D_p(f^{},f)_{}`$. We will write down exactly what $`D_1`$ and $`D_2`$ are. By definition, we have
$$D_1=\left\{(h,\alpha )\right|hH_1,\alpha \mathrm{Ker}p_2,\text{s.t.}f_1^{}(h)\underset{¯}{\alpha }=f_1(h)\}.$$
It should be clear what this means: $`\alpha `$ is 2-cell that is vertical (because it is in $`\mathrm{Ker}p_2`$) and joins the 1-cells $`f_1^{}(h)`$ to $`f_1(h)`$. The group multiplication is
$$(h,\alpha )(k,\beta )=(hk,\alpha ^{f_1^{}(k)}\beta ).$$
To determine $`D_2`$, we have to pick a 2-cell $`\beta H_2`$, and find all pointed vertical homotopies from $`f_2^{}(\beta )`$ to $`f_2(\beta )`$. A pointed homotopy from $`f_2^{}(\beta )`$ to $`f_2(\beta )`$ is a 2-cell $`\gamma G_2`$ such that $`f_2^{}(\beta )\gamma =f_2(\beta )`$. To ensure it is vertical, we need to have $`\gamma \mathrm{Ker}p_2`$. This, however, is automatic, since $`p_2f_2^{}(\beta )=p_2f_2(\beta )`$. The conclusion is that, for any $`\beta H_2`$, there is a unique vertical homotopy from $`f_2^{}(\beta )`$ to $`f_2(\beta )`$; it is given by $`\gamma =f_2^{}(\beta )^1f_2(\beta )`$. Therefore,
$$D_2=H_2.$$
The map $`D_2D_1`$ is given by
$$H_2=D_2\beta (\underset{¯}{\beta },f_2^{}(\beta )^1f_2(\beta ))D_1.$$
An element $`(h,\alpha )D_1`$ acts on $`\beta D_2=H_2`$ by sending it to $`\beta ^h`$.
There is a natural map of crossed-modules $`q:D_p(f^{},f)_{}`$ given by $`q_1(h,\alpha )=h`$ and $`q_2=\mathrm{id}`$.
### 16.3. The special case
We now consider the special case of the difference construction that is relevant to Theorem 14.6, and show that it recovers Definition 14.8 (see Proposition 16.7). Suppose we are given a group $`\mathrm{\Gamma }`$, a 2-group $`𝔾`$, and a homomorphism $`\chi :\mathrm{\Gamma }\pi _1𝔾`$. We want to study lifts of $`\chi `$ to maps $`\mathrm{\Gamma }𝔾`$. Keep in mind that here we are talking about maps in the homotopy category of 2-groups. So it would definitely be false to consider 2-group maps $`\mathrm{\Gamma }𝔾`$. One way to handle the situation is to work with weak maps. But this is not very convenient. The better way would be to pick a cofibrant replacement $``$ for $`\mathrm{\Gamma }`$ (see Example 6.5), and use the fact that a map $`\mathrm{\Gamma }𝔾`$ in the homotopy category of 2-groups can be represented by a map of 2-groups $`𝔾`$, and that the latter is unique up to pointed transformation.
Recall (Example 6.5) how we construct a cofibrant replacement for $`\mathrm{\Gamma }`$: we choose a presentation $`F/R=\mathrm{\Gamma }`$, where $`F`$ is a free group, and form the crossed-module $`=[RF]`$; the natural map $`\mathrm{\Gamma }`$ is then our cofibrant replacement. Throughout the paper, we fix such a cofibrant replacement.
Our problem is to study the difference construction for the following situation:
All maps are now honest maps of crossed-modules. (We have abused notation and denoted the induced map $`\pi _1𝔾`$ also by $`\chi `$.) As we saw in the previous subsection, this picture gives a map of crossed-modules $`D_p(f^{},f)_{}`$, which is indeed a fibration of crossed-modules in the sense of Definition 6.1. Our aim is to compare this with Definition 14.8.
Let $`(E,\rho )`$, respectively $`(E^{},\rho ^{})`$, be the object of $`(\mathrm{\Gamma },𝔾)`$ obtained from pushing out $``$ along $`f`$, respectively $`f`$; see Definition 7.6). Recall (Definition 14.8) that the difference $`D((E^{},\rho ^{}),(E,\rho ))`$ is defined to be the following exact sequence:
$$1\pi _2𝔾E^{}\times _𝔾E\mathrm{\Gamma }1.$$
We prove the following proposition.
###### Proposition 16.7.
Notation being as above, the map $`D_p(f^{},f)_{}`$ is a fibration of crossed-modules and is naturally equivalent to $`E^{}\times _𝔾E\mathrm{\Gamma }`$. More precisely, we have the following commutative square in which the horizontal arrows are equivalences of crossed-modules and the vertical arrows are fibrations:
###### Proof.
We use the explicit description of the crossed-module $`D_p(f^{},f)_{}=[D_2D_1]`$ given in $`\mathrm{\S }`$16.2:
$$D_1=\left\{(x,\alpha )\right|xF,\alpha G_2,\text{s.t.}f_1^{}(x)\underset{¯}{\alpha }=f_1(x)\}\text{and}D_2=R.$$
The map $`D_2D_1`$ is given by $`r(\underset{¯}{r},f_2^{}(r)^1f_2(r))`$, where $`\underset{¯}{r}`$ stands for the image of $`r`$ in $`F`$. Notice that this is an injection, so we can identify $`D_2`$ with a subgroup of $`D_1`$. Let us now give an explicit description of $`E^{}\times _𝔾E`$. Recall (Definition 7.6) that
$$E=F^{R,f}G_2\text{and}E^{}=F^{R,f^{}}G_2.$$
Using this, we get
$$E^{}\times _𝔾E=\left\{((x,\alpha ),(y,\beta ))\right|\overline{x}=\overline{y}\mathrm{\Gamma },f_1^{}(x)\underset{¯}{\alpha }=f_1(y)\underset{¯}{\beta }G_1\}/J,$$
where $`x,yF`$, $`\alpha ,\beta G_2`$, and $`J`$ is the subgroup generated by
$$\{((\underset{¯}{r},f_2^{}(r)^1),(1,1);rR\},\{((1,1),(\underset{¯}{s},f_2(s)^1));sR\},$$
$$\text{and}\left\{((1,\gamma ),(1,\gamma ));\gamma G_2\right\}.$$
Therefore,
$$J=\{((\underset{¯}{r},f_2^{}(r)^1\gamma ),(\underset{¯}{s},f_2(s)^1\gamma ));r,sR,\gamma G_2\}.$$
Define the map $`\mathrm{\Lambda }:D_1E^{}\times _𝔾E`$ by
$$\mathrm{\Lambda }(x,\alpha )=((x,\alpha ),(x,1)).$$
We claim that $`\mathrm{\Lambda }`$ is surjective and its kernel is $`R=D_2`$. This proves that the map $`D_p(f^{},f)_{}E^{}\times _𝔾E`$ is an equivalence of crossed-modules.
Surjectivity. Pick an element $`a=((x,\alpha ),(y,\beta ))`$ in $`E^{}\times _𝔾E`$. Since $`\overline{x}=\overline{y}`$, there is $`sR`$ such that $`y=x\underset{¯}{s}`$. Using the fact that
$$((x,\alpha ),(x\underset{¯}{s},\beta ))=((x,\alpha ),(x,f_2(s)\beta ))$$
in $`E^{}\times _𝔾E`$, we may assume that $`x=y`$, that is, $`a=((x,\alpha ),(x,\beta ))`$. On the other hand, after multiplying on the right by $`((1,\beta ^1),(1,\beta ^1))J`$, we may assume that $`\beta =1`$; that is $`a=((x,\alpha ),(x,1))`$. This is obviously in the image of $`\mathrm{\Lambda }`$.
Kernel of $`\mathrm{\Lambda }`$ is $`R`$. Easy verification.
To show that $`D_p(f^{},f)_{}`$ is a fibration, we have to show that the map $`D_1F`$ which sends $`(x,\alpha )`$ to $`x`$, and the map $`\mathrm{id}:D_2=RR`$ are surjective (Definition 6.1). The latter is obvious. The former follows from the fact that, for every $`xF`$, we have the equality $`\overline{f_1(x)}=\overline{f_1^{}(x)}`$ in $`\pi _1𝔾`$. This is true because both these elements are equal to $`\chi (x)`$.
Commutativity of the square is obvious. ∎
## 17. Appendix: 2-categories and 2-groupoids
In this appendix we quickly go over some basic facts and constructions we need about 2-categories, and fix some terminology. Most of the material in this appendix can be found in \[No1\].
For us, a 2-category means a strict 2-category. A 2-groupoid is a 2-category in which every 1-morphism and every 2-morphism has an inverse (in the strict sense). Every category (respectively, groupoid) can be thought of as a 2-category (respectively, 2-groupoid) in which all 2-morphisms are identity.
A 2-functor between 2-categories means a strict 2-functor. We sometimes refer to a 2-functor simply by a functor, or a map of 2-categories. A 2-functor between 2-groupoids is simply a 2-functor between the underlying 2-categories.
By fiber product of 2-categories we mean strict fiber product. We will not encounter homotopy fiber product of 2-categories in this paper.
The terms ‘morphism’, ‘1-morphism’ and ‘arrow’ will be used synonymously. We use multiplicative notation for elements of a groupoid, as opposed to the compositional notation (it means, $`fg`$ instead of $`gf`$).
Notation. We use the German letters $``$, $`𝔇`$,… for general 2-categories and $`𝔾`$, $``$,… for 2-groupoids. The upper case script letters $`A`$, $`B`$, $`C`$… are used for objects in such 2-categories, lower case script letters $`a`$, $`b`$, $`g`$, $`h`$… for 1-morphisms, and lower case Greek letters $`\alpha `$, $`\beta `$… for 2-morphisms. We denote the category of 2-categories by $`\mathrm{𝟐}𝐂𝐚𝐭`$ and the category of 2-groupoids by $`\mathrm{𝟐}𝐆𝐩𝐝`$.
### 17.1. 2-functors, weak 2-transformations and modifications
We recall what weak 2-transformations between strict 2-functors are. More details can be found in \[No1\]. We usually suppress the adjective weak.
Let $`P,Q:𝔇`$ be (strict) 2-functors. By a weak 2-transformation $`T:PQ`$ we mean, assignment of an arrow $`t_A`$ in $``$ to every object $`A`$ in $`𝔇`$, and a 2-morphism $`\theta _c`$ in $``$ to every arrow $`c`$ in $`𝔇`$, as in the following diagram:
We require that $`\theta _{\mathrm{id}}=\mathrm{id}`$, and that $`\theta _h`$ satisfy the obvious compatibility conditions with respect to 2-morphisms and composition of morphisms.
A transformation between two weak transformations $`T`$, $`S`$, sometimes called a modification, is a rule to assign to each object $`A𝔇`$ a 2-morphism $`\mu _A`$ in $``$ as in the following diagram:
The 2-morphisms $`\mu _A`$ should satisfy the obvious compatibility relations with $`\theta _c`$ and $`\sigma _c`$, for every arrow $`c:AB`$ in $`𝔇`$. (Here $`\sigma _c`$ are for $`S`$ what $`\theta _c`$ are for $`T`$.) This relation can be written as $`\sigma _c\mu _A=\theta _c\mu _B`$.
###### Definition 17.1.
Given 2-categories $``$ and $`𝔇`$, we define the mapping 2-category $`\underset{¯}{om}(𝔇,)`$ to be the 2-category whose objects are strict 2-functors from $`𝔇`$ to $``$, whose 1-morphisms are weak 2-transformations between 2-functors, and whose 2-morphisms are modifications. When $``$ and $`𝔇`$ are 2-groupoids, then $`\underset{¯}{om}(𝔇,)`$ is also a 2-groupoid.
Viewing $`\mathrm{𝟐}𝐆𝐩`$ as a full subcategory of $`\mathrm{𝟐}𝐆𝐩𝐝`$, we can use the same notion for 2-groups as well. In fact, in the case of 2-groups, we are more interested in the pointed versions of the above definition. Namely, a pointed 2-transformation is required to satisfy the extra condition $`t_{}=\mathrm{id}`$. A pointed modification is, by definition, the identity modification!
For 2-groups $`𝔾`$ and $``$, we denote the 2-groupoid of pointed weak maps from $``$ to $`𝔾`$ by $`\underset{¯}{om}_{}(,𝔾)`$.
### 17.2. Nerve of a 2-category
We review the nerve construction for 2-categories, and recall its basic properties \[MoSe\], \[No1\].
Let $`𝔾`$ be a 2-category . We define the nerve of $`𝔾`$, denoted by $`N𝔾`$, to be the simplicial set defined as follows. The set of of 0-simplices of $`N𝔾`$ is the set of objects of $`𝔾`$. The 1-simplices are the morphisms in $`𝔾`$. The 2-simplices are diagrams of the form
where $`\alpha :fgh`$ is a 2-morphism. The 3-simplices of $`N𝔾`$ are commutative tetrahedra of the form
Commutativity of the above tetrahedron means $`(f\gamma )(\beta )=(\alpha m)(\delta )`$. That is, the following square of transformations is commutative:
For $`n3`$, an $`n`$-simplex of $`N𝔾`$ is an $`n`$-simplex such that each of its sub 3-simplices is a commutative tetrahedron as described above. In other words, $`N𝔾`$ is the coskeleton of the 3-truncated simplicial set $`\{N𝔾_0,N𝔾_1,N𝔾_2,N𝔾_3\}`$ defined above.
The nerve gives us a functor $`N:\mathrm{𝟐}𝐂𝐚𝐭\mathrm{𝐒𝐒𝐞𝐭}`$, where $`\mathrm{𝐒𝐒𝐞𝐭}`$ is the the category of simplicial sets.
### 17.3. Moerdijk-Svensson closed model structure on 2-groupoids
We give a quick review of the Moerdijk-Svensson closed model structure on the category of 2-groupoids. The main reference is \[MoSe\].
###### Definition 17.2.
Let $``$ and $`𝔾`$ be 2-groupoids, and $`P:𝔾`$ a functor between them. We say that $`P`$ is a fibration, if it satisfies the following properties:
* For every arrow $`a:A_0A_1`$ in $`𝔾`$, and every object $`B_1`$ in $``$ such that $`P(B_1)=A_1`$, there is an object $`B_0`$ in $``$ and an arrow $`b:B_0B_1`$ such that $`P(b)=a`$.
* For every 2-morphism $`\alpha :a_0a_1`$ in $`𝔾`$ and every arrow $`b_1`$ in $``$ such that $`P(b_1)=a_1`$, there is an arrow $`b_0`$ in $``$ and a 2-morphism $`\beta :b_0b_1`$ such that $`P(\beta )=\alpha `$.
###### Definition 17.3.
Let $`𝔾`$ be a 2-groupoid, and $`A`$ an object in $`𝔾`$. We define the following.
* $`\pi _0𝔾`$ is the set of equivalence classes of objects in $`𝔾`$.
* $`\pi _1(𝔾,A)`$ is the group of 2-isomorphism classes of arrows from $`A`$ to itself. The fundamental groupoid $`\mathrm{\Pi }_1𝔾`$ is the groupoid whose objects are the same as those of $`𝔾`$ and whose morphisms are 2-isomorphism classes of 1-morphisms in $`𝔾`$.
* $`\pi _2(𝔾,A)`$ is the group of 2-automorphisms of the identity arrow $`1_A:AA`$.
These invariants are functorial with respect to 2-functors. A map $`𝔾`$ is called a (weak) equivalence of 2-groupoids if it induces a bijection on $`\pi _0`$, $`\pi _1`$ and $`\pi _2`$, for every choice of a base point.
Having defined the notions of fibration and equivalence between 2-groupoids, we define cofibrations using the left lifting property. There is a more explicit description of cofibrations which can be found in (\[MoSe\] page 194), but we skip it here.
###### Theorem 17.4 (\[MoSe\], Theorem 1.2).
With weak equivalences, fibrations and cofibrations defined as above, the category of 2-groupoids has a natural structure of a closed model category.
The nerve functor is a bridge between the homotopy theory of 2-groupoids and the homotopy theory of simplicial sets. To justify this statement, we quote the following from \[MoSe\].
###### Proposition 17.5 (see \[MoSe\], Proposition 2.1).
* The functor $`N:\mathrm{𝟐}𝐂𝐚𝐭\mathrm{𝐒𝐒𝐞𝐭}`$ is faithful, preserves fiber products, and sends transformations between 2-functors to simplicial homotopies.
* The functor $`N`$ sends a fibration between 2-groupoids (Definition 17.2) to a Kan fibration. Nerve of every 2-groupoid is a Kan complex.
* For every (pointed) 2-groupoid $`𝔾`$ we have $`\pi _i(𝔾)\pi _i(N𝔾)`$, $`i=0,1,2`$.
* A map $`f:𝔾`$ of 2-groupoids is an equivalence if and only if $`Nf:NN𝔾`$ is a weak equivalence of simplicial sets.
###### Remark 17.6.
We can think of a 2-group as a 2-groupoid with one object. This identifies $`\mathrm{𝟐}𝐆𝐩`$ with a full subcategory of $`\mathrm{𝟐}𝐆𝐩𝐝`$. So, we can talk about nerves of 2-groups. This is a functor $`N:\mathrm{𝟐}𝐆𝐩\mathrm{𝐒𝐒𝐞𝐭}_{}`$, where $`\mathrm{𝐒𝐒𝐞𝐭}_{}`$ is the category of pointed simplicial sets. The above proposition remains valid if we replace 2-groupoids by 2-groups and $`\mathrm{𝐒𝐒𝐞𝐭}`$ by $`\mathrm{𝐒𝐒𝐞𝐭}_{}`$ throughout.
The functor $`N:\mathrm{𝟐}𝐆𝐩𝐝\mathrm{𝐒𝐒𝐞𝐭}`$ has a left adjoint $`W:\mathrm{𝐒𝐒𝐞𝐭}\mathrm{𝟐}𝐆𝐩𝐝`$, called the Whitehead 2-groupoid whose definition can be found in (\[MoSe\] page 190, Example 2).
It is easy to see that $`W`$ preserves homotopy groups. In particular, it sends weak equivalences of simplicial sets to equivalences of 2-groupoids. Much less obvious is the following
###### Theorem 17.7 (\[MoSe\], $`\mathrm{\S }`$2).
The pair
$$W:\mathrm{𝐒𝐒𝐞𝐭}\mathrm{𝟐}𝐆𝐩𝐝:N$$
is a Quillen pair. It satisfies the following properties:
* Each adjoint preserves weak equivalences.
* For every 2-groupoid $`𝔾`$, the counit $`WN(𝔾)𝔾`$ is a weak equivalence
* For every simplicial set $`X`$ such that $`\pi _iX=0`$, $`i3`$, the unit of adjunction $`XNW(X)`$ is a weak equivalence.
In particular, the functor $`N:\mathrm{Ho}(\mathrm{𝟐}𝐆𝐩𝐝)\mathrm{Ho}(\mathrm{𝐒𝐒𝐞𝐭})`$ induces an equivalence of categories between $`\mathrm{Ho}(\mathrm{𝟐}𝐆𝐩𝐝)`$ and the category of homotopy 2-types. (The latter is defined to be the full subcategory of $`\mathrm{Ho}(\mathrm{𝐒𝐒𝐞𝐭})`$ consisting of all $`X`$ such that $`\pi _iX=0`$, $`i3`$.)
###### Remark 17.8.
The pointed version of the above theorem is also valid. The proof is just a minor modification of the proof of the above theorem.
The following Proposition follows from Theorem 17.7 and Remark 17.8.
###### Proposition 17.9.
The functor $`N:\mathrm{Ho}(\mathrm{𝟐}𝐆𝐩)\mathrm{Ho}(\mathrm{𝐒𝐒𝐞𝐭}_{})`$ induces an equivalence between $`\mathrm{Ho}(\mathrm{𝟐}𝐆𝐩)`$ and the full subcategory of $`\mathrm{Ho}(\mathrm{𝐒𝐒𝐞𝐭}_{})`$ consisting of connected pointed homotopy 2-types.
It is also well-known that the geometric realization functor $`||:\mathrm{𝐒𝐒𝐞𝐭}_{}\mathrm{𝐓𝐨𝐩}_{}`$ induces an equivalence of of homotopy categories. So we have the following
###### Corollary 17.10.
The functor $`|N()|:\mathrm{Ho}(\mathrm{𝟐}𝐆𝐩)\mathrm{Ho}(\mathrm{𝐓𝐨𝐩}_{})`$ induces an equivalence between $`\mathrm{Ho}(\mathrm{𝟐}𝐆𝐩)`$ and the full subcategory of $`\mathrm{Ho}(\mathrm{𝐓𝐨𝐩}_{})`$ consisting of connected pointed homotopy 2-types.
The following proposition says that derived mapping 2-groupoids have the correct homotopy type. We have denoted the category $`\{01\}`$ by $`𝐈`$.
###### Proposition 17.11.
Let $`𝔾`$ and $``$ be 2-groupoids. Then there is a natural homotopy equivalence
$$N\underset{¯}{om}(,𝔾)\mathrm{𝐇𝐨𝐦}(N,N𝔾),$$
where the left hand side is defined to be $`N\underset{¯}{om}(𝔽,𝔾)`$, where $`𝔽`$ is a cofibrant replacement for $``$, and the right hand side is the simplicial mapping space. In particular, $`N\underset{¯}{om}(𝐈,𝔾)(N𝔾)^{\mathrm{\Delta }^1}`$, that is, the nerve of the path category is naturally homotopy equivalent to the path space of the nerve.
We also have the pointed version of the above proposition.
###### Proposition 17.12.
Let $`𝔾`$ and $``$ be 2-groups. Then, there is a natural homotopy equivalence
$$N\underset{¯}{om}_{}(,𝔾)\mathrm{𝐇𝐨𝐦}_{}(N,N𝔾).$$
###### Remark 17.13.
In the definition of $`\underset{¯}{om}`$ we have used strict 2-functors as objects, but the arrows are weak 2-transformation. There is a variant of this in which the 2-functors are also weak, and this leads to different simplicial mapping spaces between 2-group(oid)s. If in Propositions 17.11 and 17.12 above we used this simplicial mapping space on the left hand side, we would not need to use a cofibrant replacement on $``$ to compute the derived mapping spaces. In other words, the derived mapping spaces would coincide with the actual mapping spaces.
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# Large time behavior of heat kernels on forms
## 1. Introduction
The goal of the present paper is to establish large time, pointwise bounds for the heat kernel on the vector bundles of forms on some noncompact manifolds.
Information on large time behavior of heat kernels on forms usually leads to interesting analytical and topological information on the manifolds. In fact heat kernel on forms contains much more information on the interplay between analysis, geometry and topology than that on functions. So far much effort has been spent on the study of short time and long time behavior of heat kernel on forms, in the case of closed manifolds, see for instance , . By contrast, the present paper is to our knowledge the first one to offer estimates for the heat kernel on one-forms on a class of non-compact Riemannian manifolds with a meaningful contents for large time, i.e. without an increasing exponential factor (see for instance ).
Let $`M`$ be a complete connected Riemannian manifold. Denote by $`d(x,y)`$ the geodesic distance between two points $`x,yM`$, and by $`B(x,r)`$ the open ball of center $`xM`$. Let $`\mu `$ be the Riemannian measure; denote also by $`|\mathrm{\Omega }|`$ the measure $`\mu (\mathrm{\Omega })`$ of a mesurable subset $`\mathrm{\Omega }`$ of $`M`$. Denote by $`\mathrm{\Delta }`$ the (non-negative) Laplace-Beltrami operator on functions. The heat semigroup on functions $`e^{t\mathrm{\Delta }}`$ will also be denoted by $`P_t`$, and the corresponding heat kernel by $`p_t(x,y)`$, $`t>0`$, $`x,yM`$. We will use $`\stackrel{}{\mathrm{\Delta }}`$ to denote the Hodge Laplacian on forms. The heat semigroup on forms $`e^{t\stackrel{}{\mathrm{\Delta }}}`$ will also be denoted by $`\stackrel{}{P}_t`$, and the corresponding heat kernel by $`\stackrel{}{p}_t(x,y)`$, $`t>0`$, $`x,yM`$.
The main question we shall address below is the following:
Given an upper estimate for the heat kernel on functions, under which additional assumptions can one deduce an upper bound for the heat kernel on forms?
We shall consider in particular the case where $`M`$ has the so-called volume doubling property and the heat kernel on functions satisfies a Gaussian upper estimate, that is
$$p_t(x,y)\frac{C}{|B(x,\sqrt{t})|}\mathrm{exp}(cd^2(x,y)/t),x,yM,t>0,$$
for some $`C,c>0`$.
For instance, when $`M`$ has non-negative Ricci curvature, it was proved in that $`p_t`$ satisfies a Gaussian upper estimate, and in that case, the answer to the above question for $`1`$-forms is straightforward by the semigroup domination theory (see (1.2) below, and also e.g. , , , ): the heat kernel on $`1`$-forms is also bounded from above by a Gaussian.
The following simple example shows that this may be false in general. Let $`M`$ be the connected sum of two copies of $`\mathrm{I}\mathrm{R}^n`$, $`n3`$. It is known that the heat kernel on functions has a Gaussian upper bound (see ). If the heat kernel on 1-forms $`\stackrel{}{p}_t`$ also had a Gaussian upper bound, then by , pp.1740–1741, it would follow that
$$|p_t(x,y)|\frac{C^{}}{\sqrt{t}|B(x,\sqrt{t})|}\mathrm{exp}(c^{}d^2(x,y)/t),x,yM,t>0.$$
A classical argument shows that $`p_t`$ would then be bounded below by a Gaussian (see for instance ). This is false as was noticed in . See for more on this example.
Another case where the behaviour of the heat kernel on $`1`$-forms is well understood is the case of the Heisenberg group and more generally stratified Lie groups, see , .
In the present paper, we are going to see that if the negative part of the Ricci curvature is small enough in some sense, then the upper bound on the heat kernel on $`1`$-forms differs from that on the heat kernel on functions at most by a certain power of time $`t`$. One can state similar results for higher degree forms by replacing in the assumptions the Ricci curvature by a suitable curvature operator (see ). For convenience, we shall however formulate our assumptions and results in the case of $`1`$-forms. We leave the formulation of the general case to the reader.
In this article, all Riemannian manifolds under consideration will be complete non-compact. Let us layout some basic assumptions to be used below.
Assumption (A). $`M`$ satisfies the volume doubling property.
$$|B(x,2r)|C|B(x,r)|$$
for all $`xM`$, $`r>0`$ and some $`C>0`$.
Assumption (B). The heat kernel $`p_t(x,y)`$ on functions satisfies a Gaussian upper bound:
$$p_t(x,y)\frac{C}{|B(x,\sqrt{t})|}\mathrm{exp}(cd^2(x,y)/t),$$
for some $`C,c>0`$, and all $`x,yM`$ and $`t>0`$.
It was proved in that Assumptions $`(A)`$ and $`(B)`$ together are equivalent to the following relative Faber-Krahn inequality:
For all $`xM`$, $`r>0`$, and every non-empty subset $`\mathrm{\Omega }B(x,r)`$,
$`(FK)`$
$$\lambda _1(\mathrm{\Omega })\frac{c}{r^2}\left(\frac{|B(x,r)|}{|\mathrm{\Omega }|}\right)^{2/\nu }.$$
Here $`\lambda _1(\mathrm{\Omega })`$ is the first Dirichlet eigenvalue of $`\mathrm{\Omega }`$ and $`c>0`$.
Note that $`(FK)`$ implies
(1.1)
$$\frac{|B(x,s)|}{|B(x,r)|}C\left(\frac{s}{r}\right)^\nu ,$$
for all $`s>r>0`$, $`xM`$, and we shall use Assumption $`(A)`$ in this form.
Assumption (C). The Ricci curvature is bounded from below by a negative constant.
It follows from Assumption $`(C)`$ and Bishop’s comparison theorem that there exists $`C>0`$ such that $`|B(x,1)|C`$ for all $`xM`$. We shall often also need the opposite inequality.
Assumption (D). Non-collapsing of the volume of balls: there exists $`c>0`$ such that $`|B(x,1)|c`$ for all $`xM`$.
It is well-known that to estimate the heat kernel acting on one-forms, it is enough to estimate the kernel of a certain Schrödinger semigroup acting on functions, whose potential is the negative part of the Ricci curvature.
Indeed, Bochner’s formula states
$$\stackrel{}{\mathrm{\Delta }}=D^{}DRic.$$
Here $`D`$ is the covariant derivative on $`1`$-forms and $`Ric`$ is the Ricci curvature.
Let $`\lambda =\lambda (x)`$ be the lowest eigenvalue of $`Ric(x)`$, $`xM`$. We will use the notation
$$V(x)=\lambda ^{}(x)=(|\lambda (x)|\lambda (x))/2.$$
Let $`P_t^V`$ be the semigroup $`e^{t(\mathrm{\Delta }V)}`$. Under Assumption $`(C)`$, $`P_t^V`$ has a kernel which we shall denote by $`p_t^V(x,y)`$.
Let us recall the semi-group domination property, which was proved in :
(1.2)
$$|\stackrel{}{p}_t(x,y)|p_t^V(x,y)$$
for all $`x,yM`$ and $`t>0`$. Here $`\stackrel{}{p}_t(x,y)`$ is a linear operator from the tangent space $`T_yM`$ to $`T_xM`$, and here $`|\stackrel{}{p}_t(x,y)|`$ denotes its operator norm with respect to the Riemannian metrics.
We can now introduce one of our main curvature assumptions. An important property of the Hodge Laplacian $`\stackrel{}{\mathrm{\Delta }}`$ is that it is a nonnegative operator (as a consequence, $`\stackrel{}{P}_t`$ is contractive on $`L^2(M,T^{}M)`$). This means that, for every smooth compactly supported $`1`$-form $`\varphi `$,
(1.3)
$$_MRic(\varphi (x),\varphi (x))d\mu (x)_M|D\varphi (x)|^2𝑑\mu (x).$$
Now, by the Kato inequality
$$||\varphi |||D\varphi |,$$
and the fact that by definition
$$Ric(\varphi (x),\varphi (x))V(x)|\varphi (x)|^2,$$
we see that condition (1.3) is implied by
$$_MV(x)f^2(x)𝑑\mu (x)_M|f(x)|^2𝑑\mu (x),fC_0^{\mathrm{}}(M),$$
which means that $`\mathrm{\Delta }V`$ is a positive operator on $`L^2(M)`$.
We shall say that $`\mathrm{\Delta }V`$ is strongly positive (strongly subcritical in the sense of ) if it satisfies the following stronger condition: there exists $`A<1`$ such that, for all $`fC_0^{\mathrm{}}(M)`$,
(1.4)
$$_MVf^2𝑑\mu A_M|f|^2𝑑\mu .$$
The above condition sometimes is referred to as the form boundedness condition, which has its origin in the Hardy type inequality, for $`fC_0^{\mathrm{}}(\mathrm{I}\mathrm{R}^n)`$, $`n3`$,
$$\frac{(n2)^2}{4}_{\mathrm{I}\mathrm{R}^n}\frac{f^2(x)}{|x|^2}𝑑x_{\mathrm{I}\mathrm{R}^n}|f(x)|^2𝑑x.$$
For generalizations of the above inequality to the manifold case, see .
Example: If the manifold $`M`$ satisfies the Euclidean Sobolev inequality of dimension $`n`$
$$\left(_M|f|^{2n/(n2)}𝑑\mu \right)^{\frac{n2}{n}}C_M|f|^2𝑑\mu $$
for all $`fC_0^{\mathrm{}}(M)`$, for some $`n>2`$, and if $`VL^{n/2}(M)`$ with sufficiently small norm, then it is easy to see by using Hölder’s inequality that (1.4) holds.
Let us now summarize our results. Under Assumptions $`(A)`$ to $`(D)`$, the function $`V`$, the negative part of the lowest eigenvalue of the Ricci curvature largely determines the upper bound of heat hernel on $`1`$-forms. If $`V`$ is sufficiently small in certain integral sense, then $`\stackrel{}{p}_t`$ has Gaussian upper bound, which has important consequences in terms of $`L^p`$ boundedness of the Riesz transform. This is explained in Section 3. Otherwise the upper bound for $`\stackrel{}{p_t}`$ is a Gaussian times a suitable power of time $`t`$, provided that the operator $`\mathrm{\Delta }V`$ is strongly positive. The proof of this fact is contained in Sections 2.1, 2.2, 2.3. We also consider the case where the Ricci curvature is nonnegative outside of a compact set. Without any other assumptions on the Ricci curvature, there may be $`L^2`$ harmonic forms, therefore one cannot expect a decay with respect to time in general, but we show that $`\stackrel{}{p_t}`$ is bounded by a Gaussian plus the product of the Green’s function of the Laplacian in both variables. This is the subject of Section 4. Finally, we treat in Section 5 the case where the heat kernel on functions has an arbitrary uniform decay.
## 2. Bounds for the heat kernel on forms and strong positivity
Our aim in this section is the following result.
###### Theorem 2.1.
Suppose $`M`$ satisfies Assumptions $`(A)`$, $`(B)`$, $`(C)`$, and $`(D)`$, and that the operator $`\mathrm{\Delta }V`$ is strongly positive with constant $`A`$. Suppose in addition that $`VL^p(M,\mu )`$ for some $`p[1,+\mathrm{})`$. Then, if $`p=1`$, for any $`0<c<1/4`$, and any $`\epsilon >0`$, there exists $`C>0`$ such that
$$|\stackrel{}{p}_t(x,y)|C\mathrm{min}\{\frac{t^{(1+\epsilon )A}}{|B(x,\sqrt{t})|},\mathrm{\hspace{0.17em}1}\}\mathrm{exp}(cd^2(x,y)/t),x,yM,t1;$$
if $`p(1,2)`$, for any $`c<1/4`$, there exists $`C,\epsilon >0`$ such that
$$|\stackrel{}{p}_t(x,y)|C\mathrm{min}\{\frac{t^{(p\epsilon )A}}{|B(x,\sqrt{t})|},\mathrm{\hspace{0.17em}1}\}\mathrm{exp}(cd^2(x,y)/t),x,yM,t1;$$
If $`p2`$, for any $`c<1/4`$ and any $`\epsilon >0`$, there exists $`C`$ such that
$$|\stackrel{}{p}_t(x,y)|C\mathrm{min}\{\frac{t^{(p1+\epsilon )A}}{|B(x,\sqrt{t})|},\mathrm{\hspace{0.17em}1}\}\mathrm{exp}(cd^2(x,y)/t),x,yM,t1.$$
Remark: Using Assumption $`(C)`$ and the Gaussian bound on $`p_t`$, one easily obtains the following small time estimate
(2.1)
$$p_t^V(x,y)\frac{C_{t_0}}{|B(x,\sqrt{t})|}\mathrm{exp}(cd^2(x,y)/t),$$
for all $`x,yM`$ and $`0<t<t_0`$.
Thanks to the domination property (1.2), Theorem 2.1 is a consequence of the following statement, which is of independent interest.
###### Theorem 2.2.
Suppose $`M`$ satisfies Assumption $`(C)`$, and that the operator $`\mathrm{\Delta }V`$ is strongly positive with constant $`A`$. Suppose in addition that $`VL^p(M,\mu )`$ for some $`p[1,+\mathrm{})`$. Then, if $`p=1`$, for any $`0<c<1/4`$ and any $`\epsilon >0`$, there exists $`C>0`$ such that
$$p_t^V(x,y)C\mathrm{min}\{\frac{t^{(1+\epsilon )A}}{|B(x,\sqrt{t})|},\mathrm{\hspace{0.17em}1}\}\mathrm{exp}(cd^2(x,y)/t),x,yM,t1;$$
if $`p(1,2)`$, for any $`c<1/4`$, there exists $`C,\epsilon >0`$ such that
$$p_t^V(x,y)C\mathrm{min}\{\frac{t^{(p\epsilon )A}}{|B(x,\sqrt{t})|},\mathrm{\hspace{0.17em}1}\}\mathrm{exp}(cd^2(x,y)/t),x,yM,t1;$$
If $`p2`$, for any $`c<1/4`$ and any $`\epsilon >0`$, there exists $`C`$ such that
$$p_t^V(x,y)C\mathrm{min}\{\frac{t^{(p1+\epsilon )A}}{|B(x,\sqrt{t})|},\mathrm{\hspace{0.17em}1}\}\mathrm{exp}(cd^2(x,y)/t),x,yM,t1.$$
We would like to mention a number of previous papers that deal with Schrödinger heat kernels on manifolds. In the paper , a fundamental gradient estimate was derived for the heat kernel. As far as long time behavior is concerned, the emphasis is on the case without potential and nonnegative Ricci curvature. The paper studied Schrödinger heat kernels with singular oscillating potentials. The papers established long time behavior for Schrödinger heat kernels on manifolds with nonnegative Ricci curvature for potentials essentially behaving as negative powers of the distance function. The case of potentials with polynomial growth and magnetic field is considered in .
Here is the plan of the proof of Theorem 2.2.
In section 2.1, we show that, given the upper bound on $`p_t`$ and the strong positivity of $`\mathrm{\Delta }V`$, a pointwise upper bound on $`p_t^V`$ follows from an adaptation of the Nash method due to Grigor’yan, provided some $`L^1`$ to $`L^1`$ estimates for $`P_t^V`$ are available. In section 2.2, we prove such estimates under the other assumptions of Theorem 2.2, and we finish the proof.
### 2.1. Pointwise estimates
Let us first prove the following preliminary estimate.
###### Proposition 2.1.
Suppose $`M`$ satisfies Assumption $`(C)`$, and that the operator $`\mathrm{\Delta }V`$ is strongly positive. Then there exist $`C,c>0`$ such that
(2.2)
$$p_t^V(x,y)C\mathrm{exp}\left(cd^2(x,y)/t\right),x,yM,t1.$$
###### Proof.
This can be proven by a standard method of using exponential weights as in . An alternative way is to use wave equation method as in or .
Fix $`y`$ and write
$$u(x,t)=p_t^V(x,y),$$
$$I(t)=_Mu^2(x,t)w(x,t)𝑑\mu (x)$$
where $`w(x,t)=e^{\frac{d^2(x,y)}{Dt}}`$ for some $`D>0`$ to be chosen later.
One has
$`{\displaystyle \frac{d}{dt}}I(t)`$ $`=`$ $`{\displaystyle \frac{d}{dt}}{\displaystyle _M}u^2w𝑑\mu `$
$`=`$ $`2{\displaystyle _M}uw(\mathrm{\Delta }u+Vu),d\mu {\displaystyle _M}{\displaystyle \frac{u^2wd^2(x,y)}{Dt^2}}𝑑\mu (x).`$
This implies, after integration by parts,
$`{\displaystyle \frac{d}{dt}}I(t)`$ $`=`$ $`2{\displaystyle _M}|u|^2w𝑑\mu 2{\displaystyle _M}uwud\mu `$
$`+2{\displaystyle _M}V(u\sqrt{w})^2𝑑\mu {\displaystyle _M}{\displaystyle \frac{u^2wd^2(x,y)}{Dt^2}}𝑑\mu (x)`$
$`=`$ $`2{\displaystyle _M}|u|^2wd\mu 2{\displaystyle _M}uw{\displaystyle \frac{2d(x,y}{Dt}})d(x,y)ud\mu (x)`$
$`+2{\displaystyle _M}V(u\sqrt{w})^2𝑑\mu {\displaystyle _M}{\displaystyle \frac{u^2wd^2(x,y)}{Dt^2}}𝑑\mu (x)`$
$``$ $`2{\displaystyle _M}|u|^2w𝑑\mu +4{\displaystyle _M}uw{\displaystyle \frac{d(x,y)}{Dt}}|u|𝑑\mu (x)`$
$`+2{\displaystyle _M}V(u\sqrt{w})^2𝑑\mu {\displaystyle _M}{\displaystyle \frac{u^2wd^2(x,y)}{Dt^2}}𝑑\mu (x)`$
$``$ $`2{\displaystyle _M}|u|^2w𝑑\mu +C_\epsilon {\displaystyle _M}{\displaystyle \frac{u^2wd^2(x,y)}{D^2t^2}}𝑑\mu (x)`$
$`+\epsilon {\displaystyle _M}|u|^2w𝑑\mu +2{\displaystyle _M}V(u\sqrt{w})^2𝑑\mu {\displaystyle _M}{\displaystyle \frac{u^2wd^2(x,y)}{Dt^2}}𝑑\mu (x)`$
$``$ $`(2\epsilon ){\displaystyle _M}|u|^2w𝑑\mu +2{\displaystyle _M}V(u\sqrt{w})^2𝑑\mu {\displaystyle _M}{\displaystyle \frac{u^2wd^2(x,y)}{2Dt^2}}𝑑\mu (x)`$
for arbitrarily small $`\epsilon >0`$, provided $`D`$ is chosen large enough. Finally, using the strong positivity of $`V`$, we obtain
$$\frac{d}{dt}I(t)(2\epsilon )_M|u|^2w𝑑\mu +2A_M|(u\sqrt{w})|^2𝑑\mu _M\frac{u^2wd^2(x,y)}{2Dt^2}𝑑\mu (x).$$
Observe that
$`|(u\sqrt{w})|^2`$ $`=`$ $`|u|^2w+uuw+{\displaystyle \frac{u^2|w|^2}{4w}}`$
$``$ $`|u|^2w+2u|u|{\displaystyle \frac{d(x,y)}{Dt}}w+{\displaystyle \frac{d^2(x,y)u^2w}{D^2t^2}}`$
$``$ $`|u|^2w+2C{\displaystyle \frac{d^2(x,y)}{D^2t^2}}u^2w+{\displaystyle \frac{2}{C}}|u|^2w+{\displaystyle \frac{d^2(x,y)u^2w}{D^2t^2}}.`$
We find, since $`A<1`$, that
$$\frac{d}{dt}I(t)=\frac{d}{dt}_Mu^2(x,t)e^{\frac{d^2(x,y)}{Dt}}𝑑\mu (x)0$$
when $`D`$ is sufficiently large.
In particular,
$$_M\left(p_t^V(x,y)\right)^2e^{\frac{d^2(x,y)}{Dt}}𝑑\mu (x)_M\left(p_1^V(x,y)\right)^2e^{\frac{d^2(x,y)}{D}}𝑑\mu (x)$$
when $`t1`$. By (2.1) and the small time Gaussian estimate for $`p_t`$ under Assumption $`(C)`$ (see ),
$$p_1^V(x,y)\frac{C}{|B(x,1)|}e^{cd^2(x,y)}.$$
Using the well-known fact that a manifold satisfying $`(C)`$ has at most exponential volume growth around any point, we have, for $`D>0`$ large enough,
$$_M\left(p_1^V(x,y)\right)^2e^{\frac{d^2(x,y)}{D}}𝑑\mu (x)C,$$
hence
$$_M\left(p_t^V(x,y)\right)^2e^{\frac{d^2(x,y)}{Dt}}𝑑\mu (x)C,t1.$$
Next, using the semigroup property
$`p_{2t}^V(x,y)`$ $`=`$ $`{\displaystyle _M}p_t^V(x,z)p_t^V(z,y)𝑑\mu (z)`$
$`=`$ $`{\displaystyle _M}e^{\frac{d^2(x,z)}{2Dt}}p_t^V(x,z)e^{\frac{d^2(z,y)}{2Dt}}p_t^V(z,y)e^{\frac{d^2(x,z)}{2Dt}\frac{d^2(z,y)}{2Dt}}𝑑\mu (z)`$
$``$ $`e^{\frac{d^2(x,y)}{4Dt}}\left[{\displaystyle _M}e^{\frac{d^2(x,z)}{Dt}}\left(p_t^V(x,z)\right)^2𝑑z\right]^{1/2}\left[{\displaystyle _M}e^{\frac{d^2(y,z)}{Dt}}\left(p_t^V(y,z)\right)^2𝑑\mu (z)\right]^{1/2}.`$
Hence
$$p_t^V(x,y)Ce^{cd^2(x,y)/t},t2.$$
This proves the claim. ∎
We can now state our main technical result, which is an adaptation of the Nash method to the case where the semigroup under consideration is not necessarily contractive on $`L^1`$. The argument is based on the one in the proof of , Theorem 1.1 (see also , Proposition 8.1), with certain modification and localization. If $`T`$ is a operator from $`L^{p_1}`$ to $`L^{p_2}`$, then $`T_{p_1,p_2}`$ will denote the operator norm $`sup_{fL^{p_1}\{0\}}\frac{Tf_{p_2}}{f_{p_1}}.`$
###### Proposition 2.2.
Let $`M`$ satisfy Assumptions (A), (B), (C), and (D). Suppose that $`\mathrm{\Delta }V`$ is strongly positive and that there exists an non-decreasing function $`F`$ such that
(2.3)
$$P_t^V_{1,1}F(t),t1.$$
Then there exists $`C>0`$ such that
$$p_t^V(x,x)C\frac{F^2(t)[\mathrm{ln}(e+tF(t))]^{\nu /2}}{|B(x,\sqrt{t})|},$$
for all $`xM`$ and $`t1`$. where $`\nu >0`$ is the constant from (1.1). If in addition $`F`$ satisfies
$$F(2t)CF(t),t>0,$$
for some $`C>0`$, then for any $`c(0,1/4)`$, there exists $`C>0`$ such that
(2.4)
$$p_t^V(x,y)C\frac{F^2(t)[\mathrm{ln}(e+tF(t))]^{\nu /2}}{|B(x,\sqrt{t})|}\mathrm{exp}(cd^2(x,y)/t),$$
for all $`x,yM`$ and $`t1`$.
Remark. If there exists $`N`$ such that $`F(t)Ct^N`$, the above estimate takes the following simpler form:
$$p_t^V(x,y)C\frac{F^2(t)[\mathrm{ln}(e+t)]^{\nu /2}}{|B(x,\sqrt{t})|}\mathrm{exp}(cd^2(x,y)/t).$$
###### Proof.
Fix $`x_0M`$, write $`u(x,t)=p_t^V(x,x_0)`$, $`t>0`$, $`xM`$, and set
$$I(t)=_Mu^2(x,t)𝑑\mu (x)=p_{2t}^V(x_0,x_0).$$
Then
$$I^{}(t)=2_Mu(x,t)(\mathrm{\Delta }uVu)(x,t)𝑑\mu (x)=2_M|u|^2𝑑\mu +2_MVu^2𝑑\mu .$$
Using assumption (1.4), we have
(2.5)
$$I^{}(t)2(1A)_M|u(x,t)|^2𝑑\mu (x).$$
Since, for any $`s>0`$,
$$u^2(us)_+^2+2su,$$
we can write
$$I(t)_{\{x|u(x,t)>s\}}(u(x,t)s)^2𝑑\mu (x)+2s_Mu(x,t)𝑑\mu (x).$$
By assumption (2.3) and the definition of $`\lambda _1`$, this yields
$$I(t)\frac{_{\{x|u(x,t)>s\}}|(u(x,t)s)|^2𝑑\mu (x)}{\lambda _1(\{x|u(x,t)>s\})}+2sF(t),$$
hence
(2.6)
$$I(t)\frac{_{\{x|u(x,t)>s\}}|u(x,t)|^2𝑑\mu (x)}{\lambda _1(\{x|u(x,t)>s\})}+2sF(t).$$
The bound (2.2) yields
$$u(x,t)C_1e^{c_2d^2(x,x_0)/t},$$
thus
$$\{x|u(x,t)>s\}\{x|e^{c_2d^2(x,x_0)/t}>s/C_1\}=\{x|d^2(x,x_0)<c_2^1t\mathrm{ln}(C_1/s)\}.$$
Thus
$$\{x|u(x,t)>s\}B(x_0,r),$$
where
$$r=\sqrt{c_2^1t(|\mathrm{ln}(C_1/s)|+1)}$$
(we choose to take $`rc\sqrt{t}`$ for later convenience). According to $`(FK)`$, we have
$$\lambda _1(\{x|u(x,t)>s\})\frac{c}{r^2}\left(\frac{|B(x_0,r)|}{|\{x|u(x,t)>s\}|}\right)^{2/\nu }.$$
On the other hand,
$$|\{x|u(x,t)>s\}|s^1_Mu(x,t)𝑑\mu (x)s^1F(t).$$
Therefore
(2.7)
$$\lambda _1(\{x|u(x,t)>s\})\frac{c}{r^2}\left(\frac{s|B(x_0,r)|}{F(t)}\right)^{2/\nu }:=m(s,t,x_0).$$
Plugging this into (2.6), we obtain
$$I(t)\frac{_{\{x|u(x,t)>s\}}|u(x,t)|^2𝑑\mu (x)}{m(s,t,x_0)}+2sF(t).$$
Hence
(2.8)
$$_{\{x|u(x,t)>s\}}|u(x,t)|^2𝑑\mu (x)\left(I(t)2sF(t)\right)m(s,t,x_0).$$
The combination of (2.8) and (2.5) yields
$$I^{}(t)2(1A)\left(I(t)2sF(t)\right)m(s,t,x_0).$$
Choosing $`s`$ so that $`sF(t)=I(t)/4`$ yields
(2.9)
$$I^{}(t)(1A)I(t)m(s,t,x_0),$$
for all $`t>0`$ and the corresponding $`s`$.
We have that
$$I(t)=p_{2t}^V(x_0,x_0)c/t^{\nu /2},$$
for $`t1`$. Indeed, since $`V0`$, by the maximum principle $`p_t^V(x,x_0)p_t(x,x_0)`$. Now it is well known (see ) that Assumptions $`(A)`$ and $`(B)`$ imply
$$p_{2t}(x_0,x_0)\frac{c}{|B(x_0,\sqrt{t})|}.$$
One concludes by using Assumptions $`(A)`$ and $`(D)`$.
Let us now estimate $`m(s,t,x_0)`$. First, for $`t1`$,
$`c\sqrt{t}r`$ $`=`$ $`\sqrt{c_2^1t(|\mathrm{ln}(C_1/s)|+1)}`$
$`=`$ $`\sqrt{c_2^1t(|\mathrm{ln}(4C_1F(t)/I(t))|+1)}`$
$``$ $`\sqrt{c_2^1t(|\mathrm{ln}(CF(t)t^{\nu /2})|+1)}`$
$``$ $`C\sqrt{t}\sqrt{\mathrm{ln}(e+tF(t))}.`$
Finally,
(2.10)
$$m(s,t,x_0)\frac{c}{t\mathrm{ln}(e+tF(t))}\left(\frac{I(t)|B(x_0,\sqrt{t})|}{F^2(t)}\right)^{2/\nu }.$$
By (2.9) and (2.10), it follows that
$$I^{}(t)c\frac{I(t)^{1+(2/\nu )}|B(x_0,\sqrt{t})|^{2/\nu }}{tF(t)^{4/\nu }\mathrm{ln}(e+tF(t))},$$
that is
(2.11)
$$\frac{I^{}(t)}{I(t)^{1+(2/\nu )}}\frac{c|B(x_0,\sqrt{t})|^{2/\nu }}{tF(t)^{4/\nu }\mathrm{ln}(e+tF(t))}.$$
Integrating (2.11) from $`t/2`$ to $`t`$ and using the monotonicity of $`F(t)`$ and $`|B(x_0,\sqrt{t})|`$, one easily obtains
(2.12)
$$p_{2t}^V(x_0,x_0)=I(t)C\frac{F^2(t)\left[\mathrm{ln}(e+tF(t))\right]^{\nu /2}}{|B(x_0,\sqrt{t})|}.$$
From this on-diagonal bound, we can derive the full bound (2.4) by either the method in or the wave equation method in .
### 2.2. The $`L^1`$ to $`L^1`$ estimates
###### Proposition 2.3.
Suppose that $`M`$ satisfies Assumptions (A), (B), (C), (D), and that the operator $`\mathrm{\Delta }V`$ is strongly positive. If $`VL^p(M,\mu )`$ for some $`p[1,+\mathrm{})`$, then there exists $`C=C(p)`$ such that
$$\{\begin{array}{cc}P_t^V_{1,1}Ct^{1/2},t1,\text{ if }p=1,\hfill & \\ P_t^V_{1,1}Ct^{(p\theta )/2},t1,\text{ if }1<p<2,\text{ for some }\theta =\theta (p)>0,\hfill & \\ P_t^V_{1,1}Ct^{(p1)/2},t1,\text{ if }p2.\hfill & \end{array}$$
Remark We shall see in Section 5 below that, if one weakens $`(A)`$, and $`(B)`$ as suggested in the remark after Proposition 2.1, one can still prove
$$P_t^V_{1,1}C_pt^{p/2},t1,1p<+\mathrm{}.$$
If now one only assumes that $`\mathrm{\Delta }V`$ is positive instead of being strongly positive, then one can still prove, for $`VL^p(M,\mu )`$ and $`t1`$,
$$P_t^V_{1,1}\{\begin{array}{cc}Ct,\text{ if }1p2;\hfill & \\ Ct^{p/2},\text{ if }p2;\hfill & \end{array}$$
Note that, according to , Theorem 3.1, the above estimate is sharp in the range $`1p2`$, which shows the role of strong positivity in the better estimate of Proposition 2.3. For $`p>2`$, similar estimates were proved in , Theorem 8.1, by extending the method of , Theorem 3, to the manifold case. This is also what we shall do to prove the more precise Proposition 2.3, taking advantage in addition of the strong positivity of $`V`$.
###### Proof.
Let us first prove that there exists $`\delta >0`$ such that
(2.13)
$$\underset{y}{sup}p_t^V(.,y)_2=P_t^V_{2,\mathrm{}}Ct^\delta ,t1.$$
The proof goes as follows.
Given $`fC_0^{\mathrm{}}(M)`$ and $`q>1`$, by an easy consequence of the Feynman-Kac formula (see for instance , p.712), one may write
(2.14)
$$|P_t^Vf(x)|[e^{t(\mathrm{\Delta }qV)}|f|(x)]^{1/q}[e^{t\mathrm{\Delta }}|f|(x)]^{(q1)/q}.$$
Since $`\mathrm{\Delta }V`$ is strongly positive, when $`q`$ is sufficiently close to $`1`$, the operator $`\mathrm{\Delta }qV`$ is also strongly positive, thus $`P_t^{qV}`$ is contractive on $`L^2(M,\mu )`$. Hence
(2.15)
$$e^{t(\mathrm{\Delta }qV)}|f|(x)P_t^{qV}_{2,\mathrm{}}f_2P_{t1}^{qV}_{2,2}P_1^{qV}_{2,\mathrm{}}f_2P_1^{qV}_{2,\mathrm{}}f_2,$$
for $`t1`$. Using the assumption that $`p_t`$ is bounded from above by a Gaussian (Assumption $`(B)`$), we have
$$e^{t\mathrm{\Delta }}|f|(x)\underset{y}{sup}p_t(y,.)_2f_2\frac{C}{\sqrt{|B(x,\sqrt{t})|}}f_2.$$
By Assumptions $`(A)`$ and $`(D)`$,
$$|B(x,\sqrt{t})|c_1t^{\nu /2}|B(x,1)|c_2t^{\nu /2}.$$
Therefore
$$e^{t\mathrm{\Delta }}|f|(x)\frac{C}{t^{\nu /4}}f_2.$$
Substituting this and (2.15) into (2.14), we have, for some $`\delta >0`$,
$$|P_t^Vf(x)|\frac{C}{t^\delta }f_2,t1,xM.$$
This proves (2.13).
Now we are ready to prove the decay estimates in $`L^1L^1`$ norm.
Case 1. Assume $`VL^1(M,\mu )`$.
Fix $`yM`$, and let $`u(x,t)=p_t^V(x,y)`$. Since $`u`$ satisfies
$$\mathrm{\Delta }uVu+u_t=0,$$
integrating on $`[1,t]\times M`$ yields
$$_Mu(x,t)𝑑\mu (x)=_Mu(x,1)𝑑\mu (x)+_1^t_MV(x)u(x,s)𝑑\mu (x)𝑑s.$$
By (2.1) and doubling (again, subexponential growth is enough, see the remark after Proposition 2.1), it follows that
(2.16)
$$_Mu(x,t)𝑑\mu (x)C+_1^t_MV(x)u(x,s)𝑑\mu (x)𝑑s.$$
From (2.16) and the strong positivity of $`V`$,
$`{\displaystyle _M}u(x,t)𝑑\mu (x)`$ $`C+{\displaystyle _1^t}{\displaystyle _M}V(x)u(x,s)𝑑\mu (x)𝑑s`$
$`C+\left({\displaystyle _1^t}{\displaystyle _M}V(x)𝑑\mu (x)𝑑s\right)^{1/2}\left({\displaystyle _1^t}{\displaystyle _M}V(x)u^2(x,s)𝑑\mu (x)𝑑s\right)^{1/2}`$
$`C+\sqrt{A}\sqrt{t}V_1^{1/2}\left({\displaystyle _1^t}{\displaystyle _M}|u|^2(x,s)𝑑\mu (x)𝑑s\right)^{1/2}.`$
Multiplying by $`u`$ the equation
$$\mathrm{\Delta }uVu+u_t=0$$
and integrating on $`[1,t]\times M`$, we obtain
$$_1^t_M|u|^2(x,s)𝑑\mu (x)𝑑s_1^t_MVu^2𝑑\mu (x)𝑑s+\frac{1}{2}_Mu^2(x,t)𝑑\mu (x)=\frac{1}{2}_Mu^2(x,1)𝑑\mu (x).$$
Since $`\mathrm{\Delta }V`$ is strongly positive, we obtain
(2.17)
$$_1^t_M|u|^2(x,s)𝑑\mu (x)𝑑s\frac{1}{2(1A)}_Mu^2(x,1)𝑑\mu (x)\frac{C}{2(1A)}.$$
Finally
$$P_t^V_{11}=_Mu(x,t)𝑑\mu (x)C\sqrt{\frac{A}{1A}}V_1^{1/2}\sqrt{t},$$
which is the claim.
Case 2. Now we assume that $`VL^p(M,\mu )`$, $`p(1,2)`$.
The decay estimate in this range of $`p`$ seems to be new even in the Euclidean case.
Using Hölder’s inequality repeatedly, one has
$`{\displaystyle _M}u(x,t)𝑑\mu (x)C+{\displaystyle _1^t}{\displaystyle _M}V(x)u(x,s)𝑑\mu (x)𝑑s`$
$`=`$ $`C+{\displaystyle _1^t}{\displaystyle _M}V(x)^{p/2}V(x)^{1(p/2)}u(x,s)𝑑\mu (x)𝑑s`$
$``$ $`C+\left({\displaystyle _1^t}{\displaystyle _M}V(x)^p𝑑\mu (x)𝑑s\right)^{1/2}\left({\displaystyle _1^t}{\displaystyle _M}V(x)^{2p}u^2(x,s)𝑑\mu (x)𝑑s\right)^{1/2}`$
$`=`$ $`C+t^{1/2}V_p^{p/2}\left({\displaystyle _1^t}{\displaystyle _M}V(x)^{2p}u(x,s)^{2(2p)}u(x,s)^{2p2}𝑑\mu (x)𝑑s\right)^{1/2}`$
$``$ $`C+t^{1/2}V_p^{p/2}\left({\displaystyle _1^t}{\displaystyle _M}[V(x)^{2p}u(x,s)^{2(2p)}]^{1/(2p)}𝑑\mu (y)𝑑s\right)^{(2p)/2}`$
$`\times \left({\displaystyle _1^t}{\displaystyle _M}u(x,s)^{(2p2)/(p1)}𝑑\mu (x)𝑑s\right)^{(p1)/2}`$
$`=`$ $`C+t^{1/2}V_p^{p/2}\left({\displaystyle _1^t}{\displaystyle _M}V(x)u^2(x,s)𝑑\mu (y)𝑑s\right)^{(2p)/2}\left({\displaystyle _1^t}{\displaystyle _M}u^2(x,s)𝑑\mu (x)𝑑s\right)^{(p1)/2}`$
$``$ $`C+A^{(2p)/2}t^{1/2}V_p^{p/2}\left({\displaystyle _1^t}{\displaystyle _M}|u(x,s)|^2𝑑\mu (y)𝑑s\right)^{(2p)/2}\left({\displaystyle _1^t}{\displaystyle _M}u^2(x,s)𝑑\mu (x)𝑑s\right)^{(p1)/2}.`$
Using (2.17), we obtain
$$_Mu(x,t)𝑑\mu (x)C+C^{}t^{1/2}\left(_1^t_Mu^2(x,s)𝑑\mu (x)𝑑s\right)^{(p1)/2}$$
From (2.13), we know that
$$_1^t_Mu^2(x,s)𝑑\mu (x)𝑑sCt^{1\delta },$$
and using (2.17), we deduce
$$_Mu(x,t)𝑑\mu (x)C+C^{}t^{1/2+(12\delta )(p1)/2}=Ct^{(p\theta )/2},$$
with $`\theta =2\delta (p1)`$.
Case 3. The only remaining case is when $`VL^p(M,\mu )`$, $`p2`$.
This case may be skipped in a first reading; indeed, if one is prepared to replace $`(p1)/2`$ by $`p/2`$ in the estimate, the simpler proof in Section 5 will do.
The reader may guess that the claimed estimate follows from the idea in , p.99, where a similar bound for the Schrödinger heat kernel was proven in the Euclidean case. However in , the authors use the boundedness of $`\mathrm{\Delta }^{1/2}`$ from some $`L^p`$ space to another. But it is known that this property is false for most open manifolds. Therefore, we have to work considerably harder. We will show that the inverse square root of the Laplacian on $`M\times \mathrm{I}\mathrm{R}^3`$ is bounded from the space $`L^{p_1}L^{p_2}`$ to $`L^2`$ for some $`p_1,p_2`$. Then we will use the idea in to get an $`L^1`$ to $`L^1`$ bound for a version of our Schrödinger semigroup acting on $`NM\times \mathrm{I}\mathrm{R}^3`$. After integrating over $`\mathrm{I}\mathrm{R}^3`$, we will reach the desired $`L^1`$ to $`L^1`$ bound for the Schrödinger semigroup acting on $`M`$.
Points in $`N`$ will be denoted by $`\stackrel{~}{x}=(x,x^{})`$ and $`\stackrel{~}{y}=(y,y^{})`$, …, where $`x,yM`$ and $`x^{},y^{}\mathrm{I}\mathrm{R}^3`$. The distance function on $`N`$ is denoted by $`d(\stackrel{~}{x},\stackrel{~}{y})=d(x,y)+|x^{}y^{}|`$, and the Riemannian measure on $`N`$ by $`d\stackrel{~}{\mu }`$. The Laplace-Beltrami operator on $`N`$ will be denoted by $`\stackrel{~}{\mathrm{\Delta }}`$. Denote by $`\stackrel{~}{P}_t`$ the heat semigroup on $`N`$ and let $`\stackrel{~}{P}_t^V=e^{t(\stackrel{~}{\mathrm{\Delta }}V)}`$, where $`V`$ is the function on $`M`$ defined in (1.1). The kernel for $`\stackrel{~}{\mathrm{\Delta }}^{1/2}`$ is
(2.18)
$$\mathrm{\Delta }^{1/2}(\stackrel{~}{x},\stackrel{~}{y})=_0^{\mathrm{}}t^{1/2}\stackrel{~}{p}_t(\stackrel{~}{x},\stackrel{~}{y})𝑑t.$$
It is clear that $`N`$ satisfies Assumptions $`(A)`$, $`(B)`$, $`(C)`$ and $`(D)`$ just like $`M`$ does. Moreover, if $`\mathrm{\Delta }V`$ is strongly positive on $`M`$, then $`\stackrel{~}{\mathrm{\Delta }}V`$ is strongly positive on $`N`$. Now suppose we can prove that
(2.19)
$$\stackrel{~}{P}_t^V_{1,1}Ct^{(p1)/2}.$$
Then, since
(2.20)
$$\stackrel{~}{p}_t^V(\stackrel{~}{x},\stackrel{~}{y})=p_t^V(x,y)\frac{c}{t^{3/2}}e^{|x^{}y^{}|^2/4t},$$
we will deduce that
$$P_t^V_{1,1}=P_t^V_\mathrm{},\mathrm{}=P_t^V1_{\mathrm{}}=\stackrel{~}{P}_t^V1_{\mathrm{}}=\stackrel{~}{P}_t^V_\mathrm{},\mathrm{}=\stackrel{~}{P}_t^V_{1,1}Ct^{(p1)/2},$$
thus finishing case 3. The rest of the section is devoted to proving (2.19). It will be divided into several steps.
Step 1. We show that there exists $`C>0`$ such that
(2.21)
$$\stackrel{~}{\mathrm{\Delta }}^{1/2}(\stackrel{~}{x},\stackrel{~}{y})C\frac{d(\stackrel{~}{x},\stackrel{~}{y})}{|B(\stackrel{~}{x},d(\stackrel{~}{x},\stackrel{~}{y}))|}.$$
From the upper bound for $`\stackrel{~}{P}_t`$, which is a consequence of the Gaussian upper bound for $`P_t`$, we have
$`I`$ $`\stackrel{~}{\mathrm{\Delta }}^{1/2}(\stackrel{~}{x},\stackrel{~}{y})`$
$`C{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{1}{\sqrt{t}}}{\displaystyle \frac{e^{cd^2(\stackrel{~}{x},\stackrel{~}{y})/t}}{|B(\stackrel{~}{x},\sqrt{t})|}}𝑑t`$
$`=C{\displaystyle _0^{d^2(\stackrel{~}{x},\stackrel{~}{y})}}{\displaystyle \frac{1}{\sqrt{t}}}{\displaystyle \frac{e^{cd^2(\stackrel{~}{x},\stackrel{~}{y})/t}}{|B(\stackrel{~}{x},\sqrt{t})|}}𝑑t+C{\displaystyle _{d^2(\stackrel{~}{x},\stackrel{~}{y})}^{\mathrm{}}}{\displaystyle \frac{1}{\sqrt{t}}}{\displaystyle \frac{e^{cd^2(\stackrel{~}{x},\stackrel{~}{y})/t}}{|B(\stackrel{~}{x},\sqrt{t})|}}𝑑t`$
$`CI_1+CI_2.`$
By the doubling condition, for $`td^2(\stackrel{~}{x},\stackrel{~}{y})`$,
$$|B(\stackrel{~}{x},\sqrt{t})|c|B(\stackrel{~}{x},d(\stackrel{~}{x},\stackrel{~}{y}))|\left(\frac{\sqrt{t}}{d(\stackrel{~}{x},\stackrel{~}{y})}\right)^\nu .$$
Hence
(2.22)
$$I_1C\frac{d(\stackrel{~}{x},\stackrel{~}{y})}{|B(\stackrel{~}{x},d(\stackrel{~}{x},\stackrel{~}{y}))|}$$
Next we estimate $`I_2`$.
By (2.20), we have
$$I_2C_{d^2(\stackrel{~}{x},\stackrel{~}{y})}^{\mathrm{}}\frac{1}{\sqrt{t}}\frac{1}{|B(x,\sqrt{t})|t^{3/2}}𝑑t\frac{C}{|B(x,d(\stackrel{~}{x},\stackrel{~}{y}))|}_{d^2(\stackrel{~}{x},\stackrel{~}{y})}^{\mathrm{}}\frac{dt}{t^2}=\frac{C}{|B(x,d(\stackrel{~}{x},\stackrel{~}{y}))|d^2(\stackrel{~}{x},\stackrel{~}{y})}.$$
Note that
$$|B(\stackrel{~}{x},d(\stackrel{~}{x},\stackrel{~}{y}))|=|B(x,d(\stackrel{~}{x},\stackrel{~}{y}))|d^3(\stackrel{~}{x},\stackrel{~}{y}).$$
The above implies
(2.23)
$$I_2\frac{Cd(\stackrel{~}{x},\stackrel{~}{y})}{|B(\stackrel{~}{x},d(\stackrel{~}{x},\stackrel{~}{y}))|}.$$
The combination of (2.23) and (2.22) implies (2.21).
Step 2. We prove that there exist $`p_1,p_2>1`$ and $`C>0`$ such that
(2.24)
$$\stackrel{~}{\mathrm{\Delta }}^{1/2}f_2C(f_{p_1}+f_{p_2}),$$
for all $`fL^{p_1}L^{p_2}.`$
From (2.21),
$`|\stackrel{~}{\mathrm{\Delta }}^{1/2}f(\stackrel{~}{x})|`$ $`C{\displaystyle _{\stackrel{~}{M}}}{\displaystyle \frac{d(\stackrel{~}{x},\stackrel{~}{y})}{|B(\stackrel{~}{x},d(\stackrel{~}{x},\stackrel{~}{y}))|}}|f(\stackrel{~}{y})|𝑑\stackrel{~}{\mu }(\stackrel{~}{y})`$
$`C{\displaystyle _{d(\stackrel{~}{x},\stackrel{~}{y})1}}\mathrm{}+C{\displaystyle _{d(\stackrel{~}{x},\stackrel{~}{y})1}}CJ_1(\stackrel{~}{x})+CJ_2(\stackrel{~}{x}).`$
By Young’s inequality
$$J_1_2Cf_{p_1}\underset{\stackrel{~}{x}}{sup}\left(_{d(\stackrel{~}{x},\stackrel{~}{y})1}\left[\frac{d(\stackrel{~}{x},\stackrel{~}{y})}{|B(\stackrel{~}{x},d(\stackrel{~}{x},\stackrel{~}{y}))|}\right]^{r_1}𝑑\stackrel{~}{\mu }(\stackrel{~}{y})\right)^{1/r_1}$$
where $`(1/p_1)+(1/r_1)=1+(1/2)`$. Hence
$`J_1_2`$ $`f_{p_1}\underset{\stackrel{~}{x}}{sup}{\displaystyle \underset{k=0}{\overset{\mathrm{}}{}}}\left({\displaystyle _{2^{(k+1)}d(\stackrel{~}{x},\stackrel{~}{y})2^k}}\left[{\displaystyle \frac{d(\stackrel{~}{x},\stackrel{~}{y})}{|B(\stackrel{~}{x},d(\stackrel{~}{x},\stackrel{~}{y}))|}}\right]^{r_1}𝑑\stackrel{~}{\mu }(\stackrel{~}{y})\right)^{1/r_1}`$
$`Cf_{p_1}\underset{\stackrel{~}{x}}{sup}{\displaystyle \underset{k=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{2^{kr_1}}{|B(\stackrel{~}{x},2^{(k+1)})|^{r_11}}}.`$
Using the doubling property and the fact that, thanks to Assumption $`(D)`$, $`|B(\stackrel{~}{x},1)|c>0`$, we have
$$|B(\stackrel{~}{x},2^{(k+1)})|C2^{k\nu }|B(\stackrel{~}{x},1)|C2^{k\nu }.$$
Hence
$$J_1_2f_{p_1}\underset{k=0}{\overset{\mathrm{}}{}}2^{kr_1}2^{k\mu (r_11)}.$$
Choosing $`r_1`$ sufficiently close to $`1`$, the above series is convergent. Therefore
(2.25)
$$J_1_2Cf_{p_1}.$$
Similarly, by Young’s inequality again,
$$J_2_2Cf_{p_2}\underset{\stackrel{~}{x}}{sup}\left(_{d(\stackrel{~}{x},\stackrel{~}{y})1}\left[\frac{d(\stackrel{~}{x},\stackrel{~}{y})}{|B(\stackrel{~}{x},d(\stackrel{~}{x},\stackrel{~}{y}))|}\right]^{r_2}𝑑\stackrel{~}{\mu }(\stackrel{~}{y})\right)^{1/r_2}$$
where $`(1/p_2)+(1/r_2)=1+(1/2)`$. Hence
$`J_2_2`$ $`Cf_{p_2}\underset{\stackrel{~}{x}}{sup}{\displaystyle \underset{k=0}{\overset{\mathrm{}}{}}}\left({\displaystyle _{2^kd(\stackrel{~}{x},\stackrel{~}{y})2^{k+1}}}\left[{\displaystyle \frac{d(\stackrel{~}{x},\stackrel{~}{y})}{|B(\stackrel{~}{x},d(\stackrel{~}{x},\stackrel{~}{y}))|}}\right]^{r_2}𝑑\stackrel{~}{\mu }(\stackrel{~}{y})\right)^{1/r_2}`$
$`Cf_{p_2}\underset{\stackrel{~}{x}}{sup}{\displaystyle \underset{k=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{2^{kr_2}}{|B(\stackrel{~}{x},2^k)|^{r_21}}}.`$
Using the fact that $`|B(\stackrel{~}{x},2^k)|c2^{3k}`$ which follows from the definition of $`N`$ and Assumption $`(D)`$, we have
$$J_2_2f_{p_2}\underset{k=0}{\overset{\mathrm{}}{}}2^{kr_23k(r_21)}.$$
Choosing $`r_2`$ sufficiently large, the above series is convergent. Therefore
(2.26)
$$J_2_2Cf_{p_2}.$$
Inequality (2.24) immediately follows from (2.26) and (2.25).
Step 3. As in , by Duhamel’s formula, one has
$`\stackrel{~}{P}_{t+1}^V_\mathrm{},\mathrm{}`$ $`=`$ $`\stackrel{~}{P}_{t+1}^V1_{\mathrm{}}`$
$``$ $`\stackrel{~}{P}_1^V1_{\mathrm{}}+{\displaystyle _1^t}\stackrel{~}{P}_{s+1}^VV_{\mathrm{}}𝑑s`$
$``$ $`C+{\displaystyle _1^t}\stackrel{~}{P}_{s+1}^VV_{\mathrm{}}𝑑s,`$
hence by interpolation,
(2.27)
$$\stackrel{~}{P}_{t+1}^V_\mathrm{},\mathrm{}C+_1^t\stackrel{~}{P}_{s+1}^V_{2,\mathrm{}}^{2/p}\stackrel{~}{P}_{s+1}^V_\mathrm{},\mathrm{}^{1(2/p)}V_p𝑑s.$$
Let us now estimate $`\stackrel{~}{P}_{s+1}^V_{2,\mathrm{}}^{2/p}`$. Using the strong positivity of $`\stackrel{~}{\mathrm{\Delta }}V`$ on $`N`$, one has
$$(\stackrel{~}{\mathrm{\Delta }}V)a^2\stackrel{~}{\mathrm{\Delta }}$$
for some $`a>0`$. Therefore, for all $`fC_0^{\mathrm{}}(N)`$,
$$(\stackrel{~}{\mathrm{\Delta }}V)^{1/2}f_2a^1\stackrel{~}{\mathrm{\Delta }}^{1/2}f_2Cf_X.$$
Here $`X=L^{p_1}L^{p_2}`$ where $`p_1,p_2`$ are given in (2.24) and
(2.28)
$$f_X=f_{p_1}+f_{p_2}.$$
Hence
$`e^{t(\stackrel{~}{\mathrm{\Delta }}V)}f_2`$ $`=t^{1/2}e^{(\stackrel{~}{\mathrm{\Delta }}V)t}((\stackrel{~}{\mathrm{\Delta }}V)t)^{1/2}(\stackrel{~}{\mathrm{\Delta }}V)^{1/2}f_2`$
$`Ct^{1/2}(\stackrel{~}{\mathrm{\Delta }}V)^{1/2}f_2\stackrel{~}{\mathrm{\Delta }}^{1/2}f_2`$
$`Ct^{1/2}f_X.`$
This shows
$$\stackrel{~}{P}_t_{2,X^{}}=\stackrel{~}{P}_t_{X,2}Ct^{1/2},$$
where $`X^{}`$ is the dual of $`X`$. Then write
(2.29)
$$\stackrel{~}{P}_{s+1}^V_{2,\mathrm{}}\stackrel{~}{P}_s^V_{2,X^{}}\stackrel{~}{P}_1^V_{X^{},\mathrm{}}.$$
It follows easily from the bound
$$\stackrel{~}{p}_1^V(\stackrel{~}{x},\stackrel{~}{y})\frac{C}{|B(\stackrel{~}{y},1)|}e^{cd^2(\stackrel{~}{x},\stackrel{~}{y})},$$
and Assumptions $`(A)`$ and $`(D)`$ that $`\stackrel{~}{P}_1^V`$ is bounded from $`L^1`$ to any $`L^p`$, $`1p+\mathrm{}`$. Therefore
(2.30)
$$\stackrel{~}{P}_1^V_{X^{},\mathrm{}}=\stackrel{~}{P}_1^V_{1,X}=\stackrel{~}{P}_1^V_{1,p_1}+\stackrel{~}{P}_1^V_{1,p_2}<+\mathrm{}.$$
This together with (2.29) implies
$$\stackrel{~}{P}_{s+1}^V_{2,\mathrm{}}Cs^{1/2}.$$
Using this and (2.27), we obtain
$$\stackrel{~}{P}_{t+1}^V_\mathrm{},\mathrm{}C+CV_p_1^ts^{1/p}\stackrel{~}{P}_{s+1}^V_\mathrm{},\mathrm{}^{1(2/p)}𝑑s.$$
By Gronwall’s lemma,
$$\stackrel{~}{P}_{t+1}^V_\mathrm{},\mathrm{}Ct^{p/21/2}.$$
Since
$$\stackrel{~}{P}_{t+1}^V_{1,1}=\stackrel{~}{P}_{t+1}^V_\mathrm{},\mathrm{},$$
the proof of Proposition 2.3 is complete. ∎
### 2.3. Proof of Theorem 2.2
Under the assumptions of Theorem 2.2, the combination of Propositions 2.2 and 2.3 yields
(2.31)
$$p_t^V(x,y)\frac{Ct\mathrm{ln}^{\nu /2}(e+t)}{|B(x,\sqrt{t})|}\mathrm{exp}(cd^2(x,y)/t),x,yM,t1$$
if $`p=1`$,
(2.32)
$$p_t^V(x,y)\frac{Ct^{p\theta }}{|B(x,\sqrt{t})|}\mathrm{exp}(cd^2(x,y)/t),x,yM,t1$$
for some $`\theta =\theta (p)>0`$, if $`1p<2`$, and
(2.33)
$$p_t^V(x,y)\frac{Ct^{p1}\mathrm{ln}^{\nu /2}(e+t)}{|B(x,\sqrt{t})|}\mathrm{exp}(cd^2(x,y)/t),x,yM,t1$$
if $`p2`$.
The above bound does not explicitely reflect the contribution of the constant $`A`$ (which measures the size of $`V`$). To remedy this, we use again (2.14), which yields
(2.34)
$$p_t^V(x,y)\left(p_t^V^{}(x,y)\right)^{1/q}\left(p_t(x,y)\right)^{(q1)/q},x,yM,t>0,$$
where $`V^{}=qV`$ and $`q>1`$.
Assume $`A>0`$, otherwise there is nothing to prove. If $`q<1/A`$, then $`V^{}`$ is obviously strongly positive with constant $`qA<1`$.
If $`1<p<2`$, (2.32) yields
(2.35)
$$p_t^V^{}(x,y)\frac{Ct^{p\theta }}{|B(x,\sqrt{t})|}\mathrm{exp}(cd^2(x,y)/t).$$
Applying (2.35) and using the Gaussian upper bound on $`p_t`$, we obtain
$$p_t^V(x,y)\frac{Ct^{(p\theta )/q}}{|B(x,\sqrt{t})|}\mathrm{exp}(cd^2(x,y)/t).$$
Taking $`q`$ sufficiently close to $`1/A`$, one can make $`(p\theta )/q<pA`$, which finishes the proof of Theorem 2.2 in the case $`1<p<2`$. The proofs when $`p=1`$ and $`p2`$ are identical using (2.31) and (2.33) instead of (2.32). ∎
## 3. Gaussian bound on the heat kernel on forms and the Riesz transform
Our next theorem provides an all time Gaussian upper bound for the heat kernel on $`1`$-forms on complete non-compact Riemannian manifolds satisfying Assumptions $`(A)`$, $`(B)`$, $`(C)`$, together with a certain condition of smallness of the Ricci curvature. Using this theorem and an argument in , one deduces a proper bound for the gradient of the heat kernel on functions. By the main result in , one obtains the $`L^p`$ boundedness of the Riesz transform on these manifolds for all $`1<p<+\mathrm{}`$. Let us point out that in , Theorem 9.1, another sufficient condition in terms of Ricci curvature is given for $`L^p`$ boundedness of the Riesz transform. However this condition seems to exclude Ricci curvature bounded from below together with non-compactness.
###### Theorem 3.1.
Let $`M`$ be a complete non-compact Riemannian manifold satisfying Assumptions (A), (B) and (C). Then there exists $`\delta >0`$ depending only on the constants in (A) and (B) such that for any $`c>1/4`$, there exists $`C>0`$ such that
$$|\stackrel{}{p}_t(x,y)|\frac{C}{|B(x,\sqrt{t})|}\mathrm{exp}\left(cd(x,y)^2/t\right),$$
for all $`x,yM`$ and $`t>0`$, provided that
(3.1)
$$K(V)\underset{xM}{sup}_0^{\mathrm{}}_M\frac{1}{|B(x,\sqrt{s})|}e^{d^2(x,y)/s}V(y)𝑑\mu (y)𝑑s<\delta .$$
As we already said, the following statement is a consequence from Theorem 3.1 and either Theorem 5.5 in , or Theorem 1.4 in together with , pp.1740-1741. This extends the class of manifolds for which one can answer a question asked by Strichartz in .
###### Corollary 1.
Let $`M`$ be a complete non-compact Riemannian manifold satisfying Assumptions (A), (B), (C), and condition (3.1) for $`\delta >0`$ small enough. Then, for all $`p(0,+\mathrm{})`$, there exist $`C_p,c_p>0`$ such that
$$c_p|f|_p\mathrm{\Delta }^{1/2}f_pC_p|f|_p,fC_0^{\mathrm{}}(M).$$
Proof of Theorem 3.1. As we have seen in the introduction,
$$|\stackrel{}{p}_t(x,y)|p_t^V(x,y)$$
under Assumptions $`(A)`$ to $`(C)`$. Now let us recall Theorem A, part (b) in . It implies that
$$p_t^V(x,y)\frac{C}{|B(x,\sqrt{t})|}e^{cd^2(x,y)/t},$$
for all $`x,yM`$ and $`t>0`$, provided that, for certain $`c_0,\epsilon _0>0`$, there holds
$$N(V)<\epsilon _0.$$
Here
$`N(V)`$ $``$ $`\underset{xM,t>0}{sup}{\displaystyle _0^t}{\displaystyle _M}{\displaystyle \frac{e^{c_0d^2(x,y)/(ts)}}{|B(x,\sqrt{ts})|}}V(y)𝑑\mu (y)𝑑s`$
$`+\underset{yM,s>0}{sup}{\displaystyle _s^{\mathrm{}}}{\displaystyle _M}{\displaystyle \frac{e^{c_0d^2(x,y)/(ts)}}{|B(x,\sqrt{ts})|}}V(x)𝑑\mu (x)𝑑t.`$
We should mention that this theorem was stated for a doubling metric in the Euclidean space and for time dependent functions $`V`$, under the extra assumption $`(D)`$. However the proof was a general one applicable verbatim to any manifold under Assumptions $`(A)`$ and $`(B)`$ only.
Changing variables and using doubling, one sees that
$$N(V)C_0K(V),$$
where $`C_0`$ only depends on the doubling constants. The conclusion follows. ∎
Remarks:
\- Suppose in addition that $`|B(x,r)|r^n`$ with $`n>2`$, uniformly in $`xM`$, for large $`r`$, then it is any easy exercise to check that the theorem holds if
$$V(x)\frac{a}{1+d(x,x_0)^{2+b}}$$
with any $`b>0`$ and $`a`$ sufficiently small, for some fixed $`x_0M`$.
\- Let $`(M,g_0)`$ be a nonparabolic manifold with nonnegative Ricci curvature and volume growth property as in the last remark. Let $`h`$ be another metric and $`\eta `$ be a smooth cut-off function on $`M`$. Then the manifold $`(M,g)`$ with $`g=g_0+\lambda \eta h`$ is covered by the theorem when $`\lambda 0`$ is sufficiently small. This is so because the constants in $`(A)`$ and $`(B)`$ are uniformly bounded when $`0\lambda 1`$ while $`V=V(x)`$ (for the metric $`g`$), being a compactly supported function is arbitrarily small when $`\lambda 0`$.
## 4. The case of non-negative Ricci curvature outside a compact set
In the next theorem, we establish an upper bound for the heat kernel on $`1`$-forms assuming Ricci curvature is nonnegative outside a compact set. The upshot of the theorem is that no other restriction on the Ricci curvature is needed. This upper bound gives a good control of the heat kernel even in the presence of harmonic forms. In general one can not expect the heat kernel on forms to decay to zero, due to the possible presence of $`L^2`$ harmonic forms. Here we are able to show that the heat kernel has certain spatial decay anyway. Using the spectral decomposition of heat kernels, one can see that the upper bound in the theorem below is quite sharp near the diagonal at least. As far as we know, this bound is new even for Schrödinger heat kernels in the Euclidean case.
Moreover the assumption that the Ricci curvature is $`0`$ outside of a compact set can be improved to assuming that the negative part of the Ricci curvature decays sufficiently fast near infinity. But we will not seek the full generality this time.
###### Theorem 4.1.
Let $`M`$ be a manifold satisfying Assumptions (A), (B), (C) and (D). In addition, we assume that the Ricci curvature of $`M`$ is nonnegative outside a compact set and the manifold is nonparabolic. Then, for a fixed $`0M`$, there exist $`C,c>0`$ such that
$$|\stackrel{}{p}_t(x,y)|C\mathrm{min}\{\mathrm{\Gamma }(x,0)\mathrm{\Gamma }(y,0),1\}e^{cd^2(x,y)/t}+\frac{C}{|B(x,\sqrt{t})|}e^{cd^2(x,y)/t}$$
for all $`x,yM`$, $`t>1`$. Here $`\mathrm{\Gamma }`$ is the Green’s function of the Laplacian $`\mathrm{\Delta }`$ on $`M`$.
###### Proof.
We divide the proof into two steps.
Step 1.
As in Section 2.1, we need a preliminary estimate.
(4.1)
$$|\stackrel{}{p}_t(x,y)|Ce^{cd^2(x,y)/t},t1,$$
for some $`C,c>0`$. Note that this estimate does not follow from Proposition 2.1 and (1.2), since we do not assume strong positivity of $`\mathrm{\Delta }V`$ any more.
However, the method is very similar to the one in Proposition 2.1. Let $`u_0`$ be a smooth compactly supported $`1`$-form. Write
$$u(x,t)=\stackrel{}{P}_tu_0(x).$$
Direct computations show, for any fixed $`yM`$ and $`D>0`$,
$`{\displaystyle \frac{d}{dt}}`$ $`{\displaystyle _M}|u|^2e^{\frac{d^2(x,y)}{Dt}}𝑑\mu (x)`$
$`=2{\displaystyle _M}e^{\frac{d^2(x,y)}{Dt}}u\stackrel{}{\mathrm{\Delta }}u𝑑\mu (x){\displaystyle _M}|u|^2e^{\frac{d^2(x,y)}{Dt}}{\displaystyle \frac{d^2(x,y)}{Dt^2}}𝑑\mu (x).`$
Noticing that $`\stackrel{}{\mathrm{\Delta }}=d^{}d+dd^{}`$, the above implies, after integration by parts,
$`{\displaystyle \frac{d}{dt}}`$ $`{\displaystyle _M}|u|^2e^{\frac{d^2(x,y)}{Dt}}𝑑\mu (x)`$
$`=2{\displaystyle _M}e^{\frac{d^2(x,y)}{Dt}}\left(d\left({\displaystyle \frac{d^2(x,y)}{Dt}}\right)u\right)𝑑u𝑑\mu (x)2{\displaystyle _M}e^{\frac{d^2(x,y)}{Dt}}𝑑u𝑑u𝑑\mu (x)`$
$`2{\displaystyle _M}e^{\frac{d^2(x,y)}{Dt}}|d^{}u|^2𝑑\mu (x){\displaystyle _M}|u|^2e^{\frac{d^2(x,y)}{Dt}}{\displaystyle \frac{d^2(x,y)}{Dt^2}}𝑑\mu (x)`$
$`C{\displaystyle _M}e^{\frac{d^2(x,y)}{Dt}}{\displaystyle \frac{d(x,y)}{Dt}}|u||du|𝑑\mu (x)2{\displaystyle _M}e^{\frac{d^2(x,y)}{Dt}}|du|^2𝑑\mu (x)`$
$`{\displaystyle _M}|u|^2e^{\frac{d^2(x,y)}{Dt}}{\displaystyle \frac{d^2(x,y)}{Dt^2}}𝑑\mu (x).`$
Using the inequality $`\frac{d(x,y)}{Dt}|u||du|\frac{1}{\epsilon }\frac{d^2(x,y)}{D^2t^2}|u|^2+\epsilon |du|^2`$, we find that
$$\frac{d}{dt}_M|u|^2e^{\frac{d^2(x,y)}{Dt}}𝑑\mu (x)0$$
when $`D`$ is sufficiently large.
Letting $`u_0`$ converge to the Dirac delta function centered at $`y`$, we obtain
$$_M|\stackrel{}{p}_t(x,y)|^2e^{\frac{d^2(x,y)}{Dt}}𝑑\mu (x)_M|\stackrel{}{p}_1(x,y)|^2e^{\frac{d^2(x,y)}{D}}𝑑\mu (x)$$
when $`t1`$. By the semigroup domination property (1.2), Assumptions $`(C)`$ and $`(B)`$,
$$|\stackrel{}{p}_1(x,y)|p_1^V(x,y)Cp_1(x,y)\frac{C}{|B(x,1)|}e^{cd^2(x,y)}.$$
Integrating and using Assumptions $`(D)`$ and $`(A)`$, we have, for a suitable $`D>0`$,
(4.2)
$$_M|\stackrel{}{p}_t(x,y)|^2e^{\frac{d^2(x,y)}{Dt}}𝑑\mu (x)C,t1.$$
Next, using the semigroup property
$`|\stackrel{}{p}_{2t}(x,y)|`$ $`=|{\displaystyle _M}\stackrel{}{p}_t(x,z)\stackrel{}{p}_t(z,y)𝑑\mu (z)|`$
$`=|{\displaystyle _M}e^{\frac{d^2(x,z)}{2Dt}}\stackrel{}{p}_t(x,z)e^{\frac{d^2(z,y)}{2Dt}}\stackrel{}{p}_t(z,y)e^{\frac{d^2(x,z)}{2Dt}\frac{d^2(z,y)}{2Dt}}𝑑\mu (z)|`$
$`e^{\frac{d^2(x,y)}{4Dt}}\left[{\displaystyle _M}e^{\frac{d^2(x,z)}{Dt}}|\stackrel{}{p}_t(x,z)|^2𝑑\mu (z)\right]^{1/2}\left[{\displaystyle _M}e^{\frac{d^2(y,z)}{Dt}}|\stackrel{}{p}_t(y,z)|^2𝑑\mu (z)\right]^{1/2}.`$
Together with (4.2), this implies (4.1).
Step 2.
We assume that the Ricci curvature is nonnegative outside of a ball $`B(0,A)`$ for a fixed $`A>0`$. Write, for any given $`x_0B(0,A)`$,
$$u(y,t)=|\stackrel{}{p}_t(x_0,y)|.$$
Then it is an immediate consequence of Bochner’s formula (see for instance , Lemma 4.1) that $`u`$ is a subsolution of the scalar heat equation in $`B^c(0,A)\times (0,\mathrm{})`$. i.e.
$$\mathrm{\Delta }u(y,t)+u_t(y,t)0.$$
Since, according to (4.1), $`u`$ is bounded from above by a constant for $`t1`$, and $`\mathrm{\Gamma }(y,0)`$ is bounded from below by a positive constant on any compact set, there exists $`C>0`$ such that
$$u(y,t)C\mathrm{\Gamma }(y,0),yB(0,A),t1.$$
Moreover, using again (4.1),
$$u(y,1)Ce^{cd^2(y,x_0)}C^{}\mathrm{\Gamma }(y,0)$$
for $`yB^c(0,A)`$. The last inequality is due to the Cheng-Yau gradient estimate from , which implies
$$\mathrm{\Gamma }(y,0)ce^{Cd(y,0)}$$
for some positive constants $`C,c>0`$ and $`yB^c(0,A)`$. Now by the maximum principle, using (4.1) again, we deduce
$$u(y,t)C\mathrm{\Gamma }(y,0)$$
for all $`t1`$ and $`yB^c(0,A)`$. Here we just used the simple observation that $`\mathrm{\Gamma }(y,0)`$ is a solution of the scalar heat equation, whereas $`u`$ is a subsolution as already observed.
We have proved that
$$|\stackrel{}{p}_t(x_0,y)|C\mathrm{\Gamma }(y,0)$$
for all $`t1`$, $`yB^c(0,A)`$ and $`x_0B(0,A)`$. Let us now to explain how to keep the same estimate while moving away $`x_0`$.
For a fixed $`yB^c(0,A)`$, define the function
$$w(x,t)=|\stackrel{}{p}_t(x,y)|.$$
Then $`w`$ is a subsolution of the scalar heat equation on $`B^c(0,A)\times (0,+\mathrm{})`$. For $`xB(0,A)`$, by the above estimate on $`u`$ we have
$$w(x,t)=|\stackrel{}{p}_t(x,y)|C\mathrm{\Gamma }(y,0).$$
Since $`\mathrm{\Gamma }(x,0)`$ is bounded away from $`0`$ for $`xB(0,A)`$, it holds
$$w(x,t)C^{}\mathrm{\Gamma }(y,0)\mathrm{\Gamma }(x,0)$$
for some $`C^{}>0`$. It is clear that the function
$$h(x,t)=_Mp_{t1}(x,z)w(z,1)𝑑\mu (z)+C\mathrm{\Gamma }(y,0)\mathrm{\Gamma }(x,0)$$
is a solution of the scalar heat equation in $`B^c(0,A)\times [1,\mathrm{})`$. Moreover, on the parabolic boundary of the region, $`h`$ dominates $`w`$. By the maximum principle again
$$w(x,t)h(x,t)=_Mp_{t1}(x,z)w(z,1)𝑑\mu (z)+C\mathrm{\Gamma }(y,0)\mathrm{\Gamma }(x,0)$$
for $`xB^c(0,A)`$ and $`t1.`$ Next we estimate the above integral term in the following way, by using the Gaussian upper bound for $`p_t`$:
$`{\displaystyle _M}p_{t1}(x,z)w(z,1)𝑑\mu (z)`$ $``$ $`{\displaystyle p_{t1}(x,z)e^{cd(y,z)^2}𝑑\mu (z)}`$
$``$ $`{\displaystyle _{d(x,z)d(x,y)/2}}\mathrm{}𝑑\mu (z)+{\displaystyle _{d(y,z)d(x,y)/2}}\mathrm{}𝑑\mu (z){\displaystyle \frac{C}{|B(x,\sqrt{t})|}}.`$
Finally, incorporating (4.1),
$$|\stackrel{}{p}_t(x,y)|=w(x,t)C\mathrm{min}\{\mathrm{\Gamma }(x,0)\mathrm{\Gamma }(y,0),1\}+\frac{C}{|B(x,\sqrt{t})|}$$
for all $`x,y`$ in $`M`$ and $`t2`$. This is the desired on-diagonal estimate.
Now the theorem follows from the standard process of going from on to off-diagonal estimate. See e.g.∎
## 5. Bounds on manifolds without doubling condition
In this final section we turn to noncompact manifolds not necessarily satisfying the volume doubling condition. This class of manifolds offers a much richer variety than the doubling ones.
We show that under reasonable conditions the on-diagonal upper bound on the heat kernel on forms differs from that on functions only by a suitable power of time $`t`$.
Let us consider a $`n`$-dimensional manifold $`M`$ with Ricci curvature bounded from below, and whose small balls do not collapse, in other words Assumptions $`(C)`$ and $`(D)`$ are satisfied. Then
(5.1)
$$p_t(x,y)Ct^{n/2}\mathrm{exp}\left(cd^2(x,y)/t\right),\mathrm{\hspace{0.17em}0}<t1,x,yM,$$
for some $`C,c>0`$ and $`p_t^V(x,y)`$, $`|\stackrel{}{p}_t(x,y)|`$ satisfy similar estimates.
Let us assume that the heat kernel on functions has a uniform rate of decay $`\gamma `$, where $`\gamma `$ is increasing, $`C^1`$ and one-to-one on $`\mathrm{I}\mathrm{R}_+`$:
(5.2)
$$\underset{xM}{sup}p_t(x,x)\frac{1}{\gamma (t)},t>0.$$
According to , this implies the following so-called uniform Faber-Krahn inequality: for any set $`\mathrm{\Omega }M`$,
$`(UFK)`$
$$\lambda _1(\mathrm{\Omega })\mathrm{\Lambda }(|\mathrm{\Omega }|),$$
where $`\mathrm{\Lambda }`$ is given by
$$\mathrm{\Lambda }(t)=\frac{\gamma ^{}(t)}{\gamma (t)},$$
i.e.
$$t=_0^{\gamma (t)}\frac{d\eta }{\eta \mathrm{\Lambda }(\eta )}.$$
Conversely, if $`\gamma `$ satisfies a mild condition, the converse is true. For more on this as well as examples where one can compute $`\mathrm{\Lambda }`$, therefore $`\gamma `$, see for instance .
We can now state our result in this setting.
###### Theorem 5.1.
Suppose $`M`$ satisfies Assumptions (C) and (D). Assume that $`VL^p(M,\mu )`$ for some $`p[1,+\mathrm{})`$, and that $`\mathrm{\Delta }V`$ is strongly positive. Finally assume that the heat kernel on functions on $`M`$ satisfies the estimate (5.2). Then there exist positive constants $`c`$ and $`C`$ such that
$$|\stackrel{}{p}_t(x,x)|\frac{Ct^p}{\gamma (ct)},t1,xM.$$
###### Proof.
We divide the proof into two parts.
Step 1. $`L^1`$ to $`L^1`$ bound.
Under the assumptions of the theorem we will prove that
(5.3)
$$P_t^V_{1,1}Ct^{p/2},t1,\text{ if }p1.$$
Comparing with the proof of Proposition 2.3, we no longer have the doubling condition. However the growth rate of $`P_t^V_{1,1}`$ here is worse. On the other hand, the proof, a simple application of the idea of , is much shorter. In case $`p>2`$, the above estimate is essentially contained in , see also . We present the proof for completeness.
Case 1. Assume $`VL^1(M)`$.
Fixing $`y`$, we write
$$u(x,t)=p_t^V(x,y).$$
In this case the proof is almost identical to that of Case 1, Proposition 2.3. The only change is that we use the small time bound (5.1) on $`p_t^V`$ and the subexponential volume growth of $`M`$ (due to Assumption $`(C)`$) to conclude that
$$_Mu(x,1)𝑑\mu (x)=_Mp_1^V(x,y)𝑑\mu (x)C.$$
The rest of the proof is identical.
Case 2. Assume $`VL^p(M,\mu )`$ with $`1<p<2`$.
This is almost identical to that of Case 2, Proposition 2.3. Indeed, from that case, we have
$$_Mu(x,t)𝑑\mu (x)C+C^{}\sqrt{t}\left(_1^t_Mu^2(x,s)𝑑\mu (x)𝑑s\right)^{(p1)/2}$$
Since $`p_t^V`$ is contractive in $`L^2`$, this implies
$$_Mu(x,t)𝑑\mu (x)Ct^{p/2},t1.$$
Case 3. Assume $`VL^p(M,\mu )`$ with $`p2`$.
Let $`u=u(x,t)`$ be as above. Then, as before,
$`{\displaystyle _M}u(x,t)𝑑\mu (x)`$ $`{\displaystyle _M}u(x,1)𝑑\mu (x)+{\displaystyle _1^t}{\displaystyle _M}V(x)u(x,s)𝑑\mu (x)𝑑s`$
$`C+{\displaystyle _1^t}V_pu(,s)_{p/(p1)}𝑑s`$
$`C+CV_p{\displaystyle _1^t}P_s^V_{1,p/(p1)}𝑑s.`$
Applying Riesz-Thorin interpolation with the parameters
$$p_2=1,q_2=p/(p1);p_0=1,q_0=2;p_1=1,q_1=1;\theta =(2q_2)/q_2,$$
we have
$$\frac{1}{p_2}=\frac{1\theta }{p_0}+\frac{\theta }{p_1},\frac{1}{q_2}=\frac{1\theta }{q_0}+\frac{\theta }{q_1}$$
and
$$P_s^V_{1,p/(p1)}=P_s^V_{p_2,q_2}P_s^V_{p_0,q_0}^{1\theta }P_s^V_{p_1,q_1}^\theta =P_s^V_{1,2}^{2/p}P_s^V_{1,1}^{1(2/p)}.$$
Therefore
$$_M|u(x,t)|𝑑\mu (x)C+V_p_1^tP_s^V_{1,2}^{2/p}P_s^V_{1,1}^{1(2/p)}𝑑s$$
Notice that
$$P_s^V_{1,2}P_s^V_{2,\mathrm{}}P_{s1}^V_{2,2}P_1^V_{2,\mathrm{}}C.$$
We obtain
$$_M|u(x,t)|𝑑\mu (x)C+V_p_1^tP_s^V_{1,1}^{1(2/p)}𝑑s.$$
i.e.
$$P_t^V_{1,1}C+V_p_1^tP_s^V_{1,1}^{1(2/p)}𝑑s.$$
From here it is easy to see that
$$P_t^V_{1,1}Ct^{p/2}.$$
This completes Step 1.
Step 2. Write
$$I(t)=_Mu^2(x,t)𝑑\mu (x).$$
As in the proof of Proposition 2.2, the strong positivity of $`\mathrm{\Delta }V`$ yields
(5.4)
$$I(t)\frac{_{\{x|u(x,t)>s\}}|(u(x,t)s)|^2𝑑\mu (x)}{\lambda _1(\{x|u(x,t)>s\})}+2sF(t).$$
Here, according to (5.3), $`F(t)=t^{p/2}`$.
Using the fact that
$$|\{x|u(x,t)>s\}|s^1_Mu(x,t)𝑑\mu (x),$$
and $`(UFK)`$, we deduce
$$I(t)\frac{_{\{x|u(x,t)>s\}}|(u(x,t)s)|^2𝑑\mu (x)}{\mathrm{\Lambda }(s^1F(t))}+2sF(t).$$
Hence
(5.5)
$$_{\{x|u(x,t)>s\}}|u(x,t)|^2𝑑\mu (x)[I(t)2sF(t)]\mathrm{\Lambda }(s^1F(t)).$$
By the strong positivity of $`\mathrm{\Delta }V`$, we have as usual
$$I^{}(t)2(1A)_M|u(x,t)|^2𝑑\mu (x),$$
thus the combination of the above inequalities yields
(5.6)
$$I^{}(t)2(1A)[I(t)2sF(t)]\mathrm{\Lambda }(s^1F(t)).$$
Take $`sF(t)=I(t)/4`$, i.e.
$$s^1=4I^1(t)F(t).$$
Then (5.6) becomes
(5.7)
$$I^{}(t)(1A)I(t)\mathrm{\Lambda }(4F^2(t)I^1(t)).$$
Hence
$$_t^{2t}\frac{I^{}(l)}{I(l)\mathrm{\Lambda }(4F^2(l)I^1(l))}𝑑l(1A)t.$$
Notice that $`\mathrm{\Lambda }`$ is a decreasing and $`F`$ is an increasing function. Therefore, for $`lt`$,
$$\mathrm{\Lambda }(4F^2(l)I^1(l))\mathrm{\Lambda }(4F^2(t)I^1(l))$$
Consequently
$$_t^{2t}\frac{I^{}(l)}{I(l)\mathrm{\Lambda }(4F^2(t)I^1(l))}𝑑l(1A)t.$$
Take $`\eta =4F^2(t)I^1(l)`$. One gets
$$_{4F^2(t)I^1(t)}^{4F^2(t)I^1(2t)}\frac{d\eta }{\eta \mathrm{\Lambda }(\eta )}(1A)t.$$
Following the definition of $`\gamma `$, i.e. $`t=_0^\gamma \frac{d\eta }{\eta \mathrm{\Lambda }(\eta )}`$, we have
$$\frac{4F^2(t)}{I(2t)}\gamma ((1A)t),$$
i.e.
$$I(t)\frac{4F^2(t)}{\gamma (ct)}.$$
From here the desired bound for $`\stackrel{}{p_t}`$ follows immediately. ∎
Let us conclude by writing a semigroup version of the last part of the proof of Theorem 5.1, in the spirit of , where the case $`F`$ bounded is treated. We leave the details to the reader.
###### Proposition 5.1.
Let $`(M,\mu )`$ a $`\sigma `$-finite measure space, and $`T_t`$ be a semigroup acting on $`L^p(M,\mu )`$, for $`1p+\mathrm{}`$, with infinitesimal generator $`A`$. Suppose that there exists a non-decreasing function $`F`$ on $`\mathrm{I}\mathrm{R}_+`$ such that
$$T_t_{11},T_t_{\mathrm{}\mathrm{}}F(t),t>0,$$
and that
$$\theta (f_2^2)\mathrm{Re}(Af,f),fD(A),f_1C,$$
fro some $`C>0`$, where $`\theta :]0,+\mathrm{}[]0,+\mathrm{}[`$ is continuous and satisfies $`^+\mathrm{}\frac{dx}{\theta (x)}<+\mathrm{}`$. Then $`T_t`$ is ultracontractive and
$$T_t_1\mathrm{}CF^2(t)m(Ct),t>0,$$
for some $`C>0`$, where $`m`$ is the solution of
$$m^{}(t)=\theta (m(t))$$
on $`]0,+\mathrm{}[`$ such that $`m(0)=+\mathrm{}`$, or alternatively the inverse function of $`p(t)=_t^+\mathrm{}\frac{dx}{\theta (x)}`$.
Acknowledgement: The second author acknowledges the support of the University of Cergy-Pontoise during the preparation of this paper. Both authors thanks Adam Sikora for nice remarks on the manuscript.
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# Gamma-spectrometric uranium age-dating using intrinsic efficiency calibration
## 1 Introduction
If illicit trafficking of nuclear materials is detected, nuclear forensic investigations are carried out to determine the possible origin of the material. One of the parameters to be determined by the methods of nuclear forensic science is the age of the sample, together with the isotopic composition of the material, its physical form, chemical composition etc. In addition, in view of the Fissile Material Cut-Off Treaty knowing the age of isotopically enriched uranium is important for identifying newly produced materials.
The purpose of the present work is to develop a non-destructive, gamma-spectrometric method for age-dating of homogeneous uranium samples, based on using intrinsic efficiency calibration. Other methods for age dating of Uranium samples are mass spectrometry , , and alpha-spectrometry, which are destructive methods. The gamma-spectrometric method described in and the method given in the present work offer simple non-destructive alternatives.
Gamma-spectrometric age determination of Uranium samples is based on measuring the activity ratio <sup>214</sup>Bi/<sup>234</sup>U by high-resolution gamma-spectrometry in low-background, assuming that the daughter nuclides have been completely removed during last separation or purification of the material . <sup>214</sup>Bi is a daughter of <sup>234</sup>U, which decays through <sup>230</sup>Th to <sup>226</sup>Ra, which in turn decays to <sup>214</sup>Bi through three short-lived nuclides. The time needed for secular equilibrium between <sup>226</sup>Ra and <sup>214</sup>Bi to be achieved is about 2 weeks, so it can be assumed that the activities of <sup>226</sup>Ra and <sup>214</sup>Bi are equal at the time of the measurement. Therefore, using the law of radioactive decay the activity ratio <sup>214</sup>Bi/<sup>234</sup>U at time $`T`$ after purification of the material may be calculated as
$`{\displaystyle \frac{A_{Bi214}(T)}{A_{U234}(T)}}={\displaystyle \frac{A_{Bi214}(T)}{A_{U234}(0)}}={\displaystyle \frac{A_{Ra226}(T)}{A_{U234}(0)}}=\lambda _2\lambda _3[{\displaystyle \frac{e^{\lambda _1T}}{(\lambda _2\lambda _1)(\lambda _3\lambda _1)}}+`$
$`{\displaystyle \frac{e^{\lambda _2T}}{(\lambda _1\lambda _2)(\lambda _3\lambda _2)}}+{\displaystyle \frac{e^{\lambda _3T}}{(\lambda _1\lambda _3)(\lambda _2\lambda _3)}}],`$ (1)
where $`A_{U234}(0)`$ denotes the activity of <sup>234</sup>U at time $`T=0`$ while $`\lambda _1`$, $`\lambda _2`$ and $`\lambda _3`$ are the decay constants of <sup>234</sup>U, <sup>230</sup>Th and <sup>226</sup>Ra, respectively. Because $`\lambda _1,\lambda _2,\lambda _3<<1`$, formula (1) can be developed into a Taylor series around $`T=0`$, so it can be approximated as
$$\frac{A_{Bi214}(T)}{A_{U234}(T)}=\frac{1}{2}\lambda _2\lambda _3T^2.$$
(2)
This equation can be used for calculating the age, $`T`$, of uranium samples after the activity ratio <sup>214</sup>Bi/<sup>234</sup>U has been determined by gamma spectroscopy.
The activity ratio <sup>214</sup>Bi/<sup>234</sup>U can be determined in several ways from the gamma spectra of the investigated sample. A very precise and reliable way would be to use a reference material of approximately same age, isotopic composition and form as the investigated sample. The samples encountered in illicit trafficking, however, are usually such that an appropriate reference material cannot be found.
Another approach is represented by the method described in , that does not require any reference materials but uses the absolute efficiency of the detector determined by “point-like” standard sources. That method is applicable whenever the sample can be approximated to be “point-like” or if it has a well defined geometrical shape, for which an efficiency calibrated geometry can be constructed and for which the self-attenuation can also be accounted for.
In the present work a third approach is followed. Here intrinsic efficiency calibration was used to determine the activity ratio <sup>214</sup>Bi/<sup>234</sup>U, without the use of any standard or reference materials. The efficiency calibration is “intrinsic” in the sense that a relative efficiency curve is determined from the same spectrum as the measured activity ratios, so the calibration is “intrinsic” to the spectrum. The aim of using intrinsic efficiency calibration is to develop a method appropriate for age-dating of homogeneous uranium samples of any arbitrary physical form and shape. The idea is to use the peaks of <sup>234m</sup>Pa, which is a short-lived daughter of <sup>238</sup>U, to construct a relative efficiency curve, and to determine the activity ratio <sup>214</sup>Bi/<sup>238</sup>U. Furthermore, the activity of <sup>234</sup>U is determined relative to <sup>235</sup>U. If the activity ratio <sup>235</sup>U/<sup>238</sup>U is also known one can calculate the activity ratio <sup>214</sup>Bi/<sup>234</sup>U and thus the age of the sample can be obtained.
In the research described in this work U<sub>3</sub>O<sub>8</sub> powder with a relative isotope abundance in <sup>235</sup>U of 36% was used for the measurements. Gamma spectra of the samples were taken in the 0-300 keV region under standard laboratory conditions and in the 0-2 MeV region in low background. From each spectrum a relative efficiency curve was constructed, using the peaks of <sup>235</sup>U in the first case, and the peaks of <sup>234m</sup>Pa in the second case.
The intensity of the peaks of <sup>234m</sup>Pa can be accurately measured within a reasonable measurement time (1-2 days) if the relative <sup>235</sup>U abundance of the material is less than approximately 90%, provided a sufficient amount of the investigated material is available (depending on detector efficiency). It should be noted that for lower <sup>235</sup>U abundances the amount of <sup>234</sup>U (and therefore of <sup>214</sup>Bi) is lower as well, so the corresponding activity is more difficult to measure and the uncertainty caused by the variation of the natural background becomes more expressed. In addition, a Compton background caused by the peaks of <sup>238</sup>U daughters is also present in the spectrum, disturbing the evaluation of the <sup>214</sup>Bi peaks. Therefore there exists a lower limit on the <sup>235</sup>U abundance of the material whose age can be determined by gamma-spectrometry, depending, of course, on the amount and the age of the material, detector efficiency and background level. For example, with our present equipment the age of 10 grams of uranium oxide powder having natural uranium isotopic composition can only be measured if it is more than about 70 years old. On the other hand, it has been estimated that with a fairly average well type detector, if 10 g of uranium oxide is used for the measurement, the lower limit for age dating could be around 30 years for natural uranium, and as low as 2 years for uranium with 90% of <sup>235</sup>U. A detailed study of the limits for uranium age dating by gamma spectroscopy is the subject of further research.
The age of the material calculated using the described method was compared to the age obtained by using the absolute efficiency of the detector. The results of the two different methods coincide with each other, within the error limits, confirming the reliability of the method presented here.
The structure of this paper is as follows. In section 2 the experimental setup and the procedures for measuring the activity ratios <sup>234</sup>U/<sup>235</sup>U and <sup>214</sup>Bi/<sup>238</sup>U are presented, together with a remark on determining the activity ratio <sup>235</sup>U/<sup>238</sup>U. The results obtained using these procedures are given in section 3, where the age of the investigated material is also calculated. In section 4 a summary of this work is given with some thoughts about directions for further development of the method.
## 2 Experimental
Four samples of the investigated U<sub>3</sub>O<sub>8</sub> powder, having different thicknesses (and therefore different masses), were prepared with masses of approximately 0.5, 2, 5 and 10 g. The four different-mass samples of the material were each held in a thin, closed polyethylene cylinder with an inner diameter of 2.9 cm. The thickness of the layer of U<sub>3</sub>O<sub>8</sub> powder on the bottom of the polyethylene cylinder varied between 0.04 and 0.92 cm, depending on the mass. For determining the activity ratio <sup>234</sup>U/<sup>235</sup>U the spectra of the four samples of the assayed material were taken by planar HPGe detector, whereas the activity ratio <sup>214</sup>Bi/<sup>238</sup>U was measured by a coaxial Germanium detector in a low-background chamber.
### 2.1 Measuring the activity ratios <sup>234</sup>U/<sup>235</sup>U and <sup>235</sup>U/<sup>238</sup>U
In order to compare different instruments, the low-energy spectra were taken by two planar high-purity Germanium detectors separately (Canberra GL2020 with active diameter of 50.5 mm, thickness of 20 mm and active surface of 2000 mm<sup>2</sup> and Ortec GLP-10180/07 with active diameter of 10 mm and length of 7 mm). Both detectors were placed in a vertical position, but the “large” (50.5 mm diameter) detector was looking upward, while the “small” (10 mm diameter) detector was positioned in a downward looking position.
The samples were placed, one after another, at a fixed distance (a few centimeters) from the detector cap and the gamma peak of <sup>234</sup>U at 120.9 keV was observed. The spectrum acquisition lasted until the statistical error of the 120.9 keV line dropped below 2 % . The relative efficiency curves of the two detectors are shown on Fig. 1.
If the sample is not “thin”<sup>1</sup><sup>1</sup>1Note that the concept of a“thin” sample depends on the matrix and density of the sample and to some extent even on the detector efficiency and measurement configuration., the activity ratio <sup>234</sup>U/<sup>235</sup>U of the material can be reliably determined using the “U235” or the “MGAU” code. For thin samples, such as the samples used in this research, however, the results of the U235 and the MGAU codes contain large systematic errors . Therefore, in this research these codes were not used for determining the <sup>234</sup>U content of the material.
In the present work the activity ratio <sup>234</sup>U/<sup>235</sup>U was determined using relative efficiency calibration, as
$$\frac{A(^{234}U)}{A(^{235}U)}=\frac{I_{120.9}/B_{120.9}}{f(120.9keV)}$$
(3)
where $`A(^{234}U)`$ and $`A(^{235}U)`$ denote the activities of <sup>234</sup>U and <sup>235</sup>U, respectively, $`I_{120.9}`$ is the intensity of the 120.9 keV line of <sup>234</sup>U, $`B_{120.9}`$ is its emission probability, while $`f(120.9keV)`$ is the value of the relative efficiency function of the detector at 120.9 keV. The function $`f(E)`$ is obtained by fitting a second order polynomial to the relative efficiencies at the 143.8, 163.3, 185.7 and 205.3 keV peaks of <sup>235</sup>U.
The activity ratio <sup>235</sup>U/<sup>238</sup>U was determined from the spectra taken by the small planar detector. It was calculated as the ratio of the relative isotope abundances of <sup>235</sup>U and <sup>238</sup>U determined using the commercial implementation of the “U235” code, which is part of ORTEC’s software package “MGA++” . The values we obtained for the <sup>235</sup>U abundance using the U235 code do not depend on the thickness of the sample, confirming the findings of , so we calculated the relative isotope abundance of <sup>235</sup>U and <sup>238</sup>U as the weighted average of the values obtained by measuring the four different-mass samples. They were found to be $`36.59\pm 0.01\%`$ and $`63.08\pm 0.01\%`$, respectively. This corresponds to an <sup>235</sup>U/<sup>238</sup>U activity ratio of $`3.664\pm 0.001`$ Bq/Bq.
### 2.2 Measuring the activity ratio <sup>214</sup>Bi/<sup>238</sup>U
The activity ratio <sup>214</sup>Bi/<sup>238</sup>U was measured in a low-background iron chamber using a 150 cm<sup>3</sup> coaxial germanium detector (“PIGC 3520” Intrinsic Coaxial Detector manufactured by PGT) having 34.1 % relative efficiency (at 1332 keV measured at 25 cm source-detector distance, relative to a 3”$`\times `$ 3” NaI(Tl) detector). The wall thickness of the iron chamber is 20 cm and its inner dimensions are $`120\times 60\times 120`$ cm (height $`\times `$ width $`\times `$ length). The detector was standing in vertical position and the samples were placed, one after another, below the detector so that the bottom of the sample holder was at 6 cm from the detector cap. A 4.2 mm lead absorber was attached to the detector cap, in order to reduce the dead-time caused by the high count rate of the low-energy peaks of <sup>235</sup>U. The presence of the lead absorber had no significant influence on the background counts.
The gamma-spectra of the samples were taken in the energy range 0-2300 keV (Fig. 2). The small insertion on Fig. 2 shows the 609.3 keV peak of <sup>214</sup>Bi. The count rates of the gamma-lines of <sup>214</sup>Bi at 609.3, 1120.3 and 1764.5 keV were determined from the spectra. Because of small branching ratios and because detector efficiency decreases with energy, the statistical errors for the counts of the last two gamma rays were significant even for long counting times. In addition, these peaks have to be deconvoluted from the peaks of <sup>234m</sup>Pa at 1120.6 keV and 1765.5 keV, which also decreases the reliability of the results obtained for the counts at these energies. Therefore, the 1120.3 and 1764.5 keV lines were not used in calculating the activity ratio <sup>214</sup>Bi/<sup>238</sup>U.
The activity ratio <sup>214</sup>Bi/<sup>238</sup>U was calculated from the 609.3 keV peak of <sup>214</sup>Bi using the relative efficiency curve obtained from the peaks of <sup>234m</sup>Pa. <sup>234m</sup>Pa is assumed to be in equilibrium with its parent <sup>238</sup>U, so its activity is equal to the activity of <sup>238</sup>U. Analogously to formula (3) one may write
$$\frac{A(^{214}Bi)}{A(^{238}U)}=\frac{I_{609.3}/B_{609.3}}{F(609.3keV)}$$
(4)
where $`A(^{214}Bi)`$ and $`A(^{238}U)`$ denote the activities of <sup>214</sup>Bi and <sup>238</sup>U, respectively, $`I_{609.3}`$ is the intensity of the 609.3 keV line of <sup>214</sup>Bi, $`B_{609.3}`$ is its emission probability, while $`F(609.3keV)`$ is the value of the relative efficiency function of the detector at 609.3 keV. The function $`F(E)`$ is obtained by fitting a second order polynomial to the relative efficiencies at the 569.30, 766.37, 1000.99, 1193.69, 1510.20, 1737.73 and 1831.36 keV peaks of <sup>234m</sup>Pa.
It can be clearly seen in Fig. 2 that in the 10 g case the count rate of the 609.3 keV peak is much higher than the background, but there is also a high Compton-continuum in the spectrum of the 10 g sample, coming from the higher energy peaks of uranium daughters. The Compton-continuum unfortunately disturbs the evaluation of the relevant peaks. In order to facilitate a more accurate evaluation of the <sup>214</sup>Bi peaks, using a detector having a higher peak-to-Compton ratio would be desirable.
## 3 Determining the age of the material
In order to calculate the age of the sample, one needs the activity ratio <sup>214</sup>Bi/<sup>234</sup>U. It is calculated from the results of the above described measurements as
$$\frac{A(^{214}Bi)}{A(^{234}U)}=\frac{A(^{214}Bi)}{A(^{238}U)}\left(\frac{A(^{235}U)}{A(^{238}U)}\frac{A(^{234}U)}{A(^{235}U)}\right)^1.$$
(5)
The age of the material was calculated from equation (2) for each sample separately, as presented in Table 1. The values $`T_{1/2}(^{230}Th)=7.538\times 10^4`$ years and $`T_{1/2}(^{226}Ra)=1600`$ years taken from reference for the half life of <sup>230</sup>Th and <sup>226</sup>Ra, respectively, were used in the calculations.
It can be seen from Table 1 that the results obtained using different amounts of the investigated material agree very well with each other, except for the case of the 0.52 g sample. In the case of the smallest sample the error of determining the count rate of <sup>214</sup>Bi is very high, because in this case the relative contribution of the background to the total count rate at 609.3 keV is very high (about 70%). That is, the error of the count rate at 609.3 keV is dominated by the fluctuations of the background. Accordingly, the results of repeated measurements with the smallest sample placed at 6 cm distance from the detector have shown a big dispersion. The influence of the background can be lowered by placing the sample closer to the detector, so that the relative contribution of the background to the total counts becomes smaller. Therefore, another measurement with the 0.52 g sample was made, but this time placing the sample in such a way that the sample holder was touching the lead shielding. In this measurement the relative contribution of the background to the total counts of the 609.3 kev line was about 50 % and the activity ratio <sup>214</sup>Bi/<sup>238</sup>U was calculated to be $`4\pm 1\times 10^4`$ Bq/Bq, resulting in $`46\pm 8`$ years for the age of the sample. Although this value is close to the value obtained by using larger amounts of the investigated material, its uncertainty is still quite high. Therefore, the measurement of the age of small-amount samples and a detailed study of background fluctuations deserve further attention.
The agreement of the results obtained using different-mass (i.e. different-thickness) samples already indicates that the applied method is not sensitive to the measurement configuration. In order to confirm that the results do not depend on the measurement geometry, the activity ratio <sup>214</sup>Bi/<sup>238</sup>U in the 10 g sample was also determined in two additional, modified measurement geometries. First the sample was placed off-center with respect to the axis of the detector, at approximately 2 cm from the detector axis. Then the cylindrical sample holder was positioned again to the detector axis, but in such a way, that the axis of the cylinder was at an angle of approximately 45 degrees with respect to the axis of the detector. The activity ratio <sup>214</sup>Bi/<sup>238</sup>U was found to be $`(3.6\pm 0.3)\times 10^4`$ Bq/Bq in the first case and $`(3.5\pm 0.3)\times 10^4`$ Bq/Bq in the second case. Both these values agree well with the values listed in Table 1, confirming the independence of the results on the measurement geometry and on the shape of the sample.
As for the activity ratio <sup>234</sup>U/<sup>238</sup>U, it can be seen that the values obtained by using different detectors are consistent with each other. If the larger detector is used, the counting times needed to achieve the desired counting statistics are consequently shorter, so this detector may be the preferred choice when only a small amount of the investigated material, producing a low count rate, is available. On the other hand, if time is not critical or if a larger amount of the material to be measured is at disposal, the smaller detector might be a better choice. Namely, the smaller detector has better energy resolution, making possible the automatized calculation of the <sup>235</sup>U and <sup>238</sup>U abundance using the computer codes U235 or MGAU, which require a good energy resolution to produce correct results. Note that the uncertainty associated with the evaluation of the relative efficiency function at 120.9 keV is relatively high, producing an uncertainty of about 9-10 % in the value of the activity ratio <sup>234</sup>U/<sup>238</sup>U (see Table 1), which is much higher than the dispersion of the measured values. Presently further studies are being carried out in order to understand better and reduce the uncertainties of determining the activity ratio <sup>234</sup>U/<sup>238</sup>U.
For comparison, the age of the investigated material was also calculated by using the absolute efficiency of the detector, as in reference . Denoting the absolute efficiency of the detector at energy $`E`$ by $`ϵ_E`$, the activity, $`A`$, of a given isotope can be calculated as
$$A=\frac{I_E}{B_Eϵ_E},$$
(6)
where $`B_E`$ and $`I_E`$ are, respectively, the emission probability and measured count rate of the observed gamma line of energy $`E`$, corresponding to the given isotope. In order to account for the self absorption of the 120.9 keV line within the sample, the procedure described in was applied. The count rate corrected for self-absorption is given in Table 2. In the case of the 609 keV line of <sup>214</sup>Bi the uncertainty caused by self-absorption was negligible compared to the statistical error of the observed peaks. Namely, the layer of U<sub>3</sub>O<sub>8</sub> powder used in the measurements is practically transparent for the gamma photons at these energies and there was no need to correct for self absorption. This was also confirmed by measurements, i.e., the count rates per unit mass for the different-mass samples were approximately the same. A much more important source of error is the random fluctuation of the background level. In the case of the 10 g sample the count rate is the highest, so in this case the relative uncertainty produced by the fluctuation of the background level is the lowest. Therefore, only the results of the measurements with the 10 g sample were used for calculating the activity of <sup>214</sup>Bi (Table 2). The age obtained using the activities given in Table 2 is $`45\pm 4`$ years, which, within the error limits, coincides with the value obtained using intrinsic efficiency calibration.
## 4 Conclusion
In this work a gamma-spectrometric method for determining the age of isotopically enriched uranium samples with medium <sup>235</sup>U abundance is presented. The method does not require the use of standard samples of known age, nor the knowledge of the absolute detector efficiency and it is applicable for determining the age of uranium samples of any arbitrary geometrical shape. The age of the investigated sample is calculated from the activity ratio <sup>214</sup>Bi/<sup>234</sup>U, using formula (2). In the present work intrinsic efficiency calibration was used to determine the activity ratios <sup>234</sup>U/<sup>235</sup>U and <sup>214</sup>Bi/<sup>238</sup>U, from which the activity ratio <sup>214</sup>Bi/<sup>234</sup>U can also be calculated, provided that the activity ratio <sup>235</sup>U/<sup>238</sup>U is known. The activity ratio <sup>235</sup>U/<sup>238</sup>U can be determined from the same spectrum as the activity ratio <sup>234</sup>U/<sup>235</sup>U, using standard methods based on intrinsic efficiency calibration, such as, e.g., “U235” and “MGAU”.
The results obtained using the described method were confirmed by a method which uses the absolute efficiency of the detector. It can be concluded that the method described in this work is a reliable alternative for determining the age of isotopically enriched uranium samples with a medium <sup>235</sup>U abundance, having any arbitrary shape.
Presently the background fluctuations are being studied in detail, in order to find a way to reduce them and thus improve the precision of evaluating the 609.3 keV peak of <sup>214</sup>Bi. Another important possibility for extending the capabilities of the method would be the use of a detector with a higher efficiency and thus higher peak-to-Compton ratio for the low-background measurements. Finally, the uncertainties associated with determining the activity ratio <sup>234</sup>U/<sup>235</sup>U by intrinsic efficiency calibration also deserve further attention.
## Acknowledgment
This research has been partially supported by the Hungarian Atomic Energy Authority under the research and development contract OAH/ÁNI-ABA-02/04. The authors wish to thank Laszlo Lakosi for checking the manuscript.
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# Interstellar 12C/13C ratios through CH⁺𝜆𝜆3957,4232 absorption in local clouds: incomplete mixing in the ISM Based on observations obtained with UVES at the ESO Very Large Telescope, Paranal, Chile (proposal No. 71.C-0367(A))
## 1 Introduction
The ratio of the <sup>12</sup>C to <sup>13</sup>C isotope is a good tracer of the amount of stellar processing in low and intermediate mass stars during the asymptotic giant branch (AGB) phase. The third dredge-up increases <sup>12</sup>C/<sup>13</sup>C to about $`300`$, while stars massive enough to undergo “hot bottom burning” bring their surface carbon isotopic ratio to CNO processing equilibrium, <sup>12</sup>C/<sup>13</sup>C$``$3. The carbon yields of AGB stars are sensitive functions of their initial masses and metallicities (e.g. Renzini & Voli ren81 (1981), or Casassus & Roche cas01 (2001) for a population synthesis approach and references to model results).
A measurement of the <sup>12</sup>C/<sup>13</sup>C ratio in the interstellar medium (ISM) gives important data for the total amount of low mass stellar evolution and subsequent enrichment of the ISM in the Galaxy. The terrestrial value of <sup>12</sup>C/<sup>13</sup>C is 90 (Rosman & Taylor ros98 (1998)). The <sup>12</sup>C/<sup>13</sup>C ratio is a cornerstone of models of the nuclear history of our galactic ISM, that is, the Galactic Chemical Evolution. Models of chemical evolution predict a decrease in the <sup>12</sup>C/<sup>13</sup>C ratio with time for a given galactocentric distance, and a decrease with galactocentric distance in the Galaxy (e.g. Palla et al. pal00 (2000), their Fig. 4 and 5.). Such models use the assumption that mixing in the ISM is complete and restricted to material at a given galactocentric distance (azimuthal mixing). A comparison of results between our Galaxy and other galaxies will give important data for the nuclear processing history (see, e.g., Tosi 2000, Prantzos 2001).
There are a large number of measurements of the <sup>12</sup>C/<sup>13</sup>C ratio from radio astronomy data (see, e.g., Wilson & Rood wil94 (1994)). However, these results may be affected by interstellar chemistry in two ways. These are chemical fractionation (which enriches molecules in <sup>13</sup>C and thus lowers the ratio) and selective dissociation (which destroys the rarer species more, and thus raises the ratio). The <sup>13</sup>C enrichment is due to fractionation of the CO molecule via the ion-molecule process,
$${}_{}{}^{13}\mathrm{C}_{}^{+}+^{12}\mathrm{CO}\stackrel{12}{\stackrel{k}{}}\mathrm{C}^++^{13}\mathrm{CO}+\mathrm{\Delta }E,$$
(1)
with $`\mathrm{\Delta }E/k=35`$K and $`k=\mathrm{2\; 10}^{10}`$cm<sup>3</sup> s<sup>-1</sup> in both directions (Watson et al. wat76 (1976)). If the reaction reaches equilibrium,
$$\frac{n(^{13}\mathrm{C}^+)n(^{12}\mathrm{CO})}{n(^{12}\mathrm{C}^+)n(^{13}\mathrm{CO})}=\mathrm{exp}\left(\frac{\mathrm{\Delta }E}{kT}\right).$$
(2)
In cold gas, with $`T35`$K, CO will be enriched in <sup>13</sup>C compared to C<sup>+</sup>. The effect is observed in CO absorption towards $`\zeta `$Oph (Sheffer et al. she92 (1992)), $`\rho `$Oph and $`\chi `$Oph (Federman et al. fed03 (2003)). Thus molecules whose synthesis involve CO at a later stage than C<sup>+</sup> will show the CO fractionation, and vice-versa (e.g. Watson wat78 (1978)). The corresponding <sup>13</sup>C depletion in C<sup>+</sup> is diluted in environments where $`n(`$C$`{}_{}{}^{+})n(`$CO), such as diffuse clouds or photo-dissocitation regions (PDRs) at $`A_\mathrm{V}<2`$ (e.g. Keene et al. kee98 (1998), their Fig. 2).
Ultra-high-resolution spectroscopy of CH<sup>+</sup> (Crawford et al. cra94 (1994)) absorption towards $`\zeta `$Oph reveals broader profiles compared to CN and CH, with an upper limit kinetic temperature of $``$2000 K, and with no velocity offset. This is consistent with CH<sup>+</sup> production through the endothermic reaction
$$\mathrm{C}^++\mathrm{H}_2+0.4\mathrm{eV}\mathrm{CH}^++\mathrm{H}.$$
(3)
The analog of Eq. 2 shows that temperatures $`>1000`$K are required for CH<sup>+</sup> production via Eq. 3. CH<sup>+</sup> thus arises in the warm halo, or photo-dissociation region, surrounding the $`\zeta `$Oph cloud. Crawford (cra95 (1995)) and Crane et al. (cran95 (1991)) extend the $`\zeta `$Oph results to five and twenty other lines of sight through translucent clouds (i.e. with $`A_\mathrm{V}<2`$). With such a kinetic temperature, it is clear that CO fractionation in the sites of CH<sup>+</sup> production is negligible.
The carbon isotope ratio as measured by <sup>12</sup>CH<sup>+</sup>/<sup>13</sup>CH<sup>+</sup> is not affected by selective dissociation. Although photodissociation of CH<sup>+</sup> contributes to PDR chemical networks (Sternberg & Dalgarno ste92 (1992)), CH<sup>+</sup> is optically thin to photodissociating UV radiation in translucent clouds. The strongest optical CH<sup>+</sup> line in $`\zeta `$Oph is $`A^1\mathrm{\Pi }X^1\mathrm{\Sigma }^+`$ at 4232 Å (Morton mor75 (1975)), which reaches maximum opacities of $``$0.3 (e.g. this work). Since no ultra-violet CH<sup>+</sup> lines have been reported, they are bound to be faint, and thus thin. The photodissociation cross-section $`\sigma (\nu )`$ calculated by Kirby et al. (kir80 (1980)) is mostly due to transitions from the ground $`X^1\mathrm{\Sigma }^+`$ state to the vibrational continua of the $`2^1\mathrm{\Sigma }^+`$, $`3^1\mathrm{\Sigma }^+`$, and $`2^1\mathrm{\Pi }`$ states, which smooths out any difference between isotopes. Some resonant absorption derives from the calculation of $`\sigma (\lambda )`$ at wavelengths of 1418Å, 1453Å, and 1490Å. But Kirby et al. (kir80 (1980)) assign oscillator strengths of less than $`10^3`$ in any of these peaks, which represents Einstein $`B`$ absorption rates one order of magnitude smaller than at 4232 Å.
The <sup>12</sup>CH<sup>+</sup>/<sup>13</sup>CH<sup>+</sup> ratio is thus a secure representation of <sup>12</sup>C/<sup>13</sup>C in translucent clouds. In addition, the absorption line from <sup>12</sup>CH<sup>+</sup> at 4232Å is separated from the line of <sup>13</sup>CH<sup>+</sup> by 0.265Å, which permits unblended measurements. These ratios, obtained from optical measurements, are restricted to regions within $``$2 kpc of the Sun because of the need for bright background stars, and of narrow velocity profiles in the intervening cloud. Eight ratios have been measured but only one line of sight, toward $`\zeta `$Oph, has an excellent signal to noise ratio. The average <sup>12</sup>C/<sup>13</sup>C ratio from CH<sup>+</sup> data agrees well with those obtained for sources near the Sun, from radio astronomy data. This gives one confidence in the ratios obtained from radio data.
However, the CH<sup>+</sup> measurements so far resulted in isotopic ratios with a large scatter between different lines of sight, and there is only one line of sight with an excellent signal to noise ratio. The previous data from CH<sup>+</sup> were taken with the ESO CAT, and 4-m or smaller aperture telescopes. The best results from the point of view of the signal to noise ratio and the reliability of the ratio for $`\zeta `$Oph were presented by 3 different groups (Stahl et al. sta89 (1989), Stahl & Wilson sta92 (1992), Crane et al. cra91 (1991), Hawkins et al. haw85 (1985), haw93 (1993)). The final ratio obtained by each group was 67$`\pm `$5. The $`\zeta `$Oph value is by far the best measurement of the <sup>12</sup>C/<sup>13</sup>C ratio for the ISM near the Sun. However, the measurements for 7 other nearby weak sources were of lower quality. The ratios for 3 of these regions, toward HD26676 (64$`\pm 6`$, Centurion & Vladilo cen91 (1991)), HD110432 ($`71\pm 11`$, Centurion et al. cen95 (1995)), and $`\mu `$Norma ($`67\pm 6`$, Centurion & Vladilo cen91 (1991) ) are in very good agreement with the $`\zeta `$Oph ratio. But there are discordant ratios: $`R=38\pm 12`$ toward HD157038 (Hawkins & Meyer 1989), and $`R=49\pm 15`$ toward $`\zeta `$Per (Hawkins et al. haw93 (1993)) while the ratio toward HD152235 is $`R=126\pm 29`$ and is $`R=98\pm 19`$ towards HD152424 (Vladilo et al. vla93 (1993)). Most of the measurements are limited by noise, but the profile toward HD152424 is complex and thus a determination of the ratio is not easy, because of CH<sup>+</sup> and <sup>13</sup>CH<sup>+</sup> line blending.
We repeated the CH<sup>+</sup> observations toward $`\zeta `$Oph, HD110432, HD152235, HD152424 and HD157038, and included the new lines of sight toward HD152236, HD154368, HD161056, HD169454, HD170740. Our goal was first to determine whether the ratios are truly different from the value for $`\zeta `$Oph, which would indicate that the mixing of the ISM is not very fast and complete, and second, we wanted to increase the sample of lines of sight with good to excellent signal to noise ratios, and from this determine an average ratio near the Sun truly representative for the local ISM.
Section 2 describes the observations, and Section 3 gives details on our procedure to measure the isotopic ratio $`R`$. Section 4 contains our results , which are then summarised and discussed in Section 5. Section 6 concludes.
## 2 Observations and reduction
There are two systems of CH<sup>+</sup> lines which have appreciable oscillator strengths. Because of atmospheric absorption, the lines at 4232 Å are easier to measure. The <sup>13</sup>CH<sup>+</sup> line is 0.26 Å ($`\mathrm{\Delta }v=19`$km s<sup>-1</sup>) shortward of the <sup>12</sup>CH<sup>+</sup> line. The other system, at 3957 Å, is more difficult to measure, but has the property that the <sup>13</sup>CH<sup>+</sup> line is on the long wavelength side of the <sup>12</sup>CH<sup>+</sup> line at a separation of 0.44 Å($`\mathrm{\Delta }v=33`$km s<sup>-1</sup>). This property has been used by Stahl & Wilson (sta92 (1992)) to check that there is no accidental overlap of a weak velocity component of <sup>12</sup>CH<sup>+</sup> at $``$18.8 km sec<sup>-1</sup> from the deepest <sup>12</sup>CH<sup>+</sup> absorption line with the <sup>13</sup>CH<sup>+</sup> line. Nearly all previous data were taken at 4232 Å.
Our observations were obtained with the UVES echelle spectrograph at the VLT unit telescope Kueyen at Cerro Paranal, Chile, in three nights between June 14/15 and June 16/17, 2003. UVES allows measuring both line systems in the same detector setting, thus providing a means to correct for line blending if it is apparent in the spectra.
The observations are difficult because of the need for both very high S/N and very high spectral resolution to acquire the profile of faint <sup>13</sup>CH<sup>+</sup>, against a very bright stellar continuum and close to the much stronger <sup>12</sup>CH<sup>+</sup> component. We need S/N ratios of at least $``$10 on <sup>13</sup>CH<sup>+</sup> for accurate profile fitting. The project requires the highest possible spectral resolution. We therefore used an image slicer to minimize flux losses. In addition, the image slicer distributes the light along the slit, which improves flat-fielding and allows for our bright targets longer integration times before saturation occurs. UVES slicer #2 was used, which reformats an entrance opening of 1$`\stackrel{}{.}`$8 $`\times `$ 2$`\stackrel{}{.}`$2 to a slit of 0$`\stackrel{}{.}`$44 width and a length of 7$`\stackrel{}{.}`$9, which is imaged on the spectrograph entrance slit of 0$`\stackrel{}{.}`$45 width. The spectral resolution in this configuration is $`\lambda `$/$`\mathrm{\Delta }\lambda `$ = 75 000. The central wavelength was set to 4 370 Å, which gives a spectral coverage from 3 730 to 5 000 Å. During our observations the seeing was typically too good to fill the entrance aperture of the images slicer. Therefore, for most spectra, only one or two slices contained most of the signal, which unfortunately decreases the expected gain of the image slicer.
Accurate flat-fielding is also important. Therefore a large number of flat-fields (200) was obtained during daytime distributed along the observing run.
In addition, a rapidly rotating unreddened early-type star ($`\epsilon `$ Cap = HD205637, spectral type B3V) was observed in order to check for possible faint terrestrial features. We confirm there are no detectable telluric features under the CH<sup>+</sup> absorption, and show in Fig. 1 the template star spectrum in the region of interest, at airmasses $`<1.08`$. Higher airmasses result in telluric absorption features, as in the case of $`\zeta `$Oph discussed below.
The brighter targets where exposed until about 50% of the maximum level allowed by the CCD detector was reached. For $`\zeta `$ Oph, this limits the exposure time to about 5 sec, and typically a few minutes for the fainter targets. Series of up to 50 exposures per night were obtained to build up the required S/N-ratio. Some targets turned out to show too complex line profiles. These were dropped from our initial target list. The observations are summarized in Table 1, which also lists the mean airmass of the target objects in the different nights. At the higher airmass values, the range in airmass differed from the mean by about $`\pm `$0.1 during the observations, and less at smaller airmass.
We also list in Table 1 the radial velocities relative to the solar barycentre of the CH<sup>+</sup> absorption. These are calculated from the shift in wavelength between the rest wavelengths of the transitions and of the average of the Gaussian centroids weighted by their equivalent widths. The uncertainty on the derived velocities is conservatively $`1`$km s<sup>-1</sup>, and depends on the accuracy of the rest wavelengths.
We used the Midas package originally developed for the ESO Feros spectrograph for the reduction of the spectra (Stahl et al., sta99 (1999)). In order to maximize the S/N-ratio of the extracted spectra, all flat-fields obtained during the run were averaged. After background subtraction and flat-fielding, the spectra were extracted with a very long slit, extracting all slices together. The wavelength calibration of the extracted spectra was done with a 2D-polynomial, fitting all echelle orders in one step. A mean ThAr-spectrum obtained during day time was used for the calibration of all spectra obtained in one night. Finally, all spectra were merged to a 1D-spectrum and all observations of each night averaged to a nightly mean spectrum. Wavelengths are reported in air and refererred to the solar barycentre.
The nightly spectra were combined in a weighted average to produce coadded spectra. The weights were taken as 1/$`\sigma ^2`$, where $`\sigma `$ is the noise in each spectra, as calculated from the root-mean-square deviations from a linear fit to a region of the spectrum devoid of conspicuous features (we chose 4204.2Å to 4205.9Å). Table 2 summarise the resulting weights.
## 3 Model line profiles
In order to use the fact that both <sup>12</sup>CH<sup>+</sup> and <sup>13</sup>CH<sup>+</sup> lines have the same opacity profile $`\tau (v)`$, we must fit the <sup>12</sup>CH<sup>+</sup> absorption with a parametrised model and scale it to <sup>13</sup>CH<sup>+</sup>.
We considered using Voigt profiles to account for the intrinsic line profiles of the CH<sup>+</sup> vibronic lines. In order to fit the low-level broad wings (e.g. Stahl et al. sta89 (1989)) towards $`\zeta `$ Oph with an hypothetical Lorentzian core, Einstein $`A`$ values of order $`10^9`$ s<sup>-1</sup> are required. Gredel. et al. (gre93 (1993)) quote an oscillator strength of $`f_{}=0.00545`$, or $`A=\mathrm{2.258\; 10}^4`$, for <sup>12</sup>CH$`{}_{}{}^{+}\lambda `$4232 $`\mathrm{\Pi }\mathrm{\Sigma }(0,0)`$ and $`f_{}=0.00331`$, or $`A=\mathrm{1.568\; 10}^4`$, for <sup>12</sup>CH$`{}_{}{}^{+}\lambda `$3957 $`\mathrm{\Pi }\mathrm{\Sigma }(1,0)`$. We therefore used simple Gaussian profiles to describe the lines.
We preferred to fit the two overtones separately, and thus obtain independent measurements with which to assess the role of systematics. A simultaneous fit of both overtones could have helped constrain the opacity profiles, which in our separate fits do not always share the same Gaussian components.
The fitting algorithm is as follows.
1. Extract a 1.5Å spectrum $`F(\lambda )`$ centred on the CH<sup>+</sup> line, $`\lambda _1<\lambda <\lambda _2`$.
2. Define a baseline $`F_c(\lambda )`$ with a $`417^{\mathrm{th}}`$ order Legendre polynomial to the data, as in Sembach & Savage (sem92 (1992)), $`F_c(\lambda )=_{i=1}^la_iP_i[(\lambda \lambda _1)/(\lambda _2\lambda _1)]`$.
3. Fit for the $`\{a_i\}_{i=1}^l`$ coefficients in a least square sense, ignoring the neighbourhood of the ISM and telluric lines, and weighting the data to improve the quality of the baselines near the interstellar absorption lines. Store the rms dispersion of the residuals as the spectrum’s noise, $`\sigma _F`$. $`\sigma _F`$ is updated after fitting the parametrised model; it is replaced by the dispersion of the residuals in Step 8.
4. Define a model spectrum with
$$F_m^{}(\lambda )=F_c(\lambda )\mathrm{exp}(\tau (\lambda )),$$
(4)
where the line absorption opacity $`\tau `$ is a superposition of $`n_\mathrm{g}=1`$ to $`6`$ Gaussians on each isotope, with <sup>13</sup>CH<sup>+</sup> components sharing the parameters of the <sup>12</sup>CH<sup>+</sup> components, except their centroids are translated by a fixed velocity shift $`\mathrm{\Delta }v_{\mathrm{iso}}`$, and their amplitude are scaled by a factor $`f=^{13}`$CH$`{}_{}{}^{+}/_{}^{12}`$CH<sup>+</sup>, which is a free parameter in our fit:
$`\tau (\lambda )={\displaystyle \underset{i=1}{\overset{n_\mathrm{g}}{}}}[\tau _i^{}\mathrm{exp}(0.5(\lambda \lambda _i^{})^2/\sigma _{}^{}{}_{i}{}^{2})+`$ (5)
$`f\tau _i^{}\mathrm{exp}(0.5(\lambda \lambda _i^{}(1+\mathrm{\Delta }v_{\mathrm{iso}}/c))^2/\sigma _{}^{}{}_{i}{}^{2})].`$
The value of $`\mathrm{\Delta }v_{\mathrm{iso}}`$ depends on the overtone and is adjusted iteratively on $`\zeta `$Oph and checked for consistency on all targets. We use $`\mathrm{\Delta }v_{\mathrm{iso}}(3957)=32.721`$km s<sup>-1</sup>, $`\mathrm{\Delta }v_{\mathrm{iso}}(4232)=18.967`$km s<sup>-1</sup>, with uncertainties of $``$0.4 km s<sup>-1</sup> (roughly 1/10 the resolution element).
5. Convolve the model spectrum with the instrument response, $`B_\lambda `$:
$$F_m(\lambda )=𝑑\lambda ^{}F_m^{}(\lambda ^{})B(\lambda \lambda _{}).$$
(6)
We estimate the instrumental response from the arc lamp spectrum. The smallest line-widths are indicative of resolving powers of $`\mathrm{83\hspace{0.17em}000}\pm 2000`$ rather than the nominal $`\mathrm{75\hspace{0.17em}000}`$ for our setup. The instrumental response is approximated to a Gaussian, with a FWHM given by $`\lambda _{}/83000`$, where $`\lambda _{}`$ is the center wavelength of each CH<sup>+</sup> overtone.
6. Calculate the goodness of fit function
$$\chi ^2=\underset{j}{}(F(\lambda _j)F_m(\lambda _j))^2/\sigma _F^2,$$
(7)
using the noise from Step 3. We use a single noise value and neglect possible variations of the noise level, as would be the case for Poisson statistics. The noise variations are important for very steep spectra, which is not our case, or under deep absorption lines. But the accuracy of the model line profiles under <sup>12</sup>CH<sup>+</sup> is generally not important to infer isotope ratios (only the wings matter).
7. Optimize the model parameters by minimizing $`\chi ^2`$. Perform an initial estimate of the global minimum with a genetic algorithm (pikaia, Charbonneau cha95 (1995)), and optimize with the downhill simplex method (amoeba, Press et al. pre86 (1986)).
8. Update the noise of the stellar spectrum by replacing $`\sigma _F`$ with the rms dispersion of the residuals (the difference between the observed and model spectra).
9. Estimate the uncertainty in model parameters by two methods:
1. search parameter space one parameter at a time to find the region enclosing 68.3% confidence for the $`\chi ^2`$ distribution with one degree of freedom, thus obtaining different upwards and downwards error bars and accounting for statistical bias,
2. estimate $`1\sigma `$ uncertainties from the curvature matrix of $`\chi ^2`$ (i.e. approximate to normal errors).
The UVES spectra are oversampled, so that the spectral datapoints are correlated. Independent datapoints are separated by roughly one resolution element, of $`\mathrm{\Delta }\lambda =\lambda /750000.05\AA `$. Since the spectra are sampled with an interval of 0.01 Å, $`N_{\mathrm{cor}}5`$ consecutive datapoints are correlated. In some cases the number of Gaussian components involves as many free-parameters as independent data points (as for HD154368 and HD169454). Some Gaussian components may thus be redundant. But we caution that our purpose is to obtain isotopic ratios, not to report on the velocity structure of the ISM in CH<sup>+</sup>. To constrain the opacity profile consistently we would have had to fit the two overtones simultaneously.
The correlation between the spectral datapoints is ignored when fitting for the parametrised model. The off-diagonal terms of the covariance matrix add terms to the $`\chi ^2`$ goodness-of-fit estimator (Eq. 7). The number of additional terms is roughly $`N_{\mathrm{cor}}`$ per diagonal element, which increase $`\chi ^2`$ by a factor $`N_{\mathrm{cor}}`$. But the noise estimate from the dispersion of the residuals also increases due to the correlation of the datapoints, the actual noise should be $`\sigma _F\times \sqrt{N_{\mathrm{cor}}}`$. Both corrections cancel out and leave $`\chi ^2`$ as in Eq. 7.
We tested for systematic noise by comparing the noise in the nightly coadded spectra, $`\sigma _{\mathrm{coadd}}`$, with that expected from the noise in each of the $`N_{\mathrm{exp}}`$ individual exposures, $`\sigma _1`$. For thermal noise it should hold that $`\sigma _{\mathrm{coadd}}=\sigma _1/\sqrt{N_{\mathrm{exp}}}`$. We define $`Q=\sigma _1/(\sqrt{N_{\mathrm{exp}}}\sigma _{\mathrm{coadd}})`$, and raise the order of the baselines until $`Q1`$. The tests were run for the case of HD170740, which presents the most ragged continuum. The reference baseline is defined on the coadded spectrum, and the noise is given by the dispersion of the residuals, excluding the spectral regions affected by the ISM lines (these regions are defined on Fig. 12). We scale the reference baseline to the line-free continua of single exposures using the ratio of their median values. Fig. 2 shows we need to reach orders of $`l`$15 for $`\lambda `$3957, and of $`l`$8 for $`\lambda `$4232. Below these orders the continua are affected by systematic noise. Unfortunately the level of systematic noise changes for each observation and setting, so that we can only use the tests on HD170740 as guidelines for the other lines of sight. In practice we chose the smallest order that is compatible with the data, without exceeding the limits in HD170740.
The isotopic ratio derived from the fits is the ratio of the column densities of each isotope, $`R=N(^{12}\mathrm{CH}^+)/N(^{13}\mathrm{CH}^+)`$. With the best fit $`f`$ we can calculate $`R=1/f`$ for a line of sight with constant $`R`$. We neglect the small difference in oscillator strength between <sup>12</sup>CH<sup>+</sup> and <sup>13</sup>CH<sup>+</sup>. The two ions have slightly different vibrational structures which will lead to small changes in the $`f_{}`$ values for individual vibrationally resolved transitions (J. Tennyson, private communication). We are thus neglecting $`f_{}`$-value differences of the order of the relative difference in reduced mass, or $`\mathrm{\Delta }f_{}/f_{}1/170`$.
In taking statistics on the values of $`R=1/f`$ in the local ISM care must be taken to assess the significance of measurements along individual sightlines. The uncertainties on individual fits do not include the systematic error involved in baseline definition. If baselines were known a priori, the expectation value and uncertainty on $`f`$ for each individual sightline would be obtained from the statistics of multiple measurements of the same quantity. Thus, in the absence of this systematic error, $`f`$ and $`\sigma (f)`$ correspond to the weighted average and quadratic sum of the weighted uncertainties, where the weights are taken as the inverse variance of each measurement, $`1/\sigma ^2`$:
$`f`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{N}{}}}f_iw_i,`$ (8)
$`\sigma _1^2(f)`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{N}{}}}(\sigma _iw_i)^2=1/{\displaystyle \underset{i=1}{\overset{N}{}}}(1/\sigma _i)^2,`$ (9)
for $`N`$ measurements on a given sightline, and where the individual errors are those derived from the curvature matrix.
To estimate the systematic error involved in baseline definition, we compare a posteriori the scatter of individual measurements with that expected from the noise level. We use the rms dispersion
$$\sigma _2(f)=\sqrt{\underset{i}{}w_i\left[(f_if)\right]^2}.$$
(10)
The uncertainty $`\sigma _1(f)`$ neglects the systematic errors. $`\sigma _2(f)`$ uncertainties will give larger error bars, but are probably too conservative considering we are in fact repeatedly measuring the same quantity. Thus the actual uncertainty on our measurement of $`f`$ for each sightline is intermediate between $`\sigma _1(f)`$ and $`\sigma _2(f)`$. For this reason we list both values. We have set $`\sigma _2=\sigma _1`$ if $`\sigma _2<\sigma _1`$, as is the case when only one ratio is available to sample the effects of systematic uncertainties in the continuum definition (or when one ratio has much higher weight than the others).
## 4 Results
An application of the fitting procedure described in Sec. 3 to our sample of 10 lines of sight gives the best fit parameters and isotopic ratios listed in Table 3.
The isotopic ratio $`f=^{13}`$C$`/^{12}`$C is better suited as free parameter than $`R=1/f`$ for the purpose of fitting individual spectra, and subsequently averaging the best fit values. For noisy data the error propagation when using $`R`$ would involve second order expansions in $`R(f)`$ when calculating averages and dispersions. In what follows we prefer to list our results in terms of $`f`$, and revert to the more common usage of $`R`$ in Section 5 when comparing with previous data.
Table 3 lists the isotope ratios and best fit parameters. It is apparent the $`\lambda `$4232 line generally gives the best estimate of $`R`$. We note from Table 3 that the normal errors are good approximations to the 68.3% confidence limits, which justifies their subsequent use in the weighted averages. We also note the weighted average of several measurements for the same line of sight are consistent, within the errors, with the fit to the coadded spectra (which are in general excluded from the averages), giving confidence in our procedure.
### 4.1 Notes on individual objects
#### 4.1.1 $`\zeta `$ Oph
From Hipparcos results, the distance to $`\zeta `$Oph gives an upper limit to the absorbing cloud of 140$`\pm `$14 pc. The spectra of $`\zeta `$Oph around the CH<sup>+</sup> lines are shown on Fig. 3, where we have also plotted the Gaussian components on the <sup>12</sup>CH<sup>+</sup> line (the individual components on the <sup>13</sup>CH<sup>+</sup> line are left out for clarity).
Residual features under the absorption lines vary on scales of $``$0.02Å . These reflect the imperfection in the model line profiles, more Gaussian components reduce these residuals.
The $`\zeta `$Oph spectra were acquired at airmasses of 1.8 and 1.7 on the 15<sup>th</sup> and 16<sup>th</sup> of June. Both spectra show absorption at 3958.15$`\AA `$ (at $`3958.5\AA `$ in the observatory’s rest frame). This is absent in the spectrum from 14<sup>th</sup> of June and in all other targets. Since all other spectra were acquired with low airmass values compared with that of $`\zeta `$Oph on the 15<sup>th</sup> and 16<sup>th</sup>, we conclude the 3958.15Å feature is telluric. We did not attempt to include the 3958.15Å feature in the fits. The telluric absorption propagates into the residuals, and contributes to the noise level. Note that the telluric feature has no effect on where we set the baseline level, because we assigned zero weights to the spectral points in its neighbourhood.
Stahl et al. (sta89 (1989)) reported $`R=77\pm 3`$<sup>1</sup><sup>1</sup>1note their uncertainty is estimated as for $`\sigma _1`$ in this work, in the case where all samples are assigned the same weight from $`\lambda `$4232, using a single Gaussian fit. But Stahl & Wilson (sta92 (1992)) used two Gaussian components to account for the broad wings detected by Crane et al. (cra91 (1991)), and thus obtained a tighter fit to the $`\zeta `$Oph CH<sup>+</sup> absorption, with $`R=71\pm 3`$ for $`\lambda `$4232 and $`R=68\pm 7`$ for $`\lambda `$3957.
We also find that using two instead of one Gaussian components significantly improves the fit, bringing reduced $`\chi ^2`$ values for the coadded $`\zeta `$Oph $`\lambda `$4232 spectrum from 2.07 for one Gaussian (using the noise reference of the two Gaussian fit), to 1.05 for two Gaussians. But the resulting average for $`R`$ is unchanged, we obtain $`R=78.22\pm 2.63(3.75)`$ with one Gaussian, $`R`$ = 80.20$`\pm 2.75(4.19)`$ with two Gaussians.
What is the effect of including more than two Gaussians? It is difficult to separate an error in the definition of the continuum baseline from a real increase in confidence level of the fit. We experimented on $`\zeta `$Oph with 3 and 4 Gaussians, keeping the noise fixed to the rms dispersion of the residuals in the 2 Gaussian fit. The reduced $`\chi ^2`$ were respectively 0.92 and 0.88 for the 3 and 4 Gaussian fits to the spectral region centred on $`\lambda `$4232 in the coadded spectrum of $`\zeta `$Oph. But some of the 4-Gaussians fits to the $`\zeta `$ Oph spectra contained very narrow features, corresponding to unphysical opacity components. Since we cannot separate the uncertainty in the baseline definition, we adopt the smallest number of Gaussians consistent with the dataset (with reduced $`\chi ^2`$ close to 1). We chose 3 Gaussians rather than 2 because Crawford et al. cra94 (1994) report a 3-Gaussians decomposition of the ultra-high resolution CH<sup>+</sup> profile toward $`\zeta `$Oph. The final value from the UVES data is thus $`R=78.47\pm 2.65(3.53)`$, with a 3-Gaussian fit.
All $`R`$ values for $`\zeta `$Oph are consistent within 2$`\sigma `$. The value we obtain here is higher than previously reported by Stahl et al. (sta92 (1992)) by 1.6 $`\sigma `$, adding the statistical errors in quadrature. For comparison Hawkins et al. (haw93 (1993)) give $`R=63\pm 8`$ for $`\lambda `$4232, $`67\pm 19`$ for $`\lambda `$3957, Crane et al. (cra91 (1991)) give $`R=67.6\pm 4.5`$ for $`\lambda `$4232, and Vanden Bout & Snell (van80 (1980)) report $`R=`$ 77$`{}_{12}{}^{}{}_{}{}^{+17}`$, also for $`\lambda `$4232.
The value reported in this article has an accuracy close to that of Stahl et al. (sta92 (1992)), of $`71\pm 3`$. Combining both measurement with identical weights gives our best value for $`\zeta `$Oph as $`R=74.7\pm 2.3`$.
In what follows we will nonetheless use the UVES value of $`R`$ for $`\zeta `$Oph to compare with the other lines of sight, $`R=78.47\pm 2.65(3.53)`$ thus relying exclusively on measurements obtained with the same instrument.
#### 4.1.2 HD110432
The profile from Crawford (cra95 (1995)) is well fit by 1 Gaussian only, but we obtain best results with 2 Gaussians.
The Hipparcos distance to HD110432 is 300$`\pm `$51 pc, which confirms that the CH<sup>+</sup> absorption most likely arises in the Coalsack, at a maximum distance of 180 pc (Franco fra89 (1989)).
The fit on the $`\lambda `$3957 line is excluded from the combined measurement of $`R`$ because the adjacent telluric absorption feature could affect the wings of the CH<sup>+</sup> absorption. Also the baseline around $`\lambda `$3957 seems to be strongly affected by systematic noise. The $`\lambda `$3957 values differ from $`\lambda `$4232, which otherwise give consistent numbers for each night.
#### 4.1.3 HD152235
HD152235 is a member of the Sco OB1 association, but as discussed by Crawford (cra95 (1995)) the absorbing CH<sup>+</sup> lies in the Lupus molecular cloud, at 170 pc (Murphy et al. mur86 (1986)).
The profile from Crawford (cra95 (1995)) is fit by 3 Gaussians. Our data, presented in Fig. 5, also required 3 Gaussians for $`\lambda `$3957.
The CH<sup>+</sup> absorption has a red tail than difficults separating both isotopes at 4232$`\AA `$, while the telluric absorption feature borders the <sup>13</sup>CH<sup>+</sup> line. It is difficult to pin down the underlying continuum at 4232$`\AA `$. We tried two approaches to define the continuum from edge to edge of the CH<sup>+</sup> absorption (from 4232$`\AA `$ to 4232.7$`\AA `$): a free fit (i.e. roughly a straight-line, with zero weights over 4232$`\AA `$ – 4232.7$`\AA `$), and the inclusion of a small non-zero weight mid-way between the two isotopes (see Fig. 5).
The spectrum from 2003 June 14th is clear of telluric absorption (it was observed with the smallest air mass, see Table 1), so both overtones should give consistent values for $`R`$. This is why we chose the baseline shown on Fig. 5. The two overtones give consistently high value for June 14th, and also on average. So we are rather confident of the ratio we report. The fits on $`\lambda `$3957 for June 15<sup>th</sup> and 16<sup>th</sup> are excluded from the combined measurement of $`R`$ because they are affected by telluric absorption under the red tail of <sup>13</sup>CH<sup>+</sup>.
Vladilo et al. (vla93 (1993)) reported $`R=126\pm 29`$. Our measurement, which has better accuracy, gives lower value, although still higher than the ISM average (see Sec. 5).
This line of sight gives a significantly higher $`R`$ value than in HD110432 at 4.2 $`\sigma `$, using the conservative uncertainties $`\sigma _2`$, from the scatter of each individual measurement of $`R`$.
#### 4.1.4 HD152236
Fig. 6 presents the CH<sup>+</sup> spectra for HD152236, along with our fits. Crawford (cra95 (1995)) used 4 Gaussians, but we needed 6.
HD152236, as HD152235, is a member of the Sco OB1 association, so that the absorption probably occurs in the Lupus cloud.
There is no manifest absorption at $`\lambda `$3957, but we give the result of the formal fit, which allows assigning a lower limit $`R`$ value. We take $`f_{\mathrm{uplimit}}=f+3\sigma _1(f)`$, and $`R>1/f_{\mathrm{uplimit}}`$, giving $`R>61.4`$ at the 3 $`\sigma `$ level.
Although the $`\lambda `$4232 line has better signal to noise, both isotopes are blended. We tried to apply the technique described below in the case of HD154368 (Sec. 4.1.6), but did not obtain significant excess absorption in the blue wing of the main line. The lower limit refers only to the $`\lambda `$3957 measurement.
#### 4.1.5 HD152424
HD152424, as HD152235, is a member of the Sco OB1 association, so that the absorption probably occurs in the Lupus cloud.
Vladilo et al. vla93 (1993) report $`R=98\pm 19`$, but this target gives us our lowest $`R`$ value, of $`55.65\pm 3.55(15.61)`$.
The two isotopes are blended at $`\lambda `$4232, so we apply the same technique as for HD154368 (Sec. 4.1.6). The residuals under the <sup>13</sup>CH$`+`$ absorption are consistent with the noise, which gives us confidence in the inferred $`R`$ value. HD152424 provides our best measurement of $`R`$ for blended lines, as obtained by keeping the opacity profile fixed to the best fit on $`\lambda `$3957. But defining consistent baselines between the two overtones is very difficult without including the baseline as free parameter. So we decided to ignore the significant residuals under the main isotope at $`\lambda `$4232. We include the noisy residuals in the noise estimate, so that the significance of the $`R`$ values derived from the blended line is lowered.
The importance of quantifying systematic uncertainties is manifest in HD152424: a visual inspection of Fig. 7 shows the fit under $`\lambda `$3957 is very good, and that the baseline seems smooth. However appeareances can be deceiving: the extrapolation of the $`\lambda `$3957 fit to $`\lambda `$4232 gives a very different value of $`R`$. Had we considered only the $`\sigma _1`$ uncertainties, the value for HD152424 would have been one of the most accurate. But because CH<sup>+</sup> towards HD152424 is so broad the conservative uncertainty $`\sigma _2`$ is much higher than $`\sigma _1`$ and reflects the systematic uncertainty involved in the baseline definition.
#### 4.1.6 HD154368
This target was selected from the survey of Gredel et al. (gre93 (1993)).
The CH<sup>+</sup> absoprtion for HD154368 (with an Hipparcos distance of $`364.8\pm 128.3`$ pc) shown on Fig. 8 has four conspicuously distinct velocity components. The two isotopes are manifestly blended for the $`\lambda `$4232 transition.
Our data cover the two overtones in the same UVES dichroic setting simultaneously, and the spectral resolution is expected to be constant with wavelength in such an instrumental setup. After checking we had no detectable trend of varying resolution by measuring the arc line widths, we attempted fitting $`\lambda `$4232 using the results of the $`\lambda `$3957 fits.
A fit of both overtones simultaneously would be counterproductive because the $`\lambda `$4232 region has better sensitivity, and would dominate the fit, placing Gaussian components of <sup>12</sup>CH<sup>+</sup> under the rarer isotope. Instead we kept the <sup>12</sup>CH<sup>+</sup> opacity profile as a function of velocity $`\tau (v)`$ as inferred from the $`\lambda `$3957 fit. We then scaled $`\tau (v)`$ to $`\tau (\lambda )`$ for the $`\lambda `$4232 region, which implies scaling component widths and separations in wavelengths. After setting baselines as in the unblended case, the optimization involved two free parameters, the scaling factors on the <sup>12</sup>CH<sup>+</sup> and <sup>13</sup>CH<sup>+</sup> opacity profiles.
The fits shown on Fig. 8 are rather poor in the case of $`\lambda `$4232. Uncertainties in the baseline determination resulted in significant residuals under the main isotope line. But these do not affect the <sup>13</sup>CH<sup>+</sup> absorption. We recompute the noise level from the residuals, including the baseline uncertainties. The uncertainties on $`R`$ tabulated in Table 3 are a result of this exaggerated noise level.
#### 4.1.7 HD157038
This line of sight also has blended absorption at $`\lambda `$4232, and we apply the same technique as for HD154368. But the results on $`\lambda `$4232 are useless because they are too sensitive on the baseline definition, so we avoided the use of $`\lambda `$4232 altogether.
Hawkins & Meyer (haw89 (1989), and references therein) give $`R=38\pm 12`$ from $`\lambda `$4232 and use a distance of 1.7 kpc to HD157038. We confirm this line of sight is enriched in <sup>13</sup>C, with a lower $`R`$ value than the average ISM value (see Section 5.1).
#### 4.1.8 HD161056
<sup>13</sup>CH<sup>+</sup> absorption towards HD161056 (selected from Gredel et al. gre93 (1993)) is very weak at $`\lambda `$3957, and yet conspicuous at $`\lambda `$4232 (see Fig. 10). We first attempted to fit the two overtones separately, and found widely discrepant values: $`f(3957)=68.9\pm 11.8(12.0)`$, and $`f(4232)=159.3\pm 5.6(19.1)`$. The two isotopes may be blended at $`\lambda `$4232. We thus applied the same technique as for HD154368, extrapolating the $`\lambda `$3957 opacity profile to $`\lambda `$4232. But the inferred isotope ratios were still very different (see Table 3). There must be additional absorption bridging the two overtones at $`\lambda `$4232, not due to CH<sup>+</sup>: the red tail of $`\lambda `$4232 is also seen at $`\lambda `$3957, while the absorption at the blue edge of $`\lambda `$4232 is not seen at $`\lambda `$3957. We cannot identify the nature of the additional absorption, so we choose to exclude $`\lambda `$4232 from the combined measurement of $`R`$ in HD161056.
The hipparcos distance to HD161056 is $`425.7\pm 145.5`$ pc.
#### 4.1.9 HD169454
This line of sight, selected from Gredel et al. (gre93 (1993)), exibits at least 4 distinct velocity components. We use 6 Gaussians to describe the profile, but did not attempt to compare reduced $`\chi ^2`$ when varying the number of Gaussians because of the uncertainties in the baseline definition.
The total number of free-parameters used to describe the main line is 6$`\times `$3, with 3 parameters per Gaussian. The zero-absorption line width for HD169454 is $``$0.8 Å, or 80 spectral datapoints. But the number of independent resolution elements under <sup>12</sup>CH<sup>+</sup> is about 16. Thus 18 free parameters may seem slightly excessive. But over-constraining the opacity profile is of no consequence to the derived isotopic ratio.
We use the same technique as for HD154368 to model the blended <sup>12</sup>CH<sup>+</sup> and <sup>13</sup>CH<sup>+</sup> absorption at $`\lambda `$4232, but do not obtain satisfactory results. The fits on $`\lambda `$4232 are shown on Fig. 11 for completeness, and are not included in the combined measurement of $`R`$ towards HD169454. An unphysical value of $`f`$ is obtained, but this value is still consistent with zero <sup>13</sup>CH<sup>+</sup> within the uncertainties. Also excluded from the combined measurement of $`R`$ is the fit to $`\lambda `$3957 from June 14<sup>th</sup> because it shows non-ISM absorption at 3957.3 $`\AA `$. The $`\lambda `$3957.3 feature is not seen in other spectra acquired at similar airmass (e.g. that of HD152235 on June 15<sup>th</sup>), and is either an instrumental artifact or the effect of passing clouds.
#### 4.1.10 HD170740
This line of sight, Selected from Gredel et al. (gre93 (1993)), is important since it has the highest $`R`$ value among those considered in this work, together with HD161056. The Hipparcos distance to HD170740 is $`211.9\pm 42.4`$ pc.
Unfortunately instrumental fringing, especially at $`\lambda `$3957, limits the accuracy of the fits. The $`\lambda `$3957 <sup>13</sup>CH<sup>+</sup> line is only marginally detected in individual nights, but we give the formal fits. The value we report relies on the nigthly spectra for $`\lambda `$4232, and on the coadded spectrum for $`\lambda `$3957.
## 5 Discussion
### 5.1 Observed scatter and comparison with other sightlines
Figure 13 summarises our results. The isotopic ratio $`R`$ is lower for HD110432 than for HD152235 by 4.3 $`\sigma `$, and lower than for HD170740 by 3.2 $`\sigma `$. In this comparison we have used the conservative uncertainties $`\sigma _2`$, derived from the scatter of $`R`$ for different nights and different overtones.
The weighted average of $`R`$ for the 9 lines of sight considered in this work with <sup>13</sup>CH<sup>+</sup> detections is $`R_{\mathrm{ISM}}=^{12}`$C/<sup>13</sup>C = $`78.27\pm 1.83`$. The weighted 1 $`\sigma `$ dispersion of the data is $`12.7`$. The hypothesis that the observed scatter is derived from a single value of $`R`$ can be tested by calculating the corresponding value of $`\chi ^2/\nu `$, where $`\nu =8`$ is the number of the degrees of freedom (the number of line of sights less one free parameter). We obtain $`\chi ^2/\nu =6.0`$, which discards the hypothesis at a confidence level of essentially 1.
The observed scatter is not due to the galactocentric abundance gradient studied by Hawkins & Meyer (haw89 (1989)), and predicted by Galactic chemical evolution models. The targets in our sample all lie within 30 deg of $`l=0`$, and at distances of at most $``$200 pc (except for the clouds toward HD157038 and HD169454, whose distances are unknown). Fig. 13 also gives longitude and distance information for each target.
Our value for $`R_{\mathrm{ISM}}`$ is significantly higher than those for the Pleiades, $`\xi `$Per and P Cyg. Vanden Bout & Snell (van80 (1980)) measured 49$`{}_{8}{}^{}{}_{}{}^{+12}`$ and 59$`{}_{13}{}^{}{}_{}{}^{+24}`$ for 20 Tau (Maia in the Pleiades) and $`\xi `$Per. Hawkins & Jura (haw87 (1987)) report $`R`$ values of 46$`\pm `$15 for $`\xi `$Per, 40$`\pm `$9 for 20 Tau, 41$`\pm `$9 for 23 Tau (also in the Pleiades), and 45$`\pm `$10 for P Cyg. Hawkins et al. (haw93 (1993)) confirm a low value of 49$`\pm `$15 for $`\xi `$Per. Even taking 3 $`\sigma `$ errors, the measurements by Hawkins et al. are below our results. One can ask whether the <sup>13</sup>C enrichment in the Pleiades is due to contamination from past stellar winds in the cluster itself.
Aside from the $`R`$ values obtained by other groups and quoted so far in this work, there are also those from Centurion & Vladilo (cen91 (1991)) for $`\mu `$Nor (67$`\pm `$6) and HD26676 (64$`\pm `$6), which differ at 2 $`\sigma `$ with our average ISM value (including its scatter).
### 5.2 Are there D/H variations in the local ISM?
From our data set, the values of $`R`$ vary by 7 $`\sigma `$ in the local translucent clouds, where $`\sigma `$ is the uncertainty in the weighted average value $`R_{\mathrm{ISM}}`$. This is strong indication that the ISM at the same galactocentric distance is not completely mixed. If this is the case in the chemical composition of the ISM in general, then one expects that the D/H ratios in the local ISM will also differ from place to place. Moos et al. (moo02 (2002)) (see also Vidal-Madjar vid02 (2002)) summarise the FUSE results towards nearby white dwarfs, located within 100 pc of the Sun. The D i/H i ratios measured with interstellar D i and H i Lyman absorption varies by less than $``$10% towards 7 white dwarfs within 100 pc of the Sun, at the limits of the Local Bubble. The fractional standard deviation of D i/H i (the ratio of the scatter in the data to its weighted average) is 12%, but cannot be distinguished from a single D i/H i value within the uncertainties. By contrast, the D i/H i ratios obtained with IMAPS show a factor of 2 scatter towards $`\gamma ^2`$ Vel, $`\zeta `$ Pup and $`\delta `$ Ori, beyond the Local Bubble, at distances of 300–500 pc (Jenkins et al. jen99 (1999), Sonneborn et al. son00 (2000)).
Is the scatter in the observed interstellar D i/H i due to variations in D/H? The primary objective of the IMAPS D i/H i observations, as stated by Sonneborn et al. (son00 (2000)), is “to determine D/H with sufficient accuracy to test for spatial inhomogeneity”. But the D/H ratio is notoriously difficult to measure, and often involves the use of several different instruments. The H i lines are $`>100`$ km s<sup>-1</sup> broad and saturated so that N(H i) is measured from the Ly $`\alpha `$ damping wings, while N(D i) is best measured from Ly $`\delta `$ or Ly $`ϵ`$ in the case of the IMAPS targets. The D/H ratio involves carefully considering the effect of the underlying stellar continua and the contribution from other interstellar H i components in the rising part of the curve of growth (with a high ratio of equivalent width to opacity compared to the main absorption line). Chemical fractionation and selective dissociation of HD affects D i/H i for lines of sight where H<sub>2</sub> is observed, and differential acceleration by radiation pressure could displace D i relative to H i in the diffuse ISM. Bruston et al. (bru81 (1981)) assign the D i/H i scatter known at the time to such effects, and estimate a value of unperturbed D/H almost a factor of two higher than the observed average D/H of $`\mathrm{1.5\; 10}^5`$.
In addition Draine (dra04 (2004)) argues for variable depletion of D onto dust grains, such that the gas phase D/H would be a function of dust processing, which is in turn a function of Galactic environment. The variations in D i/H i do not seem to be a faithful measure of D/H.
### 5.3 Inefficient mixing in the local ISM
What is the level of elemental homogeneity in the ISM? What is the power-spectrum of the dispersion in elemental abundance relative to the total metalicity? Does heterogeneity increase with size?
At least one extreme of spatial scale indicates perfect mixing: the meteoritic evidence favours mixing at the molecular level of interstellar dust from different origin, during the formation of the solar system. Zinner et al. (zin91 (1991), their Table 1) measured $`{}_{}{}^{1}2`$C/<sup>13</sup>C values from 3 to 1000 in different SiC grains found in the Murchison meteorite. At some point in the presolar cloud the SiC grains must have been mixed from a previously heterogenous distribution.
In a recent review of interstellar turbulence and mixing, Scalo & Elmegreen (sca04 (2004)) summarise current knowledge of the level of elemental heterogeneity in the ISM. The stellar return to the ISM should be spotty and poorly mixed. Scalo & Elmegreen explain that in contrast with diffusive processes, turbulence transport does not homogenize the gas at the atomic level. Yet the heterogeneity is notoriously difficult to detect. The available gas-phase diagnostics depend on local excitation, while the stellar data give upper limits only on the dispersion of elemental composition.
But the <sup>12</sup>CH<sup>+</sup>/<sup>13</sup>CH<sup>+</sup> ratio, in contrast to D i/H i, is a particularly sensitive probe of the mixing efficiency of the ISM. In terms of the fractional standard deviation used by Moos et al. (moo02 (2002)), the scatter we have measured is $`\sigma _{\mathrm{ISM}}/R_{\mathrm{ISM}}=`$16.2% on a scale of 100 pc. Considering the precise location of the absorbing material is unknown, we take the Moos et al. (moo02 (2002)) measurement of a 10% scatter of D i/H i from the mean value as a indication of inhomogeneity of the local ISM, over 100 pc scales. This value is similar to the variations reported in this work, even though D is burned in stellar interiors, rather than produced as is the case of <sup>13</sup>C. Although their absolute yields are different, the mixing process (i.e. passive scalar turbulence, Scalo & Elmegreen sca04 (2004)) should be similar for both species.
### 5.4 Is the solar value for $`R`$ comparable to ISM measurements?
The terrestrial value of $`R`$ is 90 (Rosman & Taylor ros98 (1998)), and is usually extrapolated to the Sun. However, the preliminary results of Ayres et al. (ayr05 (2005)) indicate a photospheric carbon isotopic ratio of 70. The average value of our measurements of $`R`$ in the local ISM, $`78.27\pm 1.83`$, is slightly lower than the solar value of 90 (or even higher than the photospheric value of 70). But in fact $`R_{\mathrm{ISM}}`$ is indistinguishable from the solar carbon isotope ratio given the observed scatter of $`12.7`$.
Models of galactic chemical evolution predict an enhancement of <sup>13</sup>C relative to <sup>12</sup>C with time. Thus our result that <sup>13</sup>C/<sup>12</sup>C has essentially remained constant over the past 4.5 Gyr is surprising. The inconsistency is worsened when comparing with the preliminary results of Ayres et al. (ayr05 (2005)). We may be affected by a sampling bias. Around 30 lines of sight or more are required to faithfully estimate $`R_{\mathrm{ISM}}`$.
But is the solar value truly representative of the ISM at the time of collapse of the solar nebula? The process of CO fractionation followed by condensation on cold dust grains may have increased <sup>13</sup>C in the pre-solar dust. With an enhanced dust-to-gas ratio through sedimentation in the accretion disk of the Sun, the result may have been a modification in the final $`R`$ value.
We imagine a <sup>13</sup>C-rich disk with a dust mass of 0.01 $`M_{}`$, dust to gas ratio of unity, accreting on a zero-age Sun with metalicity Z=0.01. The <sup>13</sup>C-poor atmosphere of the disk may have been blown away by the early solar wind, so that once the accretion of the disk is concluded, the Sun reaches $`Z=0.02`$ and higher <sup>13</sup>C/<sup>12</sup>C than the pre-solar cloud.
Aside from the possibility of preferential accretion of <sup>13</sup>C over <sup>12</sup>C, there is evidence that the solar nebula differed in composition from the bulk ISM. For instance the metalicity of B stars in Orion is lower than that of the Sun, even after 4.5 Gyr of Galactic chemical evolution, indicating that the Sun must have been exceptionally metal-rich. This discrepancy has been interpreted in the framework of Galactic diffusion (Wielen wie77 (1977)), by which the Sun diffused from its birthplace outwards in Galactocentric radius by $`\mathrm{\Delta }r_{\mathrm{GAL}}=1.9`$kpc. Since <sup>13</sup>C/<sup>12</sup>C is also predicted to increase with $`r_{\mathrm{GAL}}`$, Galactic diffusion may bring in agreement the ISM and solar values of $`R`$ at the time of solar birth (Wielen & Wilson wie97 (1997)).
But perhaps the most compelling evidence that isotopic ratios in the Sun cannot be taken as direct constraint on the bulk ISM are isotopic anomalies in meteorites, such as the presence of extinct radionuclides. For instance Zinner et al. (zin91 (1991)) found that <sup>26</sup>Al/<sup>27</sup>Al$``$1 in the pre-solar nebula. They measured the concentration of the decay product of <sup>26</sup>Al, <sup>26</sup>Mg, locked in SiC grains from the Murchison meteorite. But in other meteorites <sup>26</sup>Al/<sup>27</sup>Al$``$1/4000 in Al-rich minerals from refractory inclusions not representative of the pre-solar composition. Since the half-life of <sup>26</sup>Al is $`10^6`$yr, the pre-solar enrichment in <sup>26</sup>Al suggests pollution by a nearby source of <sup>26</sup>Al, such as a supernova explosion (see Zinner et al. zin91 (1991) and Clayton cla94 (1994) for thorough discussions).
## 6 Conclusions
The VLT-UVES spectra of CH<sup>+</sup> absorption in the lines of sight toward 10 bright stars have allowed us to measure the carbon isotopic ratio $`R=^{12}`$C/<sup>13</sup>C with unprecedented accuracy. We confirm previous measurements of $`R`$ toward $`\zeta `$Oph, and obtain significant scatter in the local ISM.
Our value for $`\zeta `$Oph is $`R=78.47\pm 2.65(3.53)`$, where the second uncertainty in parenthesis is that obtained a-posteriori, from independent measurements, and independent baseline and absorption fits. Combining this value with that of Stahl et al. (sta92 (1992)) gives our best value for $`\zeta `$Oph, $`R=74.7\pm 2.3`$.
Averaging our measurements for the 10 lines of sight gives a value representative of the local ISM: $`R_{\mathrm{ISM}}=`$$`78.27\pm 1.83`$, with a weighted rms dispersion of $`12.7`$. The dispersion in isotopic ratio is 7 times the uncertainty in $`R_{\mathrm{ISM}}`$ \- we detect heterogeneity at 7 $`\sigma `$.
The observed scatter in <sup>13</sup>C/<sup>12</sup>C is the first significant detection of heterogeneity in the isotopic composition of the local ISM, and was obtained with a single instrument and a homogeneous analysis, with a diagnostic independent of the local excitation and of chemical fractionation. The observed variations in $`R`$ can be extrapolated to an overall elemental heterogeneity, and reveal that mixing in the ISM is not perfect.
The solar carbon isotope ratio of $`90`$ is undistinguishable from the present-day ISM ratio, considering its intrinsic scatter, even after 4.5 Gyr of galactic evolution. This statement is not affected by instrumental uncertainties, although a sampling bias could distort our ISM value. We question whether the solar carbon isotopic ratio is equal to that of the presolar nebula.
A larger sample of stars, and new data towards the Pleiades, the Taurus cloud and P Cyg with 8 m aperture telescopes, will allow a better estimate of the average ISM value, and an improved understanding of the relationship between elemental heterogeneity and spatial scale in the ISM.
###### Acknowledgements.
We thank the referee, Roland Gredel, for critical readings and constructive comments that improved the article. S.C acknowledges support from Fondecyt grant 1030805, and from the Chilean Center for Astrophysics FONDAP 15010003.
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# Persistent current and Drude weight in mesoscopic rings
## 1 INTRODUCTION
The experimental discovery of persistent currents in mesoscopic rings pierced by a magnetic flux,<sup>1-3</sup> earlier proposed theoretically,<sup>4</sup> has revealed interesting new effects. The currents measured in metallic and semiconducting rings, either in a single ring or an array of many rings, generally exhibit an unexpectedly large amplitude, i.e., larger by at least one order of magnitude, than predicted by theoretical studies of electron models with either disorder or electron-electron interaction treated perturbatively.<sup>5,6</sup> It has been suggested that the interactions and their interplay with disorder are possibly responsible for the large currents observed, expecting that the effect of the interactions could counteract the disorder effect. However, no consensus has yet been reached on the role of the interactions. In order to gain theoretical insight, it is desirable to perform numerical calculations which allow to consider both interactions and disorder directly in systems with sizes varying from small to large. Analytical calculations usually involve approximations which mainly provide the leading behavior of the properties for large system sizes. Persistent currents in mesoscopic rings strongly depend on the system size, since they emerge from the coherence of the electrons across the entire system. Hence, it is most important to study the size dependence of the current beyond leading order in microscopic models, for a complete understanding of the experimental results. Exact diagonalization was used to calculate persistent currents in systems with very few lattice sites.<sup>7,8</sup> In this work, we use the Density Matrix Renormalization Group (DMRG) algorithm,<sup>9-11</sup> to study a simplified model incorporating interactions and a single impurity, accounting for disorder, in larger system sizes. We consider a system of interacting spinless electrons on a one-dimensional ring, with a single impurity, and penetrated by a magnetic field. We study an intermediate range of system sizes, where analytical results obtained by bosonization techniques for large system sizes, do not yet fully apply. Without impurity, and at half-filling, the system undergoes a metal-insulator transition from a Luttinger Liquid (LL)<sup>12</sup> to a Charge Density Wave (CDW)<sup>13</sup> groundstate. The persistent current of the interacting system with an impurity was studied before with the DMRG, in the LL phase.<sup>14</sup> Here we study the persistent current, and also the Drude weight characterizing the conducting properties of the system, in both the LL and the CDW phase, investigating the interplay between the impurity and the interactions in the two phases. In mesoscopic systems the separation between metallic and insulating behavior is not always obvious, since the localization length can be of the order or significantly larger than the system size. Hence, a finite Drude weight and a current can be observed in the CDW phase of a mesoscopic system. It is therefore of great interest to characterize the persistent current and the Drude weight in both the LL and the CDW phases of mesoscopic systems. Although the simple model that we consider is not the most appropriate to describe the experimental situation, we hope to obtain useful information for the understanding of the more realistic systems. Under a Jordan-Wigner transformation,<sup>15</sup> the system considered is equivalent to a spin-1/2 XXZ chain with a weakened exchange coupling. Hence, our results also provide insight into the spin transport in this type of systems.
## 2 THE MODEL
The Hamiltonian describing a system of spinless fermions on a ring pierced by a magnetic flux, with repulsive interactions and a single hopping impurity, or defect, is given by,
$$H=H_t+H_U,$$
(1)
where
$$H_t=t\underset{j=1}{\overset{N}{}}\left(e^{i\varphi /N}c_j^{}c_{j+1}+e^{i\varphi /N}c_{j+1}^{}c_j\right)+(1\rho )t\left(e^{i\varphi /N}c_N^{}c_1+e^{i\varphi /N}c_1^{}c_N\right)$$
(2)
is the hopping term, $`\varphi =2\pi \mathrm{\Phi }/\mathrm{\Phi }_o`$ contains the magnetic flux $`\mathrm{\Phi }`$ in units of the flux quantum $`\mathrm{\Phi }_o=hc/e`$, $`\rho `$ measures the strength of the defect with values between $`0`$ and $`1`$, ($`\rho =1`$ corresponding to the defectless case), and
$$H_U=U\underset{j=1}{\overset{N}{}}n_jn_{j+1}$$
(3)
is the interaction term, with $`U>0`$ representing the nearest neighbor Coulomb repulsion, and $`n_j=c_j^{}c_j`$, where $`c_j^{}`$ and $`c_j`$ are the spinless fermion operators acting on the site $`j`$ of the ring. We consider a system of $`N`$ sites, with $`N`$ even, and at half-filling, when $`M=N/2`$ particles are present. The lattice constant is set to one and periodic boundary conditions, $`c_{N+1}=c_1`$, are used.
Via the gauge transformation $`c_je^{i\varphi j/N}c_j`$, the flux can be removed from the Hamiltonian, but in the impurity term where the flux is trapped, and the quantum phase $`\varphi `$ is encoded in a twisted boundary condition $`c_{N+1}=e^{i\varphi }c_1`$. It is then clear that the energy is periodic in $`\varphi `$ with period $`2\pi `$, i.e., it is periodic in the flux $`\mathrm{\Phi }`$ threading the ring with period $`\mathrm{\Phi }_0`$.<sup>16</sup>
After a Jordan-Wigner transformation, Eqs. (2) and (3) can be rewritten, respectively, as
$$H_J=\frac{J}{2}\underset{j=1}{\overset{N}{}}\left(S_j^+S_{j+1}^{}+S_j^{}S_{j+1}^+\right)+(1\rho )\frac{J}{2}\left(e^{i\varphi }S_1^+S_N^{}+e^{i\varphi }S_1^{}S_N^+\right)$$
(4)
and
$$H_\mathrm{\Delta }=\mathrm{\Delta }\underset{j=1}{\overset{N}{}}S_j^zS_{j+1}^z$$
(5)
with $`t=J/2`$ and $`U=\mathrm{\Delta }`$, and the boundary conditions $`S_{N+1}^+=(1)^{M+1}e^{i\varphi }S_1^+`$ and $`S_{N+1}^z=S_1^z`$. Hence, the model (1) of spinless fermions is equivalent to a spin-1/2 XXZ chain with a weakened exchange coupling, and twisted boundary conditions in the transverse direction. The half-filled case corresponds to total spin projection $`S^z=0`$.
The persistent current generated on a ring pierced by a magnetic flux, at temperature $`T=0`$, can be obtained from the ground state energy $`E_o(\varphi )`$, by taking the derivative with respect to $`\varphi `$,
$$I(\varphi )=\frac{E_o(\varphi )}{\varphi }.$$
(6)
For the spinless fermion system, Eqs. (2) and (3), $`I(\varphi )`$ corresponds to the ground state value of the charge current operator $`\widehat{I}_c=it_{j=1}^N\left(c_j^{}c_{j+1}c_{j+1}^{}c_j\right)`$, while for the XXZ chain, Eqs. (4) and (5), it corresponds to the ground state value of the spin current operator $`\widehat{I}_s=i\frac{J}{2}_{j=1}^N\left(S_j^+S_{j+1}^{}S_{j+1}^+S_j^{}\right)`$. As a consequence of the periodicity of the energy, the current is also periodic in $`\varphi `$, with period $`2\pi `$. Hence, it can be expressed as a Fourier series,
$$I(\varphi )=\underset{k=1}{\overset{\mathrm{}}{}}I_k\mathrm{sin}(k\varphi ).$$
(7)
and the behavior of the current can be analyzed in terms of the coefficients $`I_k`$.<sup>14</sup>
In the noninteracting case ($`U=0`$), it has been found that for large system sizes, the current is invariant under the defect transformation $`\rho 1/\rho `$,<sup>17,18</sup> i.e.,
$$I(\rho )=I(1/\rho ).$$
(8)
We shall investigate the existence of this kind of invariance in the interacting case ($`U>0`$), both in the LL and the CDW phases.
The Drude weight was proposed by Kohn as a relevant quantity to distinguish between a metal and an insulator.<sup>19</sup> It is defined as
$$D=N\frac{^2E_o}{\varphi ^2}|{}_{\varphi =\varphi _m}{}^{},$$
(9)
where $`\varphi _m`$ is the location of the minimum of $`E_o(\varphi )`$, which depends on the parity of the number of electrons, i.e., $`\varphi _m=0`$ or $`\pi `$ for, respectively, an odd or an even number of electrons. For the spinless fermion system $`D`$ represents the charge-stiffness and measures the inverse of the effective mass of the charge carriers.<sup>20</sup> In a metallic conductor $`D`$ tends to a finite value whereas in an insulator $`D`$ vanishes with the system size $`N`$, when $`N\mathrm{}`$. In the insulating state the Drude weight decays as $`D\mathrm{exp}(N/\xi )`$, where $`\xi `$ measures the localization length. For the XXZ chain the Drude weight represents the spin-stiffness.
In a model of free fermions ($`U=0`$) with no impurity ($`\rho =1`$), it is straightforward to see that the leading behavior of the persistent current $`I(\varphi )`$ in the system size $`N`$, has a saw-tooth like shape with slope $`v_F/(\pi N)`$, where $`v_F`$ is the Fermi velocity. Thus, the amplitude of the current scales with $`1/N`$, vanishing in the limit $`N\mathrm{}`$. The discontinuity in $`I(\varphi )`$, that results from a degeneracy of energy levels associated to the translation symmetry,<sup>8</sup> appears at $`\varphi =0`$ or $`\pi `$ for, respectively, an even or an odd number of electrons. In the presence of an impurity ($`\rho 1`$), bosonization<sup>21</sup> and conformal field theory<sup>17</sup> calculations predict that the shape of the current $`I(\varphi )`$ is rounded off, and its amplitude decreases with increasing strength of the scatterer potential, still vanishing with the system size as $`1/N`$. The impurity lifts the degeneracy of the energy levels and the current then varies continuously.<sup>8</sup>
The model with interactions ($`U0`$) and without defect ($`\rho =1`$), is solvable by the Bethe ansatz for periodic boundary conditions ($`\varphi =0`$)<sup>22,23</sup> and also for twisted boundary conditions ($`\varphi 0`$).<sup>24-26</sup> At half-filling, the system exhibits a metal-insulator transition, which occurs at $`U/t=2`$. For $`U/t<2`$, the system is in a gapless LL phase, while for $`U/t>2`$ it is in a gapped CDW state.<sup>27</sup> The LL phase is characterized by a power-law decay of the correlations. Bosonization predicts that in an homogeneous LL, the leading behavior of the persistent current in the system size $`N`$, has a saw tooth like shape with slope $`v_J/(\pi N)`$, where $`v_J`$ is the velocity of current excitations.<sup>28</sup> Since translation invariance is preserved in the presence of interactions, the discontinuity in the current still exists for finite $`U`$.<sup>8</sup> A Bethe ansatz calculation shows that the Drude weight of an homogeneous LL in the thermodynamic limit, has a finite value, which decreases with increasing strength of the interaction $`U/t`$.<sup>20</sup> The LL state is strongly affected by the presence of an impurity,<sup>29-32</sup> and bosonization yields that the current then vanishes as $`IN^{1\alpha _B}`$, with $`\alpha _B>0`$.<sup>21</sup> The study of the LL phase with $`\rho 1`$, performed with the DMRG,<sup>14</sup> has in fact found this kind of behavior. The CDW phase is characterized by a localization length $`\xi `$, which is associated to the energy gap. From the work of Baxter,<sup>33</sup> the Drude weight in the gapped phase, is expected to behave as $`D\mathrm{exp}(N/\xi )`$, vanishing for an infinite system size. This behavior implies that although the system is insulating in the infinite system size limit, for a finite system, provided $`N/\xi `$ is not too large, $`D`$ is still finite and a current can be observed. The localization length can then be extracted from the size dependence of the Drude weight.
## 3 NUMERICAL METHOD
We use the DMRG to numerically calculate the groundstate energy of the spinless fermion system as a function of the magnetic flux, $`E_o(\varphi )`$, for fixed interaction $`U`$ and impurity strength $`\rho `$, in rings up to $`N=82`$ sites, keeping up to $`250`$ density matrix eigenstates per block.<sup>34</sup> The DMRG is applied to the Hamiltonian (1) after performing the gauge transformation, which removes the flux into a twisted boundary condition.<sup>11</sup> The states of the system are characterized by the quantum numbers associated to the eigenvalues of the local occupation number $`n_j`$ and the total number of particles $`M=_{j=1}^Nn_j`$ operators, which commute with the Hamiltonian (1).<sup>35</sup> For each set of $`N`$, $`\rho `$ and $`U`$, we obtained the groundstate energy $`E_o`$ for $`50`$ values of $`\varphi `$ in the periodicity interval $`\pi <\varphi \pi `$, and using Chebyshev interpolation<sup>36</sup> we determined the corresponding current (6) and Drude weight (9), by numerical differentiation. We developed a DMRG algorithm for complex Hamiltonian matrices, which allowed to calculate the detailed form of the persistent current $`I`$ as a function of $`\varphi `$,<sup>14</sup> and to obtain the Drude weight. In a previous approach the DMRG was used to calculate the so called phase sensitivity $`\mathrm{\Delta }E_0`$, which is the difference of the groundstate energy at flux $`\varphi =0`$ and $`\pi `$, and can be considered a crude measure of the persistent current.<sup>37-39</sup> Although the calculation of $`\mathrm{\Delta }E_0`$ requires considerably less computational effort than the calculation of $`E_0(\varphi )`$, because then the Hamiltonian matrix is real, the phase sensitivity does not provide information on the shape of the current and the value of the Drude weight.
## 4 RESULTS
We now present the results obtained for the persistent current and the Drude weight, where we take $`t=1`$ and the interaction $`U`$ is in units of $`t`$. Figs. 1 and 2, exhibit $`NI(\varphi )`$ plotted versus $`\varphi `$, for respectively, $`U=0.80`$ and $`U=3.00`$, which correspond, respectively, to the LL and the CDW phase, considering different impurity strengths $`\rho `$, on a fixed system size $`N=26`$. We can see that the effect of the impurities, in both phases, is to reduce the intensity of the current, and to round off the shape of $`I(\varphi )`$. The amplitude of the current decreases rapidly with increasing values of $`|\rho 1|`$, and also with increasing strength of the interaction $`U`$. Fig. 3 shows that the invariance of the current with respect to the defect, described in Eq. (8), found for large system sizes in the noninteracting case, is also observed for the interacting case, both in the LL and in the CDW phases, the system size $`N`$ required to reach the invariance being larger for larger interaction $`U`$. Fig. 4 displays the Drude weights associated to the systems with different interactions $`U`$ and impurities $`\rho `$, fixed $`N=26`$, of the currents presented in Figs. 1 and 2. As one would expect the Drude weight decreases with increasing $`|\rho 1|`$, and also increasing $`U`$. Figs. 5 and 6, present $`NI(\varphi )`$ plotted for several system sizes $`N`$, respectively, for $`U=0.80`$, in the LL phase, and $`U=3.00`$, in the CDW phase, with $`\rho `$ fixed at $`0.50`$. We observe that the current vanishes faster than $`1/N`$ in both phases, exhibiting a different behavior in each phase, $`I(\varphi )`$ vanishing much faster with $`N`$ in the CDW than in the LL phase. In order to analyze the behavior of the current in more detail we have numerically evaluated the coefficients $`I_k`$ (for $`k=1,2`$) of the Fourier expansion (7). The first ($`k=1`$) and the second ($`k=2`$) Fourier coefficients of the current, for $`U=0.80`$ and $`U=3.00`$, are shown in Fig. 7. One can clearly see that the coefficients $`I_1`$ and $`I_2`$ behave similar to each other in both phases. However, their behavior in the LL and CDW phase is distinct. In the LL phase, the Fourier coefficients show a power-law decay with $`N`$, with the second order coefficient decaying faster, i.e., with a larger exponent, than the first one. In the CDW phase, the Fourier coefficients show a dominant exponential decay with $`N`$, with the second order coefficient also decaying faster, i.e., with a smaller localization length, than the first one. We observe that for longer rings, stronger interactions and also stronger impurities, the current is increasingly more precisely described by its first Fourier component. In the LL phase, this in fact corresponds to the asymptotic behavior predicted by bosonization in the large $`N`$ limit, that is $`IN^{1\alpha _B}\mathrm{sin}\varphi `$. However, for the system sizes considered here this asymptotic regime is not reached and the current displays a more complex behavior. The current is composed of a few Fourier components with decreasing weight. Fig. 8 presents the first Fourier coefficient of the current for different values of the interaction, $`U=0.80`$, $`2.50`$, $`3.00`$, fixed $`\rho =0.50`$, from which we extract the dependence of $`I_1`$ on $`N`$, in the intermediate range of sizes considered. We observe that for $`U=0.80`$, in the LL phase, the first coefficient of the current varies as $`I_1`$ $`N^{1\alpha _1}`$, with $`\alpha _10.06`$, while for $`U=2.50`$ and $`U=3.00`$, in the CDW phase, it varies as $`I_1N^{1\delta _1}\mathrm{exp}(N/\xi _1)`$, respectively, with $`\xi _1259`$, $`\delta _1=0.11`$, and $`\xi _168`$, $`\delta _1=0.10`$. The exponent $`\alpha _1`$ is given by the slope of the straight line in Fig. 8.a, and the length $`\xi _1`$ and the exponent $`\delta _1`$ were carefully adjusted in order to obtain the best collapse of the data in Fig.8.b, on a plot of $`\mathrm{ln}(N^{1+\delta _1}I_1)`$ vs $`N/\xi _1`$. The Drude weights characterizing the systems with different interactions $`U`$, fixed $`\rho =0.50`$, are presented in Fig. 9. Fig. 10 clearly shows that the results obtained for the Drude weight confirm the conducting behavior shown by the first coefficient of the currents in Fig.8. We observe that, for $`U=0.80`$ the Drude weight varies with the system size as $`DN^\alpha `$, with $`\alpha 0.04`$, while for $`U=2.50`$ and $`U=3.00`$ it varies as $`DN^\delta \mathrm{exp}(N/\xi )`$, respectively, with $`\xi 307`$, $`\delta 0.08`$ and $`\xi 68`$, $`\delta 0.06`$. The exponent $`\alpha `$ is given by the slope of the straight line in Fig. 10.a, and the localization length $`\xi `$ and the exponent $`\delta `$ were carefully adjusted in order to obtain the best collapse of the data in Fig.10.b, on a plot of $`\mathrm{ln}(N^\delta D)`$ vs $`N/\xi `$. The exponents and localization lengths characterizing the Drude weight are a little different from those characterizing the first Fourier component of the current, as one would expect, since the Drude weight contains the contribution from the various Fourier components. One sees that the localization length $`\xi `$ and the exponent $`\delta `$ decrease with increasing strength of the interaction $`U`$. Also, concerning the impurity influence, Fig. 4 implies that the exponent $`\alpha `$ in the LL phase increases, and the localization length $`\xi `$ in the CDW phase decreases, with increasing $`\left|\rho 1\right|`$. As mentioned, in the large $`N`$ limit the current is expected to behave as its first Fourier component, which in the LL phase implies that the exponent $`\alpha _1`$ should be identified with $`\alpha _B=1/K1`$ as calculated from bosonization,<sup>21</sup> where $`K`$ is the LL parameter, calculable from the Bethe ansatz.<sup>40</sup> For $`U=0.80`$ this leads to $`\alpha _B0.27`$, which is much larger than our value of $`\alpha _1`$. We should note that the size dependence found for the first Fourier component of the current, and the Drude weight, characterizes the behavior of an intermediate and limited range of system sizes. If one would consider a larger range of systems, in the LL phase, one would most probably see the data for the larger $`N`$ bending down, crossing to an asymptotic power-law behaviour with the exponent approaching $`\alpha _B`$. This was observed in Ref. $`14`$, where the behavior of the first few Fourier components of the current in the LL phase was discussed in detail, with data taken for larger values of $`N`$ and stronger interaction and impurity strengths. Also, in the CDW phase we consider systems in an intermediate regime where the localization length is larger or near the system size.<sup>41</sup> For larger systems, the power factors that occur in the first Fourier component of the current and the Drude weight may decline,<sup>42</sup> possibly leaving a pure exponential behavior in the asymptotic regime.
From the results obtained, we observe that the system with $`U=0.80`$ and $`\rho =0.50`$, is characterized by an exponent $`\alpha >0`$, which is generated by the interplay of the electron interaction with the impurity, and $`D`$ exhibits then a power-law decay with $`N`$, which implies vanishing in the limit $`N\mathrm{}`$. On the other hand, the systems with $`U=2.50`$ and $`U=3.00`$, fixed $`\rho =0.50`$, are characterized by a localization length $`\xi `$, which decreases with increasing interaction and impurity strength, and $`D`$ exhibits now an exponential decay with $`N`$, also vanishing as $`N\mathrm{}`$. Hence, we find that both in the LL and the CDW phases, with an impurity in the system, the effect of the interaction is to decrease the current and the Drude weight. As referred before, our results also provide insight into the spin transport in a spin-1/2 XXZ chain with a weak link, and a similar behavior to the one above is implied for the spin current and stiffness. So, in the gapless XY phase, the spin stiffness decays with a power-law, vanishing in the limit $`N\mathrm{}`$, while in the gapped Ising phase, it decays exponentially, also vanishing as $`N\mathrm{}`$. Comparing our results for the persistent current in the LL phase, with those obtained in Ref. $`14`$, we have similarly found that the current vanishes faster than $`1/N`$. One observes that the model parameters strongly influence when the last asymptotic regime described by bosonization is reached. A calculation of the finite-size corrections to the spin stiffness in a pure spin-1/2 XXZ chain,<sup>42-44</sup> has revealed a size dependence in the gapped phase that has a similar form to the one found here. The result that the Drude weight in a ring in the gapless phase with an impurity drops to zero, is in agreement with a previous result obtained for a spin chain,<sup>38</sup> and with renormalization group arguments, which state that the impurity term is relevant leading to a transmission cut.<sup>31,32</sup> The renormalization group studies find that either a weak barrier or a weak link lead to an insulating state for repulsive interactions, while in the noninteracting case those are marginal perturbations. In turn, our work shows that there is an invariance of the current under the defect transformation $`\rho 1/\rho `$ in the interacting system, as for the noninteracting system, and that implies that a strong link will also reduce the current and the Drude weight. The observation that with an impurity in the system, the interaction always leads to an additional decrease of the current and the Drude weight is in agreement with previous results by other authors,<sup>7,8</sup> and can be understood as it is more difficult to move correlated electrons in a scattering potential than independent electrons.
## 5 SUMMARY
We have studied the behavior of the persistent current and the Drude weight on a mesoscopic ring pierced by a magnetic flux. We considered a model of spinless fermions with repulsive interactions and a hopping impurity, which is also equivalent to a spin-1/2 XXZ chain with a weakened exchange coupling. Using a powerful numerical method, the DMRG with complex fields, we have calculated the detailed form of the current as a function of the magnetic flux, which enabled us to investigate the corrections to the large system-size limit, and also allowed to obtain the Drude weight. We show that the system at half-filling, changes from an algebraic to an exponential behavior as the interaction increases, corresponding to a change from a LL to a CDW phase. We find that the analytical predictions of bosonization for the LL phase, are not yet fully observed in the intermediate range of system sizes considered. In addition we observe that the invariance of the current under the defect transformation $`\rho 1/\rho `$, seen in the noninteracting system, is also verified in the interacting system, in both phases. Hence, an isolated strong link is not only useless for increasing the persistent current (as might have been expected), but it rather destroys coherence and reduces the current. The behavior determined for the current is consistent with the behaviour determined for the Drude weight, the LL phase being characterized by an exponent $`\alpha >0`$, which results from the interplay of the interactions with the impurity, while the CDW phase is characterized by a localization length $`\xi `$, which decreases with increasing interaction and impurity strength. We find that, both in the LL and the CDW phase, with a defect in the system the interactions always suppress the current, and the Drude weight drops to zero in the limit $`N\mathrm{}`$. Away from half-filling there is no metal-insulator transition in the pure case, and the system is always metallic. Hence, one does not expect to observe then a change in the current and the Drude weight from an algebraic to an exponential decay. Nevertheless, one still expects to observe that in the system with an impurity, the current and the Drude weight decrease with increasing impurity and interaction strengthes. Therefore, within the model considered, the interactions cannot explain the results observed in the experiments.
Electronic address: fdias@cii.fc.ul.pt
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Figure Captions:
FIG. 1. Persistent current $`NI(\varphi )`$ vs $`\varphi `$, at fixed $`N=26`$, for $`U=0.80`$ and different $`\rho `$.
FIG. 2. The same as in Fig. 1, but for $`U=3.00`$.
FIG. 3. Symmetry of the current $`NI(\varphi )`$ in $`\rho `$ (empty symbols) vs $`1/\rho `$ (filled symbols), for $`U=0.80`$ (circles) at $`N=18`$, and $`U=3.00`$ (diamonds) at $`N=58`$.
FIG. 4. The Drude weigth $`D`$ vs $`\rho `$, at fixed $`N=26`$, for $`U=0.80`$ and $`U=3.00`$.
FIG. 5. Persistent current $`NI(\varphi )`$ vs $`\varphi `$, for $`U=0.80`$, $`\rho =0.50`$ and increasing $`N`$.
FIG. 6. The same as in Fig. 5, but for $`U=3.00`$.
FIG. 7. Fourier coefficients of the current. $`\mathrm{ln}(NI_k)`$ vs $`\mathrm{ln}(N)`$, for $`U=0.80`$ (circles) and $`U=3.00`$ (diamonds), at $`\rho =0.50`$, for $`k=1`$ (filled symbols) and $`k=2`$ (empty symbols).
FIG. 8. $`NI_1`$ for different $`U`$ and $`\rho =0.50`$. (a) $`\mathrm{ln}(NI_1)`$ vs $`\mathrm{ln}(N)`$, for $`U=0.80`$, the line represents $`I_1N^{1\alpha _1}`$. (b) $`\mathrm{ln}(NI_1)`$ vs $`N`$, for $`U=2.50`$ and $`3.00`$, the lines represent $`I_1N^{1\delta _1}\mathrm{exp}(N/\xi _1)`$, with $`\xi _1`$ and $`\delta _1`$ dependent on $`U`$.
FIG. 9. The Drude weigth $`D`$ as a function of $`1/N`$, for different $`U`$ and $`\rho =0.50`$.
FIG. 10. $`D`$ for different $`U`$ and $`\rho =0.50`$. (a) $`\mathrm{ln}(D)`$ vs $`\mathrm{ln}(N)`$, for $`U=0.80`$, the line represents $`DN^\alpha `$; (b) $`\mathrm{ln}(D)`$ vs $`N`$, for $`U=2.50`$ and $`3.00`$, the lines represent $`DN^\delta \mathrm{exp}(N/\xi )`$, with $`\xi `$ and $`\delta `$ dependent on $`U`$.
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# Coherent effects in magneto-transport through Zeeman-split levels
## I Introduction
Resonant transport through a quantum dot (or impurity) has been investigated in numerous publications. Yet no special attention has been paid to quantum interference effects in this process. These effects can generate oscillations in the resonant current through two (or more) levels of an impurity. These oscillations are similar to the well-known quantum interference effects in the two-slit experiment, and, in turn, produce a peak or a dip in the current power spectrum, depending on the relative phase of the two levels carrying the current. This feature can be experimentally observed in time-resolved measurements of transport currents . It has been previously argued that the quantum interference effect can explain modulation in the tunneling current at the Larmor frequency in scanning tunneling microscope (STM) experiments.
In this paper we investigate the interference effects in polarized magneto-transport. Conductance and I-V curves for spin dependent transport through quantum dots has recently been studied in several publications. Here we study the time dependent properties of transport currents. In particular we study the effects related to interference between different spin components of the currents.
These effects can be described schematically as follows. Consider the polarized resonant current from the left reservoir (emitter) to the right reservoir (collector) through a single level of a quantum dot (impurity) in the presence of an external magnetic field. This field would split the resonant level of the dot into a Zeeman doublet, Fig. 1. Let us assume that the polarization axis of electrons in the emitter (n) is different from that of the external magnetic field (). Then a spin-polarized electron from the emitter enters into a superposition of the “spin-up” and the “spin-down” states of the Zeeman doublet, Fig. 1. As a result, the electron wave function in the collector has two components corresponding to different energies of the doublet. Yet these components are orthogonal since they correspond to different spin components and therefore cannot interfere. If, however, there is an additional spin-flip process in the transition between the dot and the collector, the two spin components interfere in the collector current. Again, this takes place if the polarization in the collector (n) is different from that in the quantum dot (), Fig. 1. Thus the system operates as a two -path interferometer, where the phase difference between the two paths, i.e., through the upper and the lower spin states in the dot, contributes to the dynamical (or spectral) properties of the collector current.
Fig. 1. Resonant tunneling of a polarized electron through a quantum dot. Here the $`\mathrm{\Omega }`$’s denote the tunneling transition amplitudes between the reservoir states and the Zeeman doublets ($`E_{1,2}`$) of the quantum dot. $`\mu _{L,R}`$ are the chemical potentials in the left and right reservoirs. The unit vectors n, and n show the polarization axes in the emitter, quantum dot, and the collector, respectively.
These interference effects can be realized experimentally in a heterostructure with a quantum dot sandwiched between the two ferromagnetic leads with easy axes different from those in the dot. A similar set up can be implemented in other systems, such as self-assembled quantum dots, ultrasmall particles, carbon nanotubes and single molecules. These systems are likely to find prominent technological applications, including random excess memory and magnetic sensors due to giant magnetoresistance effect. The value of Zeeman splitting in the dot is controlled by the external magnetic field. The leads are assumed to be thin magnetic films, so that the magnetic field inside the leads is pinned parallel to the films and therefore the magnetization in the leads remains unaffected by the application of a relatively small external magnetic field.
The plan of this paper is as follows. In Sect. II we introduce the model and describe in general the rate equation approach for calculations of the resonant current and its noise spectrum. This approach has been obtained directly from the many-body Schrödinger equation describing the entire many-particle system and allows one to treat the magneto-transport through quantum dot systems in the most simple and precise way. In Sect. III we consider the case non-interacting (of weakly interacting) electrons. We explain in detail how the rate equations are used to calculate the polarized current and the noise spectrum both for ferromagnetic and non ferromagnetic leads. The results for the time-dependent polarized current are compared with the results of the single-particle model, presented in the Appendix. We explicitly demonstrate how the polarized current exhibits oscillations due to interference effects. In Sect. IV we concentrate on the interacting case. In particular we derive the current spectra in the presence of a Coulomb blockade in the dot. We consider separately the collector and the emitter currents, as well as the circuit current. In Sect. V we summarize our calculations and briefly discuss the potential implications of our results on the noise spectroscopy of quantum dots.
## II Many-body description.
Consider the polarized transport of non-interacting electrons through a quantum dot in the external magnetic field, Fig. 1. The polarization axis of an electron inside the dot () is different from those in the right and left reservoirs (n and n). The tunneling Hamiltonian describing this system can be written as
$`H={\displaystyle \underset{l,s}{}}E_{ls}a_{ls}^{}a_{ls}+{\displaystyle \underset{d=1,2}{}}E_da_d^{}a_d+{\displaystyle \underset{r,s^{}}{}}E_{rs^{}}a_{rs^{}}^{}a_{rs^{}}`$ (1)
$`+({\displaystyle \underset{d,l,s}{}}\mathrm{\Omega }_{dls}a_{ls}^{}a_d+{\displaystyle \underset{d,r,s^{}}{}}\mathrm{\Omega }_{drs^{}}a_{rs^{}}^{}a_d+H.c)+U_Ca_1^{}a_1a_2^{}a_2,`$ (2)
where the spin indices, $`s,s^{}=\pm 1/2`$ are related to different quantization axes (n and n’, Fig. 1), and $`E_{ls},E_{rs^{}}`$ denote the energy levels in the reservoirs. The Zeeman splitting of the dot is denoted by $`E_d`$, where $`d=1,2`$. The last two terms describe the tunneling transitions between the reservoirs and the dot states generated by both the tunneling couplings $`\mathrm{\Omega }`$ and the Coulomb repulsion of two electrons inside the dot.
All parameters of the tunneling Hamiltonian (2) are related to the initial microscopic description of the system in the configuration space ($`𝒙`$). For instance, the coupling $`\mathrm{\Omega }_{dls}`$ is given by the Bardeen formula
$$\mathrm{\Omega }_{dls}=\frac{1}{2m}_{𝒙\mathrm{\Sigma }_l}\varphi _d(𝒙)\stackrel{}{}\chi _{ls}(𝒙)d\sigma ,$$
(3)
where $`\varphi _d(𝒙)`$ and $`\chi _{ls}(𝒙)`$ are the electron wave functions inside the dot and the reservoir, respectively and $`\mathrm{\Sigma }_l`$ is a surface inside the potential barrier that separates the dot from the left reservoir. Since the spin quantization axes in the dot and in the leads differ from each other, the transition matrix elements ($`\mathrm{\Omega }`$’s) in Eq. (3) depend on the relative angles between the dot and the lead polarization axes ($`\theta _L`$ and $`\theta _R`$ for left and right leads respectively). The simplest form of the couplings that respects SU(2) symmetry is
$$\mathrm{\Omega }_{dls}=\mathrm{\Omega }_l𝖽_{s,s_d}^{(1/2)}(\theta _L)\text{and}\mathrm{\Omega }_{drs^{}}=\mathrm{\Omega }_r𝖽_{s_d,s^{}}^{(1/2)}(\theta _R),$$
(4)
where $`s_d=\pm 1/2`$ denotes the electron spin inside the dot, Fig. 1, $`𝖽^{(1/2)}(\theta )`$ is spin rotation matrix,
$$𝖽^{1/2}(\theta )=\left(\begin{array}{cc}\mathrm{cos}\frac{\theta }{2}& \mathrm{sin}\frac{\theta }{2}\\ \mathrm{sin}\frac{\theta }{2}& \mathrm{cos}\frac{\theta }{2}\end{array}\right),$$
(5)
and $`\mathrm{\Omega }_{l/r}`$ is spin-independent part of the couplings. Neglecting the energy dependence of these couplings, $`\mathrm{\Omega }_{l,r}=\mathrm{\Omega }_{L,R}`$, one can relate them to the partial widths (tunneling rates) as $`\mathrm{\Gamma }_{L,R}=2\pi \mathrm{\Omega }_{L,R}^2\rho _{L,R}`$, where $`\rho _{L,R}`$ is density of the states in the left (right) reservoir.
In the case of large bias, $`|E_{1,2}\mu _{L,R}|\mathrm{\Gamma }_{L,R}`$, the many-body Coulomb repulsion effects in the magneto-transport can be accounted for in the most simple and precise way by using the modified Bloch-type equations for the reduced density matrix. These equations can be derived from the many-body Schrödinger equation by integrating out the reservoir states in the limit of weak or strong Coulomb repulsion, $`U_C\mu _LE_1`$ or $`U_C\mu _LE_1`$, without any stochastic or other approximations. In addition these equations are very useful for evaluating of the shot-noise power spectrum.
In order to apply our method we first redefine the vacuum state, $`|\mathrm{𝟎}`$, by identifying it with the initial state of the entire system. For instance we can identify it with empty dot, with the emitter and collector filled up to the chemical potentials $`\mu _{L,R}`$, respectively. We also assume that the electrons in the emitter are polarized along the n-direction, Fig. 1. The many-body wave function, describing the entire system can be written in the most general way as
$`|𝚿(t)=[b_0(t)+{\displaystyle \underset{d,l,s}{}}b_{dls}(t)a_d^{}a_{ls}+{\displaystyle \underset{l,s,r,s^{}}{}}b_{rs^{}ls}(t)a_{rs^{}}^{}a_{ls}`$ (6)
$`+{\displaystyle \underset{l,s,\overline{l},\overline{s}}{}}b_{12ls\overline{l}\overline{s}}(t)a_1^{}a_2^{}a_{ls}a_{\overline{l}\overline{s}}+{\displaystyle \underset{d,l,s,\overline{l},\overline{s},r,s^{}}{}}b_{drs^{}ls\overline{l}\overline{s}r}(t)a_d^{}a_{rs^{}}^{}a_{ls}a_{\overline{l}\overline{s}}+\mathrm{}]|\mathrm{𝟎},`$ (7)
where $`d=\{1,2\}`$ denotes a state with one electron in the dot and $`ls(rs^{})`$ denote the electron level in the emitter (collector). The amplitudes $`b_\alpha (t)`$ of finding the entire system in the state “$`\alpha `$” are obtained from the Schrödinger equation, $`i_t|\mathrm{\Psi }(t)=|\mathrm{\Psi }(t)`$ with the initial condition $`b_\alpha (0)=\delta _{\alpha ,0}`$
Let us introduce the (reduced) density matrix $`\sigma _{jj^{}}^{m,n}(t)`$, where $`n,m`$ denote the number of electrons arriving the right reservoir with the spin components $`s^{}=\pm 1/2`$, respectively, Figs. 1,2. The lower indices, $`j,j^{}`$ denote the discrete states of the quantum dot. For instance, in the case of non-interacting (or weakly interacting) electrons $`j,j^{}=\{0,1,2,3\}`$, Fig. 2. This density matrix, $`\sigma _{jj^{}}^{m,n}(t)`$ can be easily constructed from the amplitudes $`b(t)`$, Eq. (7). For example,
$`\sigma _{00}^{0,0}(t)=|b_0(t)|^2,\sigma _{11}^{0,0}(t)={\displaystyle \underset{l}{}}|b_{1l1/2}(t)|^2,\sigma _{22}^{0,0}(t)={\displaystyle \underset{l}{}}|b_{2l1/2}(t)|^2,\sigma _{33}^{0,0}(t)={\displaystyle \underset{l,\overline{l}}{}}|b_{21l\overline{l}}(t)|^2,`$ (8)
$`\sigma _{12}^{0,0}(t)={\displaystyle \underset{l}{}}b_{1l1/2}(t)𝖻_{2l1/2}^{}(t),\sigma _{00}^{1,0}(t)={\displaystyle \underset{l,r}{}}|b_{1l1/2r1/2}(t)|^2,\mathrm{}`$ (9)
The diagonal density matrix elements, $`\sigma _{jj}^{n,m}`$, are the probabilities of finding the system in one of the states shown in Fig. 2 and the off-diagonal matrix elements (“coherencies”) describe a linear superposition of these states.
Fig. 2. Four available states of the quantum dot. The indices $`n,m`$ denote the number of electrons with the spin components $`s^{}=\pm 1/2`$ in the right reservoir.
It was demonstrated in Ref. that the Schrödinger equation for the entire system, $`i_t|𝚿(t)=|𝚿(t)`$, can be reduced to Bloch-type rate equations describing the reduced density-matrix $`\sigma _{jj^{}}^{n,m}(t)`$. This reduction takes place after partial tracing over the reservoir states. It becomes exact in the limit of large bias without the explicit use of any Markov-type or weak coupling approximations. In the general case these equations are
$`\dot{\sigma }_{jj^{}}=i(E_j^{}E_j)\sigma _{jj^{}}+i\left({\displaystyle \underset{k}{}}\sigma _{jk}\stackrel{~}{\mathrm{\Omega }}_{kj^{}}{\displaystyle \underset{k}{}}\stackrel{~}{\mathrm{\Omega }}_{jk}\sigma _{kj^{}}\right)`$ (10)
$`{\displaystyle \underset{k,k^{}}{}}𝒫_2\pi \rho (\sigma _{jk}\mathrm{\Omega }_{kk^{}}\mathrm{\Omega }_{k^{}j^{}}+\sigma _{kj^{}}\mathrm{\Omega }_{kk^{}}\mathrm{\Omega }_{k^{}j})+{\displaystyle \underset{k,k^{}}{}}𝒫_2\pi \rho (\mathrm{\Omega }_{kj}\mathrm{\Omega }_{k^{}j^{}}+\mathrm{\Omega }_{kj^{}}\mathrm{\Omega }_{k^{}j})\sigma _{kk^{}},`$ (11)
(12)
(for simplicity we have omitted the indices $`m`$ and $`n`$, which, however, can be easily restored from the conservation of the total number of electrons). Here $`\mathrm{\Omega }_{kk^{}}`$ denotes the single-electron hopping amplitude that generates the $`kk^{}`$ transition. We distinguish between the amplitudes $`\stackrel{~}{\mathrm{\Omega }}`$ describing single-electron hopping among isolated states and $`\mathrm{\Omega }`$ describing transitions among isolated and continuum states. The latter can generate transitions between the isolated states of the system, but only indirectly, via two consecutive jumps of an electron, into and out of the continuum reservoir states (with the density of states $`\rho `$). These transitions are represented by the third and the fourth terms of Eq. (12). The third term describes the transitions ($`kk^{}j`$) or ($`kk^{}j^{}`$), which cannot change the number of electrons ($`n,m`$) in the collector. The fourth term describes the transitions ($`kj`$ and $`k^{}j^{}`$) or ($`kj^{}`$ and $`k^{}j`$) which increase the number of electrons in the collector by one. These two terms of Eq. (12) are analogues of the “loss” (negative) and the “gain” (positive) terms in the classical rate equations, respectively. Yet, the sign of these terms depends on the relative sign of the corresponding couplings $`\mathrm{\Omega }`$. In our case it is determined by the sign of the spin-flip amplitude, Eq. (4). In addition, there is a (permutation) operator, $`𝒫_2=\pm 1`$, due to anti-commutation of the fermions operators, $`a_{1,2}^{}`$ in Eq. (7) (See also). The prefactor $`𝒫_2=1`$ whenever the loss or the gain terms in Eq. (12) are generated by a two-particle state of the dot. Otherwise $`𝒫_2=1`$.
## III Non-interacting electrons.
Consider first the case of no electron repulsion inside the dot, $`U_C=0`$. (In fact, the results would be the same for $`U_C\mu _LE_1`$, assuming the couplings $`\mathrm{\Omega }`$ are independent of energy). As in the previous section we choose the initial (“vacuum”) state corresponding to the polarized electrons in the left reservoir, $`s=1/2`$ (Fig. 1). In this case all four configurations shown in Fig. 3 contribute to Eqs. (12). Taking into account that there is no direct coupling between the states, $`E_{1,2}`$, i.e. $`\stackrel{~}{\mathrm{\Omega }}=0`$, one obtains the following Bloch-type rate equations for the density matrix $`\sigma _{jj^{}}^{n,m}(t)`$
$`\dot{\sigma }_{00}^{n,m}=\mathrm{\Gamma }_L\sigma _{00}^{n,m}+\mathrm{\Gamma }_R^{(1)}(\sigma _{11}^{n1,m}+\sigma _{22}^{n,m1})+\mathrm{\Gamma }_R^{(2)}(\sigma _{11}^{n,m1}+\sigma _{22}^{n1,m})+\mathrm{\Gamma }_L^{(2)}\sigma _{11}^{n,m}+\mathrm{\Gamma }_L^{(1)}\sigma _{22}^{n,m}`$ (14)
$`+\mathrm{\Gamma }_L^{(12)}(\sigma _{12}^{n,m}+\sigma _{21}^{n,m})\mathrm{\Gamma }_R^{(12)}(\sigma _{12}^{n1,m}+\sigma _{21}^{n1,m}\sigma _{12}^{n,m1}\sigma _{21}^{n,m1})`$ (15)
$`\dot{\sigma }_{11}^{n,m}=\left(\mathrm{\Gamma }_R+2\mathrm{\Gamma }_L^{(2)}\right)\sigma _{11}^{n,m}+\mathrm{\Gamma }_L^{(1)}(\sigma _{00}^{n,m}+\sigma _{33}^{n,m})\mathrm{\Gamma }_L^{(12)}(\sigma _{12}^{n,m}+\sigma _{21}^{n,m})`$ (16)
$`+\mathrm{\Gamma }_R^{(2)}\sigma _{33}^{n1,m}+\mathrm{\Gamma }_R^{(1)}\sigma _{33}^{n,m1}`$ (17)
$`\dot{\sigma }_{22}^{n,m}=\left(\mathrm{\Gamma }_R+2\mathrm{\Gamma }_L^{(1)}\right)\sigma _{22}^{n,m}+\mathrm{\Gamma }_L^{(2)}(\sigma _{00}^{n,m}+\sigma _{33}^{n,m})\mathrm{\Gamma }_L^{(12)}(\sigma _{12}^{n,m}+\sigma _{21}^{n,m})`$ (18)
$`+\mathrm{\Gamma }_R^{(1)}\sigma _{33}^{n1,m}+\mathrm{\Gamma }_R^{(2)}\sigma _{33}^{n,m1}`$ (19)
$`\dot{\sigma }_{33}^{n,m}=(2\mathrm{\Gamma }_R+\mathrm{\Gamma }_L)\sigma _{33}^{n,m}+\mathrm{\Gamma }_L^{(2)}\sigma _{11}^{n,m}+\mathrm{\Gamma }_L^{(1)}\sigma _{22}^{n,m}+\mathrm{\Gamma }_L^{(12)}(\sigma _{12}^{n,m}+\sigma _{21}^{n,m})`$ (20)
$`\dot{\sigma }_{12}^{n,m}=(iϵ+\mathrm{\Gamma })\sigma _{12}^{n,m}\mathrm{\Gamma }_L^{(12)}(\sigma _{00}^{n,m}+\sigma _{11}^{n,m}+\sigma _{22}^{n,m}+\sigma _{33}^{n,m})+\mathrm{\Gamma }_R^{(12)}(\sigma _{33}^{n1,m}\sigma _{33}^{n,m1}),`$ (21)
where $`\mathrm{\Gamma }=\mathrm{\Gamma }_L+\mathrm{\Gamma }_R`$ and $`\mathrm{\Gamma }_{L,R}^{(1)}=\mathrm{\Gamma }_{L,R}\mathrm{cos}^2(\theta _{L,R}/2)`$, $`\mathrm{\Gamma }_{L,R}^{(2)}=\mathrm{\Gamma }_{L,R}\mathrm{sin}^2(\theta _{L,R}/2)`$, $`\mathrm{\Gamma }_{L,R}^{(12)}=\mathrm{\Gamma }_{L,R}\mathrm{sin}(\theta _{L,R}/2)\mathrm{cos}(\theta _{L,R}/2)`$ are the partial tunneling widths of the levels $`E_{1,2}`$.
We now trace the origin of each term in these equations taking as an example Eq. (17), corresponding to $`j=j^{}=1`$ in Eq.(12). The first term in this equation is a “loss” term generated by the transitions $`101`$ and $`131`$ in Eq. (12). corresponding to the following processes: (a) an electron at the level $`E_1`$ (Fig. 3(1)) tunnels to the right reservoir and back to the same state, with the rate $`\mathrm{\Gamma }_R`$; (b) the same electron tunnels to the available continuum states of the left reservoir and back to the level $`E_1`$, with the rate $`\mathrm{\Gamma }_L^{(2)}=\mathrm{\Gamma }_L\mathrm{sin}^2(\theta _L/2)`$. This can proceed only via spin-flip, since the spin-up states in the left reservoir are occupied; (c) an electron from occupied states of the left reservoir tunnels to the unoccupied level $`E_2`$ of the dot and then back to the same state of the left reservoir, with the rate $`\mathrm{\Gamma }_L^{(2)}`$.
The second term in Eq. (17) is a “gain” term generated by the transitions $`01,01`$ and $`31,31`$ of an electron from the left reservoir to the level $`E_1`$ and from the level $`E_2`$ to the unoccupied (spin-down) continuum states of the left reservoir.
The third, “loss”, term in Eq. (17) is generated by the transitions $`201`$ and $`231`$ via the left reservoir. These transitions involve the following processes: (a) an electron at the level $`E_2`$ (Fig. 3(2)) tunnels to an unoccupied, spin-down state of the left reservoir, and then makes a spin-flip transition to the state $`E_1`$ of the dot. The rate of this process is $`(1/2)\mathrm{\Gamma }_L\mathrm{sin}(\theta _L)\mathrm{cos}(\theta _L/2)`$, as follows from Eq. (12); (b) an electron from one of the occupied states of the left reservoir tunnels to the state $`E_1`$ with the corresponding amplitude $`\mathrm{\Omega }_L\mathrm{cos}(\theta _L/2)`$. Then an electron with energy $`E_2`$ tunnels to the vacant state of the left reservoir with the spin-flip amplitude $`\mathrm{\Omega }_L\mathrm{sin}(\theta _L/2)`$. Since this transition proceeds via the two-electron state of the dot, the corresponding permutation prefactor, $`𝒫_2=1`$. As a result, the rate of this (loss) process is $`(1/2)\mathrm{\Gamma }_L\mathrm{sin}(\theta _L)`$. Similar transitions via the right reservoir cancel. Indeed the electron from energy level $`E_2`$ can reach the level $`E_1`$ by two ways: the first through the spin-flip hopping to the right reservoir and then to the level $`E_1`$ with no spin-flip, and the second, without spin-flip to the right reservoir, and then to the level $`E_1`$ with the spin flip. These two amplitudes are of the opposite sign.
The last two terms of Eq. (17) are “gain” terms generated by the transitions $`31,31`$ of an electron from the state in Fig. 3(3) to the spin-up or spin-down states of the right reservoir. The number of electrons in the right reservoir increases by one.
### A Resonant current in the collector
Using Eqs. (III) we can easily obtain the spin-up and spin-down currents, $`I_{1/2}(t)=_{n,m}n\dot{P}_{n,m}(t)`$ and $`I_{1/2}(t)=_{n,m}m\dot{P}_{n,m}(t)`$, where $`P_{n,m}(t)=_{j=0}^{j=3}\sigma _{jj}^{n,m}`$ is the probability of finding $`n`$ electrons with spin up and $`m`$ electrons with spin down in the right reservoir. One finds
$$I_{\pm 1/2}(t)=\mathrm{\Gamma }_R\left[\frac{1\pm \mathrm{cos}\theta _R}{2}\sigma _{11}(t)+\frac{1\mathrm{cos}\theta _R}{2}\sigma _{22}(t)+\sigma _{33}(t)\frac{\mathrm{sin}\theta _R}{2}(\sigma _{12}(t)+\sigma _{21}(t))\right]$$
(22)
where $`\sigma _{jj^{}}(t)=_{n,m}\sigma _{jj^{}}^{n,m}(t)`$. The latter can be obtained from the following matrix equation
$$\dot{X}(t)+BX(t)=0,$$
(23)
obtained by the summation of Eqs. (III) over $`n,m`$. Here $`X=\{\sigma _{00},\sigma _{11},\sigma _{22},\sigma _{33},\sigma _{12},\sigma _{21}\}`$ and $`B`$ is the corresponding $`6\times 6`$ matrix,
$$B=\left(\begin{array}{cccccc}\mathrm{\Gamma }_L& \mathrm{\Gamma }_L^{(2)}\mathrm{\Gamma }_R& \mathrm{\Gamma }_L^{(2)}\mathrm{\Gamma }_R& 0& \mathrm{\Gamma }_L^{(12)}& \mathrm{\Gamma }_L^{(12)}\\ \mathrm{\Gamma }_L^{(1)}& \mathrm{\Gamma }_R+2\mathrm{\Gamma }_L^{(2)}& 0& \mathrm{\Gamma }_L^{(1)}\mathrm{\Gamma }_R& \mathrm{\Gamma }_L^{(12)}& \mathrm{\Gamma }_L^{(12)}\\ \mathrm{\Gamma }_L^{(2)}& 0& \mathrm{\Gamma }_R+2\mathrm{\Gamma }_L^{(1)}& \mathrm{\Gamma }_L^{(2)}\mathrm{\Gamma }_R& \mathrm{\Gamma }_L^{(12)}& \mathrm{\Gamma }_L^{(12)}\\ 0& \mathrm{\Gamma }_L^{(2)}& \mathrm{\Gamma }_L^{(1)}& \mathrm{\Gamma }_L+2\mathrm{\Gamma }_R& \mathrm{\Gamma }_L^{(12)}& \mathrm{\Gamma }_L^{(12)}\\ \mathrm{\Gamma }_L^{(12)}& \mathrm{\Gamma }_L^{(12)}& \mathrm{\Gamma }_L^{(12)}& \mathrm{\Gamma }_L^{(12)}& iϵ+\mathrm{\Gamma }& 0\\ \mathrm{\Gamma }_L^{(12)}& \mathrm{\Gamma }_L^{(12)}& \mathrm{\Gamma }_L^{(12)}& \mathrm{\Gamma }_L^{(12)}& 0& iϵ+\mathrm{\Gamma }\end{array}\right)$$
(24)
Solving Eqs. (23) and substituting the result into Eqs. (22) we find the following simple expressions for the average polarized current:
$`I_{\pm 1/2}(t)={\displaystyle \frac{\mathrm{\Gamma }_L\mathrm{\Gamma }_R}{2\mathrm{\Gamma }}}(1\pm \mathrm{cos}\theta _L\mathrm{cos}\theta _R)(1e^{\mathrm{\Gamma }t})`$ (25)
$`\pm {\displaystyle \frac{\mathrm{\Gamma }_L\mathrm{\Gamma }_R\mathrm{\Gamma }\mathrm{sin}\theta _L\mathrm{sin}\theta _R}{2(ϵ^2+\mathrm{\Gamma }^2)}}\left[1e^{\mathrm{\Gamma }t}\mathrm{cos}(ϵt)+e^{\mathrm{\Gamma }t}{\displaystyle \frac{ϵ}{\mathrm{\Gamma }}}\mathrm{sin}(ϵt)\right].`$ (26)
The same result can be obtained in the framework of a single electron approach, valid for the noninteracting case. (See Appendix A.) As expected, the polarized resonant current displays damped oscillations. An example of these oscillations in $`I_{1/2}(t)`$ is shown in Fig. 3.
Fig. 3. Spin-up and total resonant currents through the Zeeman doublet as a function of time for $`\theta _L=\theta _R=\pi /2`$ and $`\mathrm{\Gamma }_L=\mathrm{\Gamma }_R=0.1ϵ`$.
These oscillations, however, disappear in the total collector current,
$$I(t)=I_{1/2}(t)+I_{1/2}(t)=\frac{\mathrm{\Gamma }_L\mathrm{\Gamma }_R}{\mathrm{\Gamma }}(1e^{\mathrm{\Gamma }t}),$$
(27)
even though electrons in the emitter are polarized.
### B Noise spectrum of the total collector current
Now we evaluate the noise-spectrum of the total current, represented by a sum of the spin-up and spin-down currents in the final state, Eq. (27). We introduce the density matrix $`\sigma _{jj^{}}^N(t)=_n\sigma ^{n,Nn}(t)`$, obtained from Eqs. (III), where $`N`$ denotes the total number of electrons which have arrived at the right reservoir by time $`t`$. In order to calculate the shot-noise spectrum we use the McDonald formula
$$S(\omega )=2e^2\omega _0^{\mathrm{}}dt\mathrm{sin}(\omega t)\frac{d}{dt}\underset{N}{}N^2P_N(t)],$$
(28)
where $`P_N(t)=_{j=0}^{j=3}\sigma _{jj}^N(t)`$. One easily finds from Eqs. (III) that
$$\underset{N}{}N^2\dot{P}_N(t)=\mathrm{\Gamma }_R\underset{N}{}(2N+1)\left[\sigma _{11}^N(t)+\sigma _{22}^N(t)+2\sigma _{33}^N(t)\right].$$
(29)
Substituting Eq. (29) into the McDonald formula, Eq. (28), we finally obtain
$$S(\omega )=2e^2\omega \mathrm{\Gamma }_R\text{Im}\left[Z_{11}(\omega )+Z_{22}(\omega )+2Z_{33}(\omega )\right],$$
(30)
where $`Z(\omega )`$ is a $`6`$-vector, $`Z=\{Z_{00},Z_{11},Z_{22},Z_{33},Z_{12}Z_{21}\}`$, defined as
$$Z_{ij}(\omega )=_0^{\mathrm{}}\underset{N}{}(2N+1)\sigma _{ij}^N(t)\mathrm{exp}(i\omega t)dt.$$
(31)
One can find $`Z_{ij}(\omega )`$ directly from Eqs. (III) by performing the corresponding summation over $`N`$. As a result one obtains
$$(Bi\omega I)Z(\omega )=\overline{X}+2\mathrm{\Gamma }_R\overline{Y}(\omega ).$$
(32)
Here $`B`$ is given by Eq. (24) and $`I`$ is the unit matrix. The 6-vector $`\overline{X}`$ corresponds to the stationary solution of Eqs. (23), $`\overline{X}=X(t\mathrm{})`$ and $`\overline{Y}(\omega )=\{Y_{11}+Y_{22},Y_{33},Y_{33},0,0,0\}`$ where $`Y(\omega )=\{Y_{00},Y_{11},Y_{22},Y_{33},Y_{33},Y_{12},Y_{21}\}`$ is given by the equation
$$(Bi\omega I)Y(\omega )=\overline{X}$$
(33)
Using Eq. (32) we calculate the ratio of the shot-noise power spectrum to the Schottky noise, $`S(\omega )/2eI`$ (Fano factor), where $`I=I(t\mathrm{})=\mathrm{\Gamma }_L\mathrm{\Gamma }_R/\mathrm{\Gamma }_t`$, Eq.(27). In particular, the result has a simple analytical form for a symmetric dot, $`\mathrm{\Gamma }_L=\mathrm{\Gamma }_R=\mathrm{\Gamma }`$. We find
$$\frac{S(\omega )}{2eI}=\frac{2\mathrm{\Gamma }^2+\omega ^2}{4\mathrm{\Gamma }^2+\omega ^2}+\frac{\mathrm{\Gamma }^2ϵ^2\mathrm{sin}^2\theta _L}{(4\mathrm{\Gamma }^2+\omega ^2)(4\mathrm{\Gamma }^2+\omega ^2)}.$$
(34)
As expected the shot-noise spectrum does not display any peak or dip at frequencies corresponding to the Zeeman splitting, since the interference effects are canceled in the total collector current. Yet the noise spectrum depends on the initial polarization of incoming electrons ($`\theta _L`$), whereas the total collector current does not (see Eq. (27)). If electrons are initially polarized along the magnetic field inside the dot (), the Fano factor is the same as in the case of resonant tunneling through a single level. With increasing $`\theta _L`$, however, the current flows through both levels of the Zeeman doublet. This leads to an additional contribution to the shot noise, described by the second term of Eq. (34).
### C Ferromagnetic reservoirs
Let us consider ferromagnetic reservoirs polarized along n and n’ directions, Fig. 1. In this case the rate equations (III) have to be modified since there are no available spin-down states in the left and right reservoirs. One easily obtains the following rate equations for the density matrix $`\sigma _{jj^{}}^n(t)`$, where $`n`$ denotes the number of electron, arriving at the collector before time $`t`$:
$`\dot{\sigma }_{00}^n`$ $`=`$ $`\mathrm{\Gamma }_L\sigma _{00}^n+\mathrm{\Gamma }_R^{(1)}\sigma _{11}^{n1}+\mathrm{\Gamma }_R^{(2)}\sigma _{22}^{n1}\mathrm{\Gamma }_R^{(12)}(\sigma _{12}^{n1}+\sigma _{21}^{n1})`$ (36)
$`\dot{\sigma }_{11}^n`$ $`=`$ $`\left(\mathrm{\Gamma }_L^{(2)}+\mathrm{\Gamma }_R^{(1)}\right)\sigma _{11}^n+\mathrm{\Gamma }_L^{(1)}\sigma _{00}^n{\displaystyle \frac{1}{2}}(\mathrm{\Gamma }_L^{(12)}\mathrm{\Gamma }_R^{(12)})(\sigma _{12}^n+\sigma _{21}^n)+\mathrm{\Gamma }_R^{(2)}\sigma _{33}^{n1}`$ (37)
$`\dot{\sigma }_{22}^n`$ $`=`$ $`\left(\mathrm{\Gamma }_L^{(1)}+\mathrm{\Gamma }_R^{(2)}\right)\sigma _{22}^n+\mathrm{\Gamma }_L^{(2)}\sigma _{00}^n{\displaystyle \frac{1}{2}}(\mathrm{\Gamma }_L^{(12)}\mathrm{\Gamma }_R^{(12)})(\sigma _{12}^n+\sigma _{21}^n)+\mathrm{\Gamma }_R^{(1)}\sigma _{33}^{n1}`$ (38)
$`\dot{\sigma }_{33}^n`$ $`=`$ $`\mathrm{\Gamma }_R\sigma _{33}^n+\mathrm{\Gamma }_L^{(2)}\sigma _{11}^n+\mathrm{\Gamma }_L^{(1)}\sigma _{22}^n+\mathrm{\Gamma }_L^{(12)}(\sigma _{12}^n+\sigma _{21}^n)`$ (39)
$`\dot{\sigma }_{12}^n`$ $`=`$ $`\left(iϵ+{\displaystyle \frac{1}{2}}\mathrm{\Gamma }\right)\sigma _{12}^n\mathrm{\Gamma }_L^{(12)}\sigma _{00}^n{\displaystyle \frac{1}{2}}(\mathrm{\Gamma }_L^{(12)}\mathrm{\Gamma }_R^{(12)})(\sigma _{11}^n+\sigma _{22}^n)+\mathrm{\Gamma }_R^{(12)}\sigma _{33}^{n1}`$ (40)
Using these equations we first evaluate the average current, $`I(t)I_{1/2}(t)`$ given by Eq. (22) with $`\sigma _{jj^{}}(t)=_n\sigma _{jj^{}}^n(t)`$. The latter quantities are obtained from a summation of Eqs. (III C) over $`n`$. As a result Eqs. (III C) are reduced to the matrix equation (23), where $`B`$ is the corresponding $`6\times 6`$ matrix of the coefficients of Eqs. (III C). Solving this equation we find the average current $`I(t)`$. For instance, in the case of $`\theta _L=\theta _R`$ one finds for the stationary current, $`I=I(\mathrm{})=\mathrm{\Gamma }_L\mathrm{\Gamma }_R/(\mathrm{\Gamma }_L+\mathrm{\Gamma }_R)`$. One obtains the same expression for the resonant tunneling of unpolarized electrons through a single level.
The time dependence of the average current, $`I(t)`$, is displayed in Fig. 4 for symmetric and asymmetric dots, $`\mathrm{\Gamma }_L=\mathrm{\Gamma }_R=0.1ϵ`$ and $`\mathrm{\Gamma }_L=ϵ`$, $`\mathrm{\Gamma }_R=0.1ϵ`$, respectively. Comparing with Fig. 2 one finds that the oscillations in the average current are more pronounced in the case of ferromagnetic reservoirs. This can be anticipated since the corresponding spin-flip transitions via the spin-down states of the reservoirs do not exist. We recall that precisely these transitions resulted in the cancelation of the interference effects in the previous case.
Fig. 4. The polarized resonant current through the Zeeman doublet with ferromagnetic reservoirs and $`\theta _L=\theta _R=\pi /2`$. The solid line corresponds to $`\mathrm{\Gamma }_L=\mathrm{\Gamma }_R=0.1ϵ`$ and the dashed line to $`\mathrm{\Gamma }_L=ϵ`$ and $`\mathrm{\Gamma }_R=0.1ϵ`$.
Now we can evaluate the shot-noise spectrum, $`S(\omega )`$ using the McDonald formula. One obtains from Eqs. (28) and (III C)
$`S(\omega )`$ $`=`$ $`2e^2\omega \text{Im}\left\{\mathrm{\Gamma }_R^{(1)}Z_{11}(\omega )+\mathrm{\Gamma }_R^{(2)}Z_{22}(\omega )+\mathrm{\Gamma }_RZ_{33}(\omega )\mathrm{\Gamma }_R^{(12)}[Z_{12}(\omega )+Z_{21}(\omega )]\right\},`$ (41)
where $`Z(\omega )`$ is given by Eq. (32) with the matrix $`B`$ corresponding to Eqs. (III C) and $`\overline{Y}=\{\overline{Y}_{00},\overline{Y}_{11},\overline{Y}_{22},0,\overline{Y}_{12},\overline{Y}_{21}\}`$. Here $`\overline{Y}_{00}=\mathrm{cos}^2\frac{\theta _R}{2}Y_{11}+\mathrm{sin}^2\frac{\theta _R}{2}Y_{22}\frac{\mathrm{sin}\theta _R}{2}(Y_{12}+Y_{21})`$, $`\overline{Y}_{11}=\mathrm{sin}^2\frac{\theta _R}{2}Y_{33}`$, $`\overline{Y}_{22}=\mathrm{cos}^2\frac{\theta _R}{2}Y_{33}`$, and $`\overline{Y}_{12}=\overline{Y}_{21}=\frac{\mathrm{sin}\theta _R}{2}Y_{33}`$, while $`Y_{jj^{}}=Y_{jj^{}}(\omega )`$ are given by Eq. (33).
The corresponding Fano factor is shown in Fig. 5 for the same parameters as in Fig. 4. It clearly displays a dip at the Zeeman frequency for a symmetric dot. It reflects the damped oscillations in the average current, shown in Fig. 4. The dip, however, almost disappears for an asymmetric dot with large $`\mathrm{\Gamma }_L`$.
Fig. 5. The Fano factor versus $`\omega `$ for a polarized electron current with ferromagnetic reservoirs and $`\theta _L,\theta _R=\pi /2`$. The solid line corresponds to $`\mathrm{\Gamma }_L=\mathrm{\Gamma }_R=0.1ϵ`$ and the dashed line to $`\mathrm{\Gamma }_L=ϵ`$ and $`\mathrm{\Gamma }_R=0.1ϵ`$.
## IV Coulomb blockade
We now introduce strong Coulomb repulsion inside the dot, $`U_C\mu _LE_1`$, so that the state (3) in Fig. 3 is not available. As a result the corresponding rate equations have an even simpler form than those found for non-interacting electrons. Consider again the case of ferromagnetic reservoirs, where the quantum interference effects are most pronounced. The corresponding rate equations for the case of Coulomb blockade can be obtained from Eqs. (III C) for non-interacting electrons, by eliminating configurations with two electrons in the dot. In the following we consider separately the electron current in the right and in the left reservoirs.
### A Collector current
The electrical current in the right reservoir and its power spectrum are obtained from the following rate equation
$`\dot{\sigma }_{00}^n`$ $`=`$ $`\mathrm{\Gamma }_L\sigma _{00}^n+\mathrm{\Gamma }_R^{(1)}\sigma _{11}^{n1}+\mathrm{\Gamma }_R^{(2)}\sigma _{22}^{n1}\mathrm{\Gamma }_R^{(12)}(\sigma _{12}^{n1}+\sigma _{21}^{n1})`$ (43)
$`\dot{\sigma }_{11}^n`$ $`=`$ $`\mathrm{\Gamma }_R^{(1)}\sigma _{11}^n+\mathrm{\Gamma }_L^{(1)}\sigma _{00}^n+{\displaystyle \frac{\mathrm{\Gamma }_R^{(12)}}{2}}(\sigma _{12}^n+\sigma _{21}^n)`$ (44)
$`\dot{\sigma }_{22}^n`$ $`=`$ $`\mathrm{\Gamma }_R^{(2)}\sigma _{22}^n+\mathrm{\Gamma }_L^{(2)}\sigma _{00}^n+{\displaystyle \frac{\mathrm{\Gamma }_R^{(12)}}{2}}(\sigma _{12}^n+\sigma _{21}^n)`$ (45)
$`\dot{\sigma }_{12}^n`$ $`=`$ $`\left(iϵ+{\displaystyle \frac{\mathrm{\Gamma }_R}{2}}\right)\sigma _{12}^n\mathrm{\Gamma }_L^{(12)}\sigma _{00}^n+{\displaystyle \frac{\mathrm{\Gamma }_R^{(12)}}{2}}(\sigma _{11}^n+\sigma _{22}^n)`$ (46)
Using these equations one finds for the average (polarized) current in the collector
$$I_R(t)=\mathrm{\Gamma }_R^{(1)}\sigma _{11}(t)+\mathrm{\Gamma }_R^{(2)}\sigma _{22}(t)\mathrm{\Gamma }_R^{(12)}[\sigma _{12}(t)+\sigma _{21}(t)]$$
(47)
where the $`\sigma _{jj^{}}(t)=_{n,m}\sigma _{jj^{}}^{n,m}(t)`$ are obtained from Eq. (23) for $`X=\{\sigma _{00},\sigma _{11},\sigma _{22},\sigma _{12},\sigma _{21}\}`$, and $`B`$ is the $`5\times 5`$ matrix obtained from the coefficients of Eqs. (IV A). Solving such a modified Eq. (23) for $`\theta _L=\theta _R`$, one finds for the stationary current,
$$I_R=I_R(\mathrm{})=\frac{\mathrm{\Gamma }_L\mathrm{\Gamma }_R}{2\mathrm{\Gamma }_L+\mathrm{\Gamma }_R}$$
(48)
This expression shows an asymmetry with respect to the widths $`\mathrm{\Gamma }_L`$ and $`\mathrm{\Gamma }_R`$, in contrast with the non-interacting case. The reason is that an electron enters the dot from the left reservoir with the rate $`2\mathrm{\Gamma }_L`$. However, it leaves it with the rate $`\mathrm{\Gamma }_R`$, since the state with two levels of the dot occupied is forbidden.
The shot-noise power spectrum for the collector current is given by
$$S_R(\omega )=2e^2\omega \text{Im}\left\{\mathrm{\Gamma }_R^{(1)}Z_{11}(\omega )+\mathrm{\Gamma }_R^{(2)}Z_{22}(\omega )\mathrm{\Gamma }_R^{(12)}[Z_{12}(\omega )+Z_{21}(\omega )]\right\}.$$
(49)
Here $`Z_{ij}(\omega )`$ are obtained from Eqs. (32),(33), where $`\overline{Y}=\{\overline{Y}_{00},0,0,0,0\}`$ and $`\overline{Y}_{00}=\mathrm{cos}^2\frac{\theta _R}{2}Y_{11}+\mathrm{sin}^2\frac{\theta _R}{2}Y_{22}\frac{\mathrm{sin}\theta _R}{2}(Y_{12}+Y_{21})`$.
The results of our calculations of $`S(\omega )`$ for symmetric and asymmetric quantum dots in the case of Coulomb blockade are shown in Fig. 6.
Fig. 6. The Fano factor versus $`\omega `$ for a polarized collector current with ferromagnetic reservoirs and Coulomb blockade and $`\theta _L,\theta _R=\pi /2`$. The solid line corresponds to $`\mathrm{\Gamma }_L=\mathrm{\Gamma }_R=0.1ϵ`$ and the dashed line to $`\mathrm{\Gamma }_L=ϵ`$ and $`\mathrm{\Gamma }_R=0.1ϵ`$.
### B Emitter current
We now consider the electric current and its power spectrum in the left reservoir. These quantities are determined from the density-matrix $`\sigma _{jj^{}}^p`$, where $`p`$ is the number of electrons that left the emitter before time $`t`$, (the number of holes in the left reservoir). The corresponding rate equations are similar to Eqs. (IV A). One finds
$`\dot{\sigma }_{00}^p`$ $`=`$ $`\mathrm{\Gamma }_L\sigma _{00}^p+\mathrm{\Gamma }_R^{(1)}\sigma _{11}^p+\mathrm{\Gamma }_R^{(2)}\sigma _{22}^p\mathrm{\Gamma }_R^{(12)}(\sigma _{12}^p+\sigma _{21}^p)`$ (51)
$`\dot{\sigma }_{11}^p`$ $`=`$ $`\mathrm{\Gamma }_R^{(1)}\sigma _{11}^p+\mathrm{\Gamma }_L^{(1)}\sigma _{00}^{p1}+{\displaystyle \frac{\mathrm{\Gamma }_R^{(12)}}{2}}(\sigma _{12}^p+\sigma _{21}^p)`$ (52)
$`\dot{\sigma }_{22}^p`$ $`=`$ $`\mathrm{\Gamma }_R^{(2)}\sigma _{22}^p+\mathrm{\Gamma }_L^{(2)}\sigma _{00}^{p1}+{\displaystyle \frac{\mathrm{\Gamma }_R^{(12)}}{2}}(\sigma _{12}^p+\sigma _{21}^p)`$ (53)
$`\dot{\sigma }_{12}^p`$ $`=`$ $`\left(iϵ+{\displaystyle \frac{\mathrm{\Gamma }_R}{2}}\right)\sigma _{12}^p\mathrm{\Gamma }_L^{(12)}\sigma _{00}^{p1}+{\displaystyle \frac{\mathrm{\Gamma }_R^{(12)}}{2}}(\sigma _{11}^p+\sigma _{22}^p)`$ (54)
The average emitter current in the left reservoir is given by $`I_L(t)=\mathrm{\Gamma }_L\sigma _{00}(t)`$, which differs from Eq. (47) describing the collector current, $`I_R(t)`$. Yet, as expected, their stationary values coincide, $`I_L(\mathrm{})=I_R(\mathrm{})`$, Eq. (48).
The shot-noise power spectrum of the emitter current is given
$$S_L(\omega )=2e^2\omega \mathrm{\Gamma }_L\text{Im}Z_{00}(\omega )$$
(55)
instead of Eq. (49) for $`S_R(\omega )`$, where $`Z_{ij}(\omega )`$ are obtained from Eqs. (32),(33). Yet, $`\overline{Y}(\omega )=\{0,\mathrm{cos}^2\frac{\theta _L}{2},\mathrm{sin}^2\frac{\theta _L}{2},\frac{\mathrm{sin}\theta _L}{2},\frac{\mathrm{sin}\theta _L}{2}\}Y_{00}(\omega )`$, in contrast with the corresponding expression for $`S_R(\omega )`$. Even though the expressions for $`S_{L,R}(\omega )`$ are quite different, one finds that the shot-noise power of the emitter current is the same as that in the collector current, $`S_L(\omega )=S_R(\omega )`$.
### C Circuit current
In general the circuit current is given by $`I_c(t)=\alpha I_L(t)+\beta I_R(t)`$, where the coefficients $`\alpha ,\beta `$ with $`\alpha +\beta =1`$ depend on the junction capacities. Using charge conservation, $`I_L=I_R+\dot{Q}`$, where $`Q`$ is charge in the dot, one finds
$$I_c(t)I_c(0)=\alpha I_L(t)I_L(0)+\beta I_R(t)I_R(0)\alpha \beta \dot{Q}(t)\dot{Q}(0).$$
(56)
Using this relation one finds a simple expression for the noise spectrum of the circuit current
$$S_c(\omega )=\alpha S_L(\omega )+\beta S_R(\omega )\alpha \beta \omega ^2S_Q(\omega ).$$
(57)
where $`S_Q(\omega )`$ is Fourier transform of the charge correlation function. This quantity can be obtained straightforwardly from the matrix equation (33), where $`\overline{X}`$ is the 5-vector $`\{0,\sigma _{11}(\mathrm{}),\sigma _{22}(\mathrm{}),0,0\}`$. Then $`S_Q(\omega )=4`$Re$`[Y_{11}(\omega )+Y_{22}(\omega )]`$.
The results of our calculations of $`S_c(\omega )`$ for $`\alpha =\beta =1/2`$ are shown in Fig. 7.
Fig. 7. The Fano factor for the circuit ($`\alpha =\beta =1/2`$) versus $`\omega `$ for a polarized electron current with ferromagnetic reservoirs and Coulomb blockade and $`\theta _L,\theta _R=\pi /2`$. The solid line corresponds to $`\mathrm{\Gamma }_L=\mathrm{\Gamma }_R=0.1ϵ`$ and the dashed line to $`\mathrm{\Gamma }_L=ϵ`$ and $`\mathrm{\Gamma }_R=0.1ϵ`$.
One finds from Figs. 6 and 7 that the Coulomb blockade modifies the current spectrum very drastically with respect to the non-interacting case, Fig. 5.
## V Conclusions
In this paper, we study the interference effects in magneto-transport through Zeeman split levels of quantum dots or impurities. We concentrated on the time-dependent properties and the power spectrum of the electric current by applying the new approach using quantum rate equations, which is mostly suitable for this type of problems. We explicitly demonstrated that our method produces the same results as a single electron approach, widely used for a description of non-interacting electron transport. Yet the quantum rate equations method is valid also for the case of interacting electrons and accounts for the Coulomb blockade in the most simple and precise way.
Our results indicate that the Coulomb blockade plays an important role in the spectral properties of the transport current. First of all, in the presence of Coulomb blockade the signal-to-noise ratio significantly is amplified, as one can observe from the results in the previous sections. This is probably a consequence of the prohibition of double occupation of the resonant level in the quantum dot. Indeed, when two electrons in the dot are present, the interference effects are suppressed due to the “randomization” of the relative phase. Interestingly, the dip in the noise spectrum for the noninteracting electrons is replaced by a peak as result of Coulomb interaction. Clearly the Coulomb interaction modifies the phase of the electrons tunneling trough the dot, which “flips” the spectral feature in the noise. The details of this very interesting phenomena must be studied in the future.
We emphasize that the coherent oscillations in the current can be observed only for polarized current and that oscillations disappear for unpolarized current. This is different from the resonant transport through two orbital levels of a quantum dot or impurity, where the quantum interference effects can be observed even in unpolarized case. Therefore it is most natural to use ferromagnetic leads for observation and utilization of quantum interference effect in the magneto-transport. Thus our calculations were mostly concentrated on this case. Our results show explicitly the appearance of peak or dip at a frequency near the Zeemann splitting frequency (Larmour frequency). We believe that this phenomenon can be useful for analyzing the noise spectroscopy of quantum dots or impurities. Indeed, the Zeeman splitting of a localized quantum dot orbital must be sensitive to local magnetic fields, and therefore one can hope that such coherent effect, if observed experimentally, may allow for detection of the local hyperfine structure of the dot/impurity. This, however, must be a subject of a separate investigation.
## VI Acknowledgement
We thank J. Brown, L. Fedichkin, M. B. Hastings, M. Hawley, and I. Martin for valuable discussions. We are especially grateful to Gary D. Doolen for important remarks and for proofreading the manuscript. The work was supported by the Department of Energy under Contract No. W-7405-ENG-36 and by DOE Office of basic Energy Sciences. D. M. was supported, in part, by the US NSF grant DMR-0121146.
## A Single-electron description
In the case of non-interacting electrons one can compare our results with those obtained using a single-electron approach. Although the latter is widely used in the literature, it is usually restricted to the time-independent (stationary) case. Here we present an extension of the single-electron approach for the non-stationary case. This would allow us to evaluate the time-dependent resonant current, Fig. 3, and to compare the results with those obtained from Eqs. (12).
Let us consider a system consisting of the reservoirs and the quantum dot filled with only a single electron. We assume that this electron is initially in the left reservoir (emitter) at the level $`E_{\overline{l}}`$ with the spin polarized along the n-direction, Fig. 1. The electron motion is described by a wave function which can be written in the most general way as:
$$|\mathrm{\Psi }(t)=\left[\underset{l,s}{}𝖻_{ls}(t)a_{ls}^{}+\underset{d=1,2}{}𝖻_d(t)a_d^{}+\underset{r,s^{}}{}𝖻_{rs^{}}(t)a_{rs^{}}^{}\right]|0,$$
(A1)
where $`𝖻_\alpha (t)`$ is the amplitude of finding the electron in the state $`\alpha `$ given by a corresponding creation operator. These amplitudes are obtained from the Schrödinger equation $`|\mathrm{\Psi }(t)`$, with the initial conditions $`𝖻_{ls}(0)=\delta _{l,\overline{l}}\delta _{s,1/2}`$ and $`𝖻_d(0)=𝖻_{rs}(0)=0`$. It is useful to use the Laplace transform, $`\stackrel{~}{𝖻}(E)=_0^{\mathrm{}}𝖻(t)\mathrm{exp}(iEt)𝑑t`$. In this case the time-dependent Schrödinger equation for the amplitudes $`\stackrel{~}{𝖻}(E)`$ becomes the following system of linear algebraic equations
$`(EE_{ls})\stackrel{~}{𝖻}_{ls}(E)\mathrm{\Omega }_l{\displaystyle \underset{d^{}=1,2}{}}𝖽_{s,s_d^{}}^{(1/2)}(\theta _L)\stackrel{~}{𝖻}_d^{}(E)=i\delta _{l,\overline{l}}\delta _{s,1/2}`$ (A3)
$`(EE_d)\stackrel{~}{𝖻}_d(E){\displaystyle \underset{l,s}{}}\mathrm{\Omega }_l𝖽_{s_d,s}^{(1/2)}(\theta _L)\stackrel{~}{𝖻}_{ls}(E){\displaystyle \underset{r,s^{}}{}}\mathrm{\Omega }_r𝖽_{s_d,s^{}}^{(1/2)}(\theta _R)\stackrel{~}{𝖻}_{rs^{}}(E)=0`$ (A4)
$`(EE_{rs^{}})\stackrel{~}{𝖻}_{rs^{}}(E)\mathrm{\Omega }_r{\displaystyle \underset{d^{}=1,2}{}}𝖽_{s^{},s_d^{}}^{(1/2)}(\theta _R)\stackrel{~}{𝖻}_d^{}(E)=0.`$ (A5)
Substituting $`\stackrel{~}{𝖻}_{ls}`$ and $`\stackrel{~}{𝖻}_{rs^{}}`$ from Eqs. (A3), (A5) into Eq. (A4) and replacing the sums on $`l`$ and $`r`$ by the integrals, we obtain
$`\left(EE_1+i{\displaystyle \frac{\mathrm{\Gamma }_L+\mathrm{\Gamma }_R}{2}}\right)\stackrel{~}{𝖻}_1(E)=i{\displaystyle \frac{\mathrm{\Omega }_L\mathrm{cos}(\theta _L/2)}{EE_{\overline{l},1/2}}}`$ (A7)
$`\left(EE_2+i{\displaystyle \frac{\mathrm{\Gamma }_L+\mathrm{\Gamma }_R}{2}}\right)\stackrel{~}{𝖻}_2(E)=i{\displaystyle \frac{\mathrm{\Omega }_L\mathrm{sin}(\theta _L/2)}{EE_{\overline{l},1/2}}}.`$ (A8)
Note that the amplitudes, $`\stackrel{~}{𝖻}_1(E)`$ and $`\stackrel{~}{𝖻}_2(E)`$, are decoupled in Eqs. (A) although the corresponding states are connected via the continuum. The reason is that the spin-flip couplings of the dot with the reservoirs are of the opposite sign for the spin-up and the spin-down states of the dot ($`E_1`$ and $`E_2`$ in Fig. 1). However, for the general case of resonant tunneling through two levels, the corresponding amplitudes are coupled via the interaction through continuum.
Using the inverse Laplace transform $`𝖻_{1,2}(t)=\stackrel{~}{𝖻}_{1,2}(E)\mathrm{exp}(iEt)𝑑E/(2\pi )`$, we obtain for the amplitudes, $`𝖻_{1,2}(t)`$, for finding the electron inside the dot
$`𝖻_1(t)`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Omega }_L\mathrm{cos}(\theta _L/2)}{E_LE_1+i\frac{\mathrm{\Gamma }}{2}}}\left(e^{iE_Lt}e^{iE_1t\frac{\mathrm{\Gamma }}{2}t}\right)`$ (A10)
$`𝖻_2(t)`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Omega }_L\mathrm{sin}(\theta _L/2)}{E_LE_2+i\frac{\mathrm{\Gamma }}{2}}}\left(e^{iE_Lt}e^{iE_2t\frac{\mathrm{\Gamma }}{2}t}\right),`$ (A11)
where $`\mathrm{\Gamma }=\mathrm{\Gamma }_L+\mathrm{\Gamma }_R`$. The probability amplitude of finding the electron inside the collector is $`\stackrel{~}{𝖻}_{rs^{}}(t)=\stackrel{~}{𝖻}_{rs^{}}(E)\mathrm{exp}(iEt)𝑑E/(2\pi )`$, where $`\stackrel{~}{𝖻}_{r,s^{}}(E)`$ is given by Eq. (A5)
$$\stackrel{~}{𝖻}_{rs^{}}(E)=\frac{\mathrm{\Omega }_R}{EE_{rs^{}}}\underset{d}{}𝖽_{s^{}s_d}^{(1/2)}(\theta _R)\stackrel{~}{𝖻}_d(E).$$
(A12)
The above equations determine the motion of a single electron placed initially in the emitter. In order to obtain the polarized current, $`I_s^{}(t)`$, in the single-electron model one has to sum over all initially occupied states $`E_{\overline{l}}`$ of the emitter and over all available states $`E_r`$ of the collector. Thus $`I_s^{}=dN_s^{}(t)/dt`$, where $`N_s^{}(t)=_{\overline{l},r}|𝖻_{rs^{}}(t)|^2`$ is the average number of electrons with spin-up and spin-down ($`s^{}=\pm 1/2`$), accumulated in the collector by the time $`t`$. Using the inverse Laplace transform and replacing $`_{\overline{l},r}\rho _L\rho _R𝑑E_L𝑑E_R`$ we obtain
$$N_s^{}(t)=\rho _L\rho _R𝑑E_L𝑑E_R\frac{dEdE^{}}{(2\pi )^2}\stackrel{~}{𝖻}_{rs^{}}(E)\stackrel{~}{𝖻}_{rs^{}}^{}(E^{})e^{i(E^{}E)t}$$
(A13)
Substituting Eq. (A12) into Eq. (A13) and integrating over $`E_{rs^{}}`$ one obtains for the polarized current
$`I_{1/2}(t)`$ $`=`$ $`\mathrm{\Gamma }_R{\displaystyle _{\mu _R}^{\mu _L}}\rho _L𝑑E_L|\mathrm{cos}(\theta _R/2)𝖻_1(t)\mathrm{sin}(\theta _R/2)𝖻_2(t)|^2`$ (A15)
$`I_{1/2}(t)`$ $`=`$ $`\mathrm{\Gamma }_R{\displaystyle _{\mu _R}^{\mu _L}}\rho _L𝑑E_L|\mathrm{sin}(\theta _R/2)𝖻_1(t)+\mathrm{cos}(\theta _R/2)𝖻_2(t)|^2`$ (A16)
where the amplitudes $`𝖻_{1,2}(t)`$ are given by Eqs. (A). Note that these amplitudes in the stationary limit, $`𝖻_{1,2}(t\mathrm{})`$, are the transmission amplitudes describing the resonance tunneling through the levels $`E_{1,2}`$ . Thus Eqs. (A) represent a generalization of the Landauer formula for the time-dependent case.
For large bias, $`\mu _L\mu _R\mathrm{\Gamma }`$, the integration over $`E_L`$ in Eqs. (A) can be performed analytically using Eqs. (A) for the amplitudes $`𝖻_{1,2}(t)`$. As a result we finally arrive at Eq. (26) obtained from Eq. (12) for the case of non-interacting electrons. This agreement with the case of non-interacting electrons is quite remarkable since our rate equations dealing with many-electron states are very different from those obtained in the single electron framework. Yet, this is not surprising since in the case of non-interacting electrons the single electron description is valid. In fact, Eqs. (A) can be mapped to Eq. (23) using $`|𝖻_i(t)|^2=\sigma _{ii}(t)+\sigma _{33}(t)`$, where $`i=1,2`$ and $`𝖻_1(t)𝖻_2^{}(t)=\sigma _{12}(t)`$.
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# An ISS Small-Gain Theorem for General Networks
#### Keywords
Interconnected systems – input-to-state stability – small-gain theorem – large-scale systems – monotone maps
#### MSC-classification:
93C10 (Primary) 34D05, 90B10, 93D09, 93D30 (Secondary)
## 1 Introduction
Stability is one of the fundamental concepts in the analysis and design of nonlinear dynamical systems. The notions of input-to-state stability (ISS) and nonlinear gains have proved to be an efficient tool for the qualitative description of stability of nonlinear input systems. There are different equivalent formulations of ISS: In terms of $`𝒦`$ and $`𝒦_{\mathrm{}}`$ functions (see below), via Lyapunov functions, as an asymptotic stability property combined with asymptotic gains, and others, see . A more quantitative but equivalent formulation, which captures the long term dynamic behavior of the system, is the notion of input-to-state dynamical stability (ISDS), see .
One of the interesting properties in the study of ISS systems is that under certain conditions input-to-state stability is preserved if ISS systems are connected in cascades or feedback loops. In this paper we generalize the existing results in this area. In particular, we obtain a general condition that guarantees input-to-state stability of a general system described as an interconnection of several ISS subsystems.
The earliest interconnection result on ISS systems states that cascades of ISS systems are again ISS, see e.g., . Furthermore, small-gain theorems for the case of two ISS systems in a feedback interconnection have been obtained in . These results state in one way or another that if the composition of the gain functions of ISS subsystems is smaller than the identity, then the whole system is ISS.
The papers use different approaches to the formulation of small-gain conditions that yield sufficient stability criteria: In the proof is based on the properties of $`𝒦`$ and $`𝒦_{\mathrm{}}`$ functions. This approach requires that the composition of the gains is smaller than the identity in a robust sense, see below for the precise statement. We show in Example 12 that within the context of this approach the robustness condition cannot be weakened. The result in that paper also covers practical ISS results, which we do not treat here. An ISS-Lyapunov function for the feedback system is constructed in as some combination of the corresponding ISS-Lyapunov functions of both subsystems. The key assumption of the proof in that paper is that the gains are already provided in terms of the Lyapunov functions, by which the authors need not resort to a robust version of the small-gain condition. The proof of the small-gain theorem in is based on the ISDS property and conditions for asymptotic stability of the feedback loop without inputs are derived. These results will turn out to be special cases of our main result.
General stability conditions for large scale interconnected systems have been obtained by various authors in other contexts. In sufficient conditions for the asymptotic stability of a composite system are stated in terms of the negative definiteness of some test matrix. This matrix is defined through the given Lyapunov functions of the interconnected subsystems. Similarly, in conditions for the stability of interconnected systems in terms of Lyapunov functions of the individual systems are obtained.
In Šiljak considers structural perturbations and their effects on the stability of composite systems using Lyapunov theory. The method is to reduce each subsystem to a one-dimensional one, such that the stability properties of the reduced aggregate representation imply the same stability properties of the original aggregate system. In some cases the aggregate representation gives rise to an interconnection matrix $`\overline{W}`$, such that quasi dominance or negative definiteness of $`\overline{W}`$ yield asymptotic stability of the composite system.
In small-gain type theorems for general interconnected systems with linear gains can be found. These results are of the form that the spectral radius of a gain matrix should be less than one to conclude stability. The result obtained here may be regarded as a nonlinear generalization in the same spirit.
In this paper we consider a system which consists of two or more ISS subsystems. We provide conditions by which the stability question of the overall system can be reduced to consideration of stability of the subsystems. We choose an approach using estimates involving $`𝒦`$ and $`𝒦_{\mathrm{}}`$ functions to prove the ISS stability result for general interconnected systems. The generalized small-gain condition we obtain is, that for some monotone operator $`\stackrel{~}{\mathrm{\Gamma }}`$ related to the gains of the individual systems the condition
$$\stackrel{~}{\mathrm{\Gamma }}(s)s$$
(1.1)
holds for all $`s0,s0`$ (in the sense of the component-wise ordering of the positive orthant). We discuss interpretations of this condition in Section 4.
Although we believe our approach to be amenable to the explicit construction of a Lyapunov function given the ISS-Lyapunov functions for the subsystems, so far we have been able to prove this only for linear gains.
While the general problem can be approached by repeated application of the cascade property and the known small-gain theorem, in general this can be cumbersome and it is by no means obvious in which order subsystems have to be chosen to proceed in such an iterative manner. Hence an extension of the known small-gain theorem to larger interconnections is needed.
In this paper we obtain this extension for the general case. Further, we show how to calculate the gain matrix for linear systems and give some interpretation of our result.
The paper is organized as follows. In Section 2 we introduce notation and necessary concepts and state the problem. In particular, we will need some basic properties of the positive orthant $`_+^n`$ interpreted as a lattice. In Section 3 we prove the main result, which generalizes the known small-gain theorem, and consider the special case of linear gains, for which we also construct an ISS-Lyapunov function. In Section 4 the small-gain condition of the main result is discussed and we show in which way it may be interpreted as an extension of the linear condition that the spectral radius of the gain matrix has to be less than one. There we also point out the connection to some induced monotone dynamical system. In Section 5 we show how the gain matrix can be found for linear systems. We conclude with Section 6.
## 2 Problem description
#### Notation
By $`x^T`$ we denote the transpose of a vector $`x^n`$. For $`x,y^n`$, we use the following notation
$$xyx_iy_i,i=1,\mathrm{},n,\text{ and }x>yx_i>y_i,i=1,\mathrm{},n.$$
(2.2)
In the following $`_+:=[0,\mathrm{})`$ and by $`_+^n`$ we denote $`\{x^n:x0\}`$. For a function $`v:_+^m`$ we define its restriction to the interval $`[s_1,s_2]`$ by
$$v_{[s_1,s_2]}(t):=\{\begin{array}{cc}v(t)\hfill & \text{ if }t[s_1,s_2],\hfill \\ 0\hfill & \text{ else.}\hfill \end{array}$$
###### Definition 1.
(i) A function $`\gamma :_+_+`$ is said to be of class $`𝒦`$ if it is continuous, increasing and $`\gamma (0)=0`$. It is of class $`𝒦_{\mathrm{}}`$ if, in addition, it is proper, i.e., unbounded.
(ii) A function $`\beta :_+\times _+_+`$ is said to be of class $`𝒦`$ if, for each fixed $`t`$, the function $`\beta (,t)`$ is of class $`𝒦`$ and, for each fixed $`s`$, the function $`\beta (s,)`$ is non-increasing and tends to zero for $`t\mathrm{}`$.
Let $`||`$ denote some norm in $`^n`$, and let in particular $`|x|_{\mathrm{max}}=\mathrm{max}_i|x_i|`$ be the maximum norm. The essential supremum norm on essentially bounded functions defined on $`_+`$ is denoted by $`_{\mathrm{}}`$.
###### Definition 2.
Consider a system
$$\dot{x}=f(x,u),x^n,u^m$$
such that for all initial values $`x_0`$ and all essentially bounded inputs $`u`$ unique solutions exist for all positive times. We denote these solutions by $`\xi (t;x_0,u)`$. The system is called input to state stable (ISS), if there exist functions $`\beta `$ of class $`𝒦`$ and $`\gamma `$ of class $`𝒦`$, such that the inequality
$$|\xi (t;x_0,u)|\beta (|x_0|,t)+\gamma (u_{\mathrm{}})$$
holds for all $`t0,x_0^n,u:_+^m`$ essentially bounded.
#### Problem statement
Consider $`n`$ interconnected control systems given by
$$\begin{array}{c}\dot{x}_1=f_1(x_1,\mathrm{},x_n,u)\\ \mathrm{}\\ \dot{x}_n=f_n(x_1,\mathrm{},x_n,u)\end{array}$$
(2.3)
where $`x_i^{N_i},u^L`$ and $`f_i:^{_{j=1}^nN_j+L}^{N_i}`$ is continuous and Lipschitz in the first $`n`$ arguments uniformly with respect to $`u`$ for $`i=1,\mathrm{},n`$. Here $`x_i`$ is the state of the $`i^{\text{th}}`$ subsystem, and $`u`$ is considered as an external control variable.
We may consider $`u`$ as partitioned $`u=(u_1,\mathrm{},u_n)`$, such that each $`u_i`$ is the input for subsystem $`i`$ only. Then each $`f_i`$ is of the form $`f_i(\mathrm{},u)=\stackrel{~}{f}_i(\mathrm{},P_i(u))=\stackrel{~}{f}_i(\mathrm{},u_i)`$ with some projection $`P_i`$. So without loss of generality we may assume to have the same input for all systems.
We call the $`i^{\text{th}}`$ subsystem of (2.3) ISS, if there exist functions $`\beta _i`$ of class $`𝒦`$ and $`\gamma _{ij},\gamma `$ of class $`𝒦`$, such that the solution $`x_i(t)`$ starting at $`x_i(0)`$ satisfies
$$|x_i(t)|\beta _i(|x_i(0)|,t)+\underset{j=1}{\overset{n}{}}\gamma _{ij}(x_{j}^{}{}_{[0,t]}{}^{}_{\mathrm{}})+\gamma (u_{\mathrm{}})$$
(2.4)
for all $`t0`$.
For notational simplicity we allow the case $`\gamma _{ij}0`$ and require $`\gamma _{ii}0`$ for all $`i`$. The functions $`\gamma _{ij}`$ and $`\gamma `$ are called (nonlinear) gains. We define $`\mathrm{\Gamma }:_+^n_+^n`$ by
$$\mathrm{\Gamma }:=(\gamma _{ij}),\mathrm{\Gamma }(s_1,\mathrm{},s_n)^T:=(\underset{j=1}{\overset{n}{}}\gamma _{1j}(s_j),\mathrm{},\underset{j=1}{\overset{n}{}}\gamma _{nj}(s_j))^T$$
(2.5)
for $`s=(s_1,\mathrm{},s_n)^T_+^n`$. We refer to $`\mathrm{\Gamma }`$ as the *gain matrix*, noting that it does not represent a linear map. Note that by the properties of $`\gamma _{ij}`$ for $`s_1,s_2_+^n`$ we have the implication
$$s_1s_2\mathrm{\Gamma }(s_1)\mathrm{\Gamma }(s_2),$$
(2.6)
so that $`\mathrm{\Gamma }`$ defines a monotone map.
Assuming each of the subsystems of (2.3) to be ISS, we are interested in conditions guaranteeing that the whole system defined by $`x=(x_1^T,\mathrm{},x_n^T)^T,f=(f_1^T,\mathrm{},f_n^T)^T`$ and
$$\dot{x}=f(x,u)$$
(2.7)
is ISS (from $`u`$ to $`x`$).
#### Additional Preliminaries
We also need some notation from lattice theory, cf. for example. Although $`(_+^n,sup,inf)`$ is a lattice, with $`inf`$ denoting infimum and $`sup`$ denoting supremum, it is not complete. But still one can define the upper limit for bounded functions $`s:_+_+^n`$ by
$$\underset{t\mathrm{}}{lim\; sup}s(t):=\underset{t0}{inf}\underset{\tau t}{sup}s(\tau ).$$
For vector functions $`x=(x_1^T,\mathrm{},x_n^T)^T:_+^{N_1+\mathrm{}+N_n}`$ such that $`x_i:_+^{N_i},i=1,\mathrm{},n`$ and times $`0t_1t_2`$ we define
$$x_{[t_1,t_2]}:=\left(\begin{array}{c}x_{1,[t_1,t_2]}_{\mathrm{}}\\ \mathrm{}\\ x_{n,[t_1,t_2]}_{\mathrm{}}\end{array}\right)_+^n.$$
We will need the following property.
###### Lemma 3.
Let $`s:_+_+^n`$ be continuous and bounded. Then (setting $`N_i1`$)
$$\underset{t\mathrm{}}{lim\; sup}s(t)=\underset{t\mathrm{}}{lim\; sup}s_{[t/2,\mathrm{})}.$$
###### Proof.
Let $`lim\; sup_t\mathrm{}s(t)=:a_+^n`$ and $`lim\; sup_t\mathrm{}||s_{[t/2,\mathrm{})}||=:b_+^n`$. For every $`\epsilon _+^n,`$ $`\epsilon >0`$ (component-wise!) there exist $`t_a,t_b0`$ such that
$`tt_a:\underset{tt_a}{sup}s(t)a+\epsilon \text{and}tt_b:\underset{tt_b}{sup}s_{[t/2,\mathrm{})}b+\epsilon .`$ (2.8)
Clearly we have
$$s(t)s_{[t/2,\mathrm{})}$$
for all $`t0`$, i.e., $`ab`$. On the other hand $`s(\tau )a+\epsilon `$ for $`\tau t`$ implies $`s_{[\tau /2,\mathrm{})}a+\epsilon `$ for $`\tau 2t`$, i.e., $`ba`$. This immediately gives $`a=b`$, and the claim is proved. ∎
Before we introduce the ISS criterion for interconnected systems let us briefly discuss an equivalent formulation of ISS. A system
$$\dot{x}=f(x,u),$$
(2.9)
with $`f:^{N+L}^N`$ continuous and Lipschitz in $`x^N`$, uniformly with respect to $`u^L`$, is said to have the *asymptotic gain property* (AG), if there exists a function $`\gamma _{AG}𝒦_{\mathrm{}}`$ such that for all initial values $`x_0^N`$ and all essentially bounded control functions $`u():_+^L`$,
$$\underset{t0}{lim\; sup}|x(t;x_0,u)|\gamma _{AG}(u_{\mathrm{}}).$$
(2.10)
The asymptotic gain property states, that every trajectory must ultimately stay not far from zero, depending on the magnitude of $`u_{\mathrm{}}`$.
The system (2.9) is said to be *globally asymptotically stable at zero* (0-GAS), if there exists a $`\beta _{GAS}𝒦`$, such that for all initial conditions $`x_0^N`$
$$|x(t;x_0,0)|\beta _{GAS}(|x_0|,t).$$
(2.11)
Thus 0-GAS holds, if, when the input $`u`$ is set to zero, the system (2.9) is globally asymptotically stable at $`x^{}=0`$.
By a result of Sontag and Wang the asymptotic gain property and global asymptotic stability at 0 together are equivalent to ISS.
## 3 Main results
In the following subsection we present a nonlinear version of the small-gain theorem for networks. In Subsection 3.2 we restate this theorem for the case when the gains are linear functions. Here we also provide a method on how to construct an ISS-Lyapunov function for the whole network system from given ISS-Lyapunov functions of the subsystems.
### 3.1 Nonlinear gains
We introduce the following notation. For $`\alpha _i𝒦_{\mathrm{}},i=1,\mathrm{},n`$ define $`D:_+^n_+^n`$ by
$$D(s_1,\mathrm{},s_n)^T:=\left(\begin{array}{c}(\text{Id}+\alpha _1)(s_1)\\ \mathrm{}\\ (\text{Id}+\alpha _n)(s_n)\end{array}\right).$$
(3.12)
###### Theorem 4 (small-gain theorem for networks).
Consider the system (2.3) and suppose that each subsystem is ISS, i.e., condition (2.4) holds for all $`i=1,\mathrm{},n`$. Let $`\mathrm{\Gamma }`$ be given by (2.5). If there exists a mapping $`D`$ as in (3.12), such that
$$(\mathrm{\Gamma }D)(s)s,s_+^n\{0\},$$
(3.13)
then the system (2.7) is ISS from $`u`$ to $`x`$.
###### Remark 5.
Although looking very complicated to handle at first sight, condition (3.13) is a straightforward extension of the ISS small-gain theorem of . It has many interesting interpretations, as we will discuss in Section 4.
The following lemma provides an essential argument in the proof of Theorem 4.
###### Lemma 6.
Let $`D`$ be as in (3.12) and suppose (3.13) holds. Then there exists a $`\phi 𝒦_{\mathrm{}}`$ such that for all $`w,v_+^n`$,
$$(\text{Id}\mathrm{\Gamma })(w)v$$
(3.14)
implies $`|w|\phi (|v|)`$.
###### Proof.
Fix $`v_+^n`$. We first show, that for those $`w_+^n`$ satisfying (3.14) at least some components have to be bounded. To this end let
$`r^{}:=(D\text{Id})^1(v)=\left(\begin{array}{c}\alpha _1^1(v_1)\\ \mathrm{}\\ \alpha _n^1(v_n)\end{array}\right)`$ (3.15)
$`\text{and }s^{}:=D(r^{})=\left(\begin{array}{c}v_1+\alpha _1^1(v_1)\\ \mathrm{}\\ v_n+\alpha _n^1(v_n)\end{array}\right).`$
We claim that $`ss^{}`$ implies that $`w=s`$ does not satisfy (3.14). So let $`ss^{}`$ be arbitrary and $`r=D^1(s)r^{}`$ (as $`D^1𝒦_{\mathrm{}}^n`$). For such $`s`$ we have
$$sD^1(s)=D(r)rD(r^{})r^{}=v,$$
where we have used that $`(D\text{Id})𝒦_{\mathrm{}}^n`$. The assumption that $`w=s`$ satisfies (3.14) leads to
$$sv+\mathrm{\Gamma }(s)sD^1(s)+\mathrm{\Gamma }(s),$$
or equivalently, $`0\mathrm{\Gamma }(s)D^1(s)`$. This implies for $`r=D^1(s)`$ that
$$r\mathrm{\Gamma }D(r),$$
in contradiction to (3.13). This shows that the set of $`w_+^n`$ satisfying (3.14) does not intersect the set
$$Z_1:=\{w_+^n|ws^{}\}.$$
Assume now that $`w_+^n`$ satisfies (3.14). Let $`s^1:=s^{}`$. If $`s^1w`$, then there exists an index set $`I_1\{1,\mathrm{},n\}`$, such that
$`w_i>s_i^1,\text{for}iI_1\text{and}w_is_i^1,`$ $`\text{for}iI_1^c:=\{1,\mathrm{},n\}I_1.`$
For index sets $`I`$ and $`J`$ denote by $`y_I`$ the restriction
$$y_I:=(y_i)_{iI}$$
for vectors $`y_+^n`$ and by $`A_{IJ}:_+^{\mathrm{\#}I}_+^{\mathrm{\#}J}`$ the restriction
$$A_{IJ}:=(a_{ij})_{iI,jJ}$$
for mappings $`A=(a_{ij})_{i,j\{1,\mathrm{},n\}}:_+^n_+^n`$.
So from (3.14) we obtain
$$\left[\begin{array}{c}w_{I_1}\\ w_{I_1^c}\end{array}\right]\left[\begin{array}{cc}\mathrm{\Gamma }_{I_1I_1}& \mathrm{\Gamma }_{I_1I_1^c}\\ \mathrm{\Gamma }_{I_1^cI_1}& \mathrm{\Gamma }_{I_1^cI_1^c}\end{array}\right]\left(\left[\begin{array}{c}w_{I_1}\\ w_{I_1^c}\end{array}\right]\right)\left[\begin{array}{c}v_{I_1}\\ v_{I_1^c}\end{array}\right].$$
Hence we have in particular
$$\begin{array}{cc}\hfill w_{I_1}& \mathrm{\Gamma }_{I_1I_1}(w_{I_1})v_{I_1}+\mathrm{\Gamma }_{I_1I_1^c}(s_{I_1^c}^1)\hfill \\ & \underset{>\text{Id}}{\underset{}{D_{I_1}(D_{I_1}\text{Id}_{I_1})^1}}(v_{I_1}+\mathrm{\Gamma }_{I_1I_1^c}(s_{I_1^c}^1))=:s_{I_1}^2.\hfill \end{array}$$
(3.16)
Note that $`\mathrm{\Gamma }_{I_1I_1}`$ satisfies (3.13) with $`D`$ replaced by $`D_{I_1}`$. Thus, arguing just as before, we obtain, that $`w_{I_1}s_{I_1}^2`$ is not possible. Hence some more components of $`w`$ must be bounded.
We proceed inductively, defining
$$I_{j+1}I_j,I_{j+1}:=\{iI_j:w_i>s_i^{j+1}\},$$
with $`I_{j+1}^c:=\{1,\mathrm{},n\}I_{j+1}`$ and
$$s_{I_j}^{j+1}:=D_{I_j}(D_{I_j}\text{Id}_{I_j})^1(v_{I_j}+\mathrm{\Gamma }_{I_jI_j^c}(s_{I_j^c}^j)).$$
Obviously this nesting will end after at most $`n1`$ steps: There exists a maximal $`kn`$, such that
$$\{1,\mathrm{},n\}I_1\mathrm{}I_k\mathrm{}$$
and all components of $`w_{I_k}`$ are bounded by the corresponding components of $`s_{I_k}^{k+1}`$. For $`i=1,\mathrm{},n`$ define
$$\zeta _i:=\mathrm{max}\{j\{1,\mathrm{},n\}:iI_j\}$$
and
$$s_\zeta :=(s_1^{\zeta _1},\mathrm{},s_n^{\zeta _n}).$$
Clearly we have
$$ws_\zeta [D(D\text{Id})^1(\text{Id}+\mathrm{\Gamma })]^n(v)$$
and the term on the very right hand side does not depend on any particular choice of nesting of the index sets. Hence every $`w`$ satisfying (3.14) also satisfies
$$w[D(D\text{Id})^1(\text{Id}+\mathrm{\Gamma })]^n\left(\begin{array}{ccc}|v|_{\mathrm{max}},& \mathrm{},& |v|_{\mathrm{max}}\end{array}\right)^T$$
and taking the $`\mathrm{max}`$-norm on both sides yields
$$|w|_{\mathrm{max}}\phi (|v|_{\mathrm{max}})$$
for some function $`\phi `$ of class $`𝒦_{\mathrm{}}`$. This completes the proof of the lemma. ∎
We proceed with the proof of Theorem 4, which is divided into two main steps. First we establish the existence of a solution of the system (2.7) for all times $`t0`$. In the second step we establish the ISS property for this system.
###### Proof.
(of Theorem 4) *Existence of a solution for (2.7) for all times:* For finite times $`t0`$ and for $`s_+^n`$ we introduce the abbreviating notation
$$x(t):=\left(\begin{array}{c}|x_1(t)|\\ \mathrm{}\\ |x_n(t)|\end{array}\right)_+^n,\gamma ^n(u_{\mathrm{}}):=\left(\begin{array}{c}\gamma (u_{\mathrm{}})\\ \mathrm{}\\ \gamma (u_{\mathrm{}})\end{array}\right)_+^n$$
(3.23)
$$\text{and}\beta (s,t):=\left(\begin{array}{c}\beta _1(s_1,t)\\ \mathrm{}\\ \beta _n(s_n,t)\end{array}\right):_+^n\times _+_+^n.$$
(3.27)
Now we can rewrite the ISS conditions (2.4) of the subsystems in a vectorized form for $`\tau 0`$ as
$$x(\tau )\beta (x(0),\tau )+\mathrm{\Gamma }(x_{[0,\tau ]})+\gamma ^n(u_{\mathrm{}})$$
(3.28)
and taking the supremum on both sides over $`\tau [0,t]`$ we obtain
$`(\text{Id}\mathrm{\Gamma })x_{[0,t]}`$ $`=`$ $`x_{[0,t]}\mathrm{\Gamma }(x_{[0,t]})`$ (3.29)
$``$ $`\beta (x(0),0)+\gamma ^n(u_{\mathrm{}})`$
where we used (2.6). Now by Lemma 6 we find
$$x_{[0,t]}_{\mathrm{}}\phi \left(\right|\beta (||x(0)||,0)+\gamma ^n(u_{\mathrm{}})\left|\right)=:s_{\mathrm{}}$$
(3.30)
for some class $`𝒦`$ function $`\phi `$ and all times $`t0`$. Hence for every initial condition and essentially bounded input $`u`$ the solution of our system (2.7) exists for all times $`t0`$, since $`s_{\mathrm{}}`$ in (3.30) does not depend on $`t`$.
*Establishing ISS:* We now utilize an idea from : Instead of estimating $`|x_i(t)|`$ with respect to $`|x_i(0)|`$ in (2.4), we can also have the point of view that our trajectory started in $`x_i(\tau )`$ at time $`0\tau t`$ and we followed it for some time $`t\tau `$ and reach $`x_i(t)`$ at time $`t`$. For $`\tau =t/2`$ this reads
$`|x_i(t)|`$ $``$ $`\beta _i(|x_i(t/2)|,t/2)+{\displaystyle \underset{ji}{}}\gamma _{ij}(x_{i,[t/2,t]}_{\mathrm{}})+\gamma (u)`$ (3.31)
$``$ $`\beta _i(s_{\mathrm{}},t/2)+{\displaystyle \underset{ji}{}}\gamma _{ij}(x_{i,[t/2,\mathrm{})}_{\mathrm{}})+\gamma (u)`$
$`=`$ $`\stackrel{~}{\beta }_i(s_{\mathrm{}},t)+{\displaystyle \underset{ji}{}}\gamma _{ij}(x_{i,[t/2,\mathrm{})}_{\mathrm{}})+\gamma (u)`$ (3.32)
where we again applied (2.6) to obtain (3.31) and defined
$$\stackrel{~}{\beta }_i(s_i,t):=\beta _i(s_i,t/2),$$
which is of class $`𝒦`$.
To write inequality (3.32) in vector form, we define
$$\stackrel{~}{\beta }(s,t):=\left(\begin{array}{c}\stackrel{~}{\beta }_1(s_1,t)\\ \mathrm{}\\ \stackrel{~}{\beta }_n(s_n,t)\end{array}\right)$$
(3.33)
for all $`s_+^n`$. Denoting by $`s_{\mathrm{}}^n:=(s_{\mathrm{}},\mathrm{},s_{\mathrm{}})^T`$ we obtain the vector formulation of (3.32) as
$`x(t)\stackrel{~}{\beta }(s_{\mathrm{}}^n,t)+\mathrm{\Gamma }x_{[t/2,\mathrm{})}+\gamma ^n(u_{\mathrm{}}).`$ (3.34)
By the boundedness of the solution we can take the upper limit on both sides of (3.34). By Lemma 3 we have
$$\underset{t\mathrm{}}{lim\; sup}||x(t)||=\underset{t\mathrm{}}{lim\; sup}||x_{[t/2,\mathrm{})}||=:l(x),$$
and it follows that
$$(\text{Id}\mathrm{\Gamma })l(x)\gamma ^n(u_{\mathrm{}})$$
since $`lim_t\mathrm{}\stackrel{~}{\beta }(s_{\mathrm{}}^n,t)=0`$. Finally, by Lemma 6 we have
$$|l(x)|\phi (|\gamma ^n(u_{\mathrm{}})|)$$
(3.35)
for some $`\phi `$ of class $`𝒦_{\mathrm{}}`$. But (3.35) is the asymptotic gain property (2.10).
Now $`0`$-GAS is established as follows: First note that for $`u0`$ the quantity $`s_{\mathrm{}}`$ in (3.30) is a $`𝒦`$ function of $`|x(0)|`$. So (3.30) shows (Lyapunov) stability of the system in the case $`u0`$. Furthermore, (3.35) shows attractivity of $`x=0`$ for the system (2.7) in the case $`u0`$. This shows global asymptotic stability of $`x=0`$.
Hence system (2.7) is AG and 0-GAS, which together were proved to be equivalent to ISS in \[15, Theorem 1\]. ∎
### 3.2 Linear gains and an ISS-Lyapunov version
Suppose the gain functions $`\gamma _{ij}`$ are all linear, hence $`\mathrm{\Gamma }`$ is a linear mapping and (2.5) is just matrix-vector multiplication. Then we have the following
###### Corollary 7.
Consider $`n`$ interconnected ISS systems as in the previous section on the problem description with a linear
gain matrix $`\mathrm{\Gamma }`$, such that for the spectral radius $`\rho `$ of $`\mathrm{\Gamma }`$ we have
$$\rho (\mathrm{\Gamma })<1.$$
(3.36)
Then the system defined by (2.7) is ISS from $`u`$ to $`x`$.
###### Remark 8.
For non-negative matrices $`\mathrm{\Gamma }`$
it is well known that (see, e.g., \[1, Theorem 2.1.1, page 26, and Theorem 2.1.11, page 28\])
1. $`\rho (\mathrm{\Gamma })`$ is an eigenvalue of $`\mathrm{\Gamma }`$ and $`\mathrm{\Gamma }`$ possesses a non-negative eigenvector corresponding to $`\rho (\mathrm{\Gamma })`$,
2. $`\alpha x\mathrm{\Gamma }x`$ holds for some $`x_+^n\{0\}`$ if and only if $`\alpha \rho (\mathrm{\Gamma })`$.
Hence $`\rho (\mathrm{\Gamma })<1`$ if and only if $`\mathrm{\Gamma }ss`$ for all $`s_+^n\{0\}`$.
Also, by continuity of the spectrum it is clear that for such $`\mathrm{\Gamma },\rho (\mathrm{\Gamma })<1,`$ there always exists a matrix $`D=\text{diag}(1+\alpha _1,\mathrm{},1+\alpha _n)`$ with $`\alpha _i>0,i=1,\mathrm{},n`$, such that $`\mathrm{\Gamma }Dss`$ for all $`s_+^n\{0\}`$.
###### Remark 9.
For the case of large-scale interconnected input-output systems a similar result exists, which can be found in a monograph by Vidyasagar, cf. \[17, p. 110\]. It also covers Corollary 7 as a special case. The condition on the spectral radius is quite the same, although it is applied to a *test matrix*, whose entries are finite gains of products of *interconnection operators* and corresponding *subsystem operators*. These gains are non-negative numbers and, roughly speaking, defined as the minimal possible slope of affine bounds on the interconnection operators.
###### Proof.
(of Corollary 7) The proof is essentially the same as of Theorem 4, but note that instead of Lemma 6 we now directly have existence of
$$(\text{Id}\mathrm{\Gamma })^1=\text{Id}+\mathrm{\Gamma }+\mathrm{\Gamma }^2+\mathrm{}$$
since $`\rho (\mathrm{\Gamma })<1`$ and from the power sum expansion it is obvious that $`(\text{Id}\mathrm{\Gamma })^1`$ is a non-decreasing mapping, i.e., for $`d_1,d_20`$ we have $`(\text{Id}\mathrm{\Gamma })^1(d_1+d_2)(\text{Id}\mathrm{\Gamma })^1(d_1)0`$.
Thus at the two places where Lemma 6 has been used we can simply apply $`(\text{Id}\mathrm{\Gamma })^1`$ to get the desired estimates. ∎
#### Construction of an ISS-Lyapunov function
There is another approach to describe the ISS property via so called ISS-Lyapunov functions, cf. .
###### Definition 10.
A smooth function $`V`$ is said to be an ISS-Lyapunov function of the system (2.9) $`\dot{x}=f(x,u),f:^{N+L}^N`$ if
1. $`V`$ is proper, positive-definite, that is, there exit functions $`\psi _1,\psi _2`$ of class $`𝒦_{\mathrm{}}`$ such that
$$\psi _1(|x|)V(x)\psi _2(|x|),x^{n_N};$$
(3.37)
2. there exists a positive-definite function $`\alpha `$, a class $`𝒦`$-function $`\chi `$, such that
$$V(x)\chi (|u|)V(x)f(x,u)\alpha (|x|).$$
(3.38)
We call the function $`\chi `$ the *Lyapunov-gain*.
In case of linear Lyapunov-gains a Lyapunov function for the interconnected system can be constructed, given ISS-Lyapunov functions of the subsystems. Note, that this time we define the gain matrix $`\mathrm{\Gamma }`$ with respect to the Lyapunov-gains $`\gamma _{ij}`$.
Let $`V_1(x_1),\mathrm{},V_n(x_n)`$ be some ISS-Lyapunov functions of the subsystems (2.3), allowing for linear Lyapunov-gains $`\gamma _{ij}`$, i.e., there are some $`𝒦_{\mathrm{}}`$ functions $`\psi _{i1},\psi _{i2}`$ such that
$$\psi _{i1}(|x_i|)V_i(x_i)\psi _{i2}(|x_i|),x_i^{N_i},$$
(3.39)
and some positive-definite functions $`\alpha _i`$ such that
$$V_i(x_i)>\mathrm{max}\{\underset{j}{\mathrm{max}}\{\gamma _{ij}V_j(x_j)\},\gamma _i(|u|)\}V_i(x_i)f_i(x_,u)\alpha _i(V_i(x_i)).$$
(3.40)
Consider the positive orthant $`_+^n`$, and let $`\mathrm{\Omega }_i`$ be the subsets of $`_+^n`$ defined by
$$\mathrm{\Omega }_i:=\{(v_1,\mathrm{},v_n)_+^n:v_i>\underset{j=1}{\overset{n}{}}\gamma _{ij}v_j\}.$$
(3.41)
Note that the boundaries $`\mathrm{\Omega }_i`$ are hyperplanes in case of linear gains. Now if (3.36) or equivalently $`\mathrm{\Gamma }ss,s_+^n,s0`$, holds for $`\mathrm{\Gamma }=(\gamma _{ij}),i,j=1,\mathrm{},n`$, then it follows that
$$\underset{i=1}{\overset{n}{}}\mathrm{\Omega }_i=_+^n\{0\}\text{and}\underset{i=1}{\overset{n}{}}\mathrm{\Omega }_i\mathrm{}.$$
(3.42)
The proof is the same as of Proposition 21, see below. Thus we may choose an $`s>0`$ with $`s_{i=1}^n\mathrm{\Omega }_i,`$ which implies that
$$s_i>\underset{j}{}\gamma _{ij}s_j,i=1,\mathrm{},n;$$
(3.43)
see Fig. 1. If $`\mathrm{\Gamma }`$ is irreducible, then using Perron-Frobenius theory we see that we may choose $`s`$ to be a (positive) eigenvector of $`\mathrm{\Gamma }`$.
###### Theorem 11.
Let $`V_i`$ be an ISS-Lyapunov function as in (3.40) of the $`i^{th}`$ subsystem from (2.3), $`i=1,\mathrm{},n,`$ and $`s`$ be a positive vector with (3.43). Then an ISS Lyapunov function of the interconnected system (2.3) is given by
$$V(x_1,\mathrm{},x_n):=\underset{i}{\mathrm{max}}\frac{V_i(x_i)}{s_i}.$$
(3.44)
###### Proof.
Let $`\gamma (|u|):=\mathrm{max}_i\gamma _i(|u|)`$ which is a $`𝒦`$ class function. In the following we show that there exists a positive definite function $`\alpha `$ such that:
$$V(x)\gamma (|u|)V(x)f(x,u)\alpha (V(x)).$$
(3.45)
Let $`M_i`$ be open domains in $`_+^n`$ defined by
$$M_i:=\{(v_1,\mathrm{},v_n)_+^n:\frac{v_i}{s_i}>\underset{ji}{\mathrm{max}}\left\{\frac{v_j}{s_j}\right\}\},$$
(3.46)
and let $`P_i`$ be the 2-dimensional planes spanned by $`s`$ and the $`i`$-th axis, i.e.,
$$P_i=\left\{v_+^n\right|\frac{v_k}{s_k}=\frac{v_j}{s_j};k,ji\}.$$
(3.47)
Note that $`V`$ defined by (3.44) is continuous in $`_+^n`$ and can only fail to be differentiable on the planes $`P_i`$.
Now take any $`\widehat{x}=(\widehat{x}_1,\mathrm{},\widehat{x}_n)^n`$ with $`(V_1(\widehat{x}_1),\mathrm{},V_n(\widehat{x}_n))M_i`$ then it follows that in some neighborhood $`U`$ of $`\widehat{x}`$ we have $`V(x)=\frac{V_i(x_i)}{s_i}`$ for all $`xU`$ and
$$V_i(x_i)>\underset{ji}{\mathrm{max}}\left\{\frac{s_i}{s_j}V_j(x_j)\right\}>\underset{ji}{\mathrm{max}}\{\gamma _{ij}V_j(x_j)\}$$
(3.48)
(the last inequality follows from(3.43)), hence by (3.40), if $`V(x)=V_i(x_i)/s_i>\gamma _i(|u|)`$, then
$$V(x)f(x,u)=\frac{1}{s_i}V_i(x_i)f_i(x,u)\frac{1}{s_i}\alpha _i(V_i(x_i))<\stackrel{~}{\alpha _i}(V(x)),$$
(3.49)
where $`\stackrel{~}{\alpha _i}`$ are positive-definite functions, since $`s_i=const>0.`$
It remains to consider $`x^n`$ such that $`(V_1(x_1),\mathrm{},V_n(x_n))\overline{M_i}\overline{M_j}`$, where $`V(x)`$ may be not differentiable.
For this purpose we use some results from . For smooth functions $`f_i,i=1,\mathrm{},n`$ it follows that $`f(x,u)=\underset{i}{\mathrm{max}}\{f_i(x,u)\}`$ is Lipschitz and Clarke’s subgradient of $`f`$ is given by
$$_{Cl}f(x)=co\left\{\underset{iM(x)}{}_xf_i(x,u)\right\},M(x)=\{i:f_i(x,u)=f(x)\},$$
(3.50)
i.e., in our case
$$_{Cl}V(x)=co\{\frac{1}{s_i}V_i(x):\frac{1}{s_i}V_i(x)=V(x)\}.$$
(3.51)
Now for every extremal point of $`_{Cl}V(x)`$ a decrease condition is satisfied by (3.49). By convexity, the same is true for every element of $`_{Cl}V(x)`$. Now Theorems 4.3.8 and 4.5.5 of show strong invariance and attractivity of the set $`\{x:V(x)\gamma (u)\}`$. It follows that $`V`$ is an ISS-Lyapunov function for the interconnection (2.3). ∎
See also Section 4.4 for some more considerations into this directions.
## 4 Interpretation of the generalized small-gain condition
In this section we wish to provide insight into the small-gain condition of Theorem 4. We first show, that the result covers the known interconnection results for cascades and feedback interconnections. We then compare the condition with the linear case.
Further we state some algebraic and graph theoretical relations and investigate some associated artificial dynamical system induced by the gain matrix $`\mathrm{\Gamma }`$. We complete this section with some geometrical considerations and an overview map of all these contiguities.
### 4.1 Connections to known results
As an easy consequence of Theorem 4 we recover, that an arbitrarily long feed forward cascade of ISS subsystems is ISS again. If the subsystems are enumerated consecutively and the gain function from subsystem $`j`$ to subsystem $`i>j`$ is denoted by $`\gamma _{ij}`$, then the resulting gain matrix has non-zero entries only below the diagonal. For arbitrary $`\alpha 𝒦_{\mathrm{}}`$ the gain matrix with entries $`\gamma _{ij}(\text{Id}__++\alpha )`$ for $`i>j`$ and $`0`$ for $`ij`$ clearly satisfies (3.13). Therefore the feed forward cascade itself is ISS.
Consider $`n=2`$ in equation (2.3), i.e., two subsystems with linear gains. Then in Corollary 7 we have
$$\mathrm{\Gamma }=\left[\begin{array}{cc}0& \gamma _{12}\\ \gamma _{21}& 0\end{array}\right],\gamma _{ij}_+$$
and $`\rho (\mathrm{\Gamma })<1`$ if and only if $`\gamma _{12}\gamma _{21}<1.`$ Hence we obtain the known small-gain theorem, cf. and .
For nonlinear gains and $`n=2`$ the condition (3.13) in Theorem 4 reads as follows: There exist $`\alpha _1,\alpha _2𝒦_{\mathrm{}}`$ such that
$$\left(\begin{array}{c}\gamma _{12}(\text{Id}+\alpha _2)(s_2)\\ \gamma _{21}(\text{Id}+\alpha _1)(s_1)\end{array}\right)\left(\begin{array}{c}s_1\\ s_2\end{array}\right),$$
for all $`(s_1,s_2)^T_+^2`$. This is easily seen to be equivalent to
$$\gamma _{12}(\text{Id}+\alpha _2)\gamma _{21}(\text{Id}+\alpha _1)(s)<s,s>0.$$
To this end it suffices to check what happens for to the vector $`[\gamma _{12}(\text{Id}+\alpha _2)(s_2),s_2]^T`$ under $`\mathrm{\Gamma }`$ along with a few similar considerations. The latter is equivalent to the condition in the small-gain theorem of , namely, that for some $`\stackrel{~}{\alpha }_1,\stackrel{~}{\alpha }_2𝒦_{\mathrm{}}`$ it should hold that
$$(\text{Id}+\stackrel{~}{\alpha }_1)\gamma _{21}(\text{Id}+\stackrel{~}{\alpha }_2)\gamma _{12}(s)s,s>0,$$
(4.52)
for all $`s_+`$, hence our theorem contains this result as a particular case.
###### Example 12.
The condition (4.52) of seems to be very similar to the small-gain condition $`\gamma _{12}\gamma _{21}(s)<s`$ of and , however those $`\gamma `$’s have some different meanings in these papers. This similarity raises the question, whether the compositions with $`(\text{Id}+\stackrel{~}{\alpha }_i),i=1,2`$ in (4.52) or more generally with $`D`$ in (3.13) is necessary. The answer is positive. Namely, there is a counterexample providing a system of two ISS subsystems with $`\gamma _{12}\gamma _{21}(s)<s`$ ($`\gamma `$’s are defined as above) which is not ISS.
Consider the equation
$$\dot{x}=x+u(1e^u),x(0)=x^0,u.$$
Integrating it follows
$$x(t)=e^tx^0+_0^te^{(t\tau )}u(\tau )(1e^{u(\tau )})𝑑\tau $$
$$e^tx^0+u_{\mathrm{}}(1e^u_{\mathrm{}})=e^tx^0+\gamma (u_{\mathrm{}}),\gamma (s)<s.$$
Then for a feedback system
$`\dot{x}_1`$ $`=`$ $`x_1+x_2(1e^{x_2})+u(t),`$ (4.53)
$`\dot{x}_2`$ $`=`$ $`x_2+x_1(1e^{x_1})+u(t)`$ (4.54)
we have ISS for each subsystem with $`x_i(t)e^tx_i^0+\gamma _i(x_i)+\eta _i(u),`$ where $`\gamma _i(s)<s`$ and hence $`\gamma _1\gamma _2(s)<s`$ for $`s>0`$, but there is a solution $`x_1=x_2=const`$, i.e.,
$$\dot{x}_1=x_2e^{x_2}+u,\text{with}u=x_2e^{x_2},$$
and $`x_1=x_2`$ can be chosen arbitrary large with $`u0`$ for $`x_1\mathrm{}.`$ Hence the condition $`\mathrm{\Gamma }(s)s`$, for all $`s_+^n\{0\}`$, or for two subsystems $`\gamma _{12}\gamma _{21}(s)<s`$, for all $`s>0`$, is not sufficient for the input-to-state stability of the composite system in the nonlinear case.
### 4.2 Algebraic Interpretation
In this subsection we first relate the network small-gain condition (3.13) to well known properties of matrices in the linear case. This gives some idea how the new condition can be understood and what subtle differences appear in the nonlinear case. Then we extend some graph theoretical results for non-negative matrices to nonlinear gain matrices. These are needed later on.
For a start, we discuss some algebraic consequences from (3.13). Recall that for a non-negative matrix $`\mathrm{\Gamma }`$ the following are equivalent:
1. $`\rho (\mathrm{\Gamma })<1`$,
2. $`s_+^n\{0\}:\mathrm{\Gamma }ss`$,
3. $`\mathrm{\Gamma }^k0`$, for $`k\mathrm{}`$,
4. there exist $`a_1,\mathrm{},a_n>0`$ such that $`s_+^n\{0\}`$:
$$\mathrm{\Gamma }(I+\mathrm{diag}(a_1,\mathrm{},a_n))ss.$$
Note that (iv) is the linear version of (3.13). As condition (i) is not useful in the nonlinear setting, we have turned to (ii), which we later strengthened to (3.13).
In the nonlinear case we find the obvious implication:
###### Proposition 13.
Condition (3.13) implies that
$$\mathrm{\Gamma }(s)s\text{for any}s_+^n\{0\}.$$
(4.55)
###### Proof.
By the monotonicity of $`\mathrm{\Gamma }`$ it is obvious that (3.13) implies (4.55). ∎
Note that the contrary is not true:
###### Example 14.
Let
$$\gamma _{12}=\text{Id}__+$$
and
$$\gamma _{21}(r)=r(1e^r).$$
Since already $`lim_r\mathrm{}(\gamma _{12}\gamma _{21}\text{Id})(r)=0`$ there are certainly no class $`𝒦_{\mathrm{}}`$ functions $`\stackrel{~}{\alpha }_i,i=1,2`$ such that (4.52) holds.
###### Remark 15.
We like to point out the connections between non-negative matrices, our gain matrix and directed graphs.
A (finite) directed graph $`G=\{V,E\}`$ consists of a set $`V`$ of vertices and a set of edges $`EV\times V`$. We may identify $`V=\{1,\mathrm{},n\}`$ in case of $`n`$ vertices. The adjacency matrix $`A_G=(a_{ij})`$ of this graph is defined by
$$a_{ij}=\{\begin{array}{cc}1\hfill & \text{if}(i,j)E,\hfill \\ 0\hfill & else.\hfill \end{array}$$
The other way round, given an $`n\times n`$-matrix $`A`$, one defines the graph $`G(A)=\{V,E\}`$ by $`V:=\{1,\mathrm{},n\}`$ and $`E=\{(i,j)V\times V:a_{ij}0\}`$.
There are several concepts and results of (non-negative) matrix theory, which are of purely graph theoretical nature. Hence the same can be done for our interconnection gain matrix $`\mathrm{\Gamma }`$. We may associate a graph $`G(\mathrm{\Gamma })`$, which represents the interconnections between the subsystems, in the same manner, as we would do for matrices.
We could also use the graph of the transpose of $`\mathrm{\Gamma }`$ here for compatibility with our previous notation ($`\gamma _{ij}`$ encodes whether or not subsystem $`j`$ influences subsystem $`i`$) and the standard notation in graph theory (edge from $`i`$ to $`j`$), then the arrows in $`G(\mathrm{\Gamma })`$ would point in the ‘right’ direction. But this does not affect the following results.
For instance, we say $`\mathrm{\Gamma }`$ is *irreducible*, if $`G(\mathrm{\Gamma })`$ is *strongly connected*, that is, for every pair of vertices $`(i,j)`$ there exists a sequence of edges (a *path*) connecting vertex $`i`$ to vertex $`j`$. Obviously $`\mathrm{\Gamma }`$ is irreducible if and only if $`\mathrm{\Gamma }^T`$ is. $`\mathrm{\Gamma }`$ is called *reducible* if it is not irreducible.
The gain matrix $`\mathrm{\Gamma }`$ is *primitive*, if its associated graph $`G_\mathrm{\Gamma }:=G(\mathrm{\Gamma })`$ is primitive, i.e., there exists a positive integer $`m`$ such that $`(A_{G_\mathrm{\Gamma }})^m`$ has only positive entries.
These definitions and the following important facts can be found in and only depend on the associated graph.
If $`\mathrm{\Gamma }`$ is reducible, then a permutation transforms it into a block upper triangular matrix. From an interconnection point of view, this splits the system into cascades of subsystems each with irreducible adjacency matrix.
###### Lemma 16.
Assume the gain matrix $`\mathrm{\Gamma }`$ is irreducible. Then there are two distinct cases:
1. The gain matrix $`\mathrm{\Gamma }=(\gamma _{ij}()),`$ where $`\gamma _{ij}()𝒦`$ or $`\gamma _{ij}=0`$, is primitive and hence there is a non-negative integer $`k_0`$ such that $`\mathrm{\Gamma }^{k_0}`$ has elements $`\gamma _{ij}^{k_0}()𝒦`$ for any $`i,j`$.
2. The gain matrix $`\mathrm{\Gamma }`$ can be transformed to
$$P\mathrm{\Gamma }P^T=\left(\begin{array}{ccccc}0& A_{12}& 0& \mathrm{}& 0\\ 0& 0& A_{23}& \mathrm{}& 0\\ \mathrm{}& & \mathrm{}& \mathrm{}& \mathrm{}\\ 0& 0& 0& \mathrm{}& A_{\nu 1,\nu }\\ A_{\nu 1}& 0& 0& \mathrm{}& 0\end{array}\right)$$
(4.56)
using some permutation matrix $`P`$, where the zero blocks on the diagonal are square and where $`\mathrm{\Gamma }^\nu `$ is of block diagonal form with square primitive blocks on the diagonal.
###### Proof.
Let $`A_{G_\mathrm{\Gamma }}`$ be the adjacency matrix corresponding to the graph associated with $`\mathrm{\Gamma }`$. This matrix is primitive if and only if $`\mathrm{\Gamma }`$ is primitive. Note that the $`(i,j)^{\text{th}}`$ entry of $`A_{G_\mathrm{\Gamma }}^k`$ is zero if and only if the $`(i,j)^{\text{th}}`$ entry of $`\mathrm{\Gamma }^k`$ is zero. Multiplication of $`\mathrm{\Gamma }`$ by a permutation matrix only rearranges the positions of the class $`𝒦`$-functions, hence this operation is well defined. From these considerations it is clear, that it is sufficient to prove the lemma for the matrix $`A:=A_{G_\mathrm{\Gamma }}`$. But for non-negative matrices this result is an aggregation of known facts from the theory of non-negative matrices, see, e.g., or . ∎
### 4.3 Asymptotic Behavior of $`\mathrm{\Gamma }^k`$
A related question to the stability of the composite system (2.7) is, whether or not the discrete positive dynamical system defined by
$$s_{k+1}=\mathrm{\Gamma }(s_k),k=1,2,\mathrm{}$$
(4.57)
with given initial state $`s_0_+^n`$ is globally asymptotically stable. Under the assumptions we made for Theorem 4 this is indeed true for irreducible $`\mathrm{\Gamma }`$.
###### Theorem 17.
Assume that $`\mathrm{\Gamma }`$ is irreducible. Then the system defined by (4.57) is globally asymptotically stable if and only if $`\mathrm{\Gamma }(s)s`$ for all $`s_+^n\{0\}`$.
The proof will make use of the following result:
###### Proposition 18.
The condition
$$\underset{k\mathrm{}}{lim}\mathrm{\Gamma }^k(s)0\text{for any fixed}s_+^n$$
(4.58)
implies (4.55). Moreover if $`\mathrm{\Gamma }`$ is irreducible, then both are equivalent.
Note that the converse implication is generally not true for reducible maps $`\mathrm{\Gamma }`$, such that (4.55) holds. See Example 19. But it is trivially true, if $`\mathrm{\Gamma }`$ is linear.
###### Proof.
Condition (4.55) follows from (4.58), since if $`\mathrm{\Gamma }(s_0)s_0`$ for some $`s_0_+^n\{0\}`$ then $`\mathrm{\Gamma }^k(s_0)\mathrm{\Gamma }^{k1}(s_0)s_0`$ for $`k=2,3,\mathrm{}`$. Hence the sequence $`\{\mathrm{\Gamma }^k(s_0)\}_{k=0}^{\mathrm{}}`$ does not converge to $`0`$.
Conversely, assume that (4.55) holds and that $`\mathrm{\Gamma }`$ is irreducible.
Step 1. First we prove that for any $`s_+^n\{0\}`$
$$\mathrm{\Gamma }^k(s)s,k.$$
(4.59)
Assume there exist some $`k>1`$ and $`s0`$ with $`\mathrm{\Gamma }^k(s)s.`$ Define $`z_+^n`$ as
$$z:=\underset{l=0,\mathrm{},k1}{\mathrm{max}}\{\mathrm{\Gamma }^l(s)\}\genfrac{}{}{0.0pt}{}{}{}\mathrm{\hspace{0.17em}0}.$$
By (2.6) and using $`\mathrm{\Gamma }^k(s)s`$ we have
$$\mathrm{\Gamma }(z)\underset{l=1,\mathrm{},k}{\mathrm{max}}\{\mathrm{\Gamma }^ls\}=\underset{l=0,\mathrm{},k}{\mathrm{max}}\{\mathrm{\Gamma }^ls\}\underset{l=0,\mathrm{},k1}{\mathrm{max}}\{\mathrm{\Gamma }^ls\}=z.$$
This contradicts (4.55).
Step 2. For any fixed $`s`$ we prove that $`lim\; sup_k\mathrm{}|\mathrm{\Gamma }^k(s)|<\mathrm{}`$. By Lemma 16 we have two cases. We only consider case a), then case b) follows with a slight modification.
Assume that $`s`$ is such that $`lim\; sup_k\mathrm{}|\mathrm{\Gamma }^k(s)|=\mathrm{}`$. For $`t_i>0`$ denote the $`i^{\text{th}}`$ column of $`\mathrm{\Gamma }^{k_0}`$ by
$$\mathrm{\Gamma }_i^{k_0}(t_i)=\left(\begin{array}{c}\gamma _{1i}^{k_0}(t_i)\\ \mathrm{}\\ \gamma _{ni}^{k_0}(t_i)\end{array}\right)$$
As $`\mathrm{\Gamma }^{k_0}`$ has no zero entries, for $`i=1,\mathrm{},n`$ there are $`T_i_+`$ such that
$$\mathrm{\Gamma }_i^{k_0}(t_i)>s\text{for any}t_i>T_i.$$
(4.60)
If $`|\mathrm{\Gamma }^k(s)|\mathrm{}`$ there exists a $`k_1`$ and an index $`i_1`$ such that
$$\mathrm{\Gamma }^{k_1}(s)_{i_1}T_{i_1}.$$
The vector $`\mathrm{\Gamma }^{k_0}\mathrm{\Gamma }^k(s),`$ seen as a sum of columns, is greater than the maximum over these columns, i.e.,
$$\begin{array}{c}\hfill \mathrm{\Gamma }^{k_0}\mathrm{\Gamma }^{k_1}(s)\underset{i}{\mathrm{max}}\mathrm{\Gamma }_i^{k_0}(\mathrm{\Gamma }^{k_1}(s)_i)\mathrm{\Gamma }_{i_1}^{k_0}(\mathrm{\Gamma }^{k_1}(s)_{i_1})\\ \hfill \mathrm{\Gamma }_{i_1}^{k_0}(T_{i_1})s.\end{array}$$
(4.61)
This contradicts Step 1.
Step 3. So $`\left\{\mathrm{\Gamma }^k(s)\right\}_{k1}`$ is bounded for any fixed $`s_+^n`$. The omega-limit set $`\omega (s)`$ is defined by
$$\begin{array}{c}\hfill \omega (s)=\left\{x\right|\text{subsequence}\{k_j\}_{j=1,2,\mathrm{}}\\ \hfill \text{such that}\mathrm{\Gamma }^{k_j}(s)\stackrel{j\mathrm{}}{}x\}.\end{array}$$
(4.62)
This set is not empty by boundedness of $`\left\{\mathrm{\Gamma }^k(s)\right\}_{k1}`$. The following properties follow from this definition and boundedness of the set $`\left\{\mathrm{\Gamma }^k(s)\right\}_{k1}:`$
$$x\omega (s)\mathrm{\Gamma }(x)\omega (s),$$
$$x\omega (s)y\omega (s):\mathrm{\Gamma }(y)=x.$$
I.e., $`\omega (s)`$ is invariant under $`\mathrm{\Gamma }.`$ The boundedness of $`\omega (s)`$ allows to define a finite vector
$$z=sup\omega (s).$$
Then for any $`x\omega (s)`$ it follows $`zx`$ and hence $`\mathrm{\Gamma }(z)\mathrm{\Gamma }(x)`$. Let $`y\omega (s)`$ be such that $`\mathrm{\Gamma }(x)=y`$. Then $`\mathrm{\Gamma }(z)y`$. By the invariance of $`\omega (s)`$ it follows that
$$\mathrm{\Gamma }(z)sup\{\mathrm{\Gamma }(x)|x\omega (s)\}=z.$$
This contradicts (4.55) if $`z0`$, i.e., $`\omega (s)=\{0\}`$. This is true for any $`s_+^n`$. Hence (4.58) is proved as a consequence of (4.55), provided that $`\mathrm{\Gamma }`$ is irreducible. ∎
###### Example 19.
Consider the map $`\mathrm{\Gamma }:_+^2_+^2`$ defined by
$$\mathrm{\Gamma }:=\left[\begin{array}{cc}\gamma _{11}& \text{id}\\ 0& \gamma _{22}\end{array}\right]$$
where for $`t_+`$
$$\gamma _{11}(t):=t(1e^t)$$
and the function $`\gamma _{22}`$ is constructed in the sequel. First note that $`\gamma _{11}𝒦_{\mathrm{}}`$ and $`\gamma _{11}(t)<t,t>0`$. Let $`\{\epsilon _k\}_{k=1}^{\mathrm{}}`$ a strictly decreasing sequence of positive real numbers, such that $`lim_k\mathrm{}\epsilon _k=0`$ and $`lim_K\mathrm{}_{k=1}^K\epsilon _k=\mathrm{}`$. For $`k=1,2,\mathrm{}`$ define
$$\gamma _{22}\left(\epsilon _k+(1+\underset{j=1}{\overset{k1}{}}\epsilon _j)e^{(1+_{j=1}^{k1}\epsilon _j)}\right):=\epsilon _{k+1}+(1+\underset{j=1}{\overset{k}{}}\epsilon _j)e^{(1+_{j=1}^k\epsilon _j)}$$
and observe that
$$\epsilon _k+(1+\underset{j=1}{\overset{k1}{}}\epsilon _j)e^{(1+_{j=1}^{k1}\epsilon _j)}>\epsilon _{k+1}+(1+\underset{j=1}{\overset{k}{}}\epsilon _j)e^{(1+_{j=1}^k\epsilon _j)},$$
since $`\epsilon _k>\epsilon _{k+1}`$ for all $`k=1,2,\mathrm{}`$ and the map $`tte^t`$ is strictly decreasing on $`(1,\mathrm{})`$.
Moreover we have by assumption, that
$$\epsilon _k+(1+\underset{j=1}{\overset{k1}{}}\epsilon _j)e^{(1+_{j=1}^{k1}\epsilon _j)}\underset{k\mathrm{}}{\overset{}{}}0.$$
These facts together imply that $`\gamma _{22}`$ may be extrapolated to some $`𝒦_{\mathrm{}}`$-function, in a way such that $`\gamma _{22}(t)<t,t>0`$ holds.
Note that by our particular construction we have $`\mathrm{\Gamma }(s)s`$ for all $`s_+^2\{0\}`$. Now define $`s^1_+^2`$ by
$$s^1:=\left[\begin{array}{c}1\\ 1+e^1\end{array}\right]$$
and for $`k=1,2,\mathrm{}`$ inductively define $`s^{k+1}:=\mathrm{\Gamma }(s^k)_+^2`$.
By induction one verifies that
$$s^{k+1}=\mathrm{\Gamma }^k(s^1)=\left[\begin{array}{c}1+_{j=1}^k\epsilon _j\\ \epsilon _{k+1}+(1+_{j=1}^k\epsilon _j)e^{(1+_{j=1}^k\epsilon _j)}\end{array}\right].$$
By our previous considerations and assumptions we easily obtain that the second component of the sequence $`\{s^k\}_{k=1}^{\mathrm{}}`$ strictly decreases and converges to zero as $`k`$ tends to infinity. But at the same time the first component strictly increases above any given bound.
Hence we established that $`\mathrm{\Gamma }(s)ss0`$ in general does not imply $`s0:\mathrm{\Gamma }^k(s)0`$ as $`k\mathrm{}`$.
###### Remark 20.
Note that we can even turn the constructed 2x2-$`\mathrm{\Gamma }`$ into the null-diagonal form, that is assumed in Theorem 4. Using the same notation for $`\gamma _{ij}`$ as in Example 19, we just define
$$\mathrm{\Gamma }:=\left[\begin{array}{cccc}0& \gamma _{11}& \text{id}& 0\\ \gamma _{11}& 0& 0& \text{id}\\ 0& 0& 0& \gamma _{22}\\ 0& 0& \gamma _{22}& 0\end{array}\right]\text{and}s^1:=\left[\begin{array}{c}1\\ 1\\ 1+e^1\\ 1+e^1\end{array}\right]$$
and easily verify that $`\mathrm{\Gamma }^k(s^1)`$ does not converge to $`0`$.
###### Proof of Theorem 17.
If (4.57) is asymptotically stable, it is in particular attracted to zero, so by Proposition 18 and irreducibility of $`\mathrm{\Gamma }`$ we establish (4.55).
Conversely assume (4.55). Clearly $`0_+^n`$ is an equilibrium point for (4.57) and by Proposition 18 it is globally attractive. It remains to prove stability, i.e., for any $`\epsilon >0`$ there exists a $`\delta >0`$ such that $`|s_0|<\delta `$ implies $`\mathrm{\Gamma }^k(s_0)<\epsilon `$ for all times $`k=0,1,2,\mathrm{}`$
Given $`\epsilon >0`$ we can choose an $`r_{i=1}^n\mathrm{\Omega }_iS_\epsilon `$ where $`S_\epsilon `$ is the sphere around $`0`$ of radius $`\epsilon `$ in $`_+^n`$. Define $`\delta `$ by
$$\delta :=sup\{d_+:s<rsB_d(0)\}.$$
Here $`B_d(0)`$ denotes the ball of radius less than $`d`$ in $`_+^n`$ around the origin with respect to the Euclidean norm. Clearly we have $`r>s_0`$ for all $`|s_0|<\delta `$. Since $`r_{i=1}^n\mathrm{\Omega }_i\mathrm{}`$ we have $`r>\mathrm{\Gamma }(r)`$ and therefore $`r>\mathrm{\Gamma }(r)\mathrm{\Gamma }^2(r)\mathrm{}`$ and even $`\mathrm{\Gamma }^k(r)\stackrel{k\mathrm{}}{}0`$ again by Proposition 18.
Hence for any $`s_0`$ such that $`|s_0|<\delta `$ we have $`\mathrm{\Gamma }^k(r)\mathrm{\Gamma }^k(s_0)`$ for all $`k=0,1,2,\mathrm{}`$ by monotonicity of $`\mathrm{\Gamma }`$ and therefrom $`\mathrm{\Gamma }^k(s_0)<\epsilon `$ for all $`k=0,1,2,\mathrm{}`$
### 4.4 Geometrical Interpretation
For the following statement let us define the open domains
$$\mathrm{\Omega }_i=\{x^N:|x_i|>\underset{ji}{}\gamma _{ij}(|x_j|)\},$$
where $`N=_{j=1}^nN_j`$ and $`x`$ is partitioned to $`(x_1,\mathrm{},x_n)`$ with $`x_i^{N_i},i=1,\mathrm{},n`$, as in (2.3).
###### Proposition 21.
Condition (4.55) is equivalent to
$$\underset{i=1}{\overset{n}{}}\mathrm{\Omega }_i=^N\{0\}\text{and}\underset{i=1}{\overset{n}{}}\mathrm{\Omega }_i\mathrm{}.$$
(4.63)
###### Proof.
Let $`s0`$. Formula (4.55) is equivalent to the existence of at least one index $`i\{1,\mathrm{},n\}`$ with $`s_i>_{ji}\gamma _{ij}(s_j).`$ This proves the first part of (4.63).
It remains to show, that (4.55) implies $`_{i=1}^n\mathrm{\Omega }_i\mathrm{}`$. We may restrict ourselves to the positive orthant in $`^n`$, and the sets
$$\stackrel{~}{\mathrm{\Omega }}_i=\{s_+^n:s_i>\underset{ji}{}\gamma _{ij}(s_j)\}$$
instead of $`\mathrm{\Omega }_i,i=1,\mathrm{},n`$.
For an index set $`I`$ we define $`E_I=\{s_+^n:s_m=0\text{for}mI\}`$. Note that points of $`E_I`$ can not be in $`\stackrel{~}{\mathrm{\Omega }}_m`$ for $`mI`$. Consider $`\stackrel{~}{\mathrm{\Omega }}_i`$ and $`\stackrel{~}{\mathrm{\Omega }}_j`$ for any $`ij`$. The intersections $`\stackrel{~}{\mathrm{\Omega }}_iE_{\{i,j\}}`$ and $`\stackrel{~}{\mathrm{\Omega }}_jE_{\{i,j\}}`$ of this two domains with the plane $`E_{\{i,j\}}`$ are nonempty. The points of $`\stackrel{~}{\mathrm{\Omega }}_i`$ lying in this plane do not belong to $`\stackrel{~}{\mathrm{\Omega }}_k`$ for any $`kj`$, hence they are in $`\stackrel{~}{\mathrm{\Omega }}_j`$. Since the domains are open it follows that the intersections $`\stackrel{~}{\mathrm{\Omega }}_i\stackrel{~}{\mathrm{\Omega }}_j\mathrm{}`$ for any $`ij.`$ Denote $`\stackrel{~}{\mathrm{\Omega }}_{ij}=\stackrel{~}{\mathrm{\Omega }}_i\stackrel{~}{\mathrm{\Omega }}_j`$ which has nonempty intersection with $`E_{\{i,j\}}`$ by construction. Take any $`ki,j`$. Consider $`E_{\{i,j,k\}}E_{\{i,j\}}`$ which has non-empty intersection with $`\stackrel{~}{\mathrm{\Omega }}_{ij}`$. Let $`x\stackrel{~}{\mathrm{\Omega }}_{ij}E_{\{i,j,k\}}.`$ There is some $`yE_{\{i,j,k\}}`$ with $`y\mathrm{\Omega }_{ij}`$ (say $`yE_{\{k\}}`$). Since $`E_{\{i,j,k\}}`$ is convex the segment $`\overline{xy}E_{\{i,j,k\}},`$ hence there is some point $`zE_{\{i,j,k\}}`$ belonging to $`\stackrel{~}{\mathrm{\Omega }}_{ij}`$, i.e., $`E_{\{i,j,k\}}\stackrel{~}{\mathrm{\Omega }}_{ij}`$ is non-empty.
The points of $`\stackrel{~}{\mathrm{\Omega }}_{ij}`$, which are not in $`\stackrel{~}{\mathrm{\Omega }}_i,\stackrel{~}{\mathrm{\Omega }}_j`$ and lying in $`E_{\{i,j,k\}}`$ can not belong to $`\stackrel{~}{\mathrm{\Omega }}_\nu ,\nu k`$. Hence they are in $`\stackrel{~}{\mathrm{\Omega }}_k`$ and it follows $`\stackrel{~}{\mathrm{\Omega }}_i\stackrel{~}{\mathrm{\Omega }}_j\stackrel{~}{\mathrm{\Omega }}_k\mathrm{}.`$ By iteration the second part of (4.63) follows. ∎
Let us briefly explain, why the overlapping condition (4.63) is interesting: From the theory of ISS-Lyapunov functions it is known, that a system of the form (2.9) is ISS if and only if there exists a smooth Lyapunov function $`V`$ with the property
$$|x|\gamma (|u|)V(x)f(x,u)<W(|x|),$$
for some $`W𝒦`$. In the case of our interconnected system this condition translates to the existence of Lyapunov functions $`V_i`$ for the subsystems $`i=1,\mathrm{},n`$ with the property
$$\begin{array}{c}\hfill |x_i|\gamma _{ij}(|x_j|)+\gamma (|u|)\\ \hfill V_i(x_i)f_i(x,u)<W_i(|x_i|),\end{array}$$
(4.64)
Now for $`u=0`$ the condition of (4.64) is simply, that $`x\mathrm{\Omega }_i`$. Thus the overlapping condition states that in each point of the state space one of the Lyapunov functions of the subsystems is decreasing. It is an interesting problem if via this an ISS-Lyapunov function for the whole system may be constructed.
A typical situation in case of three one dimensional systems ($`^3`$) is presented on the Figure 2 on a plane crossing the positive semi axis. The three sectors are the intersections of the $`\mathrm{\Omega }_i`$ with this plane.
### 4.5 Summary map of the interpretations concerning $`\mathrm{\Gamma }`$
In Figure 3 we summarize the relations between various statements about $`\mathrm{\Gamma }`$ that were proved in section 4.
## 5 Application to linear systems
An important special case is, of course, when the underlying systems are linear themselves. Consider the following setup, where in the sequel we omit the external input, formerly denoted by $`u`$, for notational simplicity. Let
$$\dot{x}_j=A_jx_j,x_j^{N_j},j=1,\mathrm{},n$$
(5.65)
describe $`n`$ globally asymptotically stable linear systems, which are interconnected by the formula
$$\dot{x}_j=A_jx_j+\underset{k=1}{\overset{n}{}}\mathrm{\Delta }_{jk}x_kj=1,\mathrm{},n,$$
(5.66)
which can be rewritten as
$$\dot{x}=(A+\mathrm{\Delta })x,$$
(5.67)
where $`A`$ is block diagonal, $`A=\text{diag}(A_j,j=1,\mathrm{},n)`$, each $`A_j`$ is Hurwitz (i.e., the spectrum of $`A_j`$ is contained in the open left half plane) and the matrix $`\mathrm{\Delta }=(\mathrm{\Delta }_{jk})`$ is also in block form and encodes the connections between the $`n`$ subsystems. We suppose that $`\mathrm{\Delta }_{jj}=0`$ for all $`j`$. Define the matrix $`R=(r_{jk}),R_+^{n\times n}`$, by $`r_{jk}:=\mathrm{\Delta }_{jk}`$. For each subsystem, there exist positive constants $`M_j,\lambda _j`$, such that $`e^{A_jt}M_je^{\lambda _jt}`$ for all $`t0`$.
Define a matrix $`D_+^{n\times n}`$ by $`D:=\text{diag}(\frac{M_j}{\lambda _j},j=1,\mathrm{},n)`$.
From the last subsection we obtain
###### Corollary 22.
If $`\rho (DR)<1`$ then (5.67) is globally asymptotically stable.
Note that this is a special case of a theorem, which can be found in Vidyasagar \[17, p. 110\], see Remark 9.
###### Proof.
Denote the initial value by $`x^0`$. Then by elementary ODE theory we have
$$x_j(t)=e^{A_jt}x_j^0+\underset{kj}{}_0^te^{A_j(ts)}\mathrm{\Delta }_{jk}x_k(s)𝑑s$$
(5.68)
and by standard estimates
$$|x_j(t)|M_je^{\lambda _jt}+\underset{kj}{}r_{jk}\frac{M_k}{\lambda _k}x_{k,[0,t]}.$$
(5.69)
As one can see from (5.69), in this case the gain matrix happens to be $`\mathrm{\Gamma }=DR`$. ∎
It is noteworthy, that this particular corollary is also a consequence of more general and precise results of a recent paper by Hinrichsen, Karow and Pritchard.
## 6 Conclusions
We considered a composite system consisting of an arbitrary number of nonlinear arbitrarily interconnected subsystems, as they arise in applications.
For this general case we derived a multisystem version of the *nonlinear small-gain theorem*. For the special case of linear interconnection gains this is a special case of a known Theorem, cf. \[17, page 110\]. We also showed how our generalized small-gain theorem for networks can be applied to linear systems.
Many interesting questions remain, for instance concerning the construction of Lyapunov functions in case of nonlinear (Lyapunov-)gain functions.
## 7 Acknowledgements
This research is funded by the German Research Foundation (DFG) as part of the Collaborative Research Centre 637 ”Autonomous Cooperating Logistic Processes: A Paradigm Shift and its Limitations” (SFB 637).
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