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# Wonderful Consequences of the Kerr Theorem
## Abstract
Kerrโs multi-particle solution is obtained on the base of the Kerr theorem. Choosing generating function of the Kerr theorem $`F`$ as a product of partial functions $`F_i`$ for spinning particles i=1,โฆk, we obtain a multi-sheeted, multi-twistorial space-time over $`M^4`$ possessing unusual properties. Twistorial structures of the i-th and j-th particles do not feel each other, forming a type of its internal space. Gravitation and electromagnetic interaction of the particles occurs via a singular twistor line which is common for twistorial structures of interacting particles. The obtained multi-particle Kerr-Newman solution turns out to be โdressedโ by singular twistor lines linked to surrounding particles. We conjecture that this structure of space-time has the relation to a stringy structure of vacuum and opens a geometrical way to quantum gravity.
Introduction. It has been mentioned long ago that the Kerr-Newman solution displays some relationships to the quantum world. It is the anomalous gyromagnetic ratio $`g=2`$, as that of the Dirac electron , stringy structures and other features allowing one to construct a semiclassical model of the extended electron which has the Compton size and possesses the wave properties .
One of the mysteries of the Kerr geometry is the existence of two sheets of space-time, $`(+)`$ and $`()`$, on which the dissimilar gravitation (and electromagnetic) fields are realized, and fields living on the $`(+)`$-sheet do not feel the fields of the $`()`$-sheet. Origin of this twofoldedness lies in the Kerr theorem, generating function $`F`$ of which for the Kerr-Newman solution has two roots which determine two different twistorial structures on the same space-time.
In this letter we describe the Kerrโs multi-particle solution. Choosing generating function $`F`$ of the Kerr theorem as a product of partial functions $`F_i`$ for spinning particles i=1,โฆk, we obtain multi-sheeted, multi-twistorial space-time over $`M^4`$ possessing unusual properties. Twistorial structures of the i-th and j-th particles do not feel each other, forming a type of its internal space. Gravitation and electromagnetic interaction of the particles occurs via a singular twistor line.
This unusual structure of space-time is the direct generalization of Kerrโs twofoldedness and we conjecture that it displays also a relation to quantum physics.
The Kerr-Newman metric can be represented in the Kerr-Schild form
$$g_{\mu \nu }=\eta _{\mu \nu }+2hk_\mu k_\nu ,$$
(1)
where $`\eta _{\mu \nu }`$ is metric of auxiliary Minkowski space-time $`M^4`$,
$$h=(mre^2/2)/(r^2+a^2\mathrm{cos}^2\theta ),$$
(2)
and $`k_\mu (x)`$ is a twisting null field, which is tangent to the Kerr principal null congruence (PNC) which is geodesic and shear-free . PNC is determined by the complex function $`Y(x)`$ via the one-form
$$e^3=du+\overline{Y}d\zeta +Yd\overline{\zeta }Y\overline{Y}dv=Pk_\mu dx^\mu $$
(3)
where $`u,v,\overline{\zeta },\zeta `$ are the null Cartesian coordinates. Here $`P`$ is a normalizing factor for $`k_\mu `$ which provide $`k_0=1`$ in the rest frame.<sup>*</sup><sup>*</sup>*We replace the factors $`P`$ from $`e^3`$ to function $`h`$, so $`h`$ in differs by factor $`P^2`$. The null rays of the Kerr congruence are twistors. Form of the Kerr PNC is shown on Fig. 1.
The Kerr theorem allows one to describe the Kerr geometry in twistor terms .
It claims that any geodesic and shear-free null congruence in Minkowski space-time is defined by a function $`Y(x)`$ which is a solution of the equation
$$F=0,$$
(4)
where $`F(Y,\lambda _1,\lambda _2)`$ is an arbitrary holomorphic function of the projective twistor coordinates
$$Y,\lambda _1=\zeta Yv,\lambda _2=u+Y\overline{\zeta }.$$
(5)
In the Kerr-Schild backgrounds the Kerr theorem acquires a broader content , allowing one to determine the normalizing function $`P`$ and complex radial distance $`\stackrel{~}{r}=r+ia\mathrm{cos}\theta ,`$
$$P=_{\lambda _1}F\overline{Y}_{\lambda _2}F,\stackrel{~}{r}=PZ^1=dF/dY$$
(6)
which are important characteristics of the corresponding solutions. The position of singular lines, caustics of PNC, corresponds to $`\stackrel{~}{r}=0`$, and is determined by the system of equations $`F=0;dF/dY=0.`$
The proof of the Kerr theorem in the extended version adapted to the Kerr-Schild formalism is given in .
In the original paper , the following generating function $`F`$ was considered
$$F\varphi (Y)+(qY+c)(d\zeta Ydv)(pY+\overline{q})(u+Y\overline{\zeta }).$$
(7)
The parameters $`p,q,\overline{q},c`$ are related to the Killing vector of the solution $`K^\mu _\mu =c_u+\overline{q}_\zeta +q_{\overline{\zeta }}p_v,`$ and $`P=pY\overline{Y}+qY+\overline{q}\overline{Y}+c.`$ For the stationary Kerr solution $`p=c=2^{1/2},q=\overline{q}=0`$ and $`\stackrel{~}{r}=r+ia\mathrm{cos}\theta `$.
It was shown in that function $`\varphi (Y)`$ in (7) has to be at most quadratic in $`Y`$ to provide singular lines to be confined in a restricted region, which corresponds to the Kerr PNC up to the Lorentz boosts, orientations of angular momenta and the shifts of origin.
In the papers another form for this function was suggested $`F=(\lambda _1\lambda _1^0)\stackrel{ห}{K}\lambda _2(\lambda _2\lambda _2^0)\stackrel{ห}{K}\lambda _1`$ which is related to the Newman-initiated complex representation of the Kerr geometry. In this case, function $`F(Y)`$ can be expressed via the set of parameters $`q`$ which determine the motion and orientation of the Kerr spinning particle and takes the form $`F(Y|q)=A(x|q)Y^2+B(x|q)Y+C(x|q)`$. The equations (4) can be resolved explicitly, leading to two roots $`Y=Y^\pm (x|q)`$ which correspond to two sheets of the Kerr space-time. The root $`Y^+(x)`$ determines via (3) the out-going congruence on the $`(+)`$-sheet, while the root $`Y^{}(x)`$ gives the in-going congruence on the $`()`$-sheet. Therefore, function $`F`$ may be represented in the form $`F(Y|q)=A(x|q)(YY^+)(YY^{}),`$ which allows one to obtain all the required functions of the Kerr solution in explicit form. The detailed form of $`Y^\pm (x|q)`$ is not important for our treatment here and may be found in .
Multi-twistorial space-time. Selecting an isolated i-th particle with parameters $`q_i`$, one can obtain the roots $`Y_i^\pm (x)`$ of the equation $`F_i(Y|q_i)=0`$ and express $`F_i`$ in the form
$$F_i(Y)=A_i(x)(YY_i^+)(YY_i^{}).$$
(8)
Then, substituting the $`(+)`$ or $`()`$ roots $`Y_i^\pm (x)`$ in the relation (3), one determines congruence $`k_\mu ^{(i)}(x)`$ and consequently, the Kerr-Schild ansatz (1) for metric
$$g_{\mu \nu }^{(i)}=\eta _{\mu \nu }+2h^{(i)}k_\mu ^{(i)}k_\nu ^{(i)},$$
(9)
and finally, the function $`h^{(i)}(x)`$ may be expressed in terms of $`\stackrel{~}{r}_i=d_YF_i,`$ (6), as follows
$$h^{(i)}=\frac{m}{2}(\frac{1}{\stackrel{~}{r}_i}+\frac{1}{\stackrel{~}{r}_i^{}})+\frac{e^2}{2|\stackrel{~}{r}_i|^2}.$$
(10)
Electromagnetic field is given by the vector potential
$$A_\mu ^{(i)}=\mathrm{}e(e/\stackrel{~}{r}_i)k_\mu ^{(i)}.$$
(11)
What happens if we have a system of $`k`$ particles? One can form the function $`F`$ as a product of the known blocks $`F_i(Y)`$,
$$F(Y)\underset{i=1}{\overset{k}{}}F_i(Y).$$
(12)
The solution of the equation $`F=0`$ acquires $`2k`$ roots $`Y_i^\pm `$, and the twistorial space turns out to be multi-sheeted.
The twistorial structure on the i-th $`(+)`$ or $`()`$ sheet is determined by the equation $`F_i=0`$ and does not depend on the other functions $`F_j,ji`$. Therefore, the particle $`i`$ does not feel the twistorial structures of other particles. Similar, the condition for singular lines $`F=0,d_YF=0`$ acquires the form
$$\underset{l=1}{\overset{k}{}}F_l=0,\underset{i=1}{\overset{k}{}}\underset{li}{\overset{k}{}}F_ld_YF_i=0$$
(13)
and splits into k independent relations
$$F_i=0,\underset{li}{\overset{k}{}}F_ld_YF_i=0.$$
(14)
One sees, that i-th particle does not feel also singular lines of other particles. The space-time splits on the independent twistorial sheets, and therefore, the twistorial structure related to the i-th particle plays the role of its โinternal spaceโ.
It looks wonderful. However, it is a direct generalization of the well known twofoldedness of the Kerr space-time which remains one of the mysteries of the Kerr solution for the very long time.
For spinning particles $`|a|>>m`$ and the Kerrโs black hole horizons disappear, there appears the old problem of the source of Kerr solution with the alternative: either to remove this twofoldedness or to give it a physical interpretation. By truncation of the negative sheet, there appears the source in the form of relativistically rotating disk , bubble or bag .
Alternative way is to retain the negative sheet, treating it as the sheet of advanced fields. In this case the source of spinning particle turns out to be the Kerr singular ring (circular string, ) with the electromagnetic excitations in the form of traveling waves which generate spin and mass of the particle (microgeon model ).
Multi-particle Kerr-Schild solution. Using the Kerr-Schild formalism with the considered above generating functions $`_{i=1}^kF_i(Y)=0,`$ one can obtain the exact asymptotically flat multi-particle solutions of the Einstein-Maxwell field equations. Since congruences are independent on the different sheets, the congruence on the i-th sheet retains to be geodesic and shear-free, and one can use the standard Kerr-Schild algorithm of the paper . One could expect that result for the i-th sheet will be in this case the same as the known solution for isolated particle. Unexpectedly, there appears a new feature having a very important consequence.
Formally, we have only to replace $`F_i`$ by $`F=_{i=1}^kF_i(Y)=\mu _iF_i(Y),`$ where $`\mu _i=_{ji}^kF_j(Y)`$ is a normalizing factor which takes into account the external particles. However, in accordance with (6) this factor $`\mu _i`$ will appear also in the function $`\stackrel{~}{r}=d_YF=\mu _id_YF_i,`$ and in the function $`P=\mu _iP_i.`$
So, we obtain the different result
$$h_i=\frac{m_i(Y)}{2\mu _i^3}(\frac{1}{\stackrel{~}{r}_i}+\frac{1}{\stackrel{~}{r}_i^{}})+\frac{(e/\mu _i)^2}{2|\stackrel{~}{r}_i|^2},$$
(15)
$$A_\mu ^{(i)}=\mathrm{}e\frac{e}{\mu _i\stackrel{~}{r}_i}k_\mu ^{(i)}$$
(16)
which looks like a renormalization of the mass $`m`$ and charge $`e`$.Function $`m_i(Y)`$ is free and satisfies the condition $`(m_i),_{\overline{Y}}=0`$. It and has to be chosen in the form $`m_i(Y)=m_0\mu _i^3`$ to provide reality of metric.
This fact turns out to be still more intriguing if we note that $`\mu _i`$ is not constant, but a function of $`Y_i`$. We can specify its form by using the known structure of blocks $`F_i`$
$$\mu _i(Y_i)=\underset{ji}{}A_j(x)(Y_iY_j^+)(Y_iY_j^{}).$$
(17)
The roots $`Y_i`$ and $`Y_j^\pm `$ may coincide for some values of $`Y_i`$, which selects a common twistor for the sheets $`i`$ and $`j`$. Assuming that we are on the i-th $`(+)`$-sheet, where congruence is out-going, this twistor line will also belong to the in-going $`()`$-sheet of the particle $`j`$ . The metric and electromagnetic field will be singular along this twistor line, because of the pole $`\mu _iA(x)(Y_i^+Y_j^{})`$. Therefore, interaction occurs along a light-like Schild string which is common for twistorial structures of both particles. The field structure of this string is similar to the well known structure of pp-wave solutions.
The equations (1), (2) and (3) give the exact multi-particle solution of the Einstein-Maxwell field equations. It follows from the fact that the equations were fully integrated out in and expressed via functions $`P`$ and $`Z`$ before (without) concretization of the form of congruence, under the only condition that it is geodesic and shear free. In the same time the Kerr theorem determines the functions $`P`$ and $`Z`$ via generating function $`F,`$ eq.(6), and the condition of reality for metric may be provided by a special choice of the free function $`m(Y)`$.
The obtained multi-particle solutions show us that, in addition to the usual Kerr-Newman solution for an isolated spinning particle, there is a series of the exact โdressedโ Kerr-Newman solutions which take into account surrounding particles and differ by the appearance of singular twistor strings connecting the selected particle to external particles. This is a new gravitational phenomena which points out on a probable stringy (twistorial) texture of vacuum and may open a geometrical way to quantum gravity.
The number of surrounding particles and number of blocks in the generating function $`F`$ may be assumed countable. In this case the multi-sheeted twistorial space-time will possess the properties of the multi-particle Fock space.
Acknowledgments. Author thanks the participants of the seminar on Quantum Field Theory at the Physical Lebedev Institute for useful discussion. This work was supported by the RFBR Grant 04-0217015-a and by the ISEP Research Grant by Jack Sarfatti.
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# Universality of the Collins-Soper-Sterman nonperturbative function in vector boson production
## Abstract
We revise the $`b_{}`$ model for the Collins-Soper-Sterman resummed form factor to improve description of the leading-power contribution at nearly nonperturbative impact parameters. This revision leads to excellent agreement of the transverse momentum resummation with the data in a global analysis of Drell-Yan lepton pair and $`Z`$ boson production. The nonperturbative contributions are found to follow universal quasi-linear dependence on the logarithm of the heavy boson invariant mass, which closely agrees with an estimate from the infrared renormalon analysis.
and
Transverse momentum distributions of heavy Drell-Yan lepton pairs, $`W`$, or $`Z`$ bosons produced in hadron-hadron collisions present an interesting example of factorization for multi-scale observables. If the transverse momentum $`q_T`$ of the electroweak boson is much smaller than its invariant mass $`Q`$, $`d\sigma /dq_T`$ at an n-th order of perturbation theory includes large contributions of the type $`\alpha _s^n\mathrm{ln}^m(q_T^2/Q^2)/q_T^2`$ ($`m=0,1\mathrm{}2n1`$), which must be summed through all orders of $`\alpha _s`$ to reliably predict the cross section . The feasibility of all-order resummation is proved by a factorization theorem, first formulated for $`e^+e^{}`$ hadroproduction , stated by Collins, Soper, and Sterman (CSS) for the Drell-Yan process , and recently proved by detailed investigation of gauge transformations of $`k_T`$-dependent parton densities .
The heavy bosons acquire non-zero $`q_T`$ mostly by recoiling against QCD radiation. The CSS formalism accounts for both the short- and long-wavelength QCD radiation by means of a Fourier-Bessel transform of a resummed form factor $`\stackrel{~}{W}(b)`$ introduced in impact parameter ($`b`$) space. The perturbative contribution, characterized by $`b0.5\text{ GeV}^1`$, dominates in $`W`$ and $`Z`$ boson production at all values of $`q_T`$. The nonperturbative contribution from $`b0.5\text{ GeV}^1`$ is not negligible at $`q_T<20`$ GeV in the precision measurements of the $`W`$ boson mass $`M_W`$ at the Tevatron and LHC . The model for the nonperturbative recoil is the major source of theoretical uncertainty in the extraction of $`M_W`$ from the experimental data. This uncertainty must be reduced in order to measure $`M_W`$ with accuracy of about 30 MeV in the Tevatron Run-2 and 15 MeV at the LHC. The nonperturbative model presented below approaches the level of accuracy desired in these measurements.
The nonperturbative component \[described by the function $`_{NP}(b,Q)`$ given in Eq. (4)\] can be constrained in a few experiments by exploiting process-independence, or universality, of $`_{NP}(b,Q)`$, just as the universal $`k_T`$-integrated parton densities are constrained with the help of inclusive scattering data. The universality of $`_{NP}(b,Q)`$ in unpolarized Drell-Yan-like processes and semi-inclusive deep-inelastic scattering (SIDIS) follows from the CSS factorization theorem . In the study presented here, we carefully investigate agreement of the universality assumption with the data in a global analysis of fixed-target Drell-Yan pair and Tevatron $`Z`$ boson production. We revise the nonperturbative model used in the previous studies and improve agreement with the data without introducing additional free parameters. Renormalization-group invariance requires $`_{NP}(b,Q)`$ to depend linearly on $`\mathrm{ln}Q`$ . With our latest revisions put in place, the global $`q_T`$ fit clearly prefers a simple function $`_{NP}(b,Q)`$ with universal $`\mathrm{ln}Q`$ dependence. The new $`_{NP}(b,Q)`$ has reduced dependence on the collision energy $`\sqrt{S}`$ comparatively to the earlier fits. The slope of the $`\mathrm{ln}Q`$ dependence found in the new fit agrees numerically with its estimate made with methods of infrared renormalon analysis .
The function $`_{NP}(b,Q)`$ primarily parametrizes the โpower-suppressedโ terms, i.e., terms proportional to positive powers of $`b`$. When assessed in a fit, $`_{NP}(b,Q)`$ also contains admixture of the leading-power terms (logarithmic in $`b`$ terms), which were not properly included in the approximate leading-power function $`\stackrel{~}{W}_{LP}(b)`$ \[cf. Eq. (4)\]. In contrast, estimates of $`_{NP}(b,Q)`$ made in the infrared renormalon analysis explicitly remove all leading-power contributions from $`_{NP}(b,Q)`$ . While the recent studies point to an approximately Gaussian form of $`_{NP}(b,Q)`$ \[$`_{NP}(b,Q)b^2`$\], they disagree on the magnitude of $`_{NP}(b,Q)`$ and its $`Q`$ dependence. The source of these differences can be traced to the varying assumptions about the form of the leading-power function $`\stackrel{~}{W}_{LP}(b)`$ at $`b<2\text{ GeV}^1`$, which is correlated in the fit with $`_{NP}(b,Q)`$. The exact behavior of $`\stackrel{~}{W}(b)`$ at $`b>2\text{ GeV}^1`$ is of reduced importance, as $`\stackrel{~}{W}(b)`$ is strongly suppressed at such $`b`$. The new improvements described here (excellent agreement of $`_{NP}(b,Q)`$ with the data and renormalon analysis) result from modifications in the model for $`\stackrel{~}{W}_{LP}(b)`$ at $`b<2\text{ GeV}^1`$. The improvements are preserved under variations of the large-$`b`$ form of $`\stackrel{~}{W}_{LP}(b)`$ in a significant range of the model parameters.
Our paper follows the notations in Ref. . The form factor $`\stackrel{~}{W}(b)`$ factorizes at all $`b`$ as
$$\stackrel{~}{W}(b)=\underset{j=q,\overline{q}}{}\frac{\sigma _j^{(0)}}{S}e^{๐ฎ(b,Q)}๐ซ_j(x_1,b)๐ซ_{\overline{j}}(x_2,b),$$
(1)
where $`\sigma _j^{(0)}/S`$ is a constant prefactor , and $`x_{1,2}e^{\pm y}Q/\sqrt{S}`$ are the Born-level momentum fractions, with $`y`$ being the rapidity of the vector boson. The $`b`$-dependent parton densities $`๐ซ_j(x,b)`$ and Sudakov function
$`๐ฎ(b,Q){\displaystyle _{b_0^2/b^2}^{Q^2}}{\displaystyle \frac{d\overline{\mu }^2}{\overline{\mu }^2}}\left[๐(\alpha _s(\overline{\mu }))\mathrm{ln}\left({\displaystyle \frac{Q^2}{\overline{\mu }^2}}\right)+(\alpha _s(\overline{\mu }))\right]`$ (2)
are universal in Drell-Yan-like processes and SIDIS . When the momentum scales $`Q`$ and $`b_0/b`$ (where $`b_02e^{\gamma _E}1.123`$ is a dimensionless constant) are much larger than 1 GeV, $`\stackrel{~}{W}(b)`$ reduces to its perturbative part $`\stackrel{~}{W}_{pert}(b)`$, i.e., its leading-power (logarithmic in $`b`$) part evaluated at a finite order of $`\alpha _s`$:
$`\stackrel{~}{W}(b)|_{Q,b_0/b1GeV}\stackrel{~}{W}_{pert}(b){\displaystyle \underset{j=q,\overline{q}}{}}{\displaystyle \frac{\sigma _j^{(0)}}{S}}e^{๐ฎ_P(b,Q)}`$
$`\times \left[๐f\right]_j(x_1,b;\mu _F)\left[๐f\right]_{\overline{j}}(x_2,b;\mu _F).`$ (3)
Here $`๐ฎ_P(b,Q)`$ and $`\left[๐f\right]_j(x,b;\mu _F)_a_x^1๐\xi /\xi ๐_{ja}(x/\xi ,\mu _Fb)f_a(\xi ,\mu _F)`$ are the finite-order approximations to the leading-power parts of $`๐ฎ(b,Q)`$ and $`๐ซ_j(x,b)`$. $`f_a(x,\mu _F)`$ is the $`k_T`$-integrated parton density, computed in our study by using the CTEQ6M parameterization . $`๐_{ja}(x,\mu _Fb)`$ is the Wilson coefficient function. We compute $`๐ฎ_P(b,Q)`$ up to $`O(\alpha _s^2)`$ and $`๐_{ja}`$ up to $`O(\alpha _s)`$.
In $`Z`$ boson production, the maximum of $`b\stackrel{~}{W}(b)`$ is located at $`b0.25\text{ GeV}^1`$, and $`\stackrel{~}{W}_{pert}(b)`$ dominates the Fourier-Bessel integral. In the examined low-$`Q`$ region, the maximum of $`b\stackrel{~}{W}(b)`$ is located at $`b1\text{ GeV}^1`$, where higher-order corrections in powers of $`\alpha _s`$ and $`b`$ must be considered. We reorganize Eq. (1) to separate the leading-power (LP) term $`\stackrel{~}{W}_{LP}(b)`$, given by the model-dependent continuation of $`\stackrel{~}{W}_{pert}(b)`$ to $`b1\text{ GeV}^1`$, and the nonperturbative exponent $`e^{_{NP}(b,Q)}`$, which absorbs the power-suppressed terms:
$$\stackrel{~}{W}(b)=\stackrel{~}{W}_{LP}(b)e^{_{NP}(b,Q)}.$$
(4)
At $`b0`$, the perturbative approximation for $`\stackrel{~}{W}(b)`$ is restored: $`\stackrel{~}{W}_{LP}\stackrel{~}{W}_{pert},`$ $`_{NP}0`$. The power-suppressed contributions are proportional to even powers of $`b`$ . Detailed expressions for some power-suppressed terms are given in Ref. . At impact parameters of order 1$`\text{ GeV}^1`$, we keep only the first power-suppressed contribution proportional to $`b^2`$:
$$_{NP}b^2\left(a_1+a_2\mathrm{ln}(Q/Q_0)+a_3\varphi (x_1)+a_3\varphi (x_2)\right)+\mathrm{},$$
(5)
where $`a_1`$, $`a_2`$, and $`a_3`$ are coefficients of magnitude less than $`1\text{ GeV}^2`$, and $`\varphi (x)`$ is a dimensionless function. The terms $`a_2\mathrm{ln}(Q/Q_0)`$ and $`a_3\varphi (x_j)`$ arise from $`๐ฎ(b,Q)`$ and $`\mathrm{ln}\left[๐ซ_j(x_j,b)\right]`$ in $`\mathrm{ln}\left[\stackrel{~}{W}(b)\right]`$, respectively. We neglect the flavor dependence of $`\varphi (x)`$ in the analyzed region dominated by scattering of light $`u`$ and $`d`$ quarks. $`_{NP}`$ is consequently a universal function within this region. The dependence of $`_{NP}`$ on $`\mathrm{ln}Q`$ follows from renormalization-group invariance of the soft-gluon radiation . The coefficient $`a_2`$ of the $`\mathrm{ln}Q`$ term has been related to the vacuum average of the Wilson loop operator and estimated within lattice QCD as $`0.19_{0.09}^{+0.12}\text{ GeV}^2`$ .
The preferred $`_{NP}`$ is correlated in the fit with the assumed large-$`b`$ behavior of $`\stackrel{~}{W}_{LP}`$. We examine this correlation in a modified version of the $`b_{}`$ model . The shape of $`\stackrel{~}{W}_{LP}`$ is varied in the $`b_{}`$ model by adjusting a single parameter $`b_{max}`$. Continuity of $`\stackrel{~}{W}`$ and its derivatives, needed for the numerical stability of the Fourier transform, is always preserved. We set $`\stackrel{~}{W}_{LP}(b)\stackrel{~}{W}_{pert}(b_{}),`$ with $`b_{}(b,b_{max})b(1+b^2/b_{max}^2)^{1/2}`$. $`\stackrel{~}{W}_{LP}(b)`$ reduces to $`\stackrel{~}{W}_{pert}(b)`$ as $`b0`$ and asymptotically approaches $`\stackrel{~}{W}_{pert}(b_{max})`$ as $`b\mathrm{}`$. The $`b_{}`$ model with a relatively low $`b_{max}=0.5\text{ GeV}^1`$ was a choice of the previous $`q_T`$ fits . However, it is natural to consider $`b_{max}`$ above 1$`\text{ GeV}^1`$ in order to avoid *ad hoc* modifications of $`\stackrel{~}{W}_{pert}(b)`$ in the region where perturbation theory is still applicable. To implement $`\stackrel{~}{W}_{pert}(b_{})`$ for $`b_{max}>1\text{ GeV}^1`$, we must choose the factorization scale $`\mu _F`$ such that it stays, at any $`b`$ and $`b_{max}`$, above the initial scale $`Q_{ini}=1.3`$ GeV of the DGLAP evolution for the CTEQ6 PDFโs $`f_a(x,\mu _F)`$. We keep the usual choice $`\mu _F=C_3/b_{}(b,b_{max})`$, where $`C_3b_0`$, for $`b_{max}b_0/Q_{ini}0.86\text{ GeV}^1`$. Such choice is not acceptable at $`b_{max}>b_0/Q_{ini}`$, as it would allow $`\mu _F<Q_{ini}`$. Instead, we choose $`\mu _F=C_3/b_{}(b,b_0/Q_{ini})`$ for $`b_{max}>b_0/Q_{ini}`$, i.e., we substitute $`b_0/Q_{ini}`$ for $`b_{max}`$ in $`\mu _F`$ to satisfy $`\mu _FQ_{ini}`$ at any $`b`$. Aside from $`f_a(x,\mu _F)`$, all terms in $`\stackrel{~}{W}_{pert}(b)`$ are known, at least formally, as explicit functions of $`\alpha _s(1/b)`$ at all $`b<1/\mathrm{\Lambda }_{QCD}`$. We show in Ref. that this prescription preserves correct resummation of the large logarithms and is numerically stable up to $`b_{max}3\text{ GeV}^1`$.
We perform a series of fits for several choices of $`b_{max}`$ by closely following the previous global $`q_T`$ analysis . We consider a total of 98 data points from production of Drell-Yan pairs in E288, E605, and R209 fixed-target experiments, as well as from observation of $`Z`$ bosons with $`q_T<10`$ GeV by CDF and Dร detectors in the Tevatron Run-1. See Ref. for the experimental references. Overall normalizations for the experimental cross sections are varied as free parameters. Our best-fit normalizations agree with the published values within the systematical errors provided by the experiments, with the exception of the CDF Run-1 normalization (rescaled down by 7%).
To test the universality of $`_{NP}`$, we individually examine each bin of $`Q`$. We choose $`_{NP}=a(Q)b^2`$ and independently fit it to each of the 5 experimental data sets to determine the most plausible normalization in each experiment. We then set the normalizations equal to their best-fit values and examine $`\chi ^2`$ at each $`Q`$ as a function of $`a(Q)`$. For $`b_{max}=12\text{ GeV}^1`$, the best-fit values of $`a(Q)`$ follow a nearly linear dependence on $`\mathrm{ln}Q`$ \[cf. Fig. 1\]. The slope $`a_2da(Q)/d(\mathrm{ln}Q)`$ is close to the renormalon analysis expectation of $`0.19\text{ GeV}^2`$ . The agreement with the universal linear $`\mathrm{ln}Q`$ dependence worsens if $`b_{max}`$ is chosen outside the region 1-2$`\text{ GeV}^1`$. Since the best-fit $`a(Q)`$ are found independently in each $`Q`$ bin, we conclude that the data support the universality of $`_{NP}`$, when $`b_{max}`$ lies in the range $`12\text{ GeV}^1`$. In another test, we find that each experimental data set individually prefers a nearly quadratic dependence on $`b`$, $`_{NP}=a(Q)b^{2\beta }`$, with $`|\beta |<0.5`$ in all five experiments.
To further explore the issue, we simultaneously fit our model to all the data. We parametrize $`a(Q)`$ as $`a(Q)a_1+a_2\mathrm{ln}\left[Q/(3.2\text{ GeV})\right]+a_3\mathrm{ln}\left[100x_1x_2\right].`$ This parametrization coincides with the BLNY form , if the parameters are renamed as $`\{g_1,g_2,g_1g_3\}\text{(BLNY)}\{a_1,a_2,a_3\}\text{(here)}`$. It agrees with the generic form of $`_{NP}(b,Q)`$ in Eq. (5), if one identifies $`\varphi (x)=\mathrm{ln}(x/0.1)`$. We carry out two sequences of fits for $`C_3=b_0`$ and $`C_3=2b_0`$ to investigate the stability of our prescription for $`\mu _F`$ and sensitivity to $`๐ช(\alpha _s^2)`$ corrections. The dependence on $`C_3`$ is relatively uniform across the whole range of $`b_{max}`$, indicating that our choice of $`\mu _F`$ for $`b_{max}>b_0/Q_{ini}`$ is numerically stable.
Fig. 2 shows the dependence of the best-fit $`\chi ^2,`$ $`a_1,`$ $`a_2`$, and $`a_3`$ on $`b_{max}`$. As $`b_{max}`$ is increased above $`0.5\text{ GeV}^1`$ assumed in the BLNY study, $`\chi ^2`$ rapidly decreases, becomes relatively flat at $`b_{max}=12\text{ GeV}^1`$, and grows again at $`b_{max}>2\text{ GeV}^1`$. The global minimum of $`\chi ^2=125(111)`$ is reached at $`b_{max}1.5`$ GeV<sup>-1</sup>, where all data sets are described equally well, without major tensions among the five experiments. The improvement in $`\chi ^2`$ mainly ensues from better agreement with the low-$`Q`$ experiments (E288, E605, and R209), while the quality of all fits to the $`Z`$ data is about the same. This is illustrated by Fig. 3, which shows the differences between the measured and theoretical cross sections, divided by the experimental errors $`\delta _{exp}`$, as well as the values of $`\chi ^2`$ in each experiment, in two representative fits for $`b_{max}=0.5\text{ GeV}^1`$, $`C_3=b_0`$ (squares) and $`b_{max}=1.5\text{ GeV}^1`$, $`C_3=2b_0`$ (triangles). The data are arranged in bins of $`Q`$ (shown by gray background stripes) and $`q_T`$, with both variables increasing from left to right. For $`b_{max}=1.5\text{ GeV}^1`$, the $`(\text{Data}\text{Theory})`$ differences are reduced on average in the entire low-$`Q`$ sample, resulting in lower $`\chi ^2`$ in three low-$`Q`$ experiments. Two outlier points in the E605 sample (the first point in the second $`Q`$ bin and fifth point in the fifth $`Q`$ bin) disagree with the other E288 and E605 data in the same $`Q`$ and $`x`$ region and contribute $`1525`$ units to $`\chi ^2`$ at any $`b_{max}`$. If the two outliers were removed, one would find $`\chi ^2/d.o.f.1`$ at the global minimum.
The magnitudes of $`a_1`$, $`a_2`$, and $`a_3`$ are reduced when $`b_{max}`$ increases from $`0.5`$ to $`1.5\text{ GeV}^1`$. In the whole range $`1b_{max}2`$ GeV<sup>-1</sup>, $`a_2`$ agrees with the renormalon analysis estimate. The coefficient $`a_3`$, which parametrizes deviations from the linear $`\mathrm{ln}Q`$ dependence, is considerably smaller ($`<0.05`$) than both $`a_1`$ and $`a_2`$ ($`0.2`$). A reasonable quality of the fit is retained if $`a_3`$ is set to zero by hand: $`\chi ^2`$ increases by $`5`$ in such a fit above its minimum in the fit with a free $`a_3`$. In contrast, $`\chi ^2`$ increases by $`>200`$ units if $`a_3=g_1g_3`$ is set to zero at $`b_{max}=0.5\text{ GeV}^1`$, as it was noticed in the BLNY study.
The preference for the values of $`b_{max}`$ between $`1`$ and $`2\text{ GeV}^1`$ indicates, first, that the data do favor the extension of the $`b`$ range where $`\stackrel{~}{W}_{LP}(b)`$ is approximated by the exact $`\stackrel{~}{W}_{pert}(b)`$. In $`Z`$ boson production, this region extends up to $`34\text{ GeV}^1`$ as a consequence of the strong suppression of the large-$`b`$ tail by the Sudakov exponent. The fit to the $`Z`$ data is actually independent of $`b_{max}`$ within the experimental uncertainties for $`b_{max}>1\text{ GeV}^1`$. The best-fit form factors $`b\stackrel{~}{W}(b)`$ for $`b_{max}=0.5`$ and $`1.5\text{ GeV}^1`$ in $`Z`$ boson production are shown in Fig. 4(a).
In the low-$`Q`$ Drell-Yan process, continuation of $`b\stackrel{~}{W}_{pert}(b)`$ far beyond $`b1\text{ GeV}^1`$ raises objections, since $`b\stackrel{~}{W}_{pert}(b)`$ has a maximum and is unstable with respect to higher-order corrections at $`b1.21.5\text{ GeV}^1`$. The dubious large contributions to $`\stackrel{~}{W}_{pert}(b)`$ in this $`b`$ region would deteriorate the quality of the fit. The $`b_{}`$ prescription with $`b_{max}<2\text{ GeV}^1`$ reduces the impact of the dubious terms on $`\stackrel{~}{W}(b)`$: for $`b_{max}`$ small enough, the maximum of $`\stackrel{~}{W}_{pert}(b_{})`$ is only reached at $`b1.2\text{ GeV}^1`$, where it is suppressed by $`e^{_{NP}(b,Q)}`$. The best-fit form factors for the E605 kinematics, divided by the best-fit normalizations of the E605 data $`N_{fit}`$, are shown in Fig. 4(b).
If a very large $`b_{max}`$ comparable to $`1/\mathrm{\Lambda }_{QCD}`$ is taken, $`\stackrel{~}{W}_{LP}(b)`$ essentially coincides with $`\stackrel{~}{W}_{pert}(b)`$, extrapolated to large $`b`$ by using the known, although not always reliable, dependence of $`\stackrel{~}{W}_{pert}(b)`$ on $`\mathrm{ln}b`$. Similar, but not identical, extrapolations of $`\stackrel{~}{W}_{pert}(b)`$ to large $`b`$ are realized in the models , which describe the $`Z`$ data well, in accord with our own findings. In $`Z`$ boson production, our best-fit $`a(M_Z)=0.85\pm 0.10\text{ GeV}^2`$ agrees with $`0.8\text{ GeV}^2`$ found in the extrapolation-based models, and it is about a third of $`2.7\text{ GeV}^2`$ predicted by the BLNY parametrization. Our results support the conjecture in that $`a_3`$ is small if the exact form of $`\stackrel{~}{W}_{pert}(b)`$ is maximally preserved. To describe the low-$`Q`$ data, the model allowed a large discontinuity in the first derivative of $`\stackrel{~}{W}(b)`$ at $`b`$ equal to the separation parameter $`b_{max}^{QZ}=0.30.5\text{ GeV}^1`$, where switching from the exact $`\stackrel{~}{W}_{pert}(b)`$ to its extrapolated form occurs \[cf. Fig. 4(b)\]. In the revised $`b_{}`$ model, such discontinuity does not happen, and $`\stackrel{~}{W}_{LP}(b)`$ is closer to the exact $`\stackrel{~}{W}_{pert}(b)`$ in a wider $`b`$ range at low $`Q`$ than in the model . The two models differ substantially at $`b1\text{ GeV}^1`$, as seen in Fig. 4(b).
To summarize, the extrapolation of $`\stackrel{~}{W}_{pert}(b)`$ to $`b>1.5\text{ GeV}^1`$ is disfavored by the low-$`Q`$ data sets, if a purely Gaussian form of $`_{NP}`$ is assumed. The Gaussian approximation is adequate, on the other hand, in the $`b_{}`$ model with $`b_{max}`$ in the range $`12\text{ GeV}^1`$. Here variations in $`b_{max}`$ are compensated well by adjustments in $`a_1`$, $`a_2`$, and $`a_3`$, and the full form factor $`b\stackrel{~}{W}(b)`$ stays approximately independent of $`b_{max}`$. The best-fit parameters in $`_{NP}`$ are quoted for $`b_{max}=1.5\text{ GeV}^1`$ as $`\{a_1,a_2,a_3\}=\{0.201\pm 0.011,`$$`0.184\pm 0.018,`$$`0.026\pm 0.007\}\text{ GeV}^2`$ for $`C_3=b_0`$, and $`\{0.247\pm 0.016,`$ $`0.158\pm 0.023,`$ $`0.049\pm 0.012\}\text{ GeV}^2`$ for $`C_3=2b_0`$. In Ref. , the experimental errors are propagated into various theory predictions with the help of the Lagrange multiplier and Hessian matrix methods, discussed, e.g., in Ref. . We find that the global fit places stricter constraints on $`_{NP}`$ at $`Q=M_Z`$ than the Tevatron Run-1 $`Z`$ data alone. Theoretical uncertainties from a variety of sources are harder to quantify, and they may be substantial in the low-$`Q`$ Drell-Yan process. In particular, $`\chi ^2`$ for the low-$`Q`$ data improves by $`14`$ units when the scale parameter $`C_3`$ in $`\mu _F`$ is increased from $`b_0`$ to $`2b_0`$, reducing the size of the finite-order $`\stackrel{~}{W}_{pert}(b)`$ at low $`Q`$. The best-fit normalizations $`N_{fit}`$ also vary with $`C_3`$. The dependence of the quality of the fit on the arbitrary factorization scale $`\mu _F`$ indicates importance of $`๐ช(\alpha _s^2)`$ corrections at low $`Q`$, but does not substantially increase uncertainties at the electroweak scale. Indeed, the $`๐ช(\alpha _s^2)`$ corrections and scale dependence are smaller in $`W`$ and $`Z`$ production. In addition, the term $`a_2\mathrm{ln}Q`$, which arises from the soft factor $`๐ฎ(b,Q)`$ and dominates $`_{NP}`$ at $`Q=M_Z`$, shows little variation with $`C_3`$ \[cf. Fig. 2c\]. Consequently, the revised $`b_{}`$ model with $`b_{max}1.5\text{ GeV}^1`$ reinforces our confidence in transverse momentum resummation at electroweak scales by exposing the soft-gluon origin and universality of the dominant nonperturbative terms at collider energies.
We thank C.-P. Yuan for his crucial contribution to the setup of the fitting program, and T. Londergan, A. Szczepaniak, S. Vigdor, and CTEQ members for the helpful discussions. This work was supported by the NSF grants PHY-0100348 and PHY-0457219, and DOE grant W-31-109-ENG-38.
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# Non-perturbative orientifold transitions at the conifold
## 1 Introduction
The space of all string compactifications with $`๐ฉ=1`$ supersymmetry in four dimensions is expected to be quite rich. A poor manโs approach to the problem is to think of $`๐ฉ=1`$ as $`๐ฉ=2`$ supersymmetry broken by branes, orientifolds, and fluxes. $`๐ฉ=2`$ compactifications have been extensively studied in the framework of Type II superstrings on Calabi-Yau manifolds, and we already have a picture of the variety of vacua. The most striking aspect of this picture is that the moduli spaces corresponding to different Calabi-Yau manifolds are connected to each other through conifold transitions , which are interpreted as the $`๐ฉ=2`$ Higgs mechanism from the viewpoint of the four-dimensional spacetime physics.
Part of the motivation for our work actually comes from the desire to understand better the stringy interior of $`๐ฉ=1`$ Calabi-Yau โmoduliโ space, where we are putting quotation marks to emphasize that these moduli will generically be lifted by potentials. In particular, one would like to know which classical configurations are connected as parameters are varied, whether they are connected smoothly or through phase transitions, and what is the physical nature of the continuation or the transition. By a classical configuration, we mean the geometrical data of internal space, orientifold action, branes and fluxes in the large volume and weakly coupled regime where string and non-perturbative corrections are small. It is expected that the conifolds again play important roles.
There are various motivations to study conifolds in $`๐ฉ=1`$ compactifications. In the context of Type IIB flux compactifications , conifolds are the key to explore models with large hierarchy of scales . Also, recent study shows that a good portion of flux vacua populates a neighborhood of conifolds . Conifolds are attractive also from cosmology, see for example . In this context, conifolds were approached from the โvector sideโ (complex structure in IIB, Kรคhler class in IIA). We would like to consider also approaching from the other, โhyper sideโ (complex structure in IIA, Kรคhler class in IIB). The approach to conifolds from hypermultiplet moduli space has been studied in $`๐ฉ=2`$ systems in .
In this paper, we shall study the local behavior of the $`๐ฉ=1`$ moduli space around the conifold loci, in Type II orientifolds. Before orientifolding, the conifold locus appears in real codimension two or more in $`๐ฉ=2`$ vector multiplet or hypermultiplet moduli spaces. The orientifold projects out a part of such closed string moduli fields (see for details). In particular, it cuts a real slice in the classical geometric portion of the hypermultiplet moduli space, and the singularity appears in real codimension one. See the left hand side of Figure 1. Since scalar fields in $`๐ฉ=1`$ chiral multiplets are complex, singularities should only be expected in complex co-dimension one. The issue is that the superpartners of the geometric fields are Ramond-Ramond (RR) moduli, and the mixing between the two involves non-perturbative effects. It is therefore far from obvious what the complexified parameter space will look like in the vicinity of the classical singularity. Two possibilities are sketched on the right of Figure 1.
For example, it was found in that the tadpole cancellation conditions are generally different on the two sides of the conifold singularity in the real slice of Calabi-Yau moduli space. In other words, the charge of the orientifold plane is changing, and any claim of a smooth interpolation between the two classical limits has to account for this jumping charge.
While these points were raised for compact models, we will here focus on the singularity and study answers to these questions in the local model involving just the conifold. Although the results may not be directly applicable to compact models, we can test various methods to study the problem.
Since the conifold has been studied from a large number of perspectives in recent years, it might appear that answers to all these questions should be known. We therefore explain the basic point which we feel has not been addressed in full detail until now. Then we shall summarize our methods and results.
### 1.1 Basic question
Consider Type IIA orientifolds of the deformed conifold
$$\underset{i=1}{\overset{4}{}}z_i^2=\mu ,$$
(1.1)
with respect to the anti-holomorphic involution
$$z_i\overline{z}_i.$$
(1.2)
Under this projection, the space of complex structures of the conifold, parameterized by $`\mu `$, is restricted to the real slice $`\mu `$. The orientifold 6-plane, given by the fixed point set of (1.2) times flat four-dimensional Minkowski space is the locus of real solutions of (1.1). When $`\mu >0`$, this leaves an $`S^3`$ worth, while if $`\mu <0`$, there is no real solution, and the O-plane is empty. The point $`\mu =0`$ is the classical conifold singularity. We will refer to the transition between $`\mu `$ positive and $`\mu `$ negative as the $`\mu `$-transition. As we will see, whether or not the $`\mu `$-transition is possible depends on the case.
As alluded to above, the real parameter $`\mu `$ is complexified by a RR field, which here arises as the period of the RR three-form around the vanishing $`S^3`$ of the deformed conifold. The fundamental question is whether or not this complexification allows the two classical branches ($`\mu >0`$ and $`\mu <0`$) to be connected in the full quantum theory. Evidently there are two more classical branches joining in at $`\mu =0`$, the resolved conifolds, and the fate of these branches should be a part of the question.
More generally, we may choose to wrap D-branes on top of the O-plane on the vanishing $`S^3`$ of the deformed conifold. These will support an $`๐ฉ=1`$ gauge theory at low energies, and one has to make sure that the proposed quantum dynamics takes account of the vacuum structure.
### 1.2 Lift to M-theory
Having exposed the problem, we now explain how we will address it. The main tool for us will be the lift to M-theory, as studied in , and many other places. Recall that D6-branes and O6-planes lift in M-theory to purely geometric configurations. In the local model, the essential idea is to identify all possible classical geometries with fixed asymptotics. It then turns out that with reasonable assumptions about the dynamics, holomorphy essentially completely determines the quantum moduli space relating the various geometries.
Our second method for analyzing the possible transitions, due to is to study the critical points of a certain superpotential $`W`$. The superpotential is computed on the branch of the resolved conifold, and in the orientifold situation is a combination of flux and crosscap contributions . We give a careful analysis of this superpotential in Section 5, and reproduce the component of the moduli space including the two resolved conifold points.
### 1.3 Main results
Consider wrapping $`N`$ D6-branes on top of the O6-plane on the $`S_>^3`$ at the bottom of the orientifolded conifold (1.1), (1.2) for $`\mu >0`$. Our convention is that the total 6-brane charge as measured at infinity is $`2N4`$ in the cover ($`N2`$ in the quotient), where $`4`$ is the contribution from the O6-plane. Clearly, if a transition to $`\mu <0`$ is possible, where there is no O-plane, the charge must be carried by $`2N4`$ D6-branes wrapped on the $`S_<^3`$ at the bottom of the conifold with $`\mu <0`$. Now notice that a D6-brane wrapped on $`S_>^3`$ preserves the opposite combination of supersymmetries to a D6-brane wrapped on $`S_<^3`$. This is because the relevant calibration is the real part of the holomorphic three-form $`\mathrm{\Omega }`$, which when restricted to the two $`S^3`$โs leads to opposite orientation, depending on whether $`\mu `$ is positive or negative. In other words, if we fix the supersymmetry preserved at infinity, we can have D6-branes wrapped on $`S_>^3`$ or anti-D6-branes<sup>1</sup><sup>1</sup>1Here and throughout the paper, we refer to objects as branes or anti-branes according to the sign of the charge measured at infinity. A more natural convention would be to refer to all objects preserving the same supersymmetry as branes, but we choose our present convention to emphasize that the supersymmetric objects sometimes carry opposite charge. wrapped on $`S_<^3`$. It is clear, therefore, that we can at best expect a transition between $`\mu >0`$ and $`\mu <0`$ in the supersymmetric parameter space if $`N`$ is non-negative (so that the D6-branes preserve the same susy as the O6-plane), and $`2N4`$ is non-positive (so that to conserve the charge, we wrap $`42N`$ anti-D6-branes). In other words, we can expect a $`\mu `$-transition if $`N=0`$, $`1`$, or $`2`$.
In the previous paragraph, we discussed the possibility of wrapping an O6<sup>-</sup>-plane on $`S_>^3`$ with $`N`$ D6-branes, which yields an $`SO(2N)`$ gauge group. Alternatively, for $`N4`$ we may wrap an O6<sup>+</sup> and $`N4`$ D6-branes, yielding an $`Sp(N4)`$ gauge group. For $`N<1`$ there is a similar choice for the action of the free orientifold on the Chan-Paton matrices corresponding to the D6-branes on $`S_<^3`$. This leads to two distinct possibilities for the low energy four-dimensional gauge group, $`SO(2(2N))`$ or $`Sp(2N)`$. Finally, for any value of $`N`$ we have two semi-classical limits corresponding to the two resolved conifolds with freely acting orientifold and $`N2`$ units of RR 2-from flux through $`^2`$.
With these observations in mind, we can identify all the possible semi-classical limits in the IIA description for each value of $`N`$. Our results are summarized in the left-hand side of Table 1.
The M-theory lifts of the various semi-classical limits are described in the right-hand side of Table 1. The problem for $`N{}_{(}{}^{}\mathrm{}_)3`$ is included in . For infinite string coupling (the size of the M-theory circle growing without bounds asymptotically), the M-theory geometries are quotients of smooth $`G_2`$ holonomy manifolds, $`X_i`$, by the dihedral group $`D_N`$. The $`X_i`$ are all isomorphic to the spin bundle over $`S^3`$, whose $`G_2`$ holonomy metric was found in . They differ in the breaking pattern of the asymptotic discrete symmetries, and discrete fluxes at the singularities. In the $`D_N`$ case, there are four semi-classical limits corresponding to the four IIA limits described above.
For $`N=3`$, the dihedral group $`D_3`$ is isomorphic to $`A_3=_4`$, and the problem is equivalent, from the M-theory perspective, to a case without orientifold. But for the appropriate identification of the M-theory circle, we still end up with an orientifold in Type IIA.
The extension to lower values of $`N`$, as well as the generalization to finite values of the string coupling, requires more complicated $`G_2`$ holonomy metrics with reduced symmetry. These metrics have been partially constructed in , and we now describe their relevance to our problem.
For $`N>2`$, if one wishes to keep the asymptotic IIA string coupling finite, the M-theory lift of the deformed conifold geometry involves a $`G_2`$ metric called $`๐น_7`$ in . Roughly, the manifold $`๐น_7`$ is a Taub-NUT manifold fibered over an $`S^3`$. Quotienting $`๐น_7`$ by the dihedral group leaves a $`D_N`$ singularity supporting an $`SO(2N)`$ or $`Sp(N4)`$ gauge group, depending on a discrete flux. The resolved conifold with flux (and finite string coupling) has an M-theory lift called $`๐ป_7`$ in . This $`G_2`$ metric is smooth and $`D_N`$ acts freely on $`๐ป_7`$.
For $`N=2`$, the M-theory lift will be
$$\frac{\text{(conifold)}\times S^1}{_2},$$
where $`_2`$ acts as an antiholomorphic involution on the conifold and reverses the M-theory $`S^1`$. As such, we know the classical M-theory geometry exactly.
For $`N=0`$ or $`N=1`$, one can guess that the M-theory lift of the deformed conifold with O-plane (namely, $`\mu >0`$) will look like an Atiyah-Hitchin manifold or Dancerโs manifold fibered over an $`S^3`$. Such $`G_2`$ metrics were not previously known but we find that they do indeed exist. We will call these manifolds $`๐ธ_7/_2`$ and $`๐ธ_7`$, respectively.
In all cases with $`N<2`$, the deformed conifold with $`\mu <0`$ (which is wrapped by $`42N`$ anti-D6-branes) and the resolved conifold with flux again lift to the quotient of $`๐น_7`$ and $`๐ป_7`$, respectively, by the dihedral group $`D_{4N}`$. The difference to the $`N>2`$ case is in the action of $`D_{4N}`$ on the space, as we will explain in more detail later. The gauge group living on the $`42N`$ anti-D6-branes in the freely acting orientifold can be $`\mathrm{๐๐}(2(2N))`$ or $`\mathrm{๐๐}(2N)`$, depending on the action on the Chan-Paton factors. This corresponds to the value of a discrete torsion. Thus, the $`๐น_7/D_{4N}`$ geometries yield two semi-classical limits with $`\mu <0`$ in each case.
No analytic expressions are known for any of the $`G_2`$ metrics $`๐ธ_7`$, $`๐น_7`$, $`๐ป_7`$ except one special point on the parameter space of $`๐น_7`$ found in , as well as the limits of zero or infinite string coupling. From symmetry requirements, one can determine the metrics up to a small number of unknown functions (of a radial coordinate) and derive a set of differential equations for the unknowns. It is not difficult to verify numerically the existence of such solutions. One can also see that they depend on two parameters, one corresponding to the radius of the M-theory circle at infinity and the other to the volume of $`S^3`$ at the center.
The next question is how these semi-classical limits with the same asymptotic flux, $`N2`$, and the same supersymmetry fit together into a complex space parameterizing supersymmetric vacua. For $`N{}_{(}{}^{}\mathrm{}_)3`$, which is the case discussed in , the parameter space is a copy of $`^1`$ with four (three for $`N=3`$) marked points corresponding to the semi-classical limits we have discussed above. Good local coordinates around each of these points correspond to the volume deficits of certain three-cycles in the asymptotic geometries together with the period of the M-theory three-form around the same three-cycles. These complex parameters are the instanton coefficients in the four-dimensional low-energy gauge theory.
The case $`N=2`$ was discussed in , where it was argued that a $`\mu `$-transition should be possible. Our analysis will confirm the expectation in this case. In addition, we will find a second branch of moduli space containing two vacua of $`SO(4)`$ gauge theory as well as the two resolved conifolds that were not treated in .
In the case $`N=1`$, there are five semi-classical limits. As can be seen from Table 1 two of them have an $`\mathrm{๐๐}(2)U(1)`$ gauge theory at low energies, one has an $`\mathrm{๐๐}(1)\mathrm{๐๐}(2)`$, while the two remaining ones have no gauge theory at all. The first two have free massless gauge bosons while the latter three points do not. Since the massless spectrum is different, one must pass through a phase transition when interpolating between the various limits. This situation is very similar to ones studied recently in , and we will be able to use these methods to deduce the structure of the quantum parameter space, confirming the naive picture we have just sketched.
When $`N=0`$, it appears at first sight that we have also five semi-classical limits: two on the resolved conifold, two from the deformed conifold with $`\mu <0`$, and one from the deformed conifold with $`\mu >0`$. However, as we will see, there are in fact two distinct semi-classical limits corresponding to just the O-plane wrapping the $`S^3`$. In M-theory, this can be simply seen from the existence of an asymptotic discrete $`_2^{}`$ symmetry that is spontaneously broken in the interior of $`๐ธ_7/_2`$. In Type IIA string theory, this symmetry corresponds to D0-brane charge modulo $`2`$, which is broken only at the bare orientifold plane, but is preserved in the presence of just a single D6-brane on top. (We will explain why this statement is not in conflict with the K-theory classification of D-brane charge.) There are therefore six semi-classical limits to consider. We will argue that the quantum parameter space consists of two disconnected branches, with a certain distribution of vacua consistent with the discrete symmetries of the problem.
For $`N<0`$, we again have four semi-classical limits, each of which has a mass gap at low energies. From the analysis of the holomorphic parameters associated with the gauge theories, as well as our later superpotential analysis, we will deduce that the space on which these four limits sit is again a copy of $`^1`$. In fact, we will see that the curve for $`N<0`$ is isomorphic to the curve for $`N^{}=4N>4`$, with the only difference being the association between two of the points and $`\mathrm{๐๐}/\mathrm{๐๐}`$ gauge group!
### 1.4 Summary
In addition to the literature that we have cited already, aspects of the problem have also been discussed elsewhere. The basic question whose solution we have presented in the previous subsection has been broached, for example, in , with a restriction to the locally tadpole canceling case with exactly two D6-branes on top of the O-plane. More recently, similar transitions have been found in Type IIB compactifications in . The main result of our paper is to explain under which conditions in the local model we can actually expect a transition, and to determine the quantum parameter space, whenever possible. Our main tool of analysis is the lift to M-theory on $`G_2`$ manifolds which are described as quotients by finite groups. Similar system have been analyzed in the past also in , for example. Conifold transitions in the $`G_2`$ context have also been studied in .
We also have a number of subsidiary and complementary results to offer, as we now summarize. We will start out with a classification of A-type orientifolds <sup>2</sup><sup>2</sup>2The construction of several orientifolds of the conifold in Type IIB has been analyzed in . These orientifolds were constructed by partially blowing up orientifolds of orbifold singularities. of the deformed conifold in Section 2. In this section, we also show that when passing through a โ$`\mu `$-transitionโ between $`\mu >0`$ and $`\mu <0`$ (independent of whether this is possible dynamically or not), the orientifold charge changes by the class of the vanishing cycle. We also study how the previous orientifolds act on the resolved conifold. We then return to the main case of interest, the orientifold (1.2). In Section 3, we explain generalities about the symmetries of the underlying $`G_2`$ holonomy metrics. The nuts and bolts of these space are assembled in Section 4. We find a new class of $`G_2`$ holonomy metrics which we label $`๐ธ_7`$. In Section 5, we derive the quantum parameter spaces. In particular, we show how a careful analysis of the Vafa superpotential produces most of the structure of the parameter space. For the cases $`N=2,1`$ and $`0`$, we need to invoke some additional information, partly from . In Section 7, we briefly discuss the problem of $`\mu `$-transition for the other classes of orientifolds. Since they break more of the geometrical symmetries, we are unable to write down explicit metrics. The Vafa superpotential gives a prediction for one other class of models. However, when the orientifold does not admit the resolved conifold, the superpotential method is not applicable. In these cases, we describe our best educated guesses for possible $`\mu `$-transitions. We finally conclude in Section 8.
## 2 Orientifolds of the conifold
In this section, we describe the possible A-type orientifolds of the deformed conifold. In particular, we will see that when crossing the conifold the class of the orientifold locus changes by the class of the vanishing cycle. We then analyze how these orientifolds act on the resolved conifold.
### 2.1 Deformed conifold
An orientifold of the deformed conifold
$$z_1^2+z_2^2+z_3^2+z_4^2=\mu $$
(2.1)
in Type IIA string theory can be obtained from an anti-holomorphic involution, which acts on the complex coordinates $`z_i`$ by complex conjugation followed by a symmetry of the quadric (2.1),
$$z_ie^{i\alpha }M_i^j\overline{z}_j.$$
(2.2)
Here $`M`$ is an orthogonal matrix with $`M^2=1`$, and $`\alpha `$ is a phase, which if we assume that $`\mu `$ is real can be set to zero.
Since all such orthogonal matrices can be diagonalized with $`\pm 1`$ on the diagonal, inequivalent anti-holomorphic involutions of the deformed conifold are classified by the number of $`+1`$ and $`1`$ eigenvalues of $`M`$.
In each of these five possibilities, the fixed point set is described by setting $`z_i=x_i`$ or $`z_i=iy_i`$, with $`x_i`$, $`y_i`$ real, depending on the corresponding sign. The fixed point locus are the orientifold 6-planes. Let us take $`\mu >0`$ (the case with $`\mu `$ negative can easily be obtained by interchanging the real and imaginary components of the $`z_i`$). We have the following inequivalent cases:
* For $`(z_1,z_2,z_3,z_4)(\overline{z}_1,\overline{z}_2,\overline{z}_3,\overline{z}_4)`$ the O6-plane is described by
$$x_1^2+x_2^2+x_3^2+x_4^2=\mu .$$
(2.3)
Since $`\mu >0`$, this is an $`S^3`$.
* If the involution takes $`(z_1,z_2,z_3,z_4)(\overline{z}_1,\overline{z}_2,\overline{z}_3,\overline{z}_4)`$ the orientifold set is non-compact,
$$y_1^2+x_2^2+x_3^2+x_4^2=\mu ,$$
(2.4)
and isomorphic to $`S^2\times `$.
* When $`(z_1,z_2,z_3,z_4)(\overline{z}_1,\overline{z}_2,\overline{z}_3,\overline{z}_4)`$ the O-plane is at
$$y_1^2y_2^2+y_3^2+y_4^2=\mu ,$$
(2.5)
which describes $`S^1\times ^2`$.
* If $`(z_1,z_2,z_3,z_4)(\overline{z}_1,\overline{z}_2,\overline{z}_3,\overline{z}_4)`$ the orientifold set is not connected: The equation
$$y_1^2y_2^2y_3^2+x_4^2=\mu $$
(2.6)
is solved by two copies of $`^3`$.
* Finally, when $`(z_1,z_2,z_3,z_4)(\overline{z}_1,\overline{z}_2,\overline{z}_3,\overline{z}_4)`$, the orientifold set is empty:
$$y_1^2y_2^2y_3^2y_4^2=\mu $$
(2.7)
has no real solutions.
Generalizing the relation between (0) and (4) discussed in the introduction, flipping the sign of $`\mu `$ maps (1) to (3) and (2) to itself.
We can introduce D6-branes wrapping supersymmetric cycles of the deformed conifold. In order to have dynamical gauge symmetry in four dimensions, the cycle must be compact, and the only possibility is the three-dimensional sphere at the center of the conifold. In the cases (1) and (3), the involution reverses the orientation of the $`S^3`$, i.e., maps branes to anti-branes. Therefore, wrapping branes on $`S^3`$ in these cases will break supersymmetry. In fact, in those cases, there is a RR tadpole which originates from the non-trivial intersection of the O-plane with the compact $`S^3`$. To get a stable configuration, we are forced to wrap branes on some other non-compact cycles. This yields quite an interesting class of models, which is briefly analyzed in Section 7. The case (0) is supersymmetric: the D6-branes wrap the same $`S^3`$ as the orientifold. In the case (4) the action of the orientifold is free, the D6-branes are then wrapping an $`^3`$. In case (2) the orientifold plane intersects the $`S^3`$ where the D6-branes are wrapping in an $`S^1`$. We have depicted these circumstances in Figure 2.
Compact orientifold models including invariant conifolds must be one of these five types. For example, let us consider the Type IIA orientifold on the mirror $`X`$ of the Fermat quintic, which is a resolution of the orbifold of
$$z_1^5+z_2^5+z_3^5+z_4^5+z_5^55\psi z_1z_2z_3z_4z_5=0$$
by the $`_5^3`$ action, $`z_i\omega _iz_i`$, $`\omega _i^5=\omega _i=1`$. For $`\psi =1`$, $`X`$ has a single conifold singularity at $`z_1=\mathrm{}=z_5=1`$. We consider the orientifold with respect to the involution $`\tau _\sigma :z_i\overline{z}_{\sigma (i)}`$, where $`\sigma `$ is an order two permutation (exchange), that fixes the conifold singularity when $`\psi =1`$. The orientifold is allowed only when $`\psi `$ is real. There are three distinct cases โ $`\sigma `$ is identity, an exchange of a pair, and an exchange of two pairs. Depending on the sign of $`ฯต=\psi 1`$, these cases are one of the five possibilities: Without exchange, $`z_i\overline{z}_i`$, it is case (0) if $`ฯต>0`$ and case (4) if $`ฯต<0`$. With an exchange of one pair, such as $`(z_1,z_2,z_3,z_4,z_5)(\overline{z}_2,\overline{z}_1,\overline{z}_3,\overline{z}_4,\overline{z}_5)`$, it is case (1) if $`ฯต>0`$ and case (3) if $`ฯต<0`$. With an exchange of two pairs, such as $`(z_1,z_2,z_3,z_4,z_5)(\overline{z}_2,\overline{z}_1,\overline{z}_4,\overline{z}_3,\overline{z}_5)`$, it is case (2) for both signs of $`ฯต`$. Note that cases (0) and (4) are mirror to Type IIB orientifold with an O9-plane (Type I), case (1) and (3) are mirror to IIB orientifold with O3/O7 planes, and case (2) is mirror to IIB orientifold with O5-planes .
### 2.2 Homology of O-planes
We can study the homology classes of the Lagrangian manifolds of the orientifold loci by computing the integral of the holomorphic three-form $`\mathrm{\Omega }`$ around them, where
$$\mathrm{\Omega }=\frac{dz_1dz_2dz_3}{z_4}$$
(2.8)
on the sheet $`z_4=\sqrt{\mu z_1^2z_2^2z_3^2}`$.
We find that
* For the case (0) the orientifold wraps a compact manifold with period:
$$\varpi ^{(0)}(\mu )=_{S^3}\mathrm{\Omega }=2_{x_1^2+x_2^2+x_3^2\mu }\frac{dx_1dx_2dx_3}{\sqrt{\mu x_1^2x_2^2x_3^2}}=8\pi \mu _0^1\frac{r^2dr}{\sqrt{1r^2}}=2\pi ^2\mu .$$
(2.9)
* The orientifold locus is non-compact, and we need to introduce a cutoff, which we put at $`|z|^2=x^2+y^2=\mathrm{\Lambda }`$.
$$\begin{array}{cc}\hfill \varpi ^{(1)}(\mu )& =_{S^2\times }\mathrm{\Omega }=2i_{\mu x_2^2+x_3^2+x_4^2\frac{\mathrm{\Lambda }+\mu }{2}}\frac{dx_2dx_3dx_4}{\sqrt{x_2^2+x_3^2+x_3^2\mu }}\hfill \\ & =8\pi i\mu _1^{\sqrt{\frac{\mathrm{\Lambda }+\mu }{2\mu }}}\frac{r^2dr}{\sqrt{r^21}}=4\pi i\mu \left[\frac{\sqrt{\mathrm{\Lambda }^2\mu ^2}}{2\mu }+\mathrm{log}\frac{\sqrt{\mathrm{\Lambda }+\mu }+\sqrt{\mathrm{\Lambda }\mu }}{\sqrt{2\mu }}\right].\hfill \end{array}$$
(2.10)
* As in the previous case, the orientifold is non-compact.
$$\begin{array}{cc}\hfill \varpi ^{(2)}(\mu )& =_{S^1\times ^2}\mathrm{\Omega }=2_{\genfrac{}{}{0.0pt}{}{\mu +y_1^2+y_2^2x_3^20}{y_1^2+y_2^2{\scriptscriptstyle \frac{\mathrm{\Lambda }\mu }{2}}}}\frac{dy_1dy_2dx_3}{\sqrt{\mu +y_1^2+y_2^2x_3^2}}\hfill \\ & =\pi ^2(\mathrm{\Lambda }\mu ).\hfill \end{array}$$
(2.11)
* We have:
$$\begin{array}{cc}\hfill \varpi ^{(3)}(\mu )& =_^3\mathrm{\Omega }=2i_{y_1^2+y_2^2+y_3^2\frac{\mathrm{\Lambda }\mu }{2}}\frac{dy_1dy_2dy_3}{\sqrt{\mu +y_1^2+y_2^2+y_3^2}}\hfill \\ & =8\pi i\mu _0^{\sqrt{\frac{\mathrm{\Lambda }\mu }{2\mu }}}\frac{r^2dr}{\sqrt{1+r^2}}=4\pi i\mu \left[\frac{\sqrt{\mathrm{\Lambda }^2\mu ^2}}{2\mu }\mathrm{log}\frac{\sqrt{\mathrm{\Lambda }+\mu }+\sqrt{\mathrm{\Lambda }\mu }}{\sqrt{2\mu }}\right].\hfill \end{array}$$
(2.12)
* This is empty.
If we evaluate case (1) for $`\mu <0`$ (following the computation in case (3)), and subtract the result from the direct analytic continuation of (1), the difference is
$$\varpi ^{(1)}(\mu )\varpi ^{(1)}(\mu )=4\pi i\mu \mathrm{log}\sqrt{1}=2\pi ^2\mu ,$$
(2.13)
which is exactly the period of the vanishing $`S^3`$. The same result also holds for case (2)
$$\varpi ^{(2)}(\mu )\varpi ^{(2)}(\mu )=\pi ^2(\mathrm{\Lambda }\mu )(\pi ^2(\mathrm{\Lambda }+\mu ))=2\pi ^2\mu ,$$
(2.14)
and trivially for (0)/(4).
We can understand this by noting that the transition does not affect the boundary of the O-plane and so we may glue the O-planes for $`\mu <0`$ and $`\mu >0`$ along their common boundary to form a compact three-cycle which must then be homologous to an integer multiple of the minimal $`S^3`$. The calculation shows that this integer is one.
It is natural to propose that this holds as a universal result, and not just for the simple conifold singularity we have studied here. We conjecture:
When crossing the conifold locus of real co-dimension one in the geometric moduli space of an orientifold model, the class of the O-plane changes by the class of the vanishing cycle.
The non-compact part of the homology of the O-planes (which contributes the $`\mu \mathrm{log}\mu `$ part in the expressions above) can be understood from the intersection with the compact three-cycle. It is easy to see that when we make the intersection between O-plane and the $`S^3`$ transversal, the $`S^1`$ in case $`(2)`$ disappears completely, while the $`S^2`$ leaves two intersection points. This follows from the fact that for a Lagrangian submanifold such as our three-cycles, the normal bundle is isomorphic to the tangent bundle via contraction with the Kรคhler form. The number of intersection points is then simply the Euler characteristic of the (non-transversal) intersection locus.
### 2.3 Orientifolds of the Resolved Conifold
Let us now discover what these involutions look like for the blown-up conifold. This space, also known as $`๐ช_^1(1)๐ช_^1(1)`$, can be written as
$$\left(\begin{array}{cc}x& u\\ v& y\end{array}\right)\left(\begin{array}{c}\lambda _1\\ \lambda _2\end{array}\right)=0,$$
(2.15)
with $`(\lambda _1,\lambda _2)^1`$ and
$$\begin{array}{cc}\hfill x& =z_1+iz_2,y=z_1iz_2,\hfill \\ \hfill u& =z_3+iz_4,v=z_3+iz_4.\hfill \end{array}$$
(2.16)
The blow up breaks the $`O(4)`$ symmetry we have used above to $`\mathrm{๐๐}(4)`$, while at the same time restoring the $`U(1)`$ phase symmetry. Thus, in (2.2), $`M`$ must have determinant $`+1`$, while $`\alpha `$ is a priori arbitrary. But note that by a change of coordinates, $`\alpha `$ can be conjugated to $`0`$. The condition that the number of $`1`$ eigenvalues of $`M`$ must be even eliminates cases (1) and (3) discussed previously for the deformed conifold, while (0) and (4) become equivalent on the blown-up side. In cases (1) and (3) the Kรคhler parameter is projected out by the orientifold action and it is not possible to blow up the singularity. This leaves two cases:
* (equivalent to (4)) maps $`(x,y,u,v)(\overline{y},\overline{x},\overline{v},\overline{u})`$. In terms of the inhomogeneous coordinate $`z=\lambda _1/\lambda _2=u/x=y/v`$ the $`^1`$ is mapped as
$$z\frac{1}{\overline{z}}.$$
(2.17)
This is a freely acting orientifold.
* The action $`(x,y,u,v)(\overline{y},\overline{x},\overline{v},\overline{u})`$ must be accompanied by
$$z\frac{1}{\overline{z}}.$$
(2.18)
The fixed point set is an $`S^1\times ^2`$.
We note that these cases coincide with the cases where the orientifold action preserves the same supersymmetry as the D6-branes wrapping the $`S^3`$, and are precisely the cases discussed by Acharya, Aganagic, Hori and Vafa in .
### 2.4 The gauge group
For future reference, it is useful to describe here which gauge theories will be living on the worldvolume of D6-branes that are wrapping this geometry.
In flat space $`N`$ dynamical D6-branes on the top of an O6<sup>-</sup>-plane yield an $`\mathrm{๐๐}(2N)`$ gauge group on the worldvolume. The O6<sup>-</sup>-plane has a RR charge $`2`$ in D6-brane units, so the total charge of the system is $`N2`$. When wrapping an $`S^3`$ the D6-branes cannot be higgsed away. The four-dimensional gauge theory is pure $`๐ฉ=1`$ $`\mathrm{๐๐}(2N)`$ super Yang-Mills. This theory is confining with $`h=2N2`$ different vacua.
A similar classical configuration is a system of $`N4`$ D6-branes on the top of an O6<sup>+</sup>-plane. The O6<sup>+</sup>-plane has RR charge $`+2`$, so the whole system also has charge $`N2`$. The low energy theory is pure $`๐ฉ=1`$ $`\mathrm{๐๐}(N4)`$ super Yang-Mills. The theory is confining and has $`h=N3`$ different vacua.
This discussion was of course standard. Slightly less familiar are D-branes wrapping in freely acting orientifolds (such as case (4) above), but it is also clear what will result. The involution acts on the $`S^3`$ as the antipodal map, $`xx`$. Locally, this simply identifies excitations at antipodal points on the sphere, via an anti-unitary transformation on the D6-brane degrees of freedom. That can be understood as an action on the Chan-Paton matrices:
$$\lambda (x)\gamma _\mathrm{\Omega }\lambda (x)^T\gamma _\mathrm{\Omega }^1,$$
(2.19)
where $`\lambda `$ are $`2M\times 2M`$ hermitian matrices for $`2M`$ D6-branes wrapping the covering $`S^3x`$. Locally on $`S^3`$, this simply yields a $`U(2M)`$ gauge group. The zero modes, however, suffer a slightly different projection. The orientifold action is an involution if $`\gamma _\mathrm{\Omega }(\gamma _\mathrm{\Omega }^T)^1=ฯต`$, with $`ฯต=\pm 1`$, i.e. $`\gamma _\mathrm{\Omega }`$ is symmetric or antisymmetric. Depending on the sign the four-dimensional theory will be pure super Yang-Mills with gauge group $`\mathrm{๐๐}(2M)`$ or $`\mathrm{๐๐}(M)`$. As before the system preserves $`๐ฉ=1`$.
## 3 Symmetries of $`G_2`$ holonomy metrics
In this section and the next, we discuss aspects of the $`G_2`$ lift of the deformed and resolved conifold with branes and fluxes. Most of this section is review , but the careful discussion of the symmetry breaking pattern will be crucial in our subsequent analysis.
### 3.1 Deformed and resolved conifold
We begin by recording the symmetry group of the conical Calabi-Yau metric on the (singular) conifold,
$$z_1^2+z_2^2+z_3^2+z_4^2=0.$$
(3.1)
These isometries must preserve (3.1) together with
$$r^2=|z_i|^2.$$
(3.2)
In full glory, the symmetry group is
$$(\mathrm{๐๐}(2)\times \mathrm{๐๐}(2)\stackrel{~}{}_2\times U(1)^{\mathrm{phase}}_2^{\mathrm{cc}})/_{_2\times _2}.$$
(3.3)
Here, the $`\mathrm{๐๐}(4)\mathrm{๐๐}(2)\times \mathrm{๐๐}(2)/_2`$ is extended by $`\stackrel{~}{}_2`$ to the $`O(4)`$ leaving the quadric invariant. $`\stackrel{~}{}_2`$ acts as
$$(z_1,z_2,z_3,z_4)(z_1,z_2,z_3,z_4).$$
The $`U(1)^{\mathrm{phase}}`$ contains the rotations
$$R_\alpha :z_ie^{i\alpha /2}z_i$$
with $`\alpha [0,4\pi ]`$, and $`\alpha =2\pi `$ corresponding to an element of $`\mathrm{๐๐}(4)`$. Finally, $`_2^{\mathrm{cc}}`$ is complex conjugation, represented by
$$c_0:z_i\overline{z}_i.$$
(3.4)
When conjugated by elements of $`U(1)^{\mathrm{phase}}`$, $`c_0`$ becomes
$$c_\alpha =R_\alpha c_0R_\alpha ^1:z_ie^{i\alpha }\overline{z}_i.$$
(3.5)
When the singular conifold is smoothed out, some of these symmetries are broken. The deformation which replaces (3.1) by
$$z_i^2=\mu $$
(3.6)
breaks $`U(1)^{\mathrm{phase}}`$ to the $`_2`$ which is already part of $`O(4)`$ (namely $`\alpha =2\pi `$). Since (3.5) is a symmetry of (3.6) for both $`\alpha =\mathrm{arg}(\mu )`$ and $`\alpha =\mathrm{arg}(\mu )+\pi `$, we have a choice of orientifold $`c_0:z_i\overline{z}_i`$ or $`c_\pi :z_i\overline{z}_i`$.
The blowup of the conifold is obtained by rewriting (3.1) as
$$xyuv=0,$$
(3.7)
and then replacing it with the two equations
$$\left(\begin{array}{cc}x& u\\ v& y\end{array}\right)\left(\begin{array}{c}\lambda _1\\ \lambda _2\end{array}\right)=0$$
(3.8)
in $`^4\times ^1`$. The blowup clearly preserves $`U(1)^{\mathrm{phase}}`$, but breaks $`\stackrel{~}{}_2`$. Indeed, transposition of $`z`$ amounts to exchanging $`u`$ and $`v`$ and is equivalent to flopping the $`^1`$.
As equation (3.5) shows, orientifolding breaks $`U(1)^{\mathrm{phase}}`$ down to the $`_2`$ subgroup which is in $`SO(4)`$. The addition of RR 2-form flux also breaks this $`U(1)^{\mathrm{phase}}`$ as will become clear after lifting to M-theory. Thus the symmetry group of the asymptotic IIA geometry relevant to our problem is:
$$\left(\mathrm{๐๐}(4)\stackrel{~}{}_2\times _2^{\mathrm{cc}}\right).$$
(3.9)
Deforming the conifold leaves these symmetries intact whilst resolving breaks $`\stackrel{~}{}_2`$.
### 3.2 Lift
Next we would like to lift to M-theory and understand the action of the symmetry group (3.9) there. The asymptotic boundary of the conifold is $`T^{1,1}=\mathrm{๐๐}(2)\times \mathrm{๐๐}(2)/U(1)`$ and in the presence of RR 2-form flux the boundary of the M-theory lift is an orbifold of $`\mathrm{๐๐}(2)\times \mathrm{๐๐}(2)`$.
This can be described a little more explicitly by introducing some coordinates on the conifold. We write quaternionically
$$z=\left(\begin{array}{cc}z_1+iz_2& z_3+iz_4\\ z_3+iz_4& z_1iz_2\end{array}\right)=x+iy,$$
(3.10)
with
$$\begin{array}{cc}\hfill x& =\left(\begin{array}{cc}x_1+ix_2& x_3+ix_4\\ x_3+ix_4& x_1ix_2\end{array}\right)=X\stackrel{~}{X}^{},\hfill \\ \hfill y& =\left(\begin{array}{cc}y_1+iy_2& y_3+iy_4\\ y_3+iy_4& y_1iy_2\end{array}\right)=X\sigma \stackrel{~}{X}^{}\hfill \end{array}$$
(3.11)
and $`\sigma =\left(\begin{array}{cc}i& 0\\ 0& i\end{array}\right)`$. Here we have used that (3.1) and (3.2) imply that (for $`r=\sqrt{2}`$) $`x`$, $`y`$ are orthogonal unit quaternions
$$detx=dety=1\mathrm{Tr}y^{}x=0.$$
(3.12)
Choosing our standard traceless $`\mathrm{๐๐}(2)`$ matrix $`\sigma `$, this equation is solved in terms of two $`\mathrm{๐๐}(2)`$ matrices $`X`$ and $`\stackrel{~}{X}`$, modulo the relation $`(X,\stackrel{~}{X})(X\mathrm{\Theta },\stackrel{~}{X}\mathrm{\Theta })`$ with $`\mathrm{\Theta }=\left(\begin{array}{cc}e^{i\psi /2}& 0\\ 0& e^{i\psi /2}\end{array}\right)`$. This exhibits the base of the conifold as either a (topologically trivial) $`S^2`$ bundle over $`S^3`$ or as $`T^{1,1}=\mathrm{๐๐}(2)\times \mathrm{๐๐}(2)/U(1)`$. (We wonโt need its presentation as $`T^{1,0}`$, which uses $`x`$ and $`X`$ instead.)
$`X`$ and $`\stackrel{~}{X}`$ are the $`S^3\times S^3`$ of the M-theory lift. The action of $`SO(4)`$ is by left multiplication of $`SU(2)\times SU(2)`$ on $`X`$ and $`\stackrel{~}{X}`$. We can choose $`\stackrel{~}{}_2`$ to be $`z_1z_1`$ which corresponds to exchanging $`X`$ and $`\stackrel{~}{X}\beta `$, where $`\beta =\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right)`$. The action of complex conjugation $`c_0`$ is given by
$$XX\beta ,\stackrel{~}{X}\stackrel{~}{X}\beta .$$
(3.13)
Adding the M-theory circle adds to the asymptotic symmetry group a factor of $`U(1)^M`$ which is given by right multiplication of $`X`$ and $`\stackrel{~}{X}`$ by $`\mathrm{\Theta }`$ as described above. Note that the orientifold action (3.13) inverts the M-theory circle as it should. Also note that the right action of $`\mathrm{\Theta }`$ and $`\beta `$ on $`(X,\stackrel{~}{X})`$ generates a larger group of symmetries which in particular contains the dihedral group $`D_N`$ for any $`N`$. We denote the group generated in this way by $`G`$.
A further useful set of coordinates is given by writing $`S^3\times S^3`$ as the quotient $`\mathrm{๐๐}(2)^3/\mathrm{๐๐}(2)`$,
$$(g_1,g_2,g_3)(g_1g^{},g_2g^{},g_3g^{})\mathrm{๐๐}(2)^3/\mathrm{๐๐}(2).$$
(3.14)
The base of the conifold is obtained by reduction along the maximal torus of $`g_1`$. The explicit identification is
$$\begin{array}{cc}\hfill x& =g_2g_3^1,\hfill \\ \hfill \stackrel{~}{X}& =g_3g_1^1,\hfill \\ \hfill X& =x\stackrel{~}{X}=g_2g_1^1.\hfill \end{array}$$
(3.15)
Note that in these coordinates the flop is described by
$$(g_1,g_2,g_3)(\beta g_1,g_3,g_2).$$
(3.16)
Topologically then, the boundary of our $`G_2`$-holonomy manifolds will be an orbifold of $`S^3\times S^3`$. One interesting metric on $`S^3\times S^3`$, which is the one underlying the $`G_2`$ metric on the spin bundle on $`S^3`$ , has $`\mathrm{๐๐}(2)^3\times \mathrm{\Sigma }_3`$ symmetry, where $`\mathrm{๐๐}(2)^3`$ acts on the left in (3.14), and $`\mathrm{\Sigma }_3`$ is the permutation of the three $`\mathrm{๐๐}(2)`$ factors . $`\mathrm{\Sigma }_3`$ is โspontaneouslyโ broken in the interior, and only $`\mathrm{๐๐}(2)^3\times _2`$ are isometries of the full $`G_2`$ holonomy metrics. These metrics are relevant to the problem at infinite IIA coupling.
The metrics that are relevant for our discussion (at finite string coupling) have asymptotic symmetry group $`(\mathrm{๐๐}(2)\times \mathrm{๐๐}(2)\stackrel{~}{}_2\times G)/__2`$, where $`G`$ was defined above as the group generated by $`U(1)^M`$ and the orientifold action (3.13). The trivial $`_2`$ is generated by the element $`(1,1,1)\mathrm{๐๐}(2)\times \mathrm{๐๐}(2)\times U(1)^M`$.
In the interior, various symmetry breaking patterns are possible. Moreover, in certain cases it happens that the symmetry group is enhanced to $`\mathrm{๐๐}(2)^3`$ in the deep interior.
### 3.3 Orientifolding
As in , we can consider dividing out by the action of a discrete group $`\mathrm{\Gamma }`$ preserving the $`G_2`$ metric. Of interest to us is the case that $`\mathrm{\Gamma }`$ is the (binary) dihedral group $`D_N`$. This group has generators $`a`$, and $`b`$ satisfying the relations
$$a^{2N4}=1,b^2=a^{N2},bab^1=a^1.$$
(3.17)
The group has a presentation
$$_{2N4}D_N_2.$$
(3.18)
The dihedral group has a standard action on the three sphere coming from its embedding as a discrete subgroup of $`\mathrm{๐๐}(2)`$, namely
$$a=\left(\begin{array}{cc}e^{\pi i/(N2)}& 0\\ 0& e^{\pi i/(N2)}\end{array}\right),b=\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right).$$
(3.19)
We will call the action of $`D_N`$ on $`S^3`$ in which $`a`$ and $`b`$ are represented in this way as $`\rho `$. The quotient of our interest is obtained by letting $`D_N`$ act as $`\rho `$ on the right on both $`X`$ and $`\stackrel{~}{X}`$, in the variables (3.15). The effect of the $`_{2N4}`$ factor in (3.18) is to reduce the length of the M-theory circle, thereby increasing the flux to $`2N4`$ units. After reduction on $`U(1)^M`$, the remaining $`_2`$ sends $`(x,y)(x,y)`$, so is indeed complex conjugation $`c_0`$. In terms of the variables $`(g_1,g_2,g_3)\mathrm{๐๐}(2)^3/\mathrm{๐๐}(2)`$, $`D_N`$ acts as $`\rho `$ on $`g_1`$ and trivially on $`g_2`$ and $`g_3`$.
We have discussed this action in detail in order to make the following point. There is another action, call it $`\stackrel{~}{\rho }`$, of $`D_N`$ on $`S^3`$ in which $`b`$ is represented by $`\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right)`$. Acting with $`D_N`$ as $`\stackrel{~}{\rho }`$ on the right of $`X`$ and as $`\rho `$ on the right of $`\stackrel{~}{X}`$ is equivalent to acting on $`g_1`$ as $`\rho `$ and on $`g_2`$ via the central action, $`a=1`$, $`b=1`$. After reduction to $`T^{1,1}`$, the action is $`(x,y)(x,y)`$, corresponding to $`c_\pi `$. We shall sometimes refer to this action as $`D_N^{}`$.
Note that the actions of $`D_N`$ and $`D_N^{}`$ can be conjugated into each other by a diffeomorphism of $`S^3\times S^3`$, sending $`(X,\stackrel{~}{X})(Xi\sigma _3,\stackrel{~}{X})`$. However, this diffeomorphism is not an isometry of the boundary metric since $`U(1)^{\mathrm{phase}}`$ is explicitly broken by the addition of flux.
We should now ask, what are the symmetries of the boundary metric which also preserve the orbifold group $`D_N`$ or $`D_N^{}`$? In each case, $`(\mathrm{๐๐}(2)\times \mathrm{๐๐}(2)\stackrel{~}{}_2\times G)/__2`$ is broken to $`\mathrm{๐๐}(4)\stackrel{~}{}_2\times _2^{}`$ where $`_2^{}`$ is the centralizer of $`D_N`$ or $`D_N^{}`$ in $`G`$ and is generated by
$$\sqrt{a}=\left(\begin{array}{cc}e^{\pi i/2(N2)}& 0\\ 0& e^{\pi i/2(N2)}\end{array}\right).$$
(3.20)
We can now make a few comments about the pattern of symmetry breaking by the various geometries in the interior. As we have already commented, we expect that $`\stackrel{~}{}_2`$ will be broken for the resolved conifold and unbroken for the deformed conifold. The two resolved conifolds are interchanged by the broken $`\stackrel{~}{}_2`$. On the other hand $`_2^{}`$ is a subgroup of $`U(1)^M`$ and will be unbroken whenever translation along the M-theory circle is a symmetry. This will fail to be true only for the deformed conifold with O6-plane on $`S_>^3`$ and $`N=0`$ or $`1`$ D6-branes. Near the O6-plane we expect the geometry to look like Atiyah-Hitchin space or its double cover, wrapped on $`S_>^3`$. The M-theory circle is no longer an isometry direction, but for $`N=1`$, corresponding to the double cover, $`_2^{}`$ is unbroken whilst for $`N=0`$, $`_2^{}`$ is broken. For $`N=0`$, this leads to two distinct geometries interchanged by the broken symmetry generator. We shall return to this point and its interpretation in Type IIA later.
## 4 M-theory geometry
In this section, we will flesh out our discussion of the M-theory geometries with some details of the explicit $`G_2`$ holonomy metrics. We shall describe in turn the three distinct classes of $`G_2`$ holonomy metrics relevant to our discussion. These are labeled $`๐น_7,๐ธ_7`$ and $`๐ป_7`$ and their various orbifolds correspond respectively to orientifolds of the deformed conifold with D6-branes, an O6-plane on $`S_>^3`$ with $`N=0`$ or 1 D6-branes and orientifolds of the resolved conifold with flux.
We present the details of these metrics in order to confirm the details of the symmetry breaking discussion of the previous section and also because the existence of the $`๐ธ_7`$ metrics were not previously known in the literature. Furthermore, we show that there exists a normalizable harmonic two-form on $`๐ธ_7`$ which leads to a massless $`U(1)`$ gauge field in the IR. We can use this to rule out a smooth transition between these backgrounds and others with mass gap in the IR.
The reader who is not interested in the details of the metrics may wish to skip to the discussion of the topology of the various spaces in Section 5.1.
### 4.1 Preliminaries
We recall the Euler angle representation of $`SU(2)`$ matrices:
$$X=\left(\begin{array}{cc}\hfill \mathrm{cos}\frac{\theta }{2}e^{\frac{i}{2}(\psi +\varphi )}& \hfill \mathrm{sin}\frac{\theta }{2}e^{\frac{i}{2}(\psi \varphi )}\\ \hfill \mathrm{sin}\frac{\theta }{2}e^{\frac{i}{2}(\psi \varphi )}& \hfill \mathrm{cos}\frac{\theta }{2}e^{\frac{i}{2}(\psi +\varphi )}\end{array}\right),$$
(4.1)
where the coordinates take values
$$0\psi <4\pi ,0\varphi <2\pi ,0\theta <\pi .$$
(4.2)
The associated Maurer-Cartan one-forms $`\sigma _a`$ are defined by $`X^1dX=\frac{i}{2}\tau _a\sigma _a`$ where $`\tau _a`$ are Pauliโs matrices, and they satisfy $`d\sigma _a=\frac{1}{2}ฯต_{abc}\sigma _b\sigma _c`$. They read
$$\sigma _1=\mathrm{cos}\psi \mathrm{sin}\theta d\varphi +\mathrm{sin}\psi d\theta ,\sigma _2=\mathrm{sin}\psi \mathrm{sin}\theta d\varphi \mathrm{cos}\psi d\theta ,\sigma _3=d\psi +\mathrm{cos}\theta d\varphi .$$
(4.3)
One can write down the metric on $`^4`$ using polar coordinates $`r,\theta ,\varphi ,\psi `$ simply by relating the Cartesian and polar coordinates as follows:
$$W=\left(\begin{array}{cc}\hfill w_1+iw_2& \hfill w_3+iw_4\\ \hfill w_3+iw_4& \hfill w_1iw_2\end{array}\right)=rX(\theta ,\varphi ,\psi ),ds^2=\frac{1}{2}\mathrm{Tr}(dWdW^{})=dr^2+\frac{r^2}{4}\sigma _a\sigma _a.$$
(4.4)
### 4.2 Deformed conifold with D6-branes and orientifold
Let us consider the M-theory lift of orientifolds of deformed conifold with D6-branes. The metrics should at least preserve the $`\mathrm{๐๐}(2)\times \mathrm{๐๐}(2)\stackrel{~}{}_2`$ isometry of the deformed conifold, so we write the metrics in terms of $`SU(2)`$ matrices $`X,\stackrel{~}{X}`$ and the associated Maurer-Cartan forms $`\sigma _a,\stackrel{~}{\sigma }_a`$. The following ansatz for a $`\stackrel{~}{}_2`$-symmetric metric was considered in :
$$ds^2=dr^2+\underset{i=1}{\overset{3}{}}A_i^2(\sigma _i\stackrel{~}{\sigma }_i)^2+\underset{i=1}{\overset{3}{}}B_i^2(\sigma _i+\stackrel{~}{\sigma }_i)^2.$$
(4.5)
This metric is of $`G_2`$ holonomy provided the metric components obey
$`{\displaystyle \frac{dA_1}{dr}}`$ $`=`$ $`{\displaystyle \frac{1}{4}}\left[{\displaystyle \frac{A_1^2A_3^2B_2^2}{A_3B_2}}+{\displaystyle \frac{A_1^2A_2^2B_3^2}{A_2B_3}}\right],`$
$`{\displaystyle \frac{dA_2}{dr}}`$ $`=`$ $`{\displaystyle \frac{1}{4}}\left[{\displaystyle \frac{A_2^2A_1^2B_3^2}{A_1B_3}}+{\displaystyle \frac{A_2^2A_3^2B_1^2}{A_3B_1}}\right],`$
$`{\displaystyle \frac{dA_3}{dr}}`$ $`=`$ $`{\displaystyle \frac{1}{4}}\left[{\displaystyle \frac{A_3^2A_2^2B_1^2}{A_2B_1}}+{\displaystyle \frac{A_3^2A_1^2B_2^2}{A_1B_2}}\right],`$
$`{\displaystyle \frac{dB_1}{dr}}`$ $`=`$ $`{\displaystyle \frac{1}{4}}\left[{\displaystyle \frac{A_2^2+A_3^2B_1^2}{A_3A_2}}+{\displaystyle \frac{B_1^2B_2^2B_3^2}{B_2B_3}}\right],`$
$`{\displaystyle \frac{dB_2}{dr}}`$ $`=`$ $`{\displaystyle \frac{1}{4}}\left[{\displaystyle \frac{A_3^2+A_1^2B_2^2}{A_1A_3}}+{\displaystyle \frac{B_2^2B_3^2B_1^2}{B_3B_1}}\right],`$
$`{\displaystyle \frac{dB_3}{dr}}`$ $`=`$ $`{\displaystyle \frac{1}{4}}\left[{\displaystyle \frac{A_1^2+A_2^2B_3^2}{A_2A_1}}+{\displaystyle \frac{B_3^2B_1^2B_2^2}{B_1B_2}}\right].`$ (4.6)
and the $`G_2`$ three-form is given by
$`\mathrm{\Phi }`$ $`=`$ $`e_1e_2e_3+e_0e_ie_{\stackrel{~}{i}}+{\displaystyle \frac{1}{2}}ฯต_{ijk}e_ie_{\stackrel{~}{j}}e_{\stackrel{~}{k}}`$ (4.7)
$`=`$ $`e_1e_2e_3+e_0e_1e_{\stackrel{~}{1}}+e_0e_2e_{\stackrel{~}{2}}+e_0e_3e_{\stackrel{~}{3}}+e_1e_{\stackrel{~}{2}}e_{\stackrel{~}{3}}+e_2e_{\stackrel{~}{3}}e_{\stackrel{~}{1}}+e_3e_{\stackrel{~}{1}}e_{\stackrel{~}{2}}.`$
The M-theory geometries $`๐น_7`$ and $`๐ธ_7`$ both take the above form. The metric for $`๐น_7`$ has an additional $`U(1)^M`$ symmetry corresponding to translation along M-theory circle, while $`๐ธ_7`$ has no such $`U(1)`$ symmetry.
Numerical analysis of the differential equations proceeds in the following way. We first perform a power series analysis at the origin $`r=0`$ to find correct initial values for $`A_i,B_i`$ that make the solution smooth there. We then let them evolve according to (4.6) and see if the metric asymptotes to that of a $`(\text{conifold})\times S^1`$ with flux,
$$(A_1,A_2,A_3,B_1,B_2,B_3)\frac{r}{6}(\sqrt{3},\sqrt{3},2,\sqrt{3},\sqrt{3},0),$$
(4.8)
or that of a $`G_2`$ cone over $`S^3\times S^3`$ (which describes the IIA theory at infinite string coupling),
$$(A_1,A_2,A_3,B_1,B_2,B_3)\frac{r}{6}(\sqrt{3},\sqrt{3},\sqrt{3},1,1,1).$$
(4.9)
We will find that initial conditions which are regular at the origin do not always lead to a sensible asymptotic behavior.
#### 4.2.1 Numerical analysis for $`๐น_7`$
The correct initial data for smooth metric with an $`S^3`$ of unit size is given by the following power series
$$A_i=1+\frac{1}{16}r^2+๐ช(r^3),B_i=\frac{r}{4}+\frac{b_i}{192}r^3+๐ช(r^4),b_1+b_2+b_3=3.$$
(4.10)
This leading-order behavior uniquely determines the solution as a power series at $`r=0`$.
The family of smooth initial data is parameterized by $`b_1,b_2`$. However, the metric asymptotes to (4.8) only when the initial data is on a half-line $`C_1:b_1=b_21`$ in the parameter space. (Owing to the permutation symmetry of $`1,2,3`$, there are three related families of initial conditions leading to metrics with sensible asymptotics. They are three half-lines $`C_1,C_2,C_3`$ meeting at $`P:b_1=b_2=b_3=1`$ as depicted in Figure 3.)
Let us focus on the family of solutions $`C_1`$. Since we have fixed the size of the minimal $`S^3`$, the value of $`b_1=b_2`$ determines the radius of the M-theory circle at infinity, which is roughly $`\underset{r\mathrm{}}{lim}B_3`$. The radius blows up as $`b_1=b_2`$ approach $`1`$, and the solution has $`SU(2)`$ enhanced isometry and asymptotes to (4.9) at infinity. The solution is nothing but the familiar asymptotically conical $`G_2`$ holonomy metric on the spin bundle over $`S^3`$ . Alternatively, if we consider the family $`C_1`$ of rescaled solutions having a fixed radius of M-theory circle, then the limit of approaching $`P`$ is the limit of vanishing $`S^3`$ at the center.
It is easy to see that, for the solutions on $`C_1`$, the equalities $`A_1=A_2,B_1=B_2`$ hold all the way along the radial evolution. So the manifold $`๐น_7`$ has an additional $`U(1)^M`$ isometry corresponding to translation along the M-theory circle. At the point $`P`$ the solution satisfies $`A_1=A_2=A_3,B_1=B_2=B_3`$ along the radial evolution, so the $`U(1)^M`$ is enhanced further to an $`SU(2)`$.
The solutions at $`r=0`$ take the form
$`ds^2`$ $``$ $`dr^2+{\displaystyle \underset{i=1}{\overset{3}{}}}(\sigma _i\stackrel{~}{\sigma }_i)^2+{\displaystyle \frac{r^2}{16}}{\displaystyle \underset{i=1}{\overset{3}{}}}(\sigma _i+\stackrel{~}{\sigma }_i)^2`$ (4.11)
$`=`$ $`dr^2+{\displaystyle \underset{i=1}{\overset{3}{}}}\stackrel{~}{\mathrm{\Sigma }}_i^2+{\displaystyle \frac{r^2}{16}}{\displaystyle \underset{i=1}{\overset{3}{}}}(2\mathrm{\Sigma }_i\stackrel{~}{\mathrm{\Sigma }}_i)^2,`$
where we introduced $`Y=\stackrel{~}{X}X^1`$ and $`\mathrm{\Sigma }=XdX^1,\stackrel{~}{\mathrm{\Sigma }}=Y^1dY`$. Recalling the metric (4.4) on $`^4`$ one finds that the geometry is a smooth $`^4`$ bundle over $`S^3`$, where $`Y`$ gives the base $`S^3`$ and $`(r,X)`$ parameterize the fiber $`^4`$.
#### 4.2.2 Orbifolding
The M-theory lifts of $`N`$ D6-branes or (O6 + $`N`$D6) are orbifolds of $`๐น_7`$ by $`A_{N1}`$ or $`D_N`$ groups $`\mathrm{\Gamma }`$. The orbifold group acts diagonally on $`X`$ and $`\stackrel{~}{X}`$ from the right.
$$g\mathrm{\Gamma }:(X,\stackrel{~}{X})(Xg,\stackrel{~}{X}g).$$
(4.12)
As we have seen, near the origin the coordinate $`Y=\stackrel{~}{X}X^1`$ describes the base $`S^3`$ of finite volume, and $`X^1`$ describes the shrinking $`S^3`$. $`\mathrm{\Gamma }`$ acts on them as
$$g\mathrm{\Gamma }:(Y,X^1)(Y,g^1X^1).$$
(4.13)
Thus orbifolding gives a geometry which is $`^2/\mathrm{\Gamma }`$ fibered over $`S^3`$.
The M-theory circle corresponds to the shift $`\psi \psi +\alpha ,\stackrel{~}{\psi }\stackrel{~}{\psi }+\alpha `$. Its radius can be read off from the metric at infinity
$`ds^2`$ $`=`$ $`\mathrm{}+A_3^2(\sigma _3\stackrel{~}{\sigma }_3)^2+B_3^2(\sigma _3+\stackrel{~}{\sigma }_3)^2,`$ (4.14)
$`=`$ $`\mathrm{}+A_3^2(d\stackrel{ห}{\psi }+\mathrm{cos}\theta d\varphi \mathrm{cos}\stackrel{~}{\theta }d\stackrel{~}{\varphi })^2+B_3^2(d\widehat{\psi }+\mathrm{cos}\theta d\varphi +\mathrm{cos}\stackrel{~}{\theta }d\stackrel{~}{\varphi })^2.`$
Here $`\widehat{\psi }=\psi +\stackrel{~}{\psi }`$ is the coordinate on M-theory circle, and $`\stackrel{ห}{\psi }=\psi \stackrel{~}{\psi }`$ is one of the coordinates of $`T^{1,1}`$. From the periodicity of $`\psi ,\stackrel{~}{\psi }`$ it follows that
$$0\stackrel{ห}{\psi }<4\pi ,0\widehat{\psi }<8\pi .$$
(4.15)
The RR charge is an integral of the field strength of the one-form potential $`A`$,
$$A=\frac{1}{4}(\mathrm{cos}\theta d\varphi +\mathrm{cos}\stackrel{~}{\theta }d\stackrel{~}{\varphi })$$
(4.16)
over the $`S^2`$ defined by $`\theta =\stackrel{~}{\theta },\varphi =\stackrel{~}{\varphi }`$ (One can check that its volume becomes zero at $`r=0`$). We normalized $`A`$ so that it appears in the M-theory metric as $`(d\theta +A)^2`$, where $`\theta `$ is of period $`2\pi `$. The D6-brane charge is defined by
$$Q_{\mathrm{D6}}=_{S^2}\frac{dA}{2\pi },$$
(4.17)
which is unity for $`๐น_7`$ without orbifolding.
The $`A_{N1}`$ orbifold group is generated by $`a=\mathrm{exp}(\frac{2\pi i\tau _3}{N})`$. This simply shortens the period of M-theory circle to $`1/N`$, and therefore increases the D6-brane charge. For $`D_N`$ orbifolds, $`\mathrm{\Gamma }`$ is generated by $`a=\mathrm{exp}(\frac{i\pi \tau _3}{N2})`$ and $`b=i\tau _2`$. The element $`b`$ acts as an antipodal map on $`S^2`$, and the D6-brane charge as counted in the covering space becomes
$$Q_{\mathrm{D6}}=2\times _^2\frac{dA}{2\pi }=2N4.$$
(4.18)
The action of $`b`$ on the coordinates of the conifold was obtained in the previous section and is the complex conjugation $`c_0`$:
$$z_i\overline{z}_i.$$
(4.19)
The M-theory lift of the free orientifold of the deformed conifold with $`2N4`$ D6-branes wrapping the $`S_<^3`$ uplifts to a different orbifold $`๐น_7/D_N^{}`$. The generators $`a,b`$ of $`D_N^{}`$ act on $`(X,\stackrel{~}{X})`$ as
$$\begin{array}{ccc}\hfill a& :& (X,\stackrel{~}{X})(X\mathrm{exp}(\frac{i\pi \tau _3}{N2}),\stackrel{~}{X}\mathrm{exp}(\frac{i\pi \tau _3}{N2})),\hfill \\ \hfill b& :& (X,\stackrel{~}{X})(Xi\tau _2,\stackrel{~}{X}i\tau _2).\hfill \end{array}$$
(4.20)
The generator $`b`$ maps $`z_i\overline{z}_i`$, and in particular acts on the minimal $`S_<^3`$ as the antipodal map.
#### 4.2.3 Numerical Analysis for $`๐ธ_7`$
Next we would like to describe the M-theory metrics corresponding to O6<sup>-</sup> and $`N=0`$ or $`1`$ D6-branes on $`S_>^3`$ of the deformed conifold. In this case the net flux from the O6<sup>-</sup> and D6 is negative and the M-theory geometry is expected to look locally like Atiyah-Hitchin space ($`N=0`$) or its double cover ($`N=1`$), wrapped on $`S_>^3`$.
We are looking for a $`G_2`$ metric which is asymptotic to (4.8) or (4.9). Furthermore, we expect the orbifold at infinity to be $`D_4^{}`$ for $`N=0`$ and $`D_3^{}`$ for $`N=1`$. Recall that $`D_3^{}`$ is generated by $`b`$ which acts on $`(X,\stackrel{~}{X})`$ as in (4.20).
By analogy with the construction of Atiyah-Hitchin space, we expect that orbifolding by $`D_3^{}`$ is necessary in order to remove a conical singularity at the origin. This will be the case if $`A_2=r+\mathrm{}`$ near $`r=0`$ so that the metric becomes
$`ds^2`$ $`=`$ $`A_1^2(\sigma _1\stackrel{~}{\sigma }_1)^2+A_3^2(\sigma _3\stackrel{~}{\sigma }_3)^2+B_1^2(\sigma _1+\stackrel{~}{\sigma }_1)^2+B_2^2(\sigma _2+\stackrel{~}{\sigma }_2)^2+B_3^2(\sigma _3+\stackrel{~}{\sigma }_3)^2`$ (4.21)
$`+dr^2+r^2(\sigma _2\stackrel{~}{\sigma }_2)^2+\mathrm{}.`$
Here $`A_1,A_3,B_1,B_2,B_3`$ should be regarded as constants near $`r=0`$. The first line gives the metric of a five-dimensional bolt which is topologically $`S^2\times S^3`$. The second line yields a conical singularity which is removed by orbifolding by $`D_3^{}`$ as in (4.20). Note that in order for the shrinking $`S^1`$ to become an isometry direction at the origin (which is needed to avoid a singularity) we require $`A_1=B_3`$ and $`B_1=A_3`$ at $`r=0`$.
The initial data for a smooth metric at $`r=0`$ with five-dimensional bolt is
$$A_1=B_3,B_1=A_3,A_2=0.$$
(4.22)
Interestingly, if we set $`(A_3,B_1,B_2)`$ all equal and much greater than $`A_1(=B_3)`$, then for small $`r`$, $`(A_3,B_1,B_2)`$ stay almost constant while $`(A_1,A_2,B_3)`$ approximately obey the equations for the Atiyah-Hitchin manifold. However, the numerical analysis shows that such solutions do not extend toward large $`r`$ and we should not try to impose these extra conditions.
The power series analysis shows that any initial values for $`(A_1,B_1,B_2)`$ uniquely determine a solution, but this will have sensible asymptotics at large $`r`$ only for some particular fine-tuned initial data. After fixing an overall scale, one gets a two-dimensional parameter space of initial conditions. We choose this to be three faces of a cube,
$$\{B_1=1,B_2,A_1[0,1]\},\{B_2=1,B_1,A_1[0,1]\},\{A_1=1,B_1,B_3[0,1]\}$$
as depicted in Figure 4. One can numerically see that the solutions for generic initial conditions do not behave nicely at infinity: either $`B_1,B_2`$ or $`B_3`$ blows up much faster than the others. The generic initial conditions are grouped into regions (I), (II), (III) shown in Figure 4 according to the asymptotics. The solution behaves nicely at large $`r`$ if the initial condition is chosen on the curves $`C_{1,2,3}`$ separating three regions.
The curve $`C_1`$ is a straight line segment corresponding to the initial data
$$A_1=B_3=B_1=A_3=1,0<B_20.917,A_2=0.$$
Generic solutions with this initial condition asymptote locally to the conifold times an $`S^1`$ with flux,
$$(A_1,A_2,A_3,B_1,B_2,B_3)(\sqrt{3},2,\sqrt{3},\sqrt{3},0,\sqrt{3}).$$
(4.23)
These solutions were discovered in and named $`_7`$. They do not have the expected asymptotics (4.8) since the M-theory circle is in the $`B_2`$ direction rather than $`B_3`$. With these asymptotics, the orbifolding by $`D_3^{}`$ acts as an orbifold on the base rather than as an orientifold. In addition, the solution has a $`U(1)`$ isometry $`A_1=A_3,B_1=B_3`$ all along the flow, corresponding to translation along the M-theory circle. The solutions which we seek with O6<sup>-</sup>-plane and $`N=0`$ or 1 D6 are not expected to have such a $`U(1)^M`$ isometry. As explained in the solution $`_7`$ should be interpreted as the uplift of IIA on a manifold with local $`^1\times ^1`$ and unit flux through each $`^1`$. This solution will have no part to play in our current analysis.
The other two curves $`C_2,C_3`$ are related by the symmetry of the differential equation (4.6).
$$(A_1,A_2,A_3,B_1,B_2,B_3)(A_3,A_2,A_1,B_3,B_2,B_1).$$
(4.24)
The solutions on $`C_2`$ flow to the correct asymptotic boundary conditions (4.8) and these turn out the solutions which we label $`๐ธ_7`$ after an appropriate orbifolding. The solutions on $`C_3`$ describe the same set of solutions in different coordinates. We do not need to take account of them separately since the asymptotic boundary conditions (4.8) partially fix our choice of coordinates and the solutions on $`C_3`$ do not respect these asymptotics.
As we have discussed, it is necessary to orbifold the solutions by $`D_3^{}`$ in order to remove a conical singularity at the origin. The resulting family of solutions $`๐ธ_7`$ is asymptotically an orientifold of the conifold with RR charge $`(1)`$ and one can identify them with the M-theory lift of the O6<sup>-</sup>+D6 system.
Further orbifolding by the right action of $`(i\tau _3)(i\tau _3)`$ on $`(X,\stackrel{~}{X})`$ halves the radius of the M-theory circle and thereby doubles the RR charge. Note that this does not lead to any orbifold singularities in the M-theory geometry $`๐ธ_7/_2`$. The enlarged orbifold group acts as $`D_4^{}`$ on the boundary and this is the solution corresponding to an O6<sup>-</sup> and no D6 branes.
#### 4.2.4 Normalizable harmonic two-form?
If there is such a two-form $`\omega `$, it will lead to the presence of $`U(1)`$ gauge dynamics in the IR. We expect that this should be the case for the M-theory lift of O6<sup>-</sup>+D6 since this will have $`SO(2)=U(1)`$ gauge group.
The two-form can be written locally as a derivative of a one-form. We assume that it takes the form
$$\omega =\underset{i}{}d\left\{f_i(r)(\sigma _i\stackrel{~}{\sigma }_i)+g_i(r)(\sigma _i+\stackrel{~}{\sigma }_i)\right\}.$$
Instead of trying to solve $`d\omega =d\omega =0`$, we try to solve the simpler equation
$$\omega =\alpha \mathrm{\Phi }\omega ,$$
(4.25)
where $`\alpha `$ is a real constant, requiring $`\omega `$ to fall into an irreducible representation of $`G_2`$. A little calculation shows that $`\alpha =1`$ and $`f_i,g_i`$ have to satisfy
$$\frac{1}{f_1}\frac{df_1}{dr}+\frac{1}{2}\left(\frac{A_1}{A_2B_3}+\frac{A_1}{B_2A_3}\right)=0,\frac{1}{g_1}\frac{dg_1}{dr}+\frac{1}{2}\left(\frac{B_1}{A_2A_3}\frac{B_1}{B_2B_3}\right)=0,\mathrm{etc}.$$
(4.26)
The equation shows the existence of six independent harmonic two-forms corresponding to $`f_i,g_i`$. The normalizability of each mode can be analyzed using the form of $`(A_i,B_i)`$ at the origin and infinity.
For $`๐ธ_7`$, the mode proportional to $`f_2`$ behaves at $`r=0`$ and $`r=\mathrm{}`$ as
$$f_2\stackrel{r0}{}1\frac{1+A_1^2}{4}r^2+๐ช(r^3),f_2\stackrel{r\mathrm{}}{}\mathrm{exp}\left(\frac{r}{2B_3}\right),$$
(4.27)
and turns out to be normalizable. Note that this mode is projected out upon orbifolding further to get $`๐ธ_7/_2`$. For $`๐น_7`$, none of these two-forms is normalizable. The absence of a normalizable harmonic two-form for $`๐น_7`$ is puzzling, as discussed in , since one would expect the gauge group on $`N`$ coincident D6-branes should be $`U(N)`$ and its $`U(1)`$ part should remain in the infrared limit.
### 4.3 Resolved conifold with flux and orientifold
For the M-theory lift of the resolved conifold with flux and orientifold, one expects the symmetry $`SU(2)\times SU(2)\times U(1)^M`$ but no $`\stackrel{~}{}_2`$ exchanging the two sets of Maurer-Cartan forms. The following ansatz for the metric was considered in
$$ds^2=dr^2+a^2\{(\sigma _1+g\stackrel{~}{\sigma }_1)^2+(\sigma _2+g\stackrel{~}{\sigma }_2)^2\}+b^2(\stackrel{~}{\sigma }_1^2+\stackrel{~}{\sigma }_2^2)+c^2(\sigma _3+g_3\stackrel{~}{\sigma }_3)^2+f^2\stackrel{~}{\sigma }_3^2$$
(4.28)
The metric is of $`G_2`$ holonomy provided $`g=\frac{af}{2bc},g_3=1+2g^2`$ and
$$\begin{array}{ccc}\hfill \dot{a}& =& \frac{c}{2a}+\frac{a^5f^2}{8b^4c^3},\hfill \\ \hfill \dot{c}& =& 1+\frac{c^2}{2a^2}+\frac{c^2}{2b^2}\frac{3a^2f^2}{8b^4},\hfill \end{array}\begin{array}{ccc}\hfill \dot{b}& =& \frac{c}{2b}\frac{a^2(a^23c^2)f^2}{8b^3c^3},\hfill \\ \hfill \dot{f}& =& \frac{a^4f^3}{4b^4c^3}.\hfill \end{array}$$
(4.29)
As initial conditions we put $`a=c=0,b=1`$ and $`f=f_0`$. The regularity at $`r=0`$ requires
$$a=\frac{r}{2}+๐ช(r^2),c=\frac{r}{2}+๐ช(r^2).$$
(4.30)
If $`f_0<1`$, numerical solutions asymptote to $`\text{(conifold)}\times S^1`$,
$$(a,b,c)r(\sqrt{1/6},\sqrt{1/6},\sqrt{1/9}),f\mathrm{const},g0,g_31.$$
(4.31)
These solutions were named $`๐ป_7`$ in . At $`f_0=1`$ the $`S^1`$ decompactifies and the solution coincides with the familiar asymptotically conical $`G_2`$ metric on the spin bundle over $`S^3`$.
Interestingly, the M-theory geometry has $`SU(2)\times SU(2)\times U(1)^M`$ symmetry and no extra $`U(1)^{\mathrm{phase}}`$. It seems that the $`U(1)^{\mathrm{phase}}`$ isometry corresponding to the phase rotation $`z_ie^{i\alpha }z_i`$ of the resolved conifold is broken in the presence of flux.
One can increase the flux or introduce an orientifold action simply by orbifolding the M-theory geometry. The orbifold group can be either $`A_{N1}`$ or $`D_N`$ groups acting on $`(X,\stackrel{~}{X})`$ in the same way from the right, or it can be the group $`D_N^{}`$ defined in (4.20). All these groups act freely. The RR charge is $`N`$ for $`A_{N1}`$ orbifolds, $`N2`$ for $`D_N`$ orbifolds and $`2N`$ for $`D_N^{}`$ orbifolds.
## 5 Quantum moduli space
We are now in a position to study the quantum moduli space of supersymmetric orientifold vacua for different values of RR charge $`Q=N2`$. We first analyze them through the behavior of membrane instanton factors $`\eta _i`$ on various classical M-theory geometries, as in . We then study them using Vafaโs exact superpotential.
From Table 1, it appears that there are six classes of cases. $`N2>1`$ and $`N<0`$ are โregularโ and can be understood completely using M-theory arguments or the Vafa superpotential. $`N=3`$ is somewhat special, but the same analysis applies. In all these cases, the moduli space turns out to be a copy of $`^1`$ with some number of marked points. When $`N=0,1`$ or $`2`$, the moduli space of vacua consists of different branches and needs to be discussed separately.
### 5.1 M-theory lift
Here we try to study the quantum moduli space within the M-theory framework, through the behavior of various membrane instanton factors as holomorphic functions on moduli space. As preparation we need some basic results about the topology of the M-theory lifts.
We would like to find good chiral parameters to describe the classical and quantum moduli space. In an M-theory compactification on a $`G_2`$ manifold, $`X`$, the chiral parameters are of the form:
$$u=\mathrm{exp}\left(_Q(k\mathrm{\Phi }_3+iC)\right),$$
(5.1)
where the integral is taken over a non-trivial three-cycle $`QH_3(X,)`$. Here, $`\mathrm{\Phi }_3`$ is the $`G_2`$ three-form and $`C`$ is the M-theory three-form potential. $`k`$ is a constant related to the membrane tension and $`u`$ agrees with the exponential of the action of a membrane instanton wrapping an associative three-cycle homologous to $`Q`$.
Following we introduce a set of chiral parameters corresponding to a basis of integral three-cycles in the boundary of the manifold $`X`$. The boundary of $`X`$ is $`Y_\mathrm{\Gamma }(S^3\times S^3)/\mathrm{\Gamma }`$ where $`\mathrm{\Gamma }`$ is the relevant orbifold group which is $`D_N`$ (when $`N>2`$) or $`D_N^{}^{}`$ (with $`N^{}:=4N`$ when $`N<2`$). The space of integral three-cycles $`H_3(Y,)`$ is thus two-dimensional. On the various semi-classical branches, one three-cycle is โfilled inโ so that $`H_3(X,)`$ is one-dimensional. Classically, the chiral parameter corresponding to the filled in cycle takes the value $`\mathrm{exp}(0)=1`$. As explained in , this is subject to quantum corrections, but the classical analysis (supplemented with knowledge of the gauge theory) is reliable near limits of moduli space where the M-theory geometry is everywhere weakly curved (except for possible orbifold singularities). This gives information about the poles and zeros of the holomorphic parameters which is sufficient to reconstruct the exact moduli space.
Our first task, then, is to understand the third homology group of the boundary $`Y_\mathrm{\Gamma }`$ for $`\mathrm{\Gamma }=D_N`$ or $`D_N^{}^{}`$ ($`N^{}=4N`$). To describe $`H_3(Y_\mathrm{\Gamma },)`$, we follow . Before modding out by $`\mathrm{\Gamma }`$, the boundary is $`Y=S^3\times S^3`$. This space can usefully be described in terms of the three $`SU(2)`$ elements $`g_1,g_2,g_3`$ subject to the equivalence relation:
$$(g_1,g_2,g_3)=(g_1h,g_2h,g_3h).$$
(5.2)
Let $`\widehat{E}_iSU(2)^3`$ be the $`i^{th}`$ copy of $`SU(2)`$ โ so $`\widehat{E}_1`$ is the set $`(g,1,1),gSU(2)`$. In $`Y`$, the $`\widehat{E}_i`$ project to cycles $`E_i`$ obeying
$$E_1+E_2+E_3=0.$$
(5.3)
Under the orbifold projection $`YY_\mathrm{\Gamma }`$ for $`\mathrm{\Gamma }=D_N`$, the $`E_i`$ are mapped to cycles in $`Y_\mathrm{\Gamma }`$ which we label $`E_i^{}`$. The map of $`E_1`$ to $`E_1^{}`$ is the $`(4N8)`$-fold cover, $`S^3S^3/D_N`$, whilst $`E_2`$ and $`E_3`$ are mapped diffeomorphically to $`E_2^{}`$ and $`E_3^{}`$. Thus we have
$$(4N8)E_1^{}+E_2^{}+E_3^{}=0.$$
(5.4)
$`Y_\mathrm{\Gamma }`$ is simply the product $`E_1^{}\times E_2^{}`$ and these cycles generate $`H_3(Y_\mathrm{\Gamma },)`$.
For $`\mathrm{\Gamma }=D_N^{}^{}`$, $`E_2`$ and $`E_3`$ are again mapped diffeomorphically to $`E_2^{}`$ and $`E_3^{}`$ in $`Y_\mathrm{\Gamma }`$. To find another cycle, instead of $`E_1`$ we consider a little different three sphere in $`Y`$;
$$(g,g^1\left(\begin{array}{cc}i& 0\\ 0& i\end{array}\right)g,1),gSU(2).$$
(5.5)
Left multiplication by $`D_N^{}`$ elements on $`g`$ corresponds to the $`D_N^{}^{}`$ action on $`Y`$. Thus, this $`S^3`$ defines a $`(4N^{}8)`$-fold cover over a cycle $`ES^3/D_N^{}`$ in $`Y_\mathrm{\Gamma }`$. Note that the cover $`S^3`$ has the same homology class as $`E_1`$. This is because the inverse $`gg^1`$ reverses the orientation of $`SU(2)`$ and therefore the $`g_2`$ part in (5.5) does not contribute in homology. Using (5.3), we find the homology relation in $`Y_\mathrm{\Gamma }`$,
$$(4N^{}8)E+E_2^{}+E_3^{}=0.$$
Denoting the orientation reversal of $`E`$ by $`E_1^{}`$, we find the same relation as (5.4) with $`N=4N^{}`$. $`Y_\mathrm{\Gamma }`$ is just the product $`E_1^{}\times E_2^{}`$ and these cycles generate $`H_3(Y_\mathrm{\Gamma },)`$.
Next we would like to describe the behavior of the holomorphic parameters
$$\eta _i=\mathrm{exp}\left(_{E_i^{}}(k\mathrm{\Phi }_3+iC)\right)$$
(5.6)
at the various semi-classical limits of moduli space. The homology relation (5.4) implies the following relation of $`\eta _i`$โs
$$\eta _1^{4N8}\eta _2\eta _3=1.$$
(5.7)
As remarked in , one must be careful about the definition of $`\mathrm{exp}(i_{E_i^{}}C)`$ โ this must include the sign factor of the fermionic determinant, and there is a potential sign error in the right hand side of (5.7). It is shown in that the error is absent for $`\mathrm{\Gamma }=D_N`$, and the same proof applies also to $`\mathrm{\Gamma }=D_N^{}^{}`$. The reason is that, just as in $`D_N`$ case, $`Y_\mathrm{\Gamma }`$ is the product of spin manifolds $`Y_1=S^3/D_N^{}`$ and $`Y_2=S^3`$, and that $`E_i^{}`$ are all transverse to a $`Spin(3)`$ sub-bundle of $`TY_\mathrm{\Gamma }TY_1TY_2`$ โ a slight rotation of $`TY_1`$ in the direction of $`TY_2`$.
Thus, the relation (5.7) exactly holds at each of the semi-classical points, and therefore by holomorphy, everywhere on the moduli space. We shall determine the orders of the poles and zeroes of the $`\eta _i`$ as functions on moduli space, following the argument in . Part of this information is obtained classically by counting the number of times which a boundary cycle $`E_i^{}`$ wraps the minimal three-cycle of the geometry. However, if the geometry has an orbifold singularity, the classical analysis is modified by strong coupling effects at low energies.
Near the deformed classical points corresponding to orbifolds of $`๐น_7`$ with large $`S^3`$, $`E_1^{}`$ is filled in whilst $`E_2^{}`$ and $`E_3^{}`$ grow large with opposite orientation. Thus
$$\begin{array}{ccccc}(N>2)\hfill & \eta _11,\hfill & \eta _20^h,\hfill & \eta _3\mathrm{}^h,\hfill & h\stackrel{ห}{h}_{SO(2N)}=2N2,\hfill \\ & \eta _11,\hfill & \eta _20^{2h^{}},\hfill & \eta _3\mathrm{}^{2h^{}},\hfill & h\stackrel{ห}{h}_{Sp(N4)}=N3,\hfill \\ (N<2)\hfill & \eta _11,\hfill & \eta _2\mathrm{}^{\stackrel{~}{h}},\hfill & \eta _30^{\stackrel{~}{h}},\hfill & \stackrel{~}{h}\stackrel{ห}{h}_{SO(42N)}=22N,\hfill \\ & \eta _11,\hfill & \eta _2\mathrm{}^{2\stackrel{~}{h}^{}},\hfill & \eta _30^{2\stackrel{~}{h}^{}},\hfill & \stackrel{~}{h}^{}\stackrel{ห}{h}_{Sp(2N)}=3N.\hfill \end{array}$$
(5.8)
Note that, as explained in , the $`\pm `$ sign for $`\eta _1`$ corresponds to a choice of discrete flux which leads to $`SO`$ or $`Sp`$ gauge theory, and the order of zeroes or poles is related to the degeneracy of vacua of the corresponding super Yang-Mills theory<sup>3</sup><sup>3</sup>3 The extra factor of 2 for the symplectic group was explained in for the $`N>2`$ cases with O6-plane at $`S^3`$, using an argument in which D6-branes are deformed away from the orientifold plane. In fact, we need this extra factor also for the $`N<2`$ cases with free orientifold with minimal $`^3`$. Derivation that applies to both cases is given in Section 5.2.. Also, we have chosen the orientations of cycles $`E_i^{}`$ so that the family of vacua preserves the same supersymmetry for all $`N`$, in accordance with Table 1 in the introduction. One can see that the orders of zeroes or poles depend linearly on $`N`$ under this choice.
On the two resolved vacua described by orbifolds of $`๐ป_7`$ we have either $`E_2^{}`$ or $`E_3^{}`$ filled in. In the former case the minimal cycle is homologous to $`E_1^{}=\frac{1}{4N8}E_3^{}`$ and
$$\eta _1\mathrm{},\eta _21,\eta _30^{4N8},$$
(5.9)
whilst in the latter case one has
$$\eta _10,\eta _2\mathrm{}^{4N8},\eta _31.$$
(5.10)
The pole turn into zero and vice versa as $`N2`$ flips sign, due to our choice of the orientations of $`E_i^{}`$.
On the classical vacua corresponding to $`๐ธ_7/_2`$ and its double cover $`๐ธ_7`$, we expect a similar behavior to the deformed classical vacua except that the large $`S^3`$ should have opposite orientation. In both cases, $`E_2^{}`$ and $`E_3^{}`$ wrap the minimal cycle once and so the relevant poles and zeroes are simple.
$$\eta _11,\eta _20,\eta _3\mathrm{}.$$
(5.11)
Now, let us glue together the local behaviors of $`\eta _i`$ to get the curve describing the moduli space.
#### 5.1.1 $`N>3`$
The quantum moduli space of the orientifold plane system has been analyzed in for the cases with $`N>3`$ D6-branes. It consists of four classical geometries connected in the quantum moduli space (see Figure 5) at the points $`\eta _1=\pm 1,0,\mathrm{}`$. The point $`\eta _1=1`$ corresponds to an O6<sup>-</sup>-plane and $`N`$ D6-branes wrapping the $`S_>^3`$ of the deformed conifold. The gauge group at this point is $`SO(2N)`$. The point $`\eta _1=1`$ corresponds to an O6<sup>+</sup>-plane and $`N4`$ D6-branes wrapping the $`S_>^3`$ of the deformed conifold that has a gauge group $`Sp(N4)`$. The points $`\eta _1=0,\mathrm{}`$ correspond to the two resolved conifolds with $`N2`$ units of RR flux.
The quantum curve can be obtained by looking at the poles and zeros of the $`\eta _i`$ functions, as explained above:
| | $`P_1`$ | $`P_1^{}`$ | $`P_2`$ | $`P_3`$ |
| --- | --- | --- | --- | --- |
| $`\eta _1`$ | $`1`$ | $`1`$ | 0 | $`\mathrm{}`$ |
| $`\eta _2`$ | $`0^{2N2}`$ | $`0^{2(N3)}`$ | $`\mathrm{}^{4N8}`$ | $`1`$ |
| $`\eta _3`$ | $`\mathrm{}^{2N2}`$ | $`\mathrm{}^{2(N3)}`$ | $`1`$ | $`0^{4N8}`$ |
From this table it is easy to deduce the quantum curve:
$$\begin{array}{cc}\hfill \eta _2& =\eta _1^{(4N8)}(\eta _11)^{(2N2)}(\eta _1+1)^{(2N6)},\hfill \\ \hfill \eta _3& =(\eta _11)^{(2N2)}(\eta _1+1)^{(2N6)}.\hfill \end{array}$$
(5.12)
Notice that for this charge there is no classical configuration involving anti-D6-branes wrapped on $`S_<^3`$. They would preserve the same supersymmetry but the charge would be negative.
#### 5.1.2 $`N=3`$
The case $`N=3`$ with orientifold planes coincides with the $`A_3`$ case (see Figure 6). The quantum moduli space has three points with classical descriptions at $`\eta _1=1,0,\mathrm{}`$. The point $`\eta _1=1`$ corresponds to an O6<sup>-</sup>-plane and three D6-branes wrapping the $`S_>^3`$ of the deformed conifold. The gauge group at this point is $`SO(6)SU(4)`$, that has $`h=4`$ vacua. The points $`\eta _1=0,\mathrm{}`$ correspond to the two resolved conifolds with $`1`$ unit of RR flux. As expected the quantum curve is the same as in the $`A_3`$ case:
$$\begin{array}{cc}\hfill \eta _2=& \eta _1^4(\eta _11)^4,\hfill \\ \hfill \eta _3=& (1\eta _1)^4.\hfill \end{array}$$
(5.13)
The poles and the zeros for the $`N=3`$ case can be summarized in the following table:
| | $`P_1`$ | $`P_2`$ | $`P_3`$ |
| --- | --- | --- | --- |
| $`\eta _1`$ | 1 | 0 | $`\mathrm{}`$ |
| $`\eta _2`$ | $`0^4`$ | $`\mathrm{}^4`$ | 1 |
| $`\eta _3`$ | $`\mathrm{}^4`$ | 1 | $`0^4`$ |
As in the previous case there are no classical configurations involving anti-D6-branes and preserving the same supersymmetries.
#### 5.1.3 $`N=2`$
This case is truly exceptional, because the boundary of the relevant M-theory geometry is not an orbifold of $`S^3\times S^3`$, but rather
$$\left(\frac{S^2\times S^1}{_2}\right)\times S^3.$$
(5.14)
We denote the first factor by $`E_1^{}`$ and the second by $`E_2^{}=E_3^{}`$, and define the associated membrane instanton factors $`\eta _i`$ as explained before.
The classical point $`P_1`$ corresponding to an O6<sup>-</sup>-plane and two D6-branes wrapped on $`S_>^3`$ supports an $`SO(4)`$ gauge theory. Although $`E_2^{}=E_3^{}`$ wraps the minimal cycle once, the corresponding instanton factors $`\eta _2,\eta _3`$ develop double zero and pole there due to the $`SO(4)`$ gauge dynamics. Here we determined the degeneracy of vacua from $`\stackrel{ห}{h}_{SO(N)}=2N2=2`$ for $`N=2`$, though $`SO(4)`$ super Yang-Mills theory actually has four degenerate vacua. At the two resolved classical points $`P_{2,3}`$ the $`\eta _1`$ has a simple zero or pole, whereas $`\eta _{2,3}`$ remain finite.
In addition, we have geometries with large $`^3S_<^3/_2`$ corresponding to free orientifold of deformed conifold. We claim that there are two distinct classical points with large $`^3`$. The two will differ, from Type IIA viewpoint, in the action on the Chan-Paton indices when one wraps some D6-branes on $`^3`$. They should also be distinguished by the discrete torsion of NSNS B-field or M-theory three-form potential. Indeed, the third homology group $`H_3(X,)`$ of the relevant M-theory geometry
$$X:^3\times S^1^3$$
is $`_2`$ and therefore has a torsion part, since the third homology group of the boundary (5.14) is $``$ and only twice the first generator is trivial in $`X`$. Therefore, one has classically the choice<sup>4</sup><sup>4</sup>4 Alternatively, we may look at the twisted second homology of the IIA reduction of $`\stackrel{~}{X}`$. We again find the torsion $`H_2(\stackrel{~}{X},\stackrel{~}{})=_2`$ as $`\stackrel{~}{X}`$ contains an $`^3`$.
$$_{E_1^{}}C=0\text{or}\pi (\text{mod}2\pi ).$$
(5.15)
We denote the corresponding two classical points by $`P_1^{}`$ and $`P_1^{\prime \prime }`$.
The table of singularities reads:
| | $`P_1`$ | $`(P_1)`$ | $`P_1^{}`$ | $`P_1^{\prime \prime }`$ | $`P_2`$ | $`P_3`$ |
| --- | --- | --- | --- | --- | --- | --- |
| $`\eta _1`$ | 1 | 1 | 1 | $`1`$ | 0 | $`\mathrm{}`$ |
| $`\eta _2`$ | $`0^2`$ | $`0^2`$ | $`\mathrm{}^2`$ | $`\mathrm{}^2`$ | $`1`$ | 1 |
| $`\eta _3`$ | $`\mathrm{}^2`$ | $`\mathrm{}^2`$ | $`0^2`$ | $`0^2`$ | 1 | $`1`$ |
Here we included in the second column the contribution of the two $`SO(4)`$ vacua that we missed by the counting based on $`\stackrel{ห}{h}_{SO(4)}=2`$.
From exact superpotential explained later, one can derive the quantum curve:
$$\begin{array}{cc}\hfill \eta _2=& (\eta _11)^2(\eta _1+1)^2,\hfill \\ \hfill \eta _3=& (\eta _11)^2(\eta _1+1)^2.\hfill \end{array}$$
(5.16)
This accounts for only four classical points $`P_1,P_1^{\prime \prime },P_2,P_3`$. It is thus expected that the moduli space consists of two branches, one of which is the curve given above whereas the other contains the two missing $`SO(4)`$ vacua as well as the classical point $`P_1^{}`$. The latter branch is most likely a cylinder $`^\times `$. If we parameterize it by $`z`$ the $`\eta _i`$โs on this branch will be given by
$$\eta _1=1,\eta _2=z^2,\eta _3=z^2.$$
(5.17)
The two branches meet at the classical point with an O6<sup>-</sup>-plane and two D6-branes on $`S^3`$. The structure of quantum moduli space is summarized in Figure 7.
#### 5.1.4 $`N=1`$
There are five classical points: the one corresponding to an O6<sup>-</sup>-plane and a D6-brane wrapping the $`S_>^3`$ which we denote by $`P_1`$, the free orientifold plus two D6-branes wrapping the $`S_<^3`$ which are denoted by $`P_1^{}`$ or $`P_1^{\prime \prime }`$ according to the gauge group being $`SO(2)`$ or $`Sp(1)`$, and the two resolved vacua ($`P_2,P_3`$). The orders of zeroes or poles of $`\eta _i`$ is summarized as follows:
| | $`P_1`$ | $`P_1^{}`$ | $`P_1^{\prime \prime }`$ | $`P_2`$ | $`P_3`$ |
| --- | --- | --- | --- | --- | --- |
| $`\eta _1`$ | 1 | 1 | $`1`$ | 0 | $`\mathrm{}`$ |
| $`\eta _2`$ | $`0`$ | $`\mathrm{}`$ | $`\mathrm{}^4`$ | $`0^4`$ | 1 |
| $`\eta _3`$ | $`\mathrm{}`$ | $`0`$ | $`0^4`$ | 1 | $`\mathrm{}^4`$ |
The first two have massless $`U(1)`$ gauge bosons at low energies whereas the other three have a mass gap. The moduli space should therefore consist of two smooth components, one for vacua with $`U(1)`$ and the other for mass-gapped vacua. From the analysis of zeroes and poles we expect that the branch of vacua with $`U(1)`$ is a cylinder $`\eta _2\eta _3=1`$, and the mass-gapped branch is given by the curve
$$\begin{array}{ccc}\hfill \eta _2& =& \eta _1^4(\eta _1+1)^4,\hfill \\ \hfill \eta _3& =& (\eta _1+1)^4.\hfill \end{array}$$
(5.18)
We propose that the two branches meet at a phase transition point, as shown in Figure 8. As will be discussed in detail in Section 6, one can get the exact branch structure by following a long chain of dualities to go to a โmirrorโ IIB theory. Let us summarize here the main points.
The idea is to consider another Type IIA limit by reduction on a different circle and then take its Type IIB mirror. This is the route found in and the technique is further developed in on which the present computation is based. The Type IIB dual is the non-compact Calabi-Yau $`(\xi ,\eta ;x,y^\times )`$
$$\xi \eta =F(x,y):=y^22sxy+x^32x^2+x,$$
parameterized by $`s^2`$, together with a D5-brane located at a line $`\eta =0`$, $`\xi `$ free. The D5-brane position is parameterized by a point $`(x,y)=(x_๐,y_๐)`$ of the Riemann surface $`F(x,y)=0`$, which is generically genus one and has three punctures $`๐`$, $`๐`$, $`๐`$ at $`(x,y)=(1,0),(0,0),(\mathrm{},\mathrm{})`$. The modulus $`s`$ is a normalizable dynamical variable but $`๐=(x_๐,y_๐)`$ is a coupling constant. The presence of the D5-brane generates a superpotential $`W(x_๐,s)`$, and the extremization $`_sW=0`$ relates $`๐`$ and $`s`$ as follows: When the curve $`F(x,y)=0`$ has genus one, the D5-position $`๐`$ is determined by $`s^2`$. This is what we call the $`g=1`$ branch. At $`s^2=4`$, the curve degenerates to genus zero and then $`๐`$ is free to move on this curve. We call this the $`g=0`$ branch.
On the $`g=1`$ branch there is a $`U(1)`$ vector multiplet which is an $`๐ฉ=2`$ superpartner of the complex structure modulus $`s`$. This corresponds to the upper branch in Figure 8. The points $`P_1`$ (O6<sup>-</sup> with $`D6`$) and $`P_1^{}`$ (free orientifold with $`SO(2)`$ D6โs on $`^3`$) correspond respectively to the large complex structure limit $`s^2=\mathrm{}`$ and the โorbifold limitโ $`s^2=0`$. The $`g=0`$ branch has no massless vector and corresponds to the lower branch in Figure 8. The point $`P_1^{\prime \prime }`$ (free orientifold with $`Sp(1)`$ D6โs on $`^3`$) corresponds to the limit where $`๐`$ approaches the marked point $`๐`$ while the points $`P_2`$ and $`P_3`$ (two resolved conifolds) correspond to $`๐๐,๐`$. This branch structure is reminiscent of the result of , where the behavior of vacua of $`๐ฉ=1`$ gauge theory with an adjoint matter was studied on the space of superpotential couplings.
Let us now focus on the transition point where the two branches meet. From the $`g=1`$ side, this is the point $`s^2=4`$ where a linear combination of the $`A`$ cycle and $`B`$ cycle of the torus degenerates, and Type IIB D3-brane wrapped on this vanishing cycle becomes massless. Such D3-brane states, which constitute a charged hypermultiplet $`(M,\stackrel{~}{M})`$, must be included in the low energy effective theory near the point $`s^2=4`$. The effective superpotential is then given by
$$W_{\mathrm{๐๐๐}}=W(x_๐,s)+(s^2+4)\stackrel{~}{M}M.$$
(5.19)
Variation with respect to the normalizable variables $`s`$, $`M`$ and $`\stackrel{~}{M}`$ yields
$$_sW(x_๐,s)+2s\stackrel{~}{M}M=0,(s^2+4)\stackrel{~}{M}=0,(s^2+4)M=0.$$
The solutions with $`M=\stackrel{~}{M}=0`$ leave the $`U(1)`$ gauge symmetry unbroken and lie in the $`g=1`$ branch. There are other solutions with $`s^2=4`$ in which $`\stackrel{~}{M}M`$ is determined by $`x_๐`$ through the first equation and its non-zero value higgses the $`U(1)`$. They constitute the $`g=0`$ branch.
In the original Type IIA or M-theory description, what are the particles that become massless at the transition point? Type IIB D3-branes wrapped on $`A`$ and $`B`$ cycles of the curve (plus two other directions) near the large complex structure limit $`s^2=\mathrm{}`$ correspond in M-theory on $`๐ธ_7`$ to a membrane wrapped on the $`S^2`$ bolt of a Dancerโs fiber and a fivebrane wrapped on the $`T^{1,1}`$ bolt of $`๐ธ_7`$. The corresponding objects in Type IIA orientifold with O6<sup>-</sup>+D6 wrapped on $`S^3`$ (point $`P_1`$) are essentially the non-BPS states discussed in : a membrane wrapped on the non-holomorphic $`S^2`$ bolt of Dancerโs manifold is a massive oscillation mode of the open string stretched between a D6 and its orientifold image, while a fivebrane wrapped on the bolt is a non-BPS threebrane whose tension is $`1/\mathrm{}_s^4`$ in the strong coupling limit. Thus, the charged particle responsible for the transition to the confining branch is the electrically charged massive open string mode on D6 or the magnetic non-BPS threebrane wrapped on $`S^3`$, or some dyonic bound state. Which one becomes massless at $`s^2=4`$, is an interesting question although somewhat ambiguous because of monodromies around the other special points in moduli space.
#### 5.1.5 $`N=0`$
There are at least five classical points: an O6<sup>-</sup>-plane on $`S_>^3`$ ($`P_1`$), the free orientifold plus two D6-branes on $`S_<^3`$ ($`P_1^{}`$ or $`P_1^{\prime \prime }`$ depending on the gauge groups $`SO(4)`$ or $`Sp(2)`$), and the two of resolved vacua $`P_{2,3}`$. The orders of zeroes and poles read
| | $`P_1`$ | $`P_1^{}`$ | $`(P_1^{})`$ | $`P_1^{\prime \prime }`$ | $`P_2`$ | $`P_3`$ |
| --- | --- | --- | --- | --- | --- | --- |
| $`\eta _1`$ | 1 | 1 | $`1`$ | $`1`$ | $`0`$ | $`\mathrm{}`$ |
| $`\eta _2`$ | $`0`$ | $`\mathrm{}^2`$ | $`\mathrm{}^2`$ | $`\mathrm{}^6`$ | $`0^8`$ | 1 |
| $`\eta _3`$ | $`\mathrm{}`$ | $`0^2`$ | $`0^2`$ | $`0^6`$ | 1 | $`\mathrm{}^8`$ |
Here we included the contribution of two $`SO(4)`$ vacua at $`P_1^{}`$ that are missed by the counting based on $`\stackrel{ห}{h}_{SO(4)}=2`$. It follows that the four points $`P_1^{},P_1^{\prime \prime },P_2,P_3`$ can live on a single Riemann sphere as described by the equations
$$\begin{array}{ccc}\hfill \eta _2& =& \eta _1^8(\eta _11)^2(\eta _1+1)^6,\hfill \\ \hfill \eta _3& =& (\eta _11)^2(\eta _1+1)^6.\hfill \end{array}$$
(5.20)
This curve actually follows also from Vafaโs exact superpotential. A natural guess then is that there is another branch of moduli space containing all these missing vacua.
As a non-trivial check for this guess, let us consider the orders of zeroes or poles of the functions $`\eta _i`$ on the new branch. At the $`SO(4)`$ classical point on the new branch one should have $`\eta _20^2,\eta _3\mathrm{}^2`$ to account for the missing vacua (third column of the table above). On the other hand, at the vacuum $`P_1`$ corresponds to an M-theory geometry $`๐ธ_7/_2`$ which has a finite $`S^3E_2^{}E_3^{}`$, so that $`\eta _2\mathrm{},\eta _30`$ at the corresponding classical point. Therefore, the numbers of poles and zeroes agree on the โnewโ branch if there are two distinct classical points corresponding to an O6<sup>-</sup>-plane on $`S_>^3`$. Remarkably, this is in agreement with the fact that the corresponding M-theory geometry $`๐ธ_7/_2`$ spontaneously breaks the $`_2^{}`$ symmetry of Section 3.3.
As we have seen in Section 3.3, the asymptotic symmetry of M-theory geometry is $`SO(4)\stackrel{~}{}_2\times _2^{}`$. In the interior of some solutions, a part of $`\stackrel{~}{}_2\times Z_2^{}`$ is broken. The two resolved vacua are permuted under $`\stackrel{~}{}_2`$ defined in (3.16), while all the deformed vacua are invariant. On the other hand, $`_2^{}`$ is the centralizer of the $`D_4^{}`$ orbifold group and is identified with a half-period shift along the M-theory circle. As such, it is broken in the solution $`๐ธ_7/_2`$ corresponding to an O6<sup>-</sup>-plane on $`S_>^3`$ while all other classical solutions are invariant. Thus, there should be two classical points corresponding to an O6<sup>-</sup>-plane on $`S_>^3`$, as claimed.
Are the two branches connected? From what we discussed above it is clear that $`\stackrel{~}{}_2`$ acts non-trivially on the branch containing resolved vacua, while $`_2^{}`$ acts non-trivially on the other branch. Conversely, $`_2^{}`$ should act trivially on the branch containing resolved vacua and $`\stackrel{~}{}_2`$ should act trivially on the other branch, under the assumption that both branches are $`^1`$โs with various marked points, (since both $`_2`$โs act holomorphically on moduli space and fix at least three points on the appropriate branches.) The two branches can therefore be connected only through points invariant under both $`_2`$โs. Since the classical points $`\eta _1=\pm 1`$ are fixed by $`\stackrel{~}{}_2`$ and they are the only fixed points on the resolved branch, one can conclude immediately that the two branches are not connected through interior points. The moduli space is thus made of two disjoint branches as depicted in Figure 9.
One might have guessed that a phase transition as in the $`N=1`$ case would connect the two branches because the classical moduli space is connected. Such a phase transition would be characterized by the emergence of massless particles. For the $`N=1`$ case, the relevant particle in the M-theory framework is either a five-brane wrapped on the $`T^{1,1}`$ bolt of $`๐ธ_7`$ or a membrane wrapped on the $`S^2`$. However, the corresponding cycles for the case $`N=0`$ are both non-orientable, so there will be no massless particles when they shrink to zero size.
##### $`_2^{}`$ in Type IIA <sup>5</sup><sup>5</sup>5Discussions with S. Hellerman were instrumental in shaping the following arguments.
The presence of the discrete symmetry in $`_2^{}`$ might at first sight appear puzzling from the Type IIA perspective: A half-period shift of the M-theory circle descends in Type IIA to D0-brane charge modulo $`2`$, which according to our discussion should be preserved in front of $`N>0`$ D6-branes wrapped on top of the O6<sup>-</sup>, and broken precisely for $`N=0`$. Naively, this appears to be in conflict with the familiar classification of D-brane charge in string theory. A perturbative analysis in flat space shows that the tachyonic ground state of the 0โ0 strings is in the symmetric representation. The tachyon is therefore not orientifolded out even for a single D0-brane, which should therefore not produce a conserved charge. This situation is T-dual to the D3-brane in Type I, and in distinction to the D($`1`$)-brane there, for which the ground state is in the anti-symmetric representation.
Wrapping on $`S^3`$ does not eliminate the tachyon, and one would therefore not expect a stable D0-brane in our backgrounds, for any $`N`$. On the other hand, the M-theory analysis solidly establishes the existence of $`_2^{}`$ for $`N>0`$, and its breaking for $`N=0`$.
To reconcile the two points of view, we note that the analog of $`_2^{}`$ can already be seen in the context of O6<sup>-</sup>/D6 systems in flat space and their M-theory lifts, with the same breaking pattern. The discrete symmetry in this case can be nicely interpreted from the perspective of D2-brane probes . (Its presence was also noticed, for example, in the D0-brane scattering analysis of .) The worldvolume theory of a D2-brane pair is a $`3d`$ $`๐ฉ=4`$ supersymmetric $`Sp(1)`$ gauge theory with $`2N`$ half-hypermultiplets in the fundamental representation. The hypers are from the $`2`$-$`6`$ strings and their masses parameterize the position of the D6-branes. Now, one may interpret the $`_2^{}`$ as the global symmetry extending $`SO(2N)`$ by $`O(2N)`$. It acts as the sign flip of just one of the $`2N`$ half-hypermultiplets. That this correspond to a half-period shift of the M-theory circle follows from the fact that, in an odd monopole background, the fermion zero mode measure is odd under this symmetry. When some of the D6-brane pairs are on top of the O6-plane, the sign flip of one of the corresponding massless half-hypermultiplets is a symmetry of the system. When all D6-branes are away from the O6, all the masses are turned on, and the symmetry is broken. Namely, $`_2^{}`$ is a symmetry only when at least a single D6-brane is exactly on top of the O6-plane, and broken otherwise. This is indeed the breaking pattern we have found in the curved background.
Returning to the interpretation of $`_2^{}`$ as โD0-brane number modulo 2โ, we note that this is a good quantum number far away from the O6-plane. Indeed, a D0-brane is then far away (in the covering space) from its anti-D0-brane image, and the ground state of the open string between them is non-tachyonic. (With more than one D0-$`\overline{\mathrm{D0}}`$ pair in the covering space, a D0 and an anti-D0 from different pairs can approach each other asymptotically, and annihilate.) In a process in which the D0 approaches an O6<sup>-</sup>-plane, the tachyon will develop and the D0/$`\overline{\mathrm{D0}}`$ can annihilate. When $`N=0`$, this is fine since $`_2^{}`$ is broken. When $`N>0`$, in which case $`_2^{}`$ is globally conserved, this charge is carried away by open strings on the D6-brane describing its separation from the O-plane. (These states are massive after wrapping on $`S^3`$, but the D0-brane is much heavier at weak string coupling.)
Let us also note that $`_2^{}`$ is preserved in front of an O6<sup>+</sup>-plane (our case $`N=4`$). This is consistent with the fact that the 0โ0 tachyon lands in the anti-symmetric representation and is orientifolded out.
#### 5.1.6 $`N<0`$
Finally we come back to the regular case. The singularities can be read from the following table:
| | $`P_1^{}`$ | $`P_1^{\prime \prime }`$ | $`P_2`$ | $`P_3`$ |
| --- | --- | --- | --- | --- |
| $`\eta _1`$ | 1 | $`1`$ | 0 | $`\mathrm{}`$ |
| $`\eta _2`$ | $`\mathrm{}^{2(\stackrel{~}{N}1)}`$ | $`\mathrm{}^{2(\stackrel{~}{N}+1)}`$ | $`0^{4\stackrel{~}{N}}`$ | 1 |
| $`\eta _3`$ | $`0^{2(\stackrel{~}{N}1)}`$ | $`0^{2(\stackrel{~}{N}+1)}`$ | 1 | $`\mathrm{}^{4\stackrel{~}{N}}`$ |
where $`\stackrel{~}{N}=2N>2`$. And the quantum curve is:
$$\begin{array}{cc}\hfill \eta _2=& \eta _1^{4\stackrel{~}{N}}(\eta _11)^{2(\stackrel{~}{N}1)}(\eta _1+1)^{2(\stackrel{~}{N}+1)},\hfill \\ \hfill \eta _3=& (\eta _11)^{2(\stackrel{~}{N}1)}(\eta _1+1)^{2(\stackrel{~}{N}+1)}.\hfill \end{array}$$
(5.21)
The moduli space consists of a single smooth component as in Figure 10.
### 5.2 Using superpotential
A part of the results of the previous subsection can be obtained also by studying the superpotential proposed by Vafa and the relevant computations in .
The superpotential is computed on the branch of the resolved conifold with flux through $`^2`$. It consists of three parts, coming from the four-form flux, the two-form flux and worldsheet instantons:
$$W=W_{4\mathrm{flux}}+W_{2\mathrm{flux}}+W_{\mathrm{crosscap}}.$$
Let $`t`$ be the complexified Kรคhler class of the base $`^1`$ of the resolved conifold. In the present orientifold, the periodicity of the parameter is doubled,
$$tt+4\pi i.$$
(5.22)
The reason is that there exist crosscap diagrams associated with the odd degree maps $`S^2^1`$ which are equivariant with respect to the involution $`\mathrm{\Omega }:w1/\overline{w}`$ on the domain and the anti-holomorphic involution (2.17) on the target. For example, the identity map is such a map and has degree 1. The path-integral weight of such a diagram include odd powers of $`e^{i\mathrm{Im}(t)/2}`$. Thus, $`z=e^{t/2}`$ is the single valued parameter of the theory. Now, let us describe each term of the superpotential. The contributions from the RR two-form flux through the $`^2`$ and worldsheet instantons are
$$W_{2\mathrm{flux}}=(N2)\frac{F_0}{t}=(N2)Li_2(z^2),$$
$$W_{\mathrm{crosscap}}=4\underset{m:\mathrm{odd}1}{}\frac{e^{mt/2}}{m^2}=2(Li_2(z)Li_2(z)),$$
where $`Li_2`$ is the dilogarithm function. For convenience, some of its properties are collected in Appendix B. The contribution from the four-form flux is given by
$$W_{4\mathrm{flux}}=F_4\omega ,$$
where $`\omega `$ is the complexified Kรคhler form of the resolved conifold. According to , $`F_4`$ has an imaginary part corresponding to the RR four-form and a real NSNS part coming from the failure of the metric to be Calabi-Yau. On a non-compact space it is natural to interpret this formula as follows:
$$W_{4\mathrm{flux}}=_Md(\mathrm{Re}(\widehat{\mathrm{\Omega }})+iC_3)\omega =_M(\mathrm{Re}(\widehat{\mathrm{\Omega }})+iC_3)\omega ,$$
where $`\widehat{\mathrm{\Omega }}`$ is a suitably normalized form of the holomorphic three-form which is the superpartner of $`C_3`$. Evaluating this for the $`_2`$ quotient of the conifold we obtain
$$W_{4\mathrm{flux}}=Yt/2,$$
where $`Y`$ is the holomorphic volume of the boundary $`S^3`$:
$$Y:=_{S_{\mathrm{}}^3}(\mathrm{Re}(\widehat{\mathrm{\Omega }})+iC_3).$$
(5.23)
Summing up the three terms, we obtain the total superpotential
$$W=Yt/2(N2)Li_2(z^2)2(Li_2(z)Li_2(z)).$$
(5.24)
Note that the parameters $`t`$ and $`Y`$ introduced here are related to the coordinates that we used in the M-theory description of the moduli space as
$$z=\mathrm{exp}(t/2)=\eta _1,\mathrm{exp}(Y)=\eta _3.$$
(5.25)
Following we can use this superpotential to find the exact form of the moduli space. In order to have a supersymmetric background we should vary the superpotential with respect to $`t`$ and find a stationary point. This gives
$$_tW=0Y=\mathrm{log}\left((z1)^{(2N2)}(z+1)^{(2N6)}\right).$$
(5.26)
Using the relation (5.25), we see that this is nothing but the equation describing the component of the moduli space which is smoothly connected to the resolved geometries.
### Comparison to 4d gauge theory
The superpotential (5.24) has two branch cuts (see Figure 11), starting at $`z=\pm 1`$. For $`z=\pm e^\epsilon `$ with small $`\epsilon `$, we have
$$W=Y\epsilon b\epsilon \mathrm{log}\epsilon +\mathrm{},b=\{\begin{array}{cc}2N2\hfill & \text{at }z=1,\hfill \\ 2N6\hfill & \text{at }z=1,\hfill \end{array}$$
(5.27)
where $`+\mathrm{}`$ is a power series in $`\epsilon `$. When $`b`$ is positive (resp. negative), by the relation (5.26), $`\mathrm{Re}(Y)`$ diverges to $`+\mathrm{}`$ (resp. $`\mathrm{}`$) as $`\epsilon 0`$. This is the classical limit where we have a large minimal three sphere $`S_>^3`$ (resp. $`S_<^3`$) on which a certain number of D6-branes (resp. anti-D6-branes) are wrapped. Below we compare this behavior of the superpotential with what we expect from the gauge theory on the sixbranes.
For this purpose, one needs to understand the precise relation of the parameter $`Y`$, which can be regarded as the membrane instanton action on a cycle homologous to $`S_{\mathrm{}}^3`$, and the holomorphic gauge coupling constant $`\frac{8\pi ^2}{g^2}i\theta `$ of the $`4d`$ gauge theory on the (anti-)D6-branes wrapped on the minimal three sphere. This was discussed in for the case with O6-plane at $`S^3`$. Here we present another argument, using the embedding of $`SO(2n)`$ or $`Sp(n)`$ into $`U(2n)`$ defined by the orientifold projection, which is applicable to the more general systems we are studying. We first note that the instanton number of $`4d`$ Yang-Mills theory of a simple gauge group $`G`$ is defined as
$$k=\frac{1}{8\pi ^2}_^4\mathrm{Tr}(F_AF_A),$$
where $`\mathrm{Tr}`$ is such that the long root has length squared 2, or the trace in the adjoint representation is given by $`\mathrm{tr}_{\mathrm{adjoint}}(XY)=2\stackrel{ห}{h}\mathrm{Tr}(XY)`$, with $`\stackrel{ห}{h}`$ being the dual Coxeter number of $`G`$. For the groups $`SO(2n)`$ and $`Sp(n)`$, it is related to the trace of the fundamental representation of $`U(2n)`$ by
$$\mathrm{Tr}(F_AF_A)=\sigma \times \mathrm{tr}_{\mathrm{fund}}(F_{\stackrel{~}{A}}F_{\stackrel{~}{A}}),\sigma =\{\begin{array}{cc}\frac{1}{2}& \text{for }SO(2n)\hfill \\ 1& \text{for }Sp(n)\hfill \end{array}$$
where $`\stackrel{~}{A}`$ is the $`U(2n)`$ gauge field which is obtained from $`A`$ by the embedding of $`SO(2n)`$ or $`Sp(n)`$ into $`U(2n)`$. Thus, one $`SO(2n)`$ instanton corresponds to two $`U(2n)`$ instantons, while one $`Sp(n)`$ instanton corresponds to one $`U(2n)`$ instanton.
First, we consider the case $`\mathrm{Re}(Y)0`$ for which the Type IIA geometry is the deformed conifold with an O6-plane at the large minimal three sphere $`S_>^3`$. Note that this $`S_>^3`$ is homologous to $`S_{\mathrm{}}^3`$, and hence $`Y`$ is the action for one D2-brane wrapped on $`S_>^3`$. Before orientifold, this D2-brane is wrapped twice on $`S^3`$ and corresponds to two $`U(2n)`$ instantons on $`2n`$ D6-branes wrapped on $`S^3`$. By the remark above, after orientifold, it corresponds to one $`SO(2n)`$ instanton or two $`Sp(n)`$ instantons. This is indeed the claim in . Thus, the relation of $`Y`$ and the holomorphic gauge coupling is
$$Y=2\sigma \left(\frac{8\pi ^2}{g^2}i\theta \right).$$
(5.28)
Note that the Chan-Paton factor of the D2-branes on $`S_>^3`$ is symplectic for O6<sup>-</sup> ($`SO(2n)`$ on $`n`$ D6), while it is orthogonal for O6<sup>+</sup> ($`Sp(n)`$ on $`n`$ D6). Thus, the number of D2-branes for O6<sup>-</sup> must be even in the double cover. For O6<sup>+</sup>, on the other hand, a โhalfโ D2-brane (one in the cover) is allowed and corresponds to one $`Sp(n)`$ Yang-Mills instanton with instanton factor $`e^{Y/2}`$. In terms of $`S=2\sigma \epsilon `$, the superpotential (5.27) can be written as
$$W=\left(\frac{8\pi ^2}{g^2}i\theta \right)S\frac{b}{2\sigma }S\mathrm{log}S+\text{power series in }S.$$
(5.29)
If we identify $`z=1`$ as the large $`S_>^3`$ limit with O6<sup>-</sup>-plane and $`z=1`$ as the large $`S_>^3`$ limit with O6<sup>+</sup>-plane, the coefficient of the $`S\mathrm{log}S`$ term is
$$\frac{b}{2\sigma }=\{\begin{array}{cc}2N2\hfill & \text{at }z=1\hfill \\ N3\hfill & \text{at }z=1,\hfill \end{array}$$
which are the dual Coxeter number of the groups $`SO(2N)`$ and $`Sp(N4)`$ respectively. Then, (5.29) is exactly the Veneziano-Yankielowicz superpotential, up to a power series in $`S`$, for the gauge group $`SO(2N)`$ for $`z=1`$ and $`Sp(N4)`$ for $`z=1`$.
Let us next consider the case $`\mathrm{Re}(Y)0`$ which corresponds to the free orientifold of the deformed conifold with large minimal three sphere $`S_<^3`$. The cycle $`S_<^3`$ is also homologous to $`S_{\mathrm{}}^3`$, and $`Y`$ corresponds to the action for D2-brane wrapped twice on $`^3=S_<^3/_2`$. This again corresponds to two $`U(2n)`$ instantons on $`2n`$ D6-branes wrapped on $`S_<^3`$ before orientifold, and thus to one $`SO(2n)`$ or two $`Sp(n)`$ instanton after orientifold. On anti-D6-branes, instantons with positive instanton numbers correspond to anti-D2-branes. Thus the relation of $`Y`$ and the holomorphic gauge coupling on the anti-D6-branes is
$$Y=2\sigma \left(\frac{8\pi ^2}{g^2}i\theta \right).$$
(5.30)
As in the $`\mathrm{Re}(Y)0`$ case, the minimal instanton factor is $`e^Y`$ for $`SO(2n)`$ and $`e^{Y/2}`$ for $`Sp(n)`$. In terms of $`S=2\sigma \epsilon `$, the superpotential is
$$W=\left(\frac{8\pi ^2}{g^2}i\theta \right)S+\frac{b}{2\sigma }S\mathrm{log}S+\text{power series in }S.$$
This is exactly the expected Veneziano-Yankielowicz superpotential if we identify $`z=1`$ (resp. $`z=1`$) as the large $`S_<^3`$ limit with $`SO`$ (resp. $`Sp`$) gauge group. Indeed, under this identification, the coefficient of the $`S\mathrm{log}S`$ term is
$$\frac{b}{2\sigma }=\{\begin{array}{cc}2N+2=2\stackrel{~}{N}2\hfill & \text{at }z=1\hfill \\ N+3=\stackrel{~}{N}+1\hfill & \text{at }z=1,\hfill \end{array}$$
which are the dual Coxeter number of the groups $`SO(2\stackrel{~}{N})`$ and $`Sp(\stackrel{~}{N})`$ respectively, where $`\stackrel{~}{N}=2N`$ is half the number of anti-D6-branes.
### Summary on the component including the resolved conifolds
By the above comparison with $`4d`$ gauge theory along with the information about the expected classical gauge groups in the various limits, we arrive at the following consistent set of rules for reading the classical configurations from the superpotential analysis depending on the RR charge and the particular point in the moduli space:
* at the point $`z=1`$ there is an $`Sp`$ group, the number of vacua (or better said the instanton counting) is $`2N6`$ . The gauge group will depend on the number of vacua:
+ if $`2N6>0`$ the gauge group is $`Sp(N4)`$. The classical description is an O6<sup>+</sup> with $`N4`$ D6-branes.
+ if $`2N6=0`$ there is no classical limit at $`z=1`$.
+ if $`2N6<0`$ the gauge group is $`Sp(2N)`$. The classical description is a free orientifold with $`42N`$ anti-D6-branes.
* At the point $`z=1`$ the structure is similar:
+ if $`2N2>0`$ the gauge group is $`SO(2N)`$. The classical description is an O6<sup>-</sup> with $`N`$ D6-branes.
+ if $`2N2=0`$ there is no classical limit at $`z=1`$.
+ if $`2N2<0`$ the gauge group is $`SO(42N)`$. The classical description is the free orientifold with $`42N`$ anti-D6-branes.
These rules allow us to classify the different possibilities depending on the RR charge:
* For $`N>3`$ there are four classical points:
O6<sup>-</sup> and $`N`$ D6, gauge group $`SO(2N)`$
O6<sup>+</sup> and $`N4`$ D6, gauge group $`Sp(N4)`$
two free orientifolds of the resolved conifold
* For $`N=3`$ there are three classical points (the curve is the same as the $`SU(4)`$ curve):
O6<sup>-</sup> and $`N`$ D6, gauge group $`SO(6)SU(4)`$
two free orientifolds of the resolved conifold
* For $`N=2`$ there are four classical points (and others sitting on the other branch):
O6<sup>-</sup> and $`N`$ D6, gauge group $`SO(4)`$
free orientifold of deformed conifold, gauge group $`Sp(0)`$ (= nothing)
two free orientifolds of the resolved conifold
* For $`N=1`$ there are 3 classical points (two others sitting on the branch with an infrared $`U(1)`$):
free orientifold and $`2`$ anti-D6 on $`^3`$, gauge group $`Sp(1)=SU(2)`$
two free orientifolds of the resolved conifold
* For $`N=0`$ there are 4 classical points (and a few others on the other branch):
free orientifold and $`4`$ anti-D6 on $`^3`$, gauge group $`SO(4)`$
free orientifold and $`4`$ anti-D6 on $`^3`$, gauge group $`Sp(2)`$
two free orientifolds of the resolved conifold
* For $`N<0`$ there are four classical points:
free orientifold and $`42N`$ anti-D6 on $`^3`$, gauge group $`SO(2(2N))`$
free orientifold and $`42N`$ anti-D6 on $`^3`$, gauge group $`Sp(2N)`$
two free orientifolds of the resolved conifold
## 6 Exact branch structure for $`N=1`$
In this section we wish to present a detailed analysis of the branch structure for the $`N=1`$ case of the previous section from various dual pictures. As advertised, an exact description of the quantum moduli space is obtained by moving to a mirror Type IIB theory.
We consider the strong coupling limit of the original Type IIA, for which the relevant M-theory geometry is asymptotically a cone over $`(S^3\times S^3)/D_3^{}`$. Here the group $`D_3^{}_4`$ is introduced in Section 3.3 and acts on the triplet of $`SU(2)`$ matrices as
$$(g_1,g_2,g_3)(i\tau _2g_1,g_2,g_3).$$
Dimensional reduction along the orbit of the $`U(1)`$ action $`(g_1,g_2,g_3)(e^{i\alpha \tau _3}g_1,g_2,g_3)`$ brings the system back to the original Type IIA. If we reduce instead along a diagonal $`S^1`$,
$$(g_1,g_2,g_3)(e^{i\alpha \tau _2}g_1,e^{i\alpha \tau _2}g_2,e^{i\alpha \tau _3}g_2),$$
the resulting Type IIA configuration is a partially blown-up orbifold $`^3/D_3^{}`$ with a D6-brane of topology $`^2\times S^1`$. Here the generator of $`D_3^{}`$ acts on the coordinates of $`^3`$ as
$$(z_1,z_2,z_3)(z_1,iz_2,iz_3).$$
(6.1)
The whole system admits a GLSM description.
The main advantage of this framework is that one can read off the quantum moduli space rather directly by moving to the mirror IIB description. This chain of dualities was used in to study the geometric transition of D6-branes wrapped on the $`S^3`$ of deformed conifold, and also applied in to study its $`_2`$ orbifold.
#### GLSM description
Consider the edge vectors for the toric fan of $`^3`$,
$$v_1=(1,0,0),v_2=(0,1,0),v_3=(0,0,1),$$
generating the lattice $`^3`$ of torus actions. The orbifolding makes the lattice finer by the inclusion of an extra generator $`\rho =(\frac{1}{2},\frac{1}{4},\frac{1}{4})`$. The fan for Calabi-Yau resolution of orbifold singularity is given by including the lattice points $`v_4=(\frac{1}{2},\frac{1}{4},\frac{1}{4})`$ and $`v_5=(0,\frac{1}{2},\frac{1}{2})`$ as new edge vectors and subdividing the fan. See Figure 12.
The GLSM consists of five chiral fields $`z_{1,\mathrm{},5}`$ and $`U(1)^2`$ gauge symmetry. The $`U(1)`$ charges which span the Mori cone are given by
$$Q^1=(1,0,0,2,1),Q^2=(0,1,1,0,2),$$
(6.2)
so that the fully blown-up phase is given by the D-term equations
$$|z_1|^22|z_4|^2+|z_5|^2=r^1>0,|z_2|^2+|z_3|^22|z_5|^2=r^2>0.$$
(6.3)
In the orbifold phase with negative FI parameters the fields $`z_4,z_5`$ acquire vev and break the gauge group down to $`_4`$, which acts $`(z_1,z_2,z_3)`$ to $`(z_1,iz_2,iz_3)`$. It is useful to draw the toric skeleton diagram describing the base polytope of $`T^3`$ fibration.
In addition to the Kรคhler parameters, one has to specify the location of D6-brane which projects to a half-line ending on a one-dimensional face of the toric polytope. As was discussed in for related models, one cannot choose the Kรคhler parameters and the location of the D6-brane independently, because of a superpotential generated by the D6-brane. Also, by looking into the asymptotics one finds that the resolution mode corresponding to $`z_5`$ should not be turned on. In other words, in the skeleton diagram the points A, Aโ should coincide.
#### Three-dimensional Five-brane Web
One can understand the effect of superpotential semiclassically by relating our GLSM picture to a Type IIB five-brane web. Under a suitable choice of basis for the charge, the partial blowup of our orbifold $`^3/_4`$ is mapped to the two-dimensional web of Figure 13. The D6-brane in GLSM picture turns into another five-brane leg carrying a new kind of charge ending on a leg of the web. Addition of such a five-brane makes the web three-dimensional. The supersymmetry condition for the resulting web constrains the allowed locations of the D6-brane endpoint for each choice of Kรคhler parameter.
The charges of each leg of the three-dimensional web can be obtained by regarding the $`G_2`$ holonomy manifolds as $`T^3`$ fibrations, and analyzing the locus of degenerate fiber in the base. Since these charges are relevant in calculating the superpotential, let us go back to M-theory and calculate them explicitly. As the $`T^3`$ we take the orbit of $`U(1)^3`$ action
$$[\alpha _1,\alpha _2,\alpha _3]:(g_1,g_2,g_3)(e^{2\pi i\alpha _1\tau _2}g_1,e^{2\pi i\alpha _2\tau _2}g_2,e^{2\pi i\alpha _3\tau _2}g_3),$$
(6.4)
modulo identification
$`[\alpha _1,\alpha _2,\alpha _3]`$ $``$ $`[\alpha _1,\alpha _2,\alpha _3]+[\frac{1}{4},\frac{1}{2},0]`$ (6.5)
$``$ $`[\alpha _1,\alpha _2,\alpha _3]+[\frac{1}{2},0,0]`$
$``$ $`[\alpha _1,\alpha _2,\alpha _3]+[\frac{1}{2},\frac{1}{2},\frac{1}{2}].`$
In the following we will also use the coordinate $`(X,\stackrel{~}{X})(g_2g_1^1,g_3g_1^1)`$ of $`S^3\times S^3`$.
The classical moduli space consists of the following branches. A free orientifold of deformed conifold with two D6-branes corresponds to $`g_1`$ filled in, and the orientifold of resolved conifold with flux correspond to either $`g_2`$ or $`g_3`$ filled in. The O6<sup>-</sup>+D6 configuration corresponds to the geometry with the following $`S^1`$ shrinking at $`r=0`$,
$$(X,\stackrel{~}{X})(Xe^{2\pi i\alpha \tau _2},\stackrel{~}{X}e^{2\pi i\alpha \tau _2}).$$
(6.6)
The cross-section of the seven-manifold at any finite $`r`$ is $`S^3\times S^3`$. It can be viewed as a $`T^3`$ fibration over $`S^3`$, and a section is given by
$$(X,\stackrel{~}{X})=(e^{i\theta \tau _1/2}e^{i\phi \tau _2},e^{i\stackrel{~}{\theta }\tau _1/2}),(0\theta \pi ,0\stackrel{~}{\theta }\pi ,0\phi 2\pi ).$$
At four special points on the base, the fiber $`T^3`$ has a vanishing one-cycle labeled by the ratio of $`[\alpha _1,\alpha _2,\alpha _3]`$,
$$\begin{array}{ccc}(\theta ,\stackrel{~}{\theta })=(0,0)\hfill & \mathrm{}& [\alpha _1,\alpha _2,\alpha _3]=[+\frac{1}{2},+\frac{1}{2},+\frac{1}{2}],\hfill \\ (\theta ,\stackrel{~}{\theta })=(0,\pi )\hfill & \mathrm{}& [\alpha _1,\alpha _2,\alpha _3]=[\frac{1}{2},\frac{1}{2},+\frac{1}{2}],\hfill \\ (\theta ,\stackrel{~}{\theta })=(\pi ,0)\hfill & \mathrm{}& [\alpha _1,\alpha _2,\alpha _3]=[\frac{1}{2},+\frac{1}{2},\frac{1}{2}],\hfill \\ (\theta ,\stackrel{~}{\theta })=(\pi ,\pi )\hfill & \mathrm{}& [\alpha _1,\alpha _2,\alpha _3]=[+\frac{1}{2},\frac{1}{2},\frac{1}{2}].\hfill \end{array}$$
(6.7)
Thus we find four half-lines of degenerate fiber. There are additional loci of degenerate fiber at $`r=0`$. For the classical branch with $`g_1`$ filled in one finds a line segment of degenerate $`[1,0,0]`$ one-cycle at $`r=0`$, and similarly for the other two branches where $`g_2`$ or $`g_3`$ are filled in. For the last branch, one finds that the vanishing $`S^1`$ given in (6.6) lies along the $`T^3`$ fiber when
$$\begin{array}{ccc}\theta =0\hfill & \mathrm{}& [\alpha _1,\alpha _2,\alpha _3]=[+\frac{1}{4},+\frac{1}{2},0],\hfill \\ \theta =\pi \hfill & \mathrm{}& [\alpha _1,\alpha _2,\alpha _3]=[+\frac{1}{4},\frac{1}{2},0],\hfill \\ \stackrel{~}{\theta }=0\hfill & \mathrm{}& [\alpha _1,\alpha _2,\alpha _3]=[\frac{1}{4},0,\frac{1}{2}],\hfill \\ \stackrel{~}{\theta }=\pi \hfill & \mathrm{}& [\alpha _1,\alpha _2,\alpha _3]=[\frac{1}{4},0,+\frac{1}{2}].\hfill \end{array}$$
(6.8)
By re-labeling the shrinking one-cycle in terms of the fundamental periods given in (6.5), one finds that the four half-lines of (6.7) are labeled by
$$(0,0,1),(2,1,1),(2,1,1),(0,2,1).$$
(6.9)
We can identify them with the charges of semi-infinite five-brane legs forming three-dimensional webs. Indeed, the first leg corresponds to the D6-brane while the other three project to the legs of two-dimensional web of Figure 13. The four families of $`G_2`$ holonomy spaces thus correspond to the three-dimensional webs summarized in Figure 14. All these four webs are rigid apart from overall rescaling.
#### IIB mirror
The mirror of this GLSM is given by a LG model with superpotential
$$F(x,y)=y^2+e^{\frac{t_1}{2}+\frac{t_2}{4}}xy+x^3+e^{\frac{t_2}{2}}x^2+x.$$
(6.10)
Here $`(x,y)`$ are $`^\times `$-valued chiral superfields, and $`t_1,t_2`$ are the Kรคhler parameters in the GLSM. The spacetime theory is a Type IIB superstring on a local Calabi-Yau manifold
$$\xi \eta =F(x,y),$$
(6.11)
with a D5-brane wrapping a holomorphic curve $`\eta =F(x,y)=0`$. We denote by $`(x_๐,y_๐)`$ its position on the curve $`\mathrm{\Sigma }:F(x,y)=0`$.
The parameter $`t_2`$ is fixed by requiring the curve $`\mathrm{\Sigma }`$ to have only three punctures. There seems to be two choices $`e^{\frac{t_2}{2}}=\pm 2`$, but they should lead to the same structure for the moduli space so we choose $`e^{\frac{t_2}{2}}=2`$. The remaining normalizable parameter $`sie^{\frac{t_1}{2}}/\sqrt{2}`$ in the curve,
$$\mathrm{\Sigma }:y^22sxy+x^32x^2+x=0,$$
(6.12)
is stabilized by the superpotential generated by the D5-brane.
$$W(s)=_๐^๐+_๐^๐\mathrm{log}yd\mathrm{log}x.(๐(1,0),๐(0,0),๐(\mathrm{},\mathrm{})).$$
(6.13)
This integration contour is chosen from the observation that the four semi-infinite five-brane legs all have non-zero โthirdโ charge, which we interpret as the presence of D5 or $`\overline{\mathrm{D5}}`$-branes at three punctures A, B, C of the curve.
Note that the moduli space of curves $`\mathrm{\Sigma }`$ is the complex $`s`$-plane modulo identification $`ss`$, since $`se^{\frac{t_1}{2}}`$. Indeed, the sign flip of $`s`$ can be absorbed by the sign flip of $`y`$. The good modulus of the curve is therefore $`s^2`$.
For generic $`s`$ the curve $`\mathrm{\Sigma }`$ is of genus one. The degeneration to genus zero occurs at $`s=0`$ or $`s=\pm 2i`$. For genus one curve, the modulus $`s`$ is related to the position of D because of the superpotential. Using $`\widehat{y}=ysx`$ we rewrite the equation for the curve as
$$\widehat{y}^2+x^3(2+s^2)x^2+x=0.$$
(6.14)
The F-term condition reads
$$_{\overline{\mathrm{D5}}=๐,๐}^{\mathrm{D5}=๐,๐}\frac{dx}{\widehat{y}}=0.$$
(6.15)
One can solve this equation by fixing $`s`$ in such a way that there is a meromorphic function on $`\mathrm{\Sigma }`$ with simple zeroes at B, D and simple poles at $`๐,๐`$.
$$\left|\begin{array}{ccc}\hfill 1& \hfill x_๐& \hfill \widehat{y}_๐\\ \hfill 1& \hfill x_๐& \hfill \widehat{y}_๐\\ \hfill 1& \hfill x_๐& \hfill \widehat{y}_๐\end{array}\right|=\left|\begin{array}{ccc}1& 1& s\\ 1& 0& 0\\ 1& x_๐& \widehat{y}_๐\end{array}\right|=sx_๐\widehat{y}_๐=0(x_๐,y_๐)=(1,2s).$$
(6.16)
The moduli space for the case $`N=1`$ is made of $`g=0`$ and $`g=1`$ branches. The $`g=1`$ branch is the moduli space of genus one curves with three punctures, and is the complex $`s^2`$-plane as explained above. It has special points $`s^2=0,4`$ and infinity. At $`s=2i`$ the curve degenerates to genus zero,
$$\mathrm{\Sigma }:y^24ixy+x^32x^2+x=0.$$
(6.17)
This curve itself is regarded as the $`g=0`$ branch. Various semi-classical points on moduli space are identified with those in the conifold picture as follows:
| branch | point | gauge group | description |
| --- | --- | --- | --- |
| $`g=0`$ | A | $`Sp(1)`$ | deformed<sub>(ฮผ\<0)</sub>, with 2D6 on $`^3`$ |
| $`g=0`$ | B,C | none | resolved, with ($`1`$) flux |
| $`g=1`$ | $`s^2=0`$ | $`O(2)`$ | deformed<sub>(ฮผ\<0)</sub>, with 2D6 on $`^3`$ |
| $`g=1`$ | $`s^2=\mathrm{}`$ | $`O(2)`$ | deformed<sub>(ฮผ\>0)</sub>, with O6<sup>-</sup>+D6 on $`S^3`$ |
By a suitable identification of the variables $`\eta _i`$ with $`(x,y)`$ one can identify the $`z`$-plane of Section 5.2 with the $`g=0`$ branch here, and also see the expected behavior of $`\eta _i`$โs on the other branch which can be identified with the $`s^2`$-plane. The key fact is that the membrane instantons wrapping on three-cycles $`\{2E_1^{},E_2^{},E_3^{}\}`$ of Section 5.1 turn into disc instantons bounded by the D6-brane. Matching of their volumes gives
$$\eta _1^2e^{(|z_2|^2|z_3|^2)/2}x,\eta _2e^{|z_1|^2|z_3|^2}x^3y^2,\eta _3e^{|z_2|^2|z_1|^2}x^1y^2,$$
(6.18)
where we used the standard mirror identification
$$e^{|z_1|^2}:e^{|z_2|^2}:e^{|z_3|^2}:e^{|z_4|^2}:e^{|z_5|^2}y^2:x^3:x:xy:x^2,$$
(6.19)
and โ$``$โ expresses the identification up to phase. Under this identification, the $`g=0`$ curve (6.17) precisely agrees with the relation (5.18) among $`\eta _i`$ on the branch of mass-gapped vacua obtained in the previous section. Note that the $`\eta _1`$ defined as
$$\eta _1=i\left(\frac{y2ix}{1+x}\right)$$
(6.20)
is single valued on $`\mathrm{\Sigma }`$ and it indeed squares to $`x`$ on $`\mathrm{\Sigma }`$. Note also that, although the $`g=0`$ curve (6.17) has a double point, the singularity is blown up on generic points of $`g=0`$ branch and the $`\eta _i`$ indeed take two different values there. The functional form of $`\eta _i`$โs on the $`g=1`$ branch is obtained simply by substituting $`(x,y)=(1,2s)`$ into (6.18):
$$\eta _1=1,\eta _2=\frac{1}{4s^2},\eta _3=4s^2.$$
(6.21)
This agrees with the expectation in Section 5.1.4 that the branch of vacua with infrared $`U(1)`$ is a cylinder $`\eta _2\eta _3=1`$.
## 7 Other cases
We wish to briefly comment on the possibility of $`\mu `$-transitions at the conifold when the orientifold action is in one of the other classes discussed in Section 2. It is clear that it will be much harder to find the associated $`G_2`$ holonomy metrics. Nevertheless, qualitative considerations similar to the ones we sometimes used above give good indications when we should expect a $`\mu `$-transition.
### 7.1 Case (1)$``$(3)
The first thing to notice in cases (1) and (3) is that the O-plane intersects the compact three-cycle of the deformed conifold. This creates flux that cannot escape to infinity, and we should cancel the flux by wrapping a fixed number of D6-branes. A simple class of cycles to wrap branes around is the fixed point set of some anti-holomorphic involution which can be different from the one used to define the orientifold, but must be in the same class to preserve supersymmetry.
In case (1), the O6-plane is the fixed point locus of the involution $`zM_O\overline{z}`$, where $`M_O`$ is the orthogonal matrix $`\mathrm{diag}(1,1,1,1)`$. The O6-plane is topologically $`S^2\times `$ and the intersection number with the $`S_>^3`$ is two. If this is a standard O6<sup>-</sup>-plane with negative twice the charge of a D6-brane (as measured in the quotient space), this means that we need to wrap a D6-brane configuration that intersects the compact cycle exactly four times (namely, an invariant configuration in the covering space intersecting $`S_>^3`$ in eight points). We know of supersymmetric cycles intersecting the $`S_>^3`$ twice: The fixed point locus of $`zM\overline{z}`$, where $`M`$ is another orthogonal matrix with eigenvalues $`(1,1,1,1)`$ (generally distinct from $`M_O`$). We can write $`M=UM_OU^T`$, where $`U`$ is an element of $`SO(4)`$, and two $`U`$โs give the same $`M`$ if they differ by an element of $`SO(3)`$. In the covering space, the possible $`M`$โs live in $`SO(4)/SO(3)S^3`$, and we need two (pairs of) such cycles to cancel the charge. Orientifolding maps the brane associated with $`U`$ to the brane associated with $`M_OUM_O`$, which corresponds to acting on $`S^3`$ as the element $`\mathrm{diag}(1,1,1,1)`$ of $`SO(4)`$. In the orientifold, the space of possible brane wrappings is therefore the symmetric product $`_{}^{(1)}๐ฎ^2()`$, where $`=S^3/_2`$ with the given action of $`_2`$.
Note that if the O-plane is an O6<sup>-</sup>, we cannot wrap branes on fixed point loci of involutions with three negative eigenvalues, since those would preserve the opposite supersymmetry, and we cannot wrap anti-branes if we are to cancel the charge.
On the other hand, if the orientifold plane is an O6<sup>+</sup> with positive twice the charge of a D6-brane, we need anti-branes to cancel the charge, and those are most conveniently wrapped on fixed point loci of $`zM\overline{z}`$, where $`M`$ is an orthogonal matrix with eigenvalues $`(1,1,1,1)`$. As we recall, this gives a cycle with two disconnected components, each of which is a copy of $`^3`$ and intersects the $`S_>^3`$ once. In that case, we have the choice of four such cycles, each of which is again parameterized by the choice of $`S^3/_2\left(SO(4)/SO(3)\right)/_2`$. Thus, $`_+^{(1)}๐ฎ^4()`$.
In case (3), the situation is reversed: If the O-plane on $`^3^3`$ is an O6<sup>-</sup>, we can wrap four D6-branes on four different copies of $`^3`$, while with an O6<sup>+</sup>, we can wrap two anti-D6-branes on $`S^2\times `$. We have $`_{}^{(3)}๐ฎ^4()`$, and $`_+^{(3)}๐ฎ^2()`$, where, importantly, $`S^3/_2`$ is the same quotient as before.
It is worthwhile to point out that the parameters associated with the brane wrappings are not on the same footing as the parameter associated with the size of the $`S^3`$. The latter parameter, although not normalizable in the non-compact geometry, is still localized and will survive embedding in a compact model. The data parameterizing the positions of the branes, on the other hand, is completely fixed at infinity. Not even their complexification can be determined before embedding in a compact model.
As a consequence, in addition to fixing the sign of the O-plane, we should fix a point $`(M_1,M_2)_{}^{(1)}`$ when we attempt to take case (1) through a $`\mu `$ transition to case (3). In so doing, we can see from the matching of the D-brane configuration at infinity that we end up with a configuration in $`_{}^{(3)}`$ of the special form $`(M_1,M_1,M_2,M_2)`$. Conversely, starting from $`(M_1,M_2,M_3,M_4)_+^{(1)}`$, we can match with a point in $`_{}^{(3)}`$ only for $`M_1=M_3`$, $`M_2=M_4`$.
These constraints allow us to predict a smooth $`\mu `$-transition only for this subset of D-brane configurations. We are not able to determine the fate of the other configurations as $`\mu `$ goes to zero.
Another point that remains unclear is whether there will be a point of enhanced gauge symmetry in the moduli space. As we have explained before, the orientifolds in the classes (1) and (3) do not admit the resolved conifold. Naively, we can say that the orientifold is projecting out the scalars of the $`๐ฉ=2`$ vectormultiplet that was associated with the $`^1`$ of the resolved conifold (whereas in cases (0), (2) and (4) we are projecting out the $`๐ฉ=1`$ vector). One possible conclusion is that the $`๐ฉ=1`$ vector half of the $`๐ฉ=2`$ vectormultiplet, which is broken by the hypermultiplet vev at a generic point, reappears at the singular conifold. However, it is not clear that such a naive argument will survive a more careful treatment.
### 7.2 Case (2)
Now the O-plane locus is $`S^1\times ^2`$ and has zero intersection with the compact three-cycle. We can then wrap supersymmetrically D6-branes on the compact three-cycle and get a dynamical gauge theory in four dimensions. This is very similar to the case (0)/(4) that we have focused on before. We can not find an M-theory lift, but since the small resolution of the conifold is allowed, we can apply the Vafa superpotential method to predict the structure of a part of the quantum parameter space.
Assume that we want to wrap $`2N`$ D6-branes on $`S_>^3`$. By the computation of the period integrals in Section 2.2, this will preserve the same supersymmetry as an anti-O6-plane wrapped on the noncompact cycle $`S^1\times ^2`$. If this O-plane is an anti-O6<sup>-</sup> (with positive charge), the gauge group is $`Sp(N)`$. This gauge group type follows because the relative codimension between O-plane and D-brane is $`4`$.
A configuration with the same supersymmetry and charge at infinity is obtained by wrapping $`2N+4`$ branes, or rather $`2\stackrel{~}{N}=2N4`$ anti-D6-branes, on $`S_<^3`$. Here, the $`4`$ comes again from the jump (2.14) in the class of the fixed point locus, multiplied by the charge of the anti-O6<sup>-</sup>-plane. Note again that this is supersymmetric and gives the gauge group $`Sp(\stackrel{~}{N})`$ in front of the anti-O6<sup>-</sup>-plane because of the lower-dimensional intersection.
We see that $`N`$ and $`\stackrel{~}{N}`$ are never both non-negative at the same time, so that we do not expect a $`\mu `$-transition to be possible for any $`N`$. In fact, when $`N=1`$, also $`\stackrel{~}{N}=1`$, so this value of the flux does not admit a deformed conifold for the fixed supersymmetry.
Meanwhile, if the O-plane is an anti-O6<sup>+</sup>, and we wrap $`2N`$ D6-branes on $`S_>^3`$, the gauge group will be $`SO(2N)`$. Going through the $`\mu `$-transition, we could also wrap $`2N4`$ D6-branes, or rather $`2\stackrel{~}{N}=2N+4`$ anti-D6-branes on $`S_<^3`$, with gauge group $`SO(2\stackrel{~}{N})`$. The jump by $`4`$ is the value familiar from the case (0)/(4), and as in that case, we expect a $`\mu `$-transition to be possible for $`N=0,1,2`$, but no other values of $`N`$.
Let us now figure out the moduli space using the superpotential.
### The moduli space
To start with, we discuss the right parameter of the moduli space near the resolved conifold points. Let $`t`$ be as before the complexified Kรคhler class of the $`^1`$ of the resolved conifold. Unlike in the cases (0) and (4) considered in Section 5.2 (see (5.22)), the periodicity in the present case is the same as the one before the orientifold:
$$tt+2\pi i,$$
(7.1)
so that the single valued coordinate of the parameter space is $`e^t`$. This is because the worldsheet diagram must always have even powers of $`\mathrm{exp}(i\mathrm{Im}(t)/2)`$. Namely, there is no odd degree smooth map of the worldsheet $`S^2`$ to the target $`^1`$ compatible with the involution $`\mathrm{\Omega }:w1/\overline{w}`$ on the domain and the involution $`\tau `$ given by (2.18) on the target. This can be shown as follows. Let $`X:S^2^1`$ be such a map, where the compatibility means $`\tau X\mathrm{\Omega }=X`$. Choosing a Kรคhler form $`\omega `$ of $`^1`$ of volume $`1`$, the degree is defined as $`d=_{S^2}X^{}\omega `$. Let us decompose $`S^2`$ as a union of the upper and lower hemi-spheres $`S^2=H_+H_{}`$ which are oriented such that $`[S^2]=H_++H_{}`$ and $`\mathrm{\Omega }(H_+)=H_{}`$. Using $`\tau ^{}\omega =\omega `$, one can express the degree as
$`d`$ $`=`$ $`{\displaystyle _{H_+}}X^{}\omega +{\displaystyle _H_{}}X^{}\omega ={\displaystyle _{H_+}}X^{}\omega {\displaystyle _{H_+}}\mathrm{\Omega }^{}X^{}\omega `$
$`=`$ $`{\displaystyle _{H_+}}(X^{}\omega +\mathrm{\Omega }^{}X^{}\tau ^{}\omega )=2{\displaystyle _{H_+}}X^{}\omega .`$
The idea is to show that $`_{H_+}X^{}\omega `$ is an integer. Note that $`d`$ is an integer and thus $`_{H_+}X^{}\omega `$ is deformation invariant, as long as $`X:H_+^1`$ extends to an equivariant map of $`S^2`$ to $`^1`$. Extension to $`S^2`$ is possible if and only if the restriction to the boundary $`X:H_+^1`$ is equivariant. It is always possible to shrink this loop to a constant map to a $`\tau `$-fixed point, keeping the equivariance all the way. Once this is done, we obtain a map $`X:H_+/H_+^1`$, which is a map between two spheres. $`_{H_+}X^{}\omega `$ is its degree and thus is an integer.
We first consider the case of an anti-O6<sup>-</sup> and $`2N0`$ D6-branes on $`S_>^3`$ or $`2\stackrel{~}{N}=2N40`$ anti-D6-branes on $`S_<^3`$. The computation is done in the orientifold of the resolved conifold with certain RR two-form flux through $`^2`$. The flux is $`N`$ since the O-plane does not contribute to the two-form flux through $`^2`$ (the O-plane is still there on the resolved side). The superpotential again has three terms corresponding to the three origins: four-form flux, two-form flux and worldsheet instantons. The result is
$`W`$ $`=`$ $`{\displaystyle F_4}\omega +N{\displaystyle \frac{F_0}{t}}4{\displaystyle \underset{m:\mathrm{even}>0}{}}{\displaystyle \frac{e^{mt/2}}{m^2}}`$ (7.2)
$`=`$ $`Yt/2NLi_2(e^t)2(Li_2(e^{t/2})+Li_2(e^{t/2}))`$
$`=`$ $`Yt/2(N+1)Li_2(e^t).`$
The parameter $`Y`$ is again defined by
$$Y:=_{S_{\mathrm{}}^3}(\mathrm{Re}\widehat{\mathrm{\Omega }}+iC_3).$$
The sum over $`m`$ in the crosscap part is over even integers because only even degree maps are compatible with the present orientifold, as remarked above. Solving $`_tW=0`$, we find the equation determining the moduli space
$$e^Y=(1e^t)^{2N+2}.$$
(7.3)
It is a copy of the Riemann sphere with three marked points โ $`(e^Y,e^t)=(1,0)`$, $`(\mathrm{},\mathrm{})`$, $`(0,1)`$ for $`N0`$, while $`(e^Y,e^t)=(1,0)`$, $`(0,\mathrm{})`$, $`(\mathrm{},1)`$ for $`N2`$. In either case, the superpotential has a branch cut at the last marked point, $`e^t=1`$, at which it behaves as follows:
$$W=Yt/2(N+1)t\mathrm{log}t+\mathrm{}.$$
(7.4)
If $`N0`$, the point $`(e^Y,e^t)=(0,1)`$ corresponds to a large minimal three sphere $`S_>^3`$ with $`2N`$ D6-branes supporting an $`Sp(N)`$ gauge field. The relation of the parameter $`Y`$ and the holomorphic gauge coupling $`(\frac{8\pi ^2}{g^2}i\theta )`$ on the D6-brane can be determined following the argument given in Section 5.2. Notice that it is similar to the case (4) in that the orientifold acts non-trivially on the three sphere. Noting that we expect symplectic gauge group, we find
$$Y=2\left(\frac{8\pi ^2}{g^2}i\theta \right).$$
(7.5)
Then, with the identification of $`t`$ as the glueball field $`S`$, (7.4) is indeed the Veneziano-Yankielowicz superpotential for the gauge group $`Sp(N)`$, up to a power series in $`S`$.
If $`N2`$, at the point $`(e^Y,e^t)=(\mathrm{},1)`$ we have a large $`S_<^3`$ with $`2\stackrel{~}{N}`$ anti-D6-branes supporting an $`Sp(\stackrel{~}{N})`$ gauge field. The relation of the parameter $`Y`$ and the holomorphic gauge coupling is the same as (7.5) up to sign. With the identification $`S=t`$, (7.4) agrees with the Veneziano-Yankielowicz superpotential for the $`Sp(\stackrel{~}{N})`$ super Yang-Mills, since $`(N+1)=\stackrel{~}{N}+1`$ is the dual Coxeter number of $`Sp(\stackrel{~}{N})`$.
One may also consider the case $`N=1`$. In the resolved side, the superpotential from the flux is exactly canceled by the contribution form the crosscap instantons. The superpotential is simply $`W=Yt/2`$ and $`_tW=0`$ requires $`Y=0`$. Indeed, we do not have any candidate classical limit with deformed conifold for this value of the charge, as we have seen. What is most interesting is that the resolved branch has no singularity in the interior. By the combined effect of the flux and worldsheet instantons, the singularity, which was present in $`๐ฉ=2`$ systems, is completely washed out!
Let us next consider the case of an anti-O6<sup>+</sup>-plane and $`2N>4`$ D6-branes on $`S_>^3`$ or $`2\stackrel{~}{N}=2N+4>4`$ anti-D6-branes on $`S_<^3`$. In this case, the crosscap contribution in (7.2) changes sign, and the equation for the moduli space is
$$e^Y=(1e^t)^{2N2}.$$
(7.6)
It is again a complex plane with three marked points. For $`N>2`$, the holomorphic gauge coupling is related to the parameter $`Y`$ by
$$Y=\frac{8\pi ^2}{g^2}i\theta ,$$
(7.7)
and we have the expected behavior of the superpotential for $`S=t/2`$ near the point $`(e^Y,e^t)=(0,1)`$, as $`2(N1)`$ is the dual Coxeter number of the gauge group $`SO(2N)`$. For $`N<0`$, the relation of $`Y`$ and the gauge coupling is opposite, but again the superpotential behaves as it should.
All these cases having been checked quite nicely, we briefly comment on the exceptional cases $`N=0,1,2`$. Here we expect a $`\mu `$-transition between vacua with $`SO(2N)`$ and $`SO(42N)`$ super Yang-Mills theories to be possible, but this will occur on a branch of moduli space that is not described by the exact superpotential. Though we do not have very powerful tools of analysis for the new branch, we can at least make a guess and check the consistency by examining the order of poles of holomorphic parameters at classical points. The expected branch structure is summarized in Figure 15.
For the case $`N=2`$ (related to the case $`N=0`$ by sign flip of $`Y`$), the superpotential accounts for only two vacua of the super Yang-Mills theory on large $`S_>^3`$. Other vacua should sit on a โnewโ branch which we expect to contain also the deformed vacuum with negative $`\mu `$. The holomorphic parameter $`e^Y`$ has double zero at $`\mu +\mathrm{}`$, whereas the complex volume of minimal three-cycle at $`\mu \mathrm{}`$ is minus $`Y/2`$ so that $`e^{Y/2}`$ has a simple pole there. The new branch will therefore be a cylinder parametrized by $`e^{Y/2}`$.
For $`N=1`$, the flux and crosscap terms in the superpotential cancel out in the same way as the $`Sp`$ case with $`N=1`$. A branch of the moduli space is therefore a cylinder interpolating two resolved classical points. We expect another branch of vacua with low energy $`U(1)`$ gauge dynamics that contains deformed classical points with either sign of $`\mu `$. It would be interesting to study how the two branches are connected.
## 8 Conclusions
We have discussed supersymmetric quantum transitions between various orientifolds of the conifold. We have constructed the possible orientifolds of the deformed and resolved conifold in Type IIA string theory. In the primary case where this is possible, we have answered our basic question by considering the M-theory lift of the various IIA orientifold configurations. We identified the corresponding $`G_2`$ holonomy manifolds, and studied the quantum moduli space connecting different configurations through their topology and also the IIA exact superpotential.
Our main results are valid for the orientifold $`z_i\overline{z}_i`$, of deformed conifold $`_iz_i^2=\mu `$. With $`\mu `$ real, this has two phases. Depending on whether $`\mu `$ is positive or negative, the orientifold fixes the $`S^3`$, or acts freely so that the minimal cycle is an $`^3`$. The transition between positive and negative $`\mu `$ is possible only for special values of RR charge. Once we fix the supersymmetry, one may either consider the O$`6^\pm `$-plane ($`\mu >0`$) with charge $`\pm 2`$ and increase the charge by adding D6-branes, or start with free orientifolds ($`\mu <0`$) with zero charge and decrease the charge by adding anti-D6-branes. Therefore, $`\mu `$ can flip the sign only when the total charge $`(N2)`$ equals $`0,1`$ or $`2`$.
On the other hand, some deformed or resolved geometries uplift to $`D_N`$-type orbifolds of the $`G_2`$ holonomy spaces $`๐น_7`$ or $`๐ป_7`$, both of which are topologically $`^4\times S^3`$. Remarkably, we found that this is true also for negatively large RR charges ($`N23`$), though the action of the dihedral group turned out to be non-standard. A careful analysis of the behavior of membrane instanton factors allows us to determine the structure of quantum moduli space unambiguously. Also, the exact IIA superpotential tells us how the resolved vacua are connected smoothly to other vacua.
We concluded that when $`N=0,1,2`$ the moduli space consists of two branches. For $`N=0`$ and $`N=2`$, the two branches meet at infinity where there is a weakly coupled $`SO(4)`$ super Yang-Mills theory, and each branch contains two of the four vacua. For $`N=1`$, the branch of mass-gapped vacua meets the branch of vacua with infrared $`U(1)`$ at a phase transition point, and we found a precise description of the transition via a mirror Type IIB picture.
The quantum moduli space of IIA orientifolded conifolds thus depends on the RR charge in an interesting way. We summarize these main results of our paper in Figure 16. We found similar results in other cases of orientifolded conifolds as discussed in Section 7.
Acknowledgments We would like to thank B. Acharya, J. Gomis, S. Hellerman, M. Kleban, J. Maldacena and E. Witten for very useful discussions. K.H., K.H. and D.P. thank Research Institute for Mathematical Sciences and Yukawa Institute for Theoretical Physics, Kyoto University, for the support and hospitality during their visits where a part of this work was done. R.R. and J.W. acknowledge the hospitality of the Fields and Perimeter Institutes where some of this work was realized. The work of K.H. and K.H. was supported by NSERC and Alfred P. Sloan Foundation. The work of D.P. is supported by PREA. The work of R.R. and J.W. was supported by the DOE under grant number DE-FG02-90ER40542.
Appendix
## Appendix A Conifold
Ricci-flat Kรคhler metrics for conifold or its small deformation or resolution were obtained in . The singular conifold is defined by $`_{i=1}^4z_i^2=0`$ or equivalently by $`\mathrm{det}W=0`$ with
$$W=\frac{1}{\sqrt{2}}\left(\begin{array}{cc}z_1+iz_2& z_3+iz_4\\ z_3+iz_4& z_1iz_2\end{array}\right).$$
(A.1)
Regarding $`z_i`$ as holomorphic coordinates and putting $`K=\rho ^{2/3}`$ with $`\rho \mathrm{Tr}WW^{}`$, one obtains a Ricci-flat Kรคhler metric which is symmetric under $`SU(2)_L\times SU(2)_R`$ acting on $`W`$ as $`WLWR^{}`$. To see the symmetry of the metric, we use the coordinates (which is slightly different from the one conventionally used)
$$W=X(\theta ,\varphi ,\psi )W_0\stackrel{~}{X}(\stackrel{~}{\theta },\stackrel{~}{\varphi },\stackrel{~}{\psi })^{},W_0=r^{3/2}\left(\begin{array}{cc}1& 0\\ 0& 0\end{array}\right)$$
(A.2)
or more explicitly
$$W=r^{3/2}\left(\begin{array}{cc}\hfill \mathrm{cos}\frac{\theta }{2}\mathrm{cos}\frac{\stackrel{~}{\theta }}{2}e^{\frac{i}{2}(\psi \stackrel{~}{\psi }+\varphi \stackrel{~}{\varphi })}& \hfill \mathrm{cos}\frac{\theta }{2}\mathrm{sin}\frac{\stackrel{~}{\theta }}{2}e^{\frac{i}{2}(\psi \stackrel{~}{\psi }+\varphi +\stackrel{~}{\varphi })}\\ \hfill \mathrm{sin}\frac{\theta }{2}\mathrm{cos}\frac{\stackrel{~}{\theta }}{2}e^{\frac{i}{2}(\psi \stackrel{~}{\psi }\varphi \stackrel{~}{\varphi })}& \hfill \mathrm{sin}\frac{\theta }{2}\mathrm{sin}\frac{\stackrel{~}{\theta }}{2}e^{\frac{i}{2}(\psi \stackrel{~}{\psi }\varphi +\stackrel{~}{\varphi })}\end{array}\right)$$
(A.3)
and get
$`ds^2`$ $`=`$ $`dr^2+r^2ds_{T^{1,1}}^2,`$
$`ds_{T^{1,1}}^2`$ $`=`$ $`{\displaystyle \frac{1}{6}}(\sigma _1^2+\sigma _2^2+\stackrel{~}{\sigma }_1^2+\stackrel{~}{\sigma }_2^2)+{\displaystyle \frac{1}{9}}(\sigma _3\stackrel{~}{\sigma }_3)^2`$ (A.4)
$`=`$ $`{\displaystyle \frac{1}{6}}(d\theta ^2+\mathrm{sin}^2\theta d\varphi ^2+d\stackrel{~}{\theta }^2+\mathrm{sin}^2\stackrel{~}{\theta }d\stackrel{~}{\varphi }^2)`$
$`+{\displaystyle \frac{1}{9}}(d(\psi \stackrel{~}{\psi })+\mathrm{cos}\theta d\varphi \mathrm{cos}\stackrel{~}{\theta }d\stackrel{~}{\varphi })^2.`$
The coordinate $`\stackrel{ห}{\psi }\psi \stackrel{~}{\psi }`$ has the period $`4\pi `$ and defines an $`S^1`$ which is fibered over $`S^2\times S^2`$ labeled by $`(\theta ,\varphi )`$ and $`(\stackrel{~}{\theta },\stackrel{~}{\varphi })`$. This metric has an $`O(4)\times U(1)`$ symmetry: $`SO(4)(SU(2)_L\times SU(2)_R)/_2`$ acts on $`W`$ as explained and in particular $`z_i`$ are transformed as a 4-vector. A parity transform in $`O(4)`$ exchanges the two $`S^2`$โs, and the $`U(1)`$ shifts $`\psi `$ or acts on $`z_i`$ as phase rotation.
Small resolution is an $`๐ช(1)๐ช(1)`$ bundle over $`^1`$. To write down the Kรคhler metric, write the matrix $`W`$ with vanishing determinant as
$$W=\left(\begin{array}{cc}u\lambda & u\\ y\lambda & y\end{array}\right)=\left(\begin{array}{cc}x& x\mu \\ v& v\mu \end{array}\right).$$
(A.5)
$`(\lambda ,\mu )`$ with $`\lambda \mu =1`$ are regarded as coordinates on $`^1`$, and from the relation between $`(u,y)`$ and $`(x,v)`$ one finds they are coordinates on the fiber. A natural ansatz for the Kรคhler potential of the resolved conifold is
$$K=K(\rho )+4a^2\mathrm{ln}(1+|\lambda |^2),$$
(A.6)
the second term yielding a Fubini-Study metric on $`^1`$. From the Ricci-flatness one finds that $`r^2\rho \frac{dK}{d\rho }`$ has to satisfy
$$r^4(r^2+6a^2)=c\rho ^2+c^{}$$
(A.7)
for some constants $`c,c^{}`$. Setting $`c=1,c^{}=0`$ one obtains
$$ds^2k^1(r)dr^2+\frac{r^2}{6}(\sigma _1^2+\sigma _2^2)+\frac{r^2+4a^2}{6}(\stackrel{~}{\sigma }_1^2+\stackrel{~}{\sigma }_2^2)+\frac{r^2k(r)}{9}(\sigma _3\stackrel{~}{\sigma }_3)^2,$$
(A.8)
with $`k(r)=\frac{r^2+6a^2}{r^2+4a^2}`$. The metric is invariant under $`SU(2)\times SU(2)\times U(1)`$, but the $`_2`$ symmetry of singular conifold is lost.
Small deformation is defined by $`_iz_i^2=ฯต^2`$ or $`\mathrm{det}W=\frac{ฯต^2}{2}`$. Hereafter we assume $`ฯต`$ to be real positive, as the Ricci-flat metric will depend only on the modulus $`|ฯต|^2`$. The $`SU(2)_L\times SU(2)_R`$ invariant metric can be found by assuming the Kรคhler potential to be a function $`K(\rho )`$ of $`\rho =\mathrm{Tr}WW^{}`$. The Ricci-flatness can be solved easily by introducing $`\rho =ฯต^2\mathrm{cosh}\tau `$ and putting
$$W_0=\frac{ฯต}{\sqrt{2}}\left(\begin{array}{cc}e^{\tau /2}& 0\\ 0& e^{\tau /2}\end{array}\right).$$
(A.9)
The Ricci-flatness amounts to $`\frac{dK}{d\rho }=\frac{(\mathrm{sinh}2\tau 2\tau )^{1/3}}{\mathrm{sinh}\tau }k(\tau )`$, and one obtains the following metric
$$ds^2k(\tau )\left\{\frac{4}{3k(\tau )^3}(d\tau ^2+(\sigma _3\stackrel{~}{\sigma }_3)^2)+\mathrm{cosh}\tau (\sigma _1^2+\stackrel{~}{\sigma }_1^2+\sigma _2^2+\stackrel{~}{\sigma }_2^2)2(\sigma _1\stackrel{~}{\sigma }_1+\sigma _2\stackrel{~}{\sigma }_2)\right\}.$$
(A.10)
The metric is invariant under $`O(4)`$ but not under $`U(1)`$.
## Appendix B Properties of $`Li_2(z)`$
We can define the Euler dilogarithm function<sup>6</sup><sup>6</sup>6Defined by Euler in 1768. in the disk $`|z|<1`$ as a convergent power series <sup>7</sup><sup>7</sup>7For more properties of this function see :
$$Li_2(z)=\underset{n=1}{\overset{\mathrm{}}{}}\frac{z^n}{n^2}.$$
(B.1)
The function can be extended to the whole complex plane as a multi valued analytical function. There is a branch cut from $`z=1`$ to $`z=\mathrm{}`$. Alternatively it can be defined as an integral:
$$Li_2(z)=_0^z\frac{\mathrm{log}(1z)}{z}๐z.$$
(B.2)
A related function often found in literature is the Rogers L-function:
$$L(z)=Li_2(z)+\frac{1}{2}\mathrm{log}z\mathrm{log}(1z).$$
(B.3)
Some functional equations that we will use in the main text are:
* The Euler identity provides an expansion around the branch point $`z=1`$:
$$Li_2(1z)=Li_2(z)+\frac{\pi ^2}{6}\mathrm{log}z\mathrm{log}(1z)$$
(B.4)
Or in terms of the Rogers L-function:
$$L(z)+L(1z)=L(1)$$
(B.5)
* The expansion around $`z=\mathrm{}`$
$$Li_2(1/z)=Li_2(z)\frac{\pi ^2}{6}+\frac{1}{2}(\mathrm{log}(z))^2$$
(B.6)
* A simple relation between the value of $`Li_2(z)`$ and $`Li_2(z^2)`$
$$Li_2(z)+Li_2(z)=\frac{1}{2}Li_2(z^2)$$
(B.7)
* The above relations can be obtained in terms of the Abel identity
$$\begin{array}{c}Li_2(x)+Li_2(y)=Li_2(xy)+Li_2\left(\frac{x(1y)}{1xy}\right)+\hfill \\ \hfill Li_2\left(\frac{y(1x)}{1xy}\right)+\mathrm{log}\left(\frac{1x}{1xy}\right)\mathrm{log}\left(\frac{1y}{1xy}\right)\end{array}$$
(B.8)
Or in terms of the Rogers L-function:
$$L(x)+L(y)=L(xy)+L\left(\frac{x(1y)}{1xy}\right)+L\left(\frac{y(1x)}{1xy}\right)$$
(B.9)
And some particular values of the dilogarithm:
$$Li_2(1)=\frac{\pi ^2}{6},Li_2(1)=\frac{\pi ^2}{12},Li_2(1/2)=\frac{\pi ^2}{12}\frac{1}{2}(\mathrm{log}(2))^2,$$
(B.10)
$$L(1)=\frac{\pi ^2}{6},L(0)=0.$$
(B.11)
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# Detection and characterization of multipartite entanglement in optical lattices
## I Introduction
Multipartite entanglement is an essential ingredient for complex quantum information processing (QIP) tasks, such as quantum error-correction QEC , multi-party quantum communication GHZ and universal quantum computation in the one-way model Cluster . It has been generated between photons 4photon and in a controlled way both in ion traps IonTraps and in experiments with Mott insulating states of neutral bosonic atoms in optical lattices BlochEnt ; BECinOL . For photon and ion trap experiments full state tomography 4photon ; MLEtomography impressively proved the existence of multipartite entanglement with few particles. However, its unambiguous detection in optical lattice setups has so far not been possible. The measurements implemented in the actual experiments only provided information on the average density operator of each individual atom (representing a qubit), but not on the density operator of the composite system BlochEnt . Full state tomography in optical lattices seems to be infeasible due to the large number of involved atoms and thus other methods are required to prove the existence of entanglement. Recently, a quantum network that can detect multipartite entanglement in bosons, and can be realized in optical lattices, was proposed Carolina2004 .
This quantum network is made up of pairwise beam-splitters (BS) acting on two identical copies of a multipartite entangled state $`\rho _{12\mathrm{}n}`$ of $`n`$ atoms, and allows the determination of the purity of $`\rho _{12\mathrm{}n}`$ and of all its reduced subsystems. The purities provide us with a separability criterion based on entropic inequalities Carolina2004 that, if violated, indicate that the state $`\rho _{12\mathrm{}n}`$ is entangled. This nonlinear entanglement test can be implemented experimentally with a number of copies of $`\rho _{12\mathrm{}n}`$ independent of $`n`$, rather than the exponentially large number of copies needed for a full state tomography, and is more powerful than the other common experimental methods, such as Bell inequalities Bell64 and entanglement witnesses Witness . In fact, the nonlinear inequalities are known to be strictly stronger than the Bell-CHSH inequalities Entropies . Furthermore, if it may be assumed that the entangled state is characterized by a few parameters only it is often possible to determine them from the results obtained from the entanglement detection network.
The main aim of this paper is to analyze the operation of the entanglement detection network under realistic experimental conditions. After introducing the ideal scheme we consider the particular realization of the network in an optical lattice. We discuss three main sources of imperfections: limited or no spatial resolution, errors in the BS operations occurring with a probability $`q`$, and failure to detect an atom with probability $`p`$. In current experiments none of these errors can be suppressed completely and thus it is important to find out whether their presence still allows the unambiguous detection and characterization of multipartite entangled states. Our main results will be that (i) even without spatial resolution it is possible to unambiguously detect and characterize entanglement in multipartite states like cluster states and macroscopic superposition states, (ii) the effects of BS and detection errors can be made arbitrarily small by increasing the number of runs of the entanglement detection network, and (iii) the required number of experimental runs for achieving this is reasonably small if $`kq1`$ and $`kp1`$, where $`kn`$ is the number of particles in the subsystem whose reduced purity is being measured. Finally we will find that (iv) single site resolution is necessary to achieve significant improvement compared to having no spatial resolution at all.
The paper is organized as follows: in Sec. II we introduce the entropic inequalities and discuss how multipartite entanglement can be detected and characterized by them. Then we describe the entanglement detection network and show how it can be realized in optical lattices and how the main sources of imperfections arise. In Sec. III we investigate the detection and characterization of multipartite entangled states if no spatial resolution is available, and in Sec. IV we analyze the effect of the dominant experimental errors. We also discuss the case of limited spatial resolution in the measurements. Finally we summarize our results in Sec. V.
## II Entanglement Detection and Characterization in Optical Lattices
In this section we first present the entropic inequalities which are used to detect multipartite entanglement. We show how they can be utilized to characterize multipartite entanglement if the entangled state is described by few parameters only. Then we introduce our entanglement detection network and discuss its implementation in optical lattices. Finally we study the main types of imperfections and resulting errors that we expect to be present in this experimental realization.
### II.1 Entropic inequalities
Our network uses the entropic inequalities introduced in Carolina2004 for multipartite entanglement detection. These inequalities provide a set of necessary conditions for separability in multipartite states. Consider a state $`\rho _{123\mathrm{}n}`$ of $`n`$ subsystems. If $`\rho _{123\mathrm{}n}`$ is separable, then
$`\rho _{123\mathrm{}n}={\displaystyle \underset{\mathrm{}}{}}C_{\mathrm{}}\rho _1^{\mathrm{}}\rho _2^{\mathrm{}}\rho _3^{\mathrm{}}\mathrm{}\rho _n^{\mathrm{}},`$ (1)
where $`\rho _j^{\mathrm{}}`$ is a state of subsystem $`j`$, and $`_{\mathrm{}}C_{\mathrm{}}=1`$. The purity $`\text{Tr}\left\{\rho _B^2\right\}`$ of the reduced density operator $`\rho _B`$ with $`B\{1,2,\mathrm{},n\}`$ is less than or equal to the purity of any of its reduced density operators:
$$\text{Tr}\left\{\rho _A^2\right\}\text{Tr}\left\{\rho _B^2\right\}\text{ for all }AB.$$
(2)
We can use these inequalities to give a relation between the average purities of all reduced density operators of a given number $`k`$ of subsystems defined by
$$\overline{\text{Tr}\left\{\rho _{\left(k\right)}^2\right\}}=\left[\left(\genfrac{}{}{0pt}{}{n}{k}\right)\right]^1\underset{|B|=k}{}\text{Tr}\left\{\rho _B^2\right\},$$
(3)
where $`B`$ is summed over all combinations of $`k`$ subsystems. For completeness we define $`\overline{\text{Tr}\left\{\rho _{\left(0\right)}^2\right\}}1`$. From Eq. (2) we find that for any separable state
$$\overline{\text{Tr}\left\{\rho _{\left(k\right)}^2\right\}}\overline{\text{Tr}\left\{\rho _{\left(k^{}\right)}^2\right\}}\text{ for all }kk^{}.$$
(4)
We will now consider how Eqs. (2, 4) can be used to detect entanglement and to characterize entangled states.
#### II.1.1 Entanglement detection
Any state $`\rho `$ that violates any of the inequalities Eqs. (2, 4) is entangled. If $`\rho _{123\mathrm{}n}`$ is separable and pure, $`\text{Tr}\left\{\rho _{123\mathrm{}n}^2\right\}=1`$ and all the above inequalities become equalities; since a state with all $`\text{Tr}\left\{\rho _B^2\right\}=1`$ is necessarily a product state, all pure entangled states violate Eqs. (2, 4). All pure entangled states can hence be detected by comparing $`\text{Tr}\left\{\rho _{123\mathrm{}n}^2\right\}`$ with any other $`\text{Tr}\left\{\rho _B^2\right\}`$ or $`\overline{\text{Tr}\left\{\rho _{\left(k\right)}^2\right\}}`$ unlike entanglement witnesses or Bell inequalities where different states often require different witnesses. From the Schmidt decomposition of any pure state into two disjunct subsystems $`A`$ and $`B`$ with $`AB=\{1,2\mathrm{}n\}`$ we find $`\text{Tr}\left\{\rho _A^2\right\}=\text{Tr}\left\{\rho _B^2\right\}`$ and hence $`\overline{\text{Tr}\left\{\rho _{\left(k\right)}^2\right\}}=\overline{\text{Tr}\left\{\rho _{\left(nk\right)}^2\right\}}`$ for all pure $`\rho `$.
Mixed entangled states do not always violate Eqs. (2, 4), but since $`\text{Tr}\left\{\rho _B^2\right\}`$ is continuous in $`\rho `$ a sufficiently small amount of noise added to a pure entangled state will leave it still violating the inequalities. In the examples studied in this paper the noise level at which the inequalities no longer detect entanglement is a large fraction of the level at which entanglement ceases to be present. The permissible range values for the purities of any state $`\rho `$ is
$$2^k\text{Tr}\left\{\rho _B^2\right\}1\text{or}2^k\overline{\text{Tr}\left\{\rho _{\left(k\right)}^2\right\}}1,$$
(5)
where $`k=|B|`$. The minimum is attained by the maximally mixed state $`\rho =2^n๐`$ (with $`๐`$ the identity operator) and the maximum by any pure product state.
Our test can often verify that a state is truly $`n`$-particle entangled and not just a collection of small entangled subsystems, because $`\text{Tr}\left\{\rho _B^2\right\}<\text{Tr}\left\{\rho ^2\right\}`$ always indicates that $`B`$ is entangled with the rest of the system and cannot be caused by entanglement within $`B`$ itself. Furthermore our test is insensitive to local unitary operations, which cannot alter the entanglement structure, since they do not change the purities. It is, however, in most cases sensitive to local measurements destroying entanglement as these generally do change the purities. If the tested state $`\rho `$ is of a known form and characterized by few parameters only the kind and degree of violation of Eqs. (2, 4) can be used to determine these parameters and thus characterize the state.
#### II.1.2 Characterization of entanglement: Macroscopic superposition states
Macroscopic superposition states $`|\gamma _n`$ defined by
$$|\gamma _n=\frac{|\mathrm{\hspace{0.17em}0}^n+(\gamma |\mathrm{\hspace{0.17em}0}+\sqrt{1|\gamma |^2}|\mathrm{\hspace{0.17em}1})^n}{\sqrt{2+\gamma ^n+\overline{\gamma }^n}},$$
(6)
with a single complex parameter $`\gamma `$ satisfying $`|\gamma |1`$ arise naturally in several systems such as BECs cat-bec and superconductors cat-squid . The quantity $`n(1|\gamma |^2)`$ has been suggested as a measure of the effective size of the state, as in some respects $`|\gamma _n`$ is equivalent to a maximally entangled state of $`n(1|\gamma |^2)`$ particles Cat . The purities of $`|\gamma _n`$ are given by
$`\text{Tr}\left\{\rho _B^2\right\}`$ $`=`$ $`{\displaystyle \frac{2+2\gamma ^k\overline{\gamma }^k+2\gamma ^{nk}\overline{\gamma }^{nk}}{(2+\gamma ^n+\overline{\gamma }^n)^2}}+`$ (7)
$`{\displaystyle \frac{4\gamma ^n+4\overline{\gamma }^n+\gamma ^{2n}+\overline{\gamma }^{2n}}{(2+\gamma ^n+\overline{\gamma }^n)^2}},`$
where $`k=|B|`$. These purities violate Eq. (2), showing that $`|\gamma _n`$ is entangled, for all $`|\gamma |<1`$ and measuring the purities allows determining $`\gamma `$. For $`\gamma =0`$ the state $`|\gamma _n`$ is a maximally entangled state which is pure ($`\text{Tr}\left\{\rho ^2\right\}=1`$) while all its subsets have purity $`\text{Tr}\left\{\rho _B^2\right\}=1/2`$. Up to local unitaries it is the only state with these purities, so can be unequivocally identified by our test.
We demonstrate the effect of noise on our method by studying individual dephasing of each qubit. Such dephasing noise maps a state according to
$$\rho \underset{A}{}\left(1\frac{d}{2}\right)^{n|A|}\left(\frac{d}{2}\right)^{|A|}\left(\underset{jA}{}\sigma _{z,j}\right)\rho \left(\underset{jA}{}\sigma _{z,j}\right),$$
(8)
where $`\sigma _{z,j}`$ is a phase flip applied to qubit $`j`$, $`A`$ is the set of qubits which undergo a phase flip and $`0d1`$ is the decoherence parameter. Applying the map Eq. (8) to $`|(\gamma =0)_n`$ we find for the resulting density operator
$`\rho `$ $`=`$ $`{\displaystyle \frac{1}{2}}\left(|\mathrm{\hspace{0.17em}0}^n0|^n+|\mathrm{\hspace{0.17em}1}^n1|^n\right)`$ (9)
$`+{\displaystyle \frac{(1d)^n}{2}}\left(|\mathrm{\hspace{0.17em}0}^n1|^n+|\mathrm{\hspace{0.17em}1}^n0|^n\right),`$
and thus $`\text{Tr}\left\{\rho ^2\right\}=(1+(1d)^{2n})/2`$ while $`\text{Tr}\left\{\rho _B^2\right\}=1/2`$ for all subsets $`B`$. Therefore entanglement is detected whenever it is present. This is not generally the case as can easily be seen by looking at a Werner state $`\rho =(1d)|(\gamma =0)_n(\gamma =0)_n|+2^nd๐`$ for which entanglement is detected if $`d<1(2^{n1}+1)^{1/2}`$. However, e.g. in the case $`n=2`$ this state is entangled iff $`d<2/3`$ Werner , while our test works only for $`d<11/\sqrt{3}`$.
#### II.1.3 Characterization of entanglement: Cluster states
Cluster and graph states are multipartite entangled states which form the basic building blocks of the one-way quantum computer Cluster . We consider cluster-like states $`|\varphi _n`$ defined by
$$|\varphi _n=\frac{1}{\sqrt{2^n}}\underset{j=1}{\overset{n}{}}(|\mathrm{\hspace{0.17em}0}_je^{i\varphi \sigma _{z,j+1}}+|\mathrm{\hspace{0.17em}1}_j)=\frac{1}{\sqrt{2^n}}\underset{x=0}{\overset{2^n1}{}}e^{i\varphi c(x)}|x,$$
(10)
where $`c(x)`$ is the number of occurrences of the sequence $`01`$ in the $`n`$-bit binary number $`x`$. These states have already been created in optical lattices BlochEnt and represent a cluster state for $`\varphi =\pi `$. Current methods for detecting them essentially perform a tomographic measurement of the average single particle density matrix, which goes from pure $`(|\mathrm{\hspace{0.17em}0}+|\mathrm{\hspace{0.17em}1})(0|+1|)/2`$ at $`\varphi =0`$ to maximally mixed $`๐/2`$ at $`\varphi =\pi `$ and back again at $`\varphi =2\pi `$. This method thus yields one measurement $`\overline{\text{Tr}\left\{\rho _{\left(1\right)}^2\right\}}`$ relating to entanglement and two measurements relating to local unitaries.
Our network permits the measurement of further correlations and it does not need the assumption that all atoms have the same single particle density matrix. It therefore allows a better characterization of $`|\varphi _n`$. For any $`\varphi `$ which is not an integer multiple of $`2\pi `$, $`|\varphi _n`$ is a pure state with no separable subsystems, and hence for any subset $`B`$ we have $`\text{Tr}\left\{\rho _B^2\right\}<1=\text{Tr}\left\{\rho ^2\right\}`$ as can be seen from Fig. 1. For creating the states $`|\varphi _n`$ each qubit only needs to interact with their two nearest neighbors except for the two extremal atoms 1 and $`n`$ which will interact with only their one neighbor. Because of this the reduced purity of a subset $`B`$ is determined by the boundary between it and the rest of the system. Subsets of different sizes but with the same boundary structure have the same purity. Several examples of these purities are shown in Fig. 1 as a function of $`\varphi `$. For example, all sets of two or more adjacent atoms, located anywhere in the row that do not include either extremal atom have the same purity $`(1+\mathrm{cos}^2(\varphi /2))^2/4`$ (dash dotted curve in Fig. 1) which is independent of $`n`$. The degree of violation of Eq. (2) varies smoothly with $`\varphi `$ and measuring the various different purities allows to determine $`\varphi `$ (up to its sign). We will now introduce a quantum network which detects violations of Eq. (2) and later, in Sec. III, show how violations of Eq. (4) can be detected even without achieving spatial resolution of the different subsystems.
### II.2 Multipartite Entanglement Detection Network
A family of quantum interferometric networks that directly estimates $`\text{Tr}\{\rho ^s\}`$, $`s=2,3,4,\mathrm{}`$ for any $`\rho `$ from $`s`$ copies of $`\rho `$ was introduced in Carolina2002-3 . These networks rely on the controlled-shift operation $`CV^{(s)}`$ between the different copies, where $`V^{(s)}|\varphi _1|\varphi _2\mathrm{}|\varphi _s=|\varphi _s|\varphi _1\mathrm{}|\varphi _{s1}`$ for any $`|\varphi _i`$ with $`i=1,2\mathrm{}s`$. For the particular case of $`s=2`$ the value of $`\text{Tr}\left\{\rho ^2\right\}`$ is directly related to the probability of projecting $`\rho \rho `$ into its symmetric and antisymmetric subspaces. The values of the $`2^n`$ different purities associated with $`\rho _{123\mathrm{}n}`$ can hence be determined from the expectation value of the symmetric and antisymmetric projectors, on each different pair of reduced states $`\rho _B=\text{Tr}_A(\rho _{123\mathrm{}n})`$, where $`AB=\{1,2,\mathrm{}n\}`$. These expectation values can be measured without resorting to the implementation of the three-qubit C-Swap gate $`CV^{(2)}`$ if two identically prepared 1D rows of $`n`$ qubits which are represented by bosonic particles are coupled via pairwise BS, as shown in Fig. 2 Carolina2004 .
The BS in the $`j`$-th column projects the symmetric (antisymmetric) part of the density operator $`\rho _j`$ onto doubly (singly) occupied sites Bellstateanalyzer (for details see Appendix A). Two qubits in column $`j`$ and in state $`\rho _j\rho _j`$ will thus end up in the same site (+) or in different sites (-) with probabilities
$$P_\pm ^{(j)}=\frac{1}{2}\text{Tr}\{(๐\pm V^{(2)})\rho _j\rho _j\}=\frac{1}{2}\pm \frac{1}{2}\text{Tr}\left\{\rho _j^2\right\}.$$
(11)
Here $`S_\pm =(๐\pm V^{(2)})/2`$ is the symmetric / antisymmetric projector. By distinguishing doubly occupied sites from singly occupied ones we can thus determine the purity of $`\rho _j`$.
Extending this two-qubit scenario to the general case of two copies of a state of $`n`$ qubits undergoing pairwise BS (see Fig. 2) we obtain Carolina2004
$$P_{\pm _1\pm _2\mathrm{}\pm _n}=\text{Tr}\{\underset{i=1}{\overset{n}{}}\frac{๐\pm _iV_i^{(2)}}{2}\rho _{12\mathrm{}n}\rho _{12\mathrm{}n}\}.$$
(12)
Inverting the linear equation Eq. (12) the impurity of any subset of atoms $`B`$ is given by twice the probability of having an odd number $`j_B`$ of antisymmetric projections in subset $`B`$
$$\text{Tr}\left\{\rho _B^2\right\}=P(j_B\text{ even})P(j_B\text{ odd})=12P(j_B\text{ odd}).$$
(13)
For example, for $`n=3`$:
$`\text{Tr}\left\{\rho _{123}^2\right\}`$ $`=`$ $`P_{+++}+P_++P_++P_+`$
$`P_{}P_{++}P_{++}P_{++},`$
$`\text{Tr}\left\{\rho _{12}^2\right\}`$ $`=`$ $`P_{+++}+P_{++}+P_++P_{}`$
$`P_{++}P_+P_{++}P_+,`$
$`\text{Tr}\left\{\rho _1^2\right\}`$ $`=`$ $`P_{+++}+P_{++}+P_{++}+P_+`$
$`P_{++}P_+P_+P_{}.`$
### II.3 Realization in Optical lattices
We consider one sheet in the $`xy`$ plane of an ultracold two-species bosonic gas confined in a 3D optical lattice sufficiently deep so that the system is in a Mott insulating state with exactly one atom per lattice site Jaksch98 . Two long lived internal states of each atom $`|a`$ and $`|b`$ represent a qubit, i.e. for a site with row index $`l`$ and column index $`j`$ we define two basis states $`(a_l^{(j)})^{}|\mathrm{vac}|\mathrm{\hspace{0.17em}0}_l^j`$ and $`(b_l^{(j)})^{}|\mathrm{vac}|\mathrm{\hspace{0.17em}1}_l^j`$. Here $`|\mathrm{vac}`$ is the vacuum state and $`a`$ ($`b`$) is the bosonic destruction operator for an atom in internal state $`|a`$ ($`|b`$). We assume that starting from this Mott state identical multipartite entangled states $`\rho _{123\mathrm{}n}`$ are created in each row while different rows remain uncorrelated. This can e.g. be achieved by state selective cold controlled collisions between atoms in neighboring columns BlochEnt ; lattice-review . The above entanglement detection network can be realized in this setup to study the entanglement properties of $`\rho _{123\mathrm{}n}`$. We will first discuss the implementation of the BS by coupling pairs of rows along the $`x`$-direction and then investigate methods to distinguish doubly from singly occupied sites. Note that the intrinsic parallelism of the 3D optical lattice will allow to run many copies of these networks in the lattice at once. By exploiting this parallelism one can obtain an estimate of the desired projection probabilities and test the violation of Eqs. (2, 4) in a single or few experimental runs only.
#### II.3.1 The pairwise Beam Splitters
The dynamics of the atoms in the optical lattice is governed by the Bose-Hubbard Hamiltonian (BHM) Jaksch98 ; lattice-review . We assume that any hopping along the $`y`$ and $`z`$ directions is suppressed by a sufficient depth of the lattice in these directions. The hopping in $`x`$-direction is controlled by dynamically varying the corresponding lattice depth $`V_{0x}`$. Since $`V_{0x}`$ is proportional to the laser intensity it can easily be changed in the experiment. The pairwise BS requires that rows be coupled pairwise (this can be achieved by using a superlattice of twice the period) and thus we only need to consider one such pair labelling the two rows by $`I`$ and $`II`$, respectively, as shown in Fig. 2.
The Hamiltonian describing the dynamics of the atoms in these two rows can be written as a sum $`H_{\mathrm{BHM}}=H_{\mathrm{BS}}+H_U`$ where $`H_{\mathrm{BS}}`$ is due to hopping along the $`x`$ direction and $`H_U`$ is due to the repulsion between two atoms occupying the same lattice site Jaksch98 ; lattice-review . The two contributions are given by
$`H_{\mathrm{BS}}`$ $`=`$ $`J{\displaystyle \underset{j=1}{\overset{n}{}}}(a_I^{(j)}a_{II}^{(j)}+b_I^{(j)}b_{II}^{(j)}+\mathrm{h}.\mathrm{c}.),`$ (17)
$`H_U`$ $`=`$ $`{\displaystyle \underset{l=I,II}{}}{\displaystyle \underset{j=1}{\overset{N}{}}}{\displaystyle \frac{U}{2}}a_l^{(j)}a_l^{(j)}a_l^{(j)}a_l^{(j)}+`$
$`{\displaystyle \frac{U}{2}}b_l^{(j)}b_l^{(j)}b_l^{(j)}b_l^{(j)}+Ub_l^{(j)}b_l^{(j)}a_l^{(j)}a_l^{(j)},`$
where $`J`$ is the hopping matrix element and $`U`$ is the onsite interaction energy. The parameters $`U`$ and $`J`$ depend on the parameters of the trapping lasers, and their ratio can be varied over a wide range, on time scales much smaller than the decoherence time of the system, by dynamically changing the depth $`V_0`$ of the optical lattice Jaksch98 ; lattice-review .
The BS dynamics is perfectly realized in the noninteracting limit $`U=0`$ by applying $`H_{\mathrm{BS}}`$ for a time $`t_{bs}=\pi /(4J)`$ (for details see Appendix A). However, in practice it is impossible either to control $`J`$ perfectly accurately or to completely turn off the interaction $`U`$, and these imperfections cause the symmetric component to have a non-zero probability $`q_{bs}`$ of failing to bunch which is given by
$$q_{bs}\frac{\pi ^2}{8}\left(\frac{\delta J}{J}\right)^2+\frac{1}{16}\left(\frac{U}{J}\right)^2.$$
(18)
Here $`\delta J`$ describe the fluctuations in $`J`$ (for details see Appendix A). If the fluctuations $`\delta J`$ occur from run to run rather than from site to site $`q_{bs}`$ can be interpreted as a statistical random variable. We will discuss how to correct this error in Sec. IV.
#### II.3.2 Measuring the lattice site occupation
Recently, a method that uses atom-atom interactions to distinguish between singly and doubly occupied sites was demonstrated experimentally Measure . However, a simplified alternative method where rapid same-site two-atom loss is induced via a Feshbach resonance and the remaining singly-occupied sites are detected suffices for our purpose. The detection of the remaining singly-occupied sites is achieved by measuring the atomic density profile after the atoms are released from the lattice. A single Feshbach resonance will cause the loss of either $`|aa`$, $`|ab+|ba`$, or $`|bb`$. Hence, in order to empty doubly-occupied sites in all three states we can either turn on consecutively three separate Feshbach resonances or change the internal state of the pairs of atoms during the loss process using an appropriate sequence of laser pulses. Even if three resonances are experimentally accessible the latter might yield better results as the resonance with the best ratio of two-atom to single-atom loss can be exploited. One suitable sequence Foot1 which does not require precise control is to apply a large number of pulses of random relative phase and approximate area $`\pi /2`$ each at equal intervals. Each initial state then spends $`1/3`$ of the time in the resonant state, and hence has probability $`q_l=\mathrm{exp}(t_l/3\tau _d)`$ of failing to lose a pair of atoms occupying the same site where $`t_l`$ is the total duration of the sequence and $`\tau _d`$ the two-atom loss time constant. This method will not be perfect since single particles are also lost from the system with some time constant $`\tau _s`$ and hence $`t_l`$ cannot be chosen arbitrarily large. The probability of losing a single particle is $`p_l=1\mathrm{exp}(t_l/\tau _s)`$ where $`\tau _s\tau _d`$. Both $`p_l`$ and $`q_l`$ are error probabilities, and their sum $`p_l+q_l`$ is minimized by choosing $`t_l=\mathrm{ln}(\tau _s/(3\tau _d))/(1/3\tau _d1/\tau _s)`$. An experiment recently performed by Widera et. al Measure measured $`\tau _d=1.3\mathrm{ms}`$ and no detectable loss of single atoms for resonance times up to $`100\mathrm{m}\mathrm{s}`$. If we take $`\tau _s=500\mathrm{m}\mathrm{s}`$ the above gives $`t_l=19\mathrm{m}\mathrm{s}`$ and error probabilities of $`q_l=0.8\%`$ and $`p_l=3.7\%`$.
In summary, the optical lattice realization of our entanglement detection network contains the following stages at which experimental errors are likely to occur: (i) $`q_{bs}`$ from the implementation of $`H_{BS}`$; (ii) $`p_l`$ and $`q_l`$ occurring during the loss stage; (iii) detector errors $`p_d`$ in counting the number of singly occupied lattice sites. In addition (iv) the setup might lack spatial resolution in atom counting. We consider how to correct errors (i) - (iii) in Sec. IV and next study the case (iv) of no spatial resolution in an otherwise perfect experimental setup.
## III Entanglement detection and characterization without spatial resolution
We assume that the measurement of the total number of singly/doubly occupied sites is accurate but that we cannot know their locations. We first show that this information is sufficient to detect a violation of Eq. (4). Then we study how various different experimentally realizable multipartite entangled states might be characterized using such measurements.
### III.1 Entanglement detection
The probabilities $`P(j)`$ of measuring $`j`$ singly occupied sites in one row ($`2j`$ single atoms in total) are given by
$`P(j)`$ $`=`$ $`{\displaystyle \underset{|A|=j}{}}{\displaystyle \frac{1}{2^n}}{\displaystyle \underset{B}{}}(1)^{|AB|}\text{Tr}\left\{\rho _B^2\right\}`$
$`=`$ $`{\displaystyle \frac{1}{2^n}}{\displaystyle \underset{k=0}{\overset{n}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{n}{k}}\right)\overline{\text{Tr}\left\{\rho _{\left(k\right)}^2\right\}}{\displaystyle \underset{l}{}}\left({\displaystyle \genfrac{}{}{0pt}{}{k}{l}}\right)\left({\displaystyle \genfrac{}{}{0pt}{}{nk}{jl}}\right)(1)^l,`$
where the summation indices are $`k=|B|`$, $`l=|AB|`$, $`j=|A|`$, and $`A`$ is the set of singly occupied (antisymmetric) sites. We form the generating function
$`{\displaystyle \underset{j=0}{\overset{n}{}}}x^jP(j)`$ $`=`$ $`{\displaystyle \frac{1}{2^n}}{\displaystyle \underset{k=0}{\overset{n}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{n}{k}}\right)\overline{\text{Tr}\left\{\rho _{\left(k\right)}^2\right\}}(1x)^k(1+x)^{nk}`$
$`=`$ $`\left({\displaystyle \frac{1+x}{2}}\right)^n{\displaystyle \underset{k=0}{\overset{n}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{n}{k}}\right)\overline{\text{Tr}\left\{\rho _{\left(k\right)}^2\right\}}\left({\displaystyle \frac{1x}{1+x}}\right)^k`$
and let $`y=(1x)/(1+x)`$, to obtain
$$\underset{j=0}{\overset{n}{}}(1y)^j(1+y)^{nj}P(j)=\underset{k=0}{\overset{n}{}}y^k\left(\genfrac{}{}{0pt}{}{n}{k}\right)\overline{\text{Tr}\left\{\rho _{\left(k\right)}^2\right\}},$$
(21)
from which we find
$$\overline{\text{Tr}\left\{\rho _{\left(k\right)}^2\right\}}=\left[\left(\genfrac{}{}{0pt}{}{n}{k}\right)\right]^1\underset{j=0}{\overset{n}{}}P(j)\underset{l}{}\left(\genfrac{}{}{0pt}{}{j}{l}\right)\left(\genfrac{}{}{0pt}{}{nj}{kl}\right)(1)^l.$$
(22)
Therefore, although we cannot determine the purity of a given subset of the row of atoms we can still determine average purities associated with subsets of atoms of a given size by measuring $`P(j)`$. We will prove later (see Eq. (30)) that the accuracy in finding $`P(j)`$ required for obtaining a given accuracy of $`\overline{\text{Tr}\left\{\rho _{\left(k\right)}^2\right\}}`$ has an upper bound independent of $`n`$ and $`k`$ if no errors are present. The network is thus efficient in detecting the presence of entanglement in all pure (and some mixed) entangled states via the violation of Eq. (4).
In Fig. 3 we show the probabilities $`P(j)`$ and the resulting average purities for a variety of different states. For a classically correlated state shown in Fig. 3a the values of $`\overline{\text{Tr}\left\{\rho _{\left(k\right)}^2\right\}}`$ are monotonically decreasing with $`k`$ showing that Eq. (4) is not violated. The maximally entangled state shown in Fig. 3b has the characteristic that $`\overline{\text{Tr}\left\{\rho _{\left(k\right)}^2\right\}}=1/2`$ for $`0<k<n`$ while $`\overline{\text{Tr}\left\{\rho _{\left(n\right)}^2\right\}}=1`$ and thus the inequalities are violated in this case. The cluster state shown in Fig. 3c violates the inequalities for all $`j>n/2`$ and therefore its entanglement is detected. Finally, in Fig. 3d we show a noisy cluster state which was affected by phase noise acting independently on each atom. It can clearly be seen that decoherence reduces the violation of the inequalities but that entanglement is detectable for small amounts of noise.
### III.2 Characterization of entanglement
The $`n`$ measurable quantities $`\overline{\text{Tr}\left\{\rho _{\left(k\right)}^2\right\}}`$, $`k=1,\mathrm{},n`$ do not provide us with enough information to determine an arbitrary state $`\rho `$. However, if it may be assumed that the state $`\rho `$ is of a known form with less than $`n`$ unknown parameters, it is often possible to determine these parameters from $`\overline{\text{Tr}\left\{\rho _{\left(k\right)}^2\right\}}`$. We demonstrate this by considering macroscopic superposition states $`|\gamma _n`$ and cluster-like states $`|\varphi _j`$ introduced in Sec. II.1.3. Finally, we will also look at product states of states of subsystems containing several atoms.
#### III.2.1 Macroscopic superposition states
Because the state $`|\gamma _n`$ is totally symmetric the individual purities given in Eq. (7) only depend on the size of the subsystem $`k=|B|`$ and thus $`\overline{\text{Tr}\left\{\rho _{\left(k\right)}^2\right\}}=\text{Tr}\left\{\rho _B^2\right\}`$. Hence, from the knowledge of $`\overline{\text{Tr}\left\{\rho _{\left(k\right)}^2\right\}}`$ we can determine the value of $`\gamma `$ which in principle only requires two of the purities. The remaining equations allow a partial check of the assumption that the measured state indeed has the form $`|\gamma _n`$.
#### III.2.2 Cluster states
The states $`|\varphi _n`$ are parameterized by the entangling phase $`\varphi `$. The average purities as a function of $`\varphi `$ are depicted in Fig. 4. For any value of $`0<\varphi <2\pi `$ the states violate the inequalities Eq. (4). The degree of violation increases with $`\varphi `$ until $`\varphi =\pi `$ where the state is a cluster state and the degree of violation is a maximum. If $`\varphi `$ is increased further the state again approaches a product state and the degree of violation of the inequalities correspondingly decreases. Hence, from experimentally measured $`\overline{\text{Tr}\left\{\rho _{\left(k\right)}^2\right\}}`$ one can determine the phase $`\varphi `$ up to its sign. Again the over determined system ($`n`$ equations for one unknown $`\varphi `$) provides a check on how well the state fits the assumed form $`|\varphi _n`$.
The effect of dephasing according to the map Eq. (8) on a cluster state is shown in Fig. 5. The average purities decrease with increasing noise level $`d`$. Entanglement is certainly present and in principle detectable by our method as long as the $`\overline{\text{Tr}\left\{\rho _{\left(k\right)}^2\right\}}`$ are not in descending order, in the case shown in Fig. 5 up to $`d0.45`$. Again, the parameter $`d`$ can in principle be determined from measuring the average purities.
#### III.2.3 Products of entangled subsystem states
Finally we give an example of a class of states where even though $`\rho `$ is characterized by more than $`n`$ parameters, the associated average purities $`\overline{\text{Tr}\left\{\rho _{\left(k\right)}^2\right\}}`$ only depend on $`n`$ parameters. Consider the case where $`\rho `$ is a product state of $`L`$ subsystems, $`\rho =_{i=1}^L\rho _i`$, with each subsystem $`\rho _i`$ composed of a known number $`n_i`$ of atoms. In this case we have
$$\overline{\text{Tr}\left\{\rho _{\left(k\right)}^2\right\}}=\underset{\stackrel{\{k_1,\mathrm{},k_j\}}{_ik_i=k}}{}\underset{i}{}\left(\genfrac{}{}{0pt}{}{n_i}{k_i}\right)\overline{\text{Tr}\left\{\rho _{i\left(k_i\right)}^2\right\}}.$$
(23)
Since there are $`_in_i=n`$ different $`\overline{\text{Tr}\left\{\rho _{i\left(k_i\right)}^2\right\}}`$ in total, $`\overline{\text{Tr}\left\{\rho _{\left(k\right)}^2\right\}}`$ provides us with enough information to determine the average purities $`\overline{\text{Tr}\left\{\rho _{i\left(k_i\right)}^2\right\}}`$ of every subsystem. In particular, $`L=n/2`$ and all $`n_i=2`$ is the case of only pairwise entanglement. Since Eq. (5) holds for all $`\rho _i`$ taking $`n_i=2`$ provides a test for multi-particle as opposed to two-particle correlations in a given state. We note that unless $`\rho `$ is pure classical multipartite correlations will also be detected by our network.
We will now study how the working of the network is affected by experimental errors. In particular we will estimate how many runs are necessary to obtain the probabilities $`P(j)`$ with sufficient accuracy in the presence of errors.
## IV Effects of experimental error
The errors introduced in Sec. II.3 affect the ability to find the purities $`\text{Tr}\left\{\rho _B^2\right\}`$ as well as the average purities $`\overline{\text{Tr}\left\{\rho _{\left(k\right)}^2\right\}}`$ associated with $`\rho _{12\mathrm{}n}`$. All of these errors are of one of two mathematical kinds: extra *pairs* of atoms and missing *single* atoms. We call these two errors โbeam-splitterโ and โdetectorโ error respectively. Their respective probabilities $`q=q_{bs}+q_l`$ and $`p=p_d+p_l`$ are understood to include also errors occurring while particles are lost from doubly occupied sites. The relationship Eq. (13) between the purities of $`\rho _B`$ and the probabilities of an even/odd number $`j_B`$ of singly occupied sites in $`B`$ indicates that an experimental error, occurring with probability $`p`$ per site, changes the result $`\text{Tr}\left\{\rho _B^2\right\}`$ by $`O(|B|p)`$ if we do not attempt to correct for it. This renders the measured results totally meaningless as soon as $`|B|p1`$, because the purity of any given state $`\rho _B`$ is smaller or equal to one. However certain types of error, including the BS and detector errors, can be corrected by a suitable modification of the formulas Eqs. (13, 22) yielding $`\text{Tr}\left\{\rho _B^2\right\}`$ and $`\overline{\text{Tr}\left\{\rho _{\left(k\right)}^2\right\}}`$. This correction eliminates systematic errors, making $`\text{Tr}\left\{\rho _B^2\right\}`$ and $`\overline{\text{Tr}\left\{\rho _{\left(k\right)}^2\right\}}`$ correct on average, but tends to amplify the random errors that are inevitable in measuring probabilities using a finite number of experimental runs. These random errors can in principle be made arbitrarily small for any $`p,q<1`$ by increasing the number of runs, but in practice there is a limit because, as we will show, the number of runs required scales approximately exponentially in $`|B|p`$. We will now investigate the effects of these errors on the performance of the entanglement detection network both without and with spatial resolution.
### IV.1 Without spatial resolution
We assume the probabilities $`p`$ and $`q`$ to be the same for all $`2n`$ lattice sites and in the case of $`q`$ for all symmetric atom pair states $`|aa`$, $`|bb`$ and $`|ab+|ba`$. Errors at different lattice sites are assumed to be uncorrelated.
#### IV.1.1 Beam splitter error
Let $`P_{\mathrm{exp}}(i)`$ be the probability of detecting $`i`$ atoms in an experimental run. If only BS errors are present this probability is given by
$$P_{\mathrm{exp}}(2i)=\underset{j=0}{\overset{i}{}}P(j)\left(\genfrac{}{}{0pt}{}{nj}{ni}\right)q^{ij}(1q)^{ni},$$
(24)
where the factor two in $`P_{\mathrm{exp}}(2i)`$ accounts for $`P(j)`$ being the probability of having $`j`$ antisymmetric pairs. We can use generating functions to invert Eq. (24)
$$\underset{i=0}{\overset{n}{}}x^iP_{\mathrm{exp}}(2i)=\underset{j=0}{\overset{n}{}}P(j)x^j(1q+xq)^{nj},$$
(25)
leading to
$$P(j)=\underset{i=0}{\overset{j}{}}\left(\genfrac{}{}{0pt}{}{ni}{nj}\right)\frac{(q)^{ji}}{(1q)^{ni}}P_{\mathrm{exp}}(2i).$$
(26)
We now apply Eq. (26) to a subsystem $`B`$ and substitute this into Eq. (13), giving
$$\text{Tr}\left\{\rho _B^2\right\}=\underset{i_B=0}{\overset{k}{}}\left(\frac{1+q}{1q}\right)^{ki_B}(1)^{i_B}P_{\mathrm{exp}}(2i_B),$$
(27)
where $`i_B`$ refers to the number of atoms detected in $`B`$. This expression is then averaged over all $`B`$ of size $`|B|=k`$ to give
$$\overline{\text{Tr}\left\{\rho _{\left(k\right)}^2\right\}}=\underset{i=0}{\overset{n}{}}A_{ki}P_{\mathrm{exp}}(2i),$$
(28)
with
$$A_{ki}=\left[\left(\genfrac{}{}{0pt}{}{n}{k}\right)\right]^1\underset{l=0}{\overset{k}{}}(1)^l\left(\frac{1+q}{1q}\right)^{kl}\left(\genfrac{}{}{0pt}{}{i}{l}\right)\left(\genfrac{}{}{0pt}{}{ni}{kl}\right).$$
(29)
Hence, using Eq. (29) instead of Eq. (22) corrects all the *systematic* error caused by an imperfect BS. We are still left with the inherent *random* error associated with the measurement of the probabilities $`P_{\mathrm{exp}}(2i)`$, which is reduced by increasing the number of experimental runs.
Because of this random error the estimate of $`\overline{\text{Tr}\left\{\rho _{\left(k\right)}^2\right\}}`$ obtained from $`N`$ experimental runs (each using one pair $`\rho \rho `$) has the correct mean but a nonzero standard deviation $`\sqrt{V_k/N}`$, where this defines $`V_k(p,q,\rho )`$; hence $`O(V_k)`$ runs are necessary for meaningful results. Note that in general $`V_k>0`$ even when $`p=q=0`$, as it includes the inherent quantum uncertainty as well as that added by experimental error. In the case of BS error
$`V_k`$ $`=`$ $`{\displaystyle \underset{i}{}}P_{\mathrm{exp}}(i)A_{ki}^2\left(\overline{\text{Tr}\left\{\rho _{\left(k\right)}^2\right\}}\right)^2<\mathrm{max}_i\left(A_{ki}^2\right)`$ (30)
$``$ $`\left({\displaystyle \frac{1+q}{1q}}\right)^{2k}e^{4kq},`$
where the approximation is valid for $`k1`$, $`q1`$. The bound Eq. (30) proves that the number of runs required to obtain meaningful estimates of $`\overline{\text{Tr}\left\{\rho _{\left(k\right)}^2\right\}}`$ is reasonable for $`k1/q`$, however large $`n`$ is.
We numerically computed the worst case by maximizing $`V_k`$ with respect to $`\overline{\text{Tr}\left\{\rho _{\left(k\right)}^2\right\}}`$ subject only to Eq. (5) and compare it to cluster states $`|(\varphi =\pi )_n`$ in Fig. 6. The results confirm the analytically found exponential increase of $`V_k`$ with $`q`$ in the worst case. For the cluster state $`V_k`$ increases only slowly with $`q`$ for small $`k1/q`$ while for $`k1/q`$ we find approximately exponential growth of $`V_k`$ with $`q`$. We also computed the variances for maximally entangled states and found that they are quite close to the worst case shown in Fig. 7a. Therefore one may in an experiment generally not expect the variances $`V_k`$ to be much smaller than our worst case results. Thus only BS errors up to $`q=1/k`$ are acceptable and yield reliable results in a reasonable number of runs for all $`0kn`$. However, as shown in Fig. 7 for $`q=1/k`$, the average purities with $`kn`$ will be determined much more accurately than those with $`kn`$ and should thus be preferentially used for determining parameters characterizing the measured state.
We finally note that if $`J`$ is fluctuating from run to run and $`q_{bs}`$ hence becomes a random variable error correction is still possible. In this case we have to replace Eq. (24) by
$`P_{\mathrm{exp}}(2i)`$ $`=`$ $`{\displaystyle \underset{j=0}{\overset{i}{}}}P(j)\left({\displaystyle \genfrac{}{}{0pt}{}{nj}{ni}}\right)\times `$ (31)
$`{\displaystyle dJf(J)q(J)^{ij}(1q(J))^{ni}},`$
where $`f(J)`$ is the probability density function of $`J`$ and $`q(J)`$ denotes the BS error as a function of $`J`$. The resulting system of linear equations Eq. (31) can be treated using the methods introduced above.
#### IV.1.2 Detector error
We now assume that only detector errors are present and each atom has a probability $`p`$ of failing to be detected. In this case $`P_{\mathrm{exp}}(i)`$ is related to $`P(j)`$ via
$$P_{\mathrm{exp}}(i)=\underset{j=i/2}{\overset{n}{}}P(j)\left(\genfrac{}{}{0pt}{}{2j}{i}\right)p^{2ji}(1p)^i.$$
(32)
We can again solve this equation by methods similar to those used in Sec. IV.1.1 (though this time it is not a true inverse because the system is over-determined Foot2 ) to obtain
$$P(j)=\underset{i=2j}{\overset{n}{}}\left(\genfrac{}{}{0pt}{}{i}{2j}\right)\frac{(p)^{i2j}}{(1p)^i}P_{\mathrm{exp}}(i),$$
(33)
where non-integer values of $`j`$ are discarded. By combining Eqs. (22, 33) we obtain $`\overline{\text{Tr}\left\{\rho _{\left(k\right)}^2\right\}}`$ in terms of $`P_{\mathrm{exp}}(i)`$. For the remaining random error measured by $`V_k`$ we obtain the upper bound
$$V_k\left(\frac{1+p}{1p}\right)^{4n}e^{8np}.$$
(34)
Numerical calculations for $`n15`$ confirm this exponential growth of $`V_k`$ with $`p`$ at $`np1`$. The results are shown in Fig. 8. This time, however, $`V_k`$ is typically much smaller than the analytic bound, e.g. for the $`n=15`$ cluster state, fitting $`V_k\mathrm{exp}(\beta np)`$ gives $`\beta 2`$. However, the exponential growth with $`n`$ implies a practical limit of $`n1/p`$ for any $`k`$ contrary to the case of BS errors. This scaling can be improved by using the least squares method to handle the over-determined linear system of equations Eq. (32). We do not have an analytic bound analogous to Eq. (34) for the least squares method but numerical calculations for the cluster state and the worst case result in significantly lower values for $`V_k`$ than those obtained from Eq. (33). Most importantly as shown in Fig. 9 it appears that the scaling of $`V_k`$ becomes exponential in $`kp`$ rather than $`np`$ similarly to the case of BS errors. This implies that an error of $`p1/k`$ is acceptable for obtaining meaningful estimates of $`\overline{\text{Tr}\left\{\rho _{\left(k\right)}^2\right\}}`$ in a reasonable number of experimental runs.
In an actual experiment both BS and detector errors will be present and first correcting for the detector error $`p`$ using Eq. (33), then substituting the resulting $`P(j)`$ for $`P_{\mathrm{exp}}(2j)`$ in Eq. (29) to correct for the BS error $`q`$ yields a combined analytical error bound of
$$V_k\left(\frac{1+p}{1p}\right)^{4n}\left(\frac{1+q}{1q}\right)^{2k}e^{8np+4kq}.$$
(35)
According to our numerical results using the least squares method this bound can be improved requiring only $`p,q1/k`$ for obtaining sufficiently small errors.
### IV.2 With spatial resolution
If spatial resolution is available then $`\text{Tr}\left\{\rho _B^2\right\}`$, not just its average $`\overline{\text{Tr}\left\{\rho _{\left(k\right)}^2\right\}}`$ over all subsets $`B`$ of size $`k`$, becomes accessible to measurement. This allows us to do some state characterizations which would otherwise be impossible and also introduces extra redundancy. As we will show below this does not affect the tolerance to BS errors but we will find that the detector error tolerance improves to $`p^21/|B|`$. Imperfections in the spatial resolution, however, will lead to additional errors in determining $`\text{Tr}\left\{\rho _B^2\right\}`$.
#### IV.2.1 Beam splitter error
The variance of $`\text{Tr}\left\{\rho _B^2\right\}`$ due to BS errors can be directly inferred from Eq. (27), where only the atoms in $`B`$ are counted. Using the same methods as in Sec. IV.1.1 we find an upper bound for $`V_B`$ (defined analogously to $`V_k`$) given by $`V_B((1+q)/(1q))^{2|B|}\mathrm{exp}(4|B|q)`$. The BS error tolerance (for fixed $`k=|B|`$) is thus independent of whether spatial resolution is available or not.
#### IV.2.2 Detector error
The situation is different for the detector error. Each antisymmetric pair contains two atoms only one of which needs to be detected to know that it was antisymmetric. Therefore the effective error probability becomes $`p^2`$. The resulting formula is
$$\text{Tr}\left\{\rho _B^2\right\}=\underset{i=0}{\overset{|B|}{}}(1)^i\left(\frac{1+p^2}{1p^2}\right)^iP_{\mathrm{exp}}^B(i),$$
(36)
where $`P_{\mathrm{exp}}^B(i)`$ is the probability of measuring $`i`$ antisymmetric sites in subset $`B`$. The variance bound is given by $`V_B((1+p^2)/(1p^2))^{2|B|}\mathrm{exp}(4|B|p^2)`$.
#### IV.2.3 Imperfect spatial resolution
There is a new type of error to consider as the spatial resolution itself will not in practice be perfect. If we let $`f(x,y)`$ be the probability of finding at position $`x`$ a particle which is actually at position $`y`$, then
$$P_{\mathrm{exp}}(A)=\underset{iB}{}\underset{\varsigma }{}\frac{1}{s(A)}\underset{i=0}{\overset{k}{}}f(A_{\varsigma (i+k)},B_i)f(A_{\varsigma (i)},B_i)P(B),$$
(37)
where $`A`$ is the โsetโ (its elements are not necessarily distinct) of observed atom positions, and $`P_{\mathrm{exp}}(A)`$ is the experimental probability of observing atoms exactly at these positions $`A`$. The set $`B`$ denotes the antisymmetric sites and the probability of having antisymmetric sites at positions $`B`$ is $`P(B)`$, $`k=|B|=|A|/2`$, and $`\varsigma `$ runs over all $`(2k)!`$ permutations of the $`2k`$ atoms in $`A`$. The $`i`$-th element of $`B`$ and $`A`$ are written as $`B_i`$ and $`A_i`$, respectively. The symmetry factor $`s(A)`$ stands for the number of permutations $`\varsigma `$ which leave the ordered lists of atoms $`\{A_i\}`$ invariant (e.g., $`s(\{1,1\})=2`$, $`s(\{1,2\})=1`$), and is needed because our summation runs over different ordered lists $`\{A_i\}`$ of the same set $`A`$. This is an over-determined linear system, and just as in the case of detector error, we can either explicitly solve it by discarding some of the equations or apply the least squares method. As an explicit solution we can e.g. use
$$P(B)=\underset{\{A_i\}}{}\frac{s(A)}{2^k}\left(\underset{i=1}{\overset{2k}{}}f^1(B_i,A_i)\right)P_{\mathrm{exp}}(A),$$
(38)
where we define $`B_{k+i}B_i`$. The sum runs over all ordered lists $`\{A_i\}`$ and $`f^1`$ is the matrix inverse of $`f`$, that is $`f^1(y,x)f(x,z)=\delta _{xz}`$.
As before by performing numerical calculations using the least squares method we find much lower variances than with Eq. (38). An example is shown in Fig. 10 where we plot the variance of $`\text{Tr}\left\{\rho _B^2\right\}`$ due to a Gaussian position error of the form
$$f(x,y)=\frac{1}{\sqrt{2\pi }\sigma }\underset{x\lambda /4}{\overset{x+\lambda /4}{}}dz\mathrm{exp}[(zy)^2/(2\sigma ^2)]$$
(39)
where $`\sigma `$ is the standard deviation and $`\lambda `$ is the wave length of the laser creating the optical lattice. The corresponding lattice spacing is $`\lambda /2`$. The results obtained from Eq. (38) shown in Fig. 10a, c require a resolution of $`\sigma \lambda /2`$ whereas the least squares method shown in Fig. 10b, d yields reasonable variances $`V_B`$ for spatial resolutions up to $`\sigma 3\lambda /2`$. However, due to the exponentially large number of possibilities for $`A`$ the least squares calculation becomes intractable for large $`n`$.
## V Conclusion
We discussed the detection and characterization of multipartite entanglement in optical lattices with the quantum network introduced in Carolina2004 under different experimental conditions. We first described how the network can be implemented in ideal experimental conditions and showed that it allows to characterize a number of important classes of states like cluster-like states and macroscopic superposition states. We investigated the experimental realization of the network in an optical lattice and identified lack of spatial resolution, errors in the BS operation and imperfect atom detection as the main sources of error. We showed that even in the absence of spatial resolution entanglement can be detected. In cases where the entangled state is characterized by a few parameters only we found that these can often be determined from the measurement results. We also studied the influence of BS errors occurring with probability $`q`$ and detection errors with probability $`p`$ and concluded that with small numbers of experimental runs entanglement can be detected between $`k`$ atoms as long as $`kp1`$ and $`kq1`$. Finally we showed that for obtaining purities $`\text{Tr}\left\{\rho _B^2\right\}`$ of subsets $`B`$ rather than average purities with reasonable experimental effort a spatial resolution of $`\sigma \lambda `$ is necessary.
The results obtained in this work show that unambiguous multipartite entanglement detection in optical lattices is possible with current technology. This has not yet been achieved experimentally BlochEnt . Furthermore it will even be possible to determine some of the characteristics of entangled states created in these experiments without the requirement of performing spatially resolved measurements. Our network is thus a viable alternative to detecting entanglement via witnesses or full quantum state tomography.
###### Acknowledgements.
This work was supported by EPSRC through the QIP IRC (www.qipirc.org) GR/S82176/01 and project EP/C51933/1, and by the EU network OLAQUI. C.M.A. thanks Artur Ekert for useful discussions and is supported by the Fundaรงรฃo para a Ciรชncia e Tecnologia (Portugal).
## Appendix A The beam splitter operation
Since the BS only couples two lattice sites in rows I and II of each column (see Fig. 2) we consider a single such pair and omit the column superscript $`j`$ in this section. For $`U=0`$, we obtain from the Heisenberg equations for the operators $`\alpha =a,b`$
$`\alpha _I(t)`$ $`=`$ $`\mathrm{cos}(Jt)\alpha _Ii\mathrm{sin}(Jt)\alpha _{II},`$
$`\alpha _{II}(t)`$ $`=`$ $`\mathrm{cos}(Jt)\alpha _{II}i\mathrm{sin}(Jt)\alpha _I.`$ (40)
Hence, applying $`H_{\mathrm{BS}}`$ for a time $`t_{bs}=\pi /(4J)`$ implements a perfect pairwise BS.
Initially the atom pair is in a state of the form $`\rho \rho `$, where $`\rho `$ is a single qubit state and hence has a spectral decomposition of the form $`\rho =\lambda _1|cc\left|+\lambda _2\right|dd|`$, with $`\lambda _1+\lambda _2=1`$, $`|c=c^{}|\mathrm{vac}`$, and $`|d=d^{}|\mathrm{vac}`$. Here $`c^{}`$, $`d^{}`$ are linear superpositions of $`a^{}`$, $`b^{}`$ with coefficients depending on $`\rho `$. Therefore we can write
$`\rho \rho `$ $`=`$ $`\lambda _1^2|c_Ic_{II}c_Ic_{II}\left|+\lambda _2^2\right|d_Id_{II}d_Id_{II}|`$
$`+{\displaystyle \frac{\lambda _1\lambda _2}{2}}(|c_Id_{II}+|d_Ic_{II})(c_Id_{II}|+d_Ic_{II}|)`$
$`+{\displaystyle \frac{\lambda _1\lambda _2}{2}}(|c_Id_{II}|d_Ic_{II})(c_Id_{II}|d_Ic_{II}|),`$
which is a classical mixture of a symmetric state with probability $`P_+=1\lambda _1\lambda _2`$ and an antisymmetric state with probability $`P_{}=\lambda _1\lambda _2`$. After the BS the resulting state $`\rho ^{}=\mathrm{exp}(iH_{BS}t)\rho \rho \mathrm{exp}(iH_{BS}t)`$ is given by
$`\rho ^{}`$ $`=`$ $`\lambda _1^2|\mathrm{\Phi }_1\mathrm{\Phi }_1\left|+\lambda _2^2\right|\mathrm{\Phi }_2\mathrm{\Phi }_2|`$
$`+\lambda _1\lambda _2|\mathrm{\Phi }_3\mathrm{\Phi }_3\left|+\lambda _1\lambda _2\right|\mathrm{\Phi }_4\mathrm{\Phi }_4|,`$
where $`|\mathrm{\Phi }_1=(c_I^{}c_I^{}+c_{II}^{}c_{II}^{})|\mathrm{vac}/2`$, $`|\mathrm{\Phi }_2=(d_I^{}d_I^{}+d_{II}^{}d_{II}^{})|\mathrm{vac}/2`$, $`|\mathrm{\Phi }_3=(c_I^{}d_I^{}+c_{II}^{}d_{II}^{})|\mathrm{vac}/\sqrt{2}`$ are states with double occupancy in one row and an empty site in the other row while $`|\mathrm{\Phi }_4=(c_I^{}d_{II}^{}c_{II}^{}d_I^{})|\mathrm{vac}/\sqrt{2}`$ is a state with a singly occupied site in each row. Hence, after the BS we will find a doubly occupied site with probability $`1\lambda _1\lambda _2=P_+`$ while two singly occupied sites result with probability $`\lambda _1\lambda _2=P_{}`$.
However, in practice it is impossible to completely turn off the interaction $`U`$ which may result in a symmetric state failing to bunch. We consider any symmetric state $`|\mathrm{\Psi }_s`$ and because $`H_{BS}+H_U`$ acts only on the row indices, not the internal states, and is symmetric between the two rows the probability $`q_{bs}`$ of failure to bunch is given by
$`q_{bs}`$ $`=`$ $`|\mathrm{\Psi }_s\left|e^{i(H_{BS}+H_U)t}\right|\mathrm{\Psi }_s|^2`$
$`=`$ $`{\displaystyle \frac{16J^2}{16J^2+U^2}}\mathrm{cos}^2(\sqrt{16J^2+U^2}{\displaystyle \frac{t_{bs}}{2}})`$
$`+{\displaystyle \frac{U^2}{16J^2+U^2}}.`$
The optimal choice for the BS time is $`t_{bs}=\pi /\sqrt{16J^2+U^2}`$ for which $`q_{bs}=U^2/(16J^2+U^2)`$. If the hopping term is not controlled perfectly accurately but fluctuates by $`\delta J`$ around a mean $`J`$ we set $`t_{bs}=\pi /\sqrt{16J^2+U^2}`$ and obtain Eq. (18).
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# Tangled magnetic fields and CMBR signal from reionization epoch
## I Introduction
Coherent magnetic field of micro-gauss strength are observed in galaxies and clusters of galaxies (parker ,zeldovich , for a recent review see e.g. widrow ). Observational evidence exist for even larger scale magnetic fields kim . The origin of these observed magnetic fields however is not well understood. The observed magnetic fields could have arisen from dynamo amplification of small seed ($`\stackrel{<}{}10^{20}\mathrm{G}`$) magnetic fields (see e.g. ruzmaikin ; shukurov ) which originated from various astrophysical processes in the early universe (harrison , subramanian1 , kulsrud , grasso , widrow , gopal1 , matarrese ) . Alternatively the magnetic fields of nano-Gauss strength could have originated from some early universe process like electroweak phase transition or during inflation (e.g. turner , ratra , see grasso , giovannini for reviews). In this scenario, the observed micro-gauss magnetic fields then result from adiabatic compression of this primordial magnetic field.
The existence of primordial magnetic fields of nano-Gauss strength can influence the large scale structure formation in the universe (wasserman , kim1 , subramanian2 , sethi , gopal , sethi1 ). Also these magnetic fields could leave observable signatures in the CMBR anisotropies (barrow , subramanian3 , subramanian4 , durrer , seshadri , mack , lewis ).
In recent years, the study of CMBR anisotropies has proved to be the best probe of the theories of structure formation in the universe (see e.g. hu1 for a recent review). The simplest model of scalar, adiabatic perturbations, generated during the inflationary era, appear to be in good agreement with both the CMBR anisotropy measurements and the distribution of matter at the present epoch (see e.g. spergel , tegmark ). Tensor perturbations could have been sourced by primordial gravitational waves during the inflationary epoch. There is no definitive evidence of the existence of tensor perturbations in the CMBR anisotropy data; the WMAP experiment, from temperature anisotropy data, obtained upper limits on the amplitude of tensor perturbations spergel . Vector perturbations are generally not considered in the standard analysis as the primordial vector perturbations would have decayed by the epoch of recombination in the absence of a source. An indisputable signal of vector and tensor modes is that unlike scalar modes these perturbations generate $`B`$-type CMBR polarization anisotropies (see e.g. hu4 and references therein). At present, only upper limits exist on this polarization mode kovac . However, the on-going CMBR probe WMAP and the upcoming experiment Planck surveyor have the capability of unravelling the effects of vector and tensor perturbations.
Recent WMAP results suggest that the universe underwent an epoch of re-ionization at $`z15`$; in particular WMAP analysis concluded that the optical depth to the last reionization surface is $`\tau _{\mathrm{reion}}=0.17\pm 0.04`$ kogut ; which means that nearly 20% of CMBR photons re-scattered during the period of reionization. The secondary anisotropies generated during this re-scattering leave interesting signatures especially in CMBR polarization anisotropies (see e.g. zaldarriaga ), as is evidenced by the recent WMAP results kogut .
Primordial magnetic fields source all three kinds of perturbations. In this paper we study the secondary CMBR anisotropies, generated during the epoch of reionization, from vector, tensor, and scalar modes, in the presence of primordial tangled magnetic fields. Recently, Lewis lewis computed fully-numerically CMBR vector and tensor temperature and polarization anisotropies in the presence of magnetic fields including the effects of reionization. Seshadri and Subramanian seshadri1 calculated the secondary temperature anisotropies from vector modes owing to reionization. Our approach is to compute the secondary temperature and polarization anisotropies semi-analytically by identifying the main sources of anisotropies in each case; we compute the anisotropies by using the formalism of hu4 . We also compute the tensor primary signal to compare with the already existing analytical results for tensor anisotropies mack .
In the next section, we set up the preliminaries by discussing the models for primordial magnetic fields and the process of reionization. In ยง3, ยง4, and ยง5, we consider vector, tensor, and scalar modes. In ยง6 the detectability of the signal is discussed. In ยง7, we present and summarize our conclusions. While presenting numerical results in this paper, we use the currently-favoured FRW model: spatially flat with $`\mathrm{\Omega }_m=0.3`$ and $`\mathrm{\Omega }_\mathrm{\Lambda }=0.7`$ (spergel , perlmutter , riess ) with $`\mathrm{\Omega }_bh^2=0.024`$ (spergel , tytler ) and $`h=0.7`$ (freedman ).
## II Primordial magnetic fields, reionization, and CMBR anisotropies
Assuming that the tangled magnetic fields are generated by some process in the early universe, e.g. during inflationary epoch, magnetic fields at large scales ($`\stackrel{>}{}0.1\mathrm{Mpc}`$) are not affected appreciably by different processes in either the pre-recombination or the post-recombination universe (sethi1 , jedamzik , subramanian2 ). In this regime, the magnetic field decays as $`1/a^2`$ from the expansion of the universe. This allows us to express: $`๐(๐ฑ,\eta )=\stackrel{~}{๐}(๐ฑ)/a^2`$; here $`๐ฑ`$ is the comoving coordinate. We further assume tangled magnetic fields, $`\stackrel{~}{๐}`$, present in the early universe, to be an isotropic, homogeneous, and Gaussian random process. This allows one to write, in Fourier space (see .e.g. landau ):
$$\stackrel{~}{B}_i(๐ช)\stackrel{~}{B}_j^{}(๐ค)=\delta _D^3(๐ช๐ค)\left(\delta _{ij}k_ik_j/k^2\right)M(k)$$
(1)
Here $`M(k)`$ is the magnetic field power spectrum and $`k=|๐ค|`$ is the comoving wavenumber. We assume a power-law magnetic field power spectrum here: $`M(k)=Ak^n`$. We consider the range of scales between $`k_{\mathrm{min}}`$ taken to be zero here and small scale cut-off at $`k=k_{\mathrm{max}}`$; $`k_{\mathrm{max}}`$ is determined by the effects of damping by radiative viscosity before recombination. Following Jedamzik et al. jedamzik , $`k_{\mathrm{max}}60\mathrm{Mpc}^1(B_0/(3\times 10^9\mathrm{G})`$; $`B_0`$ is the RMS of magnetic field fluctuations at the present epoch. $`A`$ can be calculated by fixing the value of the RMS of the magnetic field, $`B_0`$, smoothed at a given scale, $`k_\mathrm{c}`$. Using a sharp $`k`$-space filter, we get,
$$A=\frac{\pi ^2(3+n)}{k_c^{(3+n)}}B_0^2$$
(2)
We take $`k_c=1\mathrm{M}\mathrm{p}\mathrm{c}^1`$ throughout this paper. For $`n3`$, the spectral indices of interest in this paper, the RMS filtered at any scale has weak dependence on the scale of filtering.
Recent WMAP observations showed that the universe might have got ionized at redshifts $`z15`$. However the details of the ionization history of the universe during the reionization era are still unknown; for instance the universe might have got reionized at $`z=15`$ and remained fully ionized till the present or the universe might have got partially reionized with ionized fraction $`x_e\stackrel{<}{}0.3`$ at $`z30`$ and became fully ionized for $`z\stackrel{<}{}10`$. Both these ionization histories are compatible with the WMAP results kogut . Given this lack of knowledge we model the reionization history by assuming the following visibility function, which gives the normalized probability that the photon last scattered between epoch $`\eta `$ and $`\eta +d\eta `$, to model the period of reionization:
$$g(\eta ,\eta _0)\dot{\tau }\mathrm{exp}(\tau )=\frac{(1\mathrm{exp}(\tau _{\mathrm{reion}}))}{\sqrt{\pi }\mathrm{\Delta }\eta _{\mathrm{reion}}}\mathrm{exp}\left[(\eta \eta _{\mathrm{reion}})^2/\mathrm{\Delta }\eta _{\mathrm{reion}}^2\right]$$
(3)
Here $`\tau (\eta ,\eta _0)=_{\eta _0}^\eta n_e\sigma _t๐t`$ is the optical depth from Thompson scattering; $`\tau _{\mathrm{reion}}`$ is the optical depth to the epoch of reionization; for compatibility with WMAP results, we use $`\tau _{\mathrm{reion}}=0.17`$ throughout. $`\eta _{\mathrm{reion}}`$ and $`\mathrm{\Delta }\eta _{\mathrm{reion}}`$ are the epoch of reionization and the width of reionization phase, respectively; we take $`\eta _{\mathrm{reion}}`$ corresponding to $`z_{\mathrm{reion}}=15`$ and $`\mathrm{\Delta }\eta _{\mathrm{reion}}=0.25\eta _{\mathrm{reion}}`$. Notice that the visibility function is normalized to $`\tau _{\mathrm{reion}}`$ for $`\tau _{\mathrm{reion}}<<1`$.
## III CMBR anisotropies from vector modes
From a given wave number $`k`$ of vector perturbations, the contribution to CMBR temperature and polarization anisotropies to a given angular mode $`\mathrm{}`$ can be expressed as (see e.g. hu4 ):
$`{\displaystyle \frac{\mathrm{\Theta }_T\mathrm{}^\mathrm{v}(\eta _0,k)}{(2\mathrm{}+1)}}`$ $`=`$ $`{\displaystyle _0^{\eta _0}}๐\eta \mathrm{exp}(\tau )\left[\dot{\tau }(v_\mathrm{b}^\mathrm{v}V)j_{\mathrm{}}^{(11)}[k(\eta _0\eta )]+(\dot{\tau }P^\mathrm{v}(\eta )+{\displaystyle \frac{1}{\sqrt{3}}}kV)j_{\mathrm{}}^{(21)}[k(\eta _o\eta )]\right]`$ (4)
$`{\displaystyle \frac{\mathrm{\Theta }_E\mathrm{}^\mathrm{v}(\eta _0,k)}{(2\mathrm{}+1)}}`$ $`=`$ $`\sqrt{6}{\displaystyle _0^{\eta _0}}๐\eta \mathrm{exp}(\tau )\dot{\tau }P^\mathrm{v}(\eta )ฯต_{\mathrm{}}^\mathrm{v}[k(\eta _0\eta )]`$ (5)
$`{\displaystyle \frac{\mathrm{\Theta }_B\mathrm{}^\mathrm{v}(\eta _0,k)}{(2\mathrm{}+1)}}`$ $`=`$ $`\sqrt{6}{\displaystyle _0^{\eta _0}}๐\eta \mathrm{exp}(\tau )\dot{\tau }P^\mathrm{v}(\eta )\beta _{\mathrm{}}^\mathrm{v}[k(\eta _0\eta )]`$ (6)
Here $`v_\mathrm{b}^\mathrm{v}`$ and $`V`$ are the line-of-sight components of the vortical component of the baryon velocity and the vector metric perturbation. $`P^\mathrm{v}(\eta )=1/10[\mathrm{\Theta }_{T2}^\mathrm{v}\sqrt{6}\mathrm{\Theta }_{E2}^\mathrm{v}]`$ and the Bessel functions, $`j_{\mathrm{}}`$, $`ฯต_{\mathrm{}}`$ and $`\beta _{\mathrm{}}`$ that give radial projection for a given mode are given in reference hu4 . The evolution of vector metric perturbations, $`V_i(๐ค,\eta )`$ is determined from Einsteinโs equations (e.g. hu4 , mack ):
$`\dot{V}_i+2{\displaystyle \frac{\dot{a}}{a}}V_i`$ $`=`$ $`{\displaystyle \frac{16\pi Ga^2S_i(๐ค,\eta )}{k}}`$ (7)
$`k^2V_i`$ $`=`$ $`16\pi Ga^2{\displaystyle \underset{j}{}}(\rho _j+p_j)(v_{ij}^\mathrm{v}V_i)`$ (8)
Here $`S_i`$, the source of vector perturbations, is determined by primordial tangled magnetic field in this paper. The index $`j`$ correspond to baryonic, photons and dark matter vortical component of velocities. For tangled magnetic fields, the vortical velocity component of the dark matter doesnโt couple to the source of vector perturbations to linear order and $`\mathrm{\Omega }_i=v_i^\mathrm{v}V_i`$ decays as $`1/a`$ for dark matter (see e.g. mack ); and hence the dark matter contribution can be dropped from the Einsteinโs equations. The photons couple to baryons through Thompson scattering. In the pre-recombination epoch, the photons are tightly coupled to the baryons as the time scale of Thompson scattering is short as compared to the expansion rate; besides the photon density is comparable to baryon density at the epoch of recombination. In the reionized models we consider here, neither the photons are tightly coupled to baryons nor are they dynamically important. Therefore photon contribution can also be dropped from Eq. (8). Eq. (8) then simplifies to:
$$k^2V_i=16\pi Ga^2\rho _B\mathrm{\Omega }_B^\mathrm{v}$$
(9)
with $`\mathrm{\Omega }_B^v=(v_\mathrm{b}^\mathrm{v}V_i)`$. The quantity of interest is the angular power spectrum of the CMBR anisotropies which is obtained from squaring Eqs. (4), (5), and (6), taking ensemble average, and integrating over all $`k`$:
$$C_{\mathrm{}T,E,B}=\frac{4}{\pi }๐kk^2\left[\frac{\mathrm{\Theta }_{T,E,B\mathrm{}}(k,\eta _0)}{2\mathrm{}+1}\right]^2$$
(10)
This expression is valid for both vector and tensor perturbations; for scalar perturbation the prefactor is $`2/\pi `$.
For primordial magnetic field, the sources $`S_i(๐ค,\eta )`$ of vector perturbation (Eq. (7)) is the vortical component of the Lorentz force:
$$S_i(๐ค,\eta )=\frac{1}{a^44\pi }\widehat{k}\mathrm{xF}.\mathrm{T}.[\stackrel{~}{๐}(๐ฑ)\mathrm{x}(\mathrm{x}\stackrel{~}{๐}(๐ฑ))]S_i(๐ค)\frac{1}{a^4}$$
(11)
It can be checked that this Newtonian expression for $`S_i`$ is the same as the more rigorously defined $`\mathrm{\Pi }_i^V`$ in Appendix A (Eq. (28)).
### III.1 Temperature anisotropies from vector modes
As seen from Eq. (4), there are three sources of temperature anisotropies. The most important contribution comes from vorticity $`\mathrm{\Omega }_B^\mathrm{v}`$. For the reionized models, using Eqs. (7) and (8), it can be expressed as:
$$\mathrm{\Omega }_B^\mathrm{v}(๐ค,\eta )=\frac{kS_i(๐ค)\eta }{a\rho _{b0}}$$
(12)
Here $`\rho _{b0}`$ is the baryon density at the present epoch. The other major contribution is from temperature quadrupole $`\mathrm{\Theta }_{T2}^\mathrm{v}`$. For reionized models, the quadrupole at the epoch of reionization is dominated by the free-streaming of the dipole from the last scattering surface (see discussion below, Eq. (15)). This contribution is generally small but in this case can be comparable to the vorticity effects at small values of $`\mathrm{}`$. This is owing to the fact that the vorticity is decaying and therefore during reionization epoch its contribution is smaller as compared to the epoch of recombination. The quadrupole term on the other hand gets its contribution from the vorticity computed at the epoch of recombination (Eq. (15)). This, as we shall discuss below, is not the case for scalar and tensor anisotropies, as the dominant source of anisotropy is either constant (metric perturbations for tensor perturbations) or is increasing (compressional velocity mode for scalar perturbation) as the universe evolves. The third source of temperature anisotropies is metric vector perturbation $`V`$; this term can be comparable to the other terms only at super-horizon scales. We drop this term in this paper.
In Figure 1 we show the secondary temperature anisotropies generated during the epoch of reionization from vector modes. It is seen that the quadrupole term has significant contribution only for $`\mathrm{}\stackrel{<}{}20`$. The dominant contribution at larger $`\mathrm{}`$ is from the vorticity during reionization. The vorticity source contribution can be approximated as:
$$\frac{\mathrm{\Theta }_T\mathrm{}^\mathrm{v}(\eta _0,k)}{(2\mathrm{}+1)}\frac{S(k)}{4\pi \rho _{\mathrm{b0}}}\frac{\eta _0^2}{\eta _{\mathrm{reion}}}j_{\mathrm{}}^{(11)}[k(\eta _0\eta _{\mathrm{reion}})]\tau _{\mathrm{reion}}$$
(13)
for $`\mathrm{}20`$ and
$$\frac{\mathrm{\Theta }_T\mathrm{}^\mathrm{v}(\eta _0,k)}{(2\mathrm{}+1)}\frac{S(k)}{k}\frac{1}{8\pi \rho _{\mathrm{b0}}}\frac{\eta _0^2g(\eta _0\mathrm{}/k,\eta _0)}{\eta _{\mathrm{reion}}}\sqrt{\frac{\pi }{\mathrm{}}}$$
(14)
for $`\mathrm{}\stackrel{>}{}50`$. The temperature angular power spectrum from vorticity increases roughly as $`\mathrm{}^{2.4}`$ for $`\mathrm{}\stackrel{>}{}50`$, with the signal reaching a value roughly $`0.3\mu \mathrm{k}`$ at $`\mathrm{}10^4`$. This is in agreement with the results of seshadri1 .
### III.2 Polarization anisotropies from vector modes
The main source of the polarization anisotropies is the temperature quadrupole $`\mathrm{\Theta }_{T2}^\mathrm{v}`$. One contribution to the temperature quadrupole at the epoch of reionization is from the free-streaming of the dipole from the last scattering surface. The dipole at the last scattering surface can be obtained from the tight-coupling solutions to the temperature anisotropies mack . The quadrupole from the free-streaming of the dipole at the epoch of recombination is:
$$\mathrm{\Theta }_{T2}^\mathrm{v}(๐ค,\eta )=5\mathrm{\Omega }_B^\mathrm{v}(๐ค,\eta _{\mathrm{rec}})j_2^{(11)}(k(\eta \eta _{\mathrm{rec}}))$$
(15)
Here $`\eta _{\mathrm{rec}}`$ corresponds to the epoch of recombination. As is the case for scalar perturbation-induced polarization in the reionized model (e.g. zaldarriaga ) this quadrupole doesnโt suffer the suppression as the quadrupole prior to the epoch of recombination when the photons and baryons are tightly coupled. The structure of anisotropies generated by the quadrupole is determined by $`j_2^{(11)}(k\eta )`$ around the epoch of reionization. This typically gives a peak in anisotropies at $`\mathrm{}2\eta _0/\eta _{\mathrm{reion}}`$. This source dominates the contribution to polarization anisotropies for $`\mathrm{}\stackrel{<}{}10`$. Another contribution to the temperature quadrupole at the epoch of reionization comes from the secondary temperature anisotropies generated at the epoch of reionization. The approximate value of this quadrupole can be got from retaining the first term in Eq. (4):
$$\frac{\mathrm{\Theta }_{T2}^\mathrm{v}(\eta ,k)}{(2\mathrm{}+1)}=_0^\eta ๐\eta ^{}g(\eta ,\eta ^{})\mathrm{\Omega }_B^\mathrm{v}(\eta ^{})j_2^{(11)}[k(\eta \eta ^{})]$$
(16)
This contribution is generically smaller than the first contribution. Firstly, this depends on the vorticity evaluated close to the epoch of reionization as opposed to the first contribution which is proportional to the vorticity at the epoch of recombination. As the vorticity decays as $`a^{1/2}`$ in the matter-dominated era (Eq. (12)), the latter contribution is suppressed by nearly a factor a $`100`$ in the angular power spectrum. Second, as only a small fraction of photons re-scatter (nearly 20%), this contribution is further suppressed by a factor of $`\tau _{\mathrm{reion}}^2`$. However, this contribution is not suppressed at small angular scales and, therefore, might dominate the polarization anisotropies at large values of $`\mathrm{}`$.
In Figure 2, we show the $`E`$ and $`B`$ polarization angular power spectrum from the sources given by Eqs. (15) and (16). As discussed above, the secondary polarization anisotropies are dominated by the quadrupole generated by free-streaming of dipole at the last scattering surface. As expected for vector modes (hu4 ), the $`B`$-mode signal is larger than the $`E`$-mode signal; the signal strength reaches $`10^3\mu \mathrm{k}`$ at $`\mathrm{}10`$ in both cases. This dominates the primary signal for $`\mathrm{}\stackrel{<}{}10`$ as also seen in the numerical results of Lewis lewis . The contribution from the quadrupole generated at the epoch of reionization is seen to be completely sub-dominant
## IV CMBR anisotropies from tensor modes
The energy-momentum tensor for magnetic fields has a non-vanishing traceless, transverse component which sources the corresponding tensor metric perturbation. This in turn affects the propagation of radiation from the last scattering surface to the present and hence gets manifested as additional anisotropies. In this section we calculate the effect of reionization on the resultant anisotropies. For the temperature anisotropies, we study this effect, by calculating the power spectra separately for the standard recombination (no-reionization) and reionized scenario whereas for the polarization anisotropies we compute the secondary anisotropies by using the visibility function given by Eq (3).
### IV.1 Tensor temperature anisotropies
The line-of-sight integral solution for temperature anisotropies, for tensor perturbations is given by (hu4 ):
$$\frac{\mathrm{\Theta }_\mathrm{}T^T(k,\eta _0)}{2\mathrm{}+1}=_0^{\eta _0}๐\eta e^\tau [\dot{\tau }P^{(T)}\dot{h}]j_{\mathrm{}}^{(22)}[k(\eta _0\eta )]$$
(17)
Here, $`P^T(\eta )=1/10[\mathrm{\Theta }_{T2}^\mathrm{T}\sqrt{6}\mathrm{\Theta }_{E2}^\mathrm{T}]`$ is the tensor polarization source and $`\dot{h}`$ is the gravitational wave contribution whose evolution is detailed in Appendix B. The polarization source is modulated by the visibility function and hence is localized to the last-scattering surface. In the tight coupling limit before recombination, $`P^T\dot{h}/(3\dot{\tau })`$ (mack ); for a more detailed derivation of $`P^T`$ in the tight-coupling regime see Appendix B. In the post-recombination epoch, $`P^T`$ is determined by the free-streaming of quadrupole generated at the last scattering surface. However, the visibility function is very small at epochs prior to reionization. Therefore the main contribution of this term comes only from epochs prior to recombination. The gravitational wave source on the other hand being modulated by the cumulative visibility $`\mathrm{exp}(\tau )`$ contributes at all epochs. As a result, the $`P^T`$ contributes negligibly to temperature anisotropies at all multipoles for the case of standard recombination. In the reionized model, this term gets additional contribution from epochs close to reionization redshift but continues to be sub-dominant to the other term. We have also checked this numerically. Hence we can neglect the first term in the above solution and using the matter-dominated solution for $`\dot{h}`$ (Appendix B) we arrive at the following expression for the angular power spectrum:
$$C_\mathrm{}T^T=\frac{4}{\pi }\left(\frac{9R_\gamma }{\rho _\gamma }\right)^2\left(\frac{8(l+2)!}{3(l2)!}\right)๐kk^2\mathrm{\Pi }_T^2(k)\left(_{x_d}^{x_0}๐x\mathrm{exp}(\tau )\frac{j_2(x)}{x}\frac{j_l(x_0x)}{(x_0x)^2}\right)^2$$
(18)
Here, $`xk\eta `$, $`x_0k\eta _0`$ and $`x_dk\eta _{\mathrm{rec}}`$. The above expression is evaluated numerically for the two different ionization histories: standard recombination with and without reionization which are essentially characterized by the different behaviour of the cumulative visibility $`\mathrm{exp}(\tau )`$. The temperature power spectra are shown in Figure 3. As seen in the figure, the temperature power spectrum in both cases shows similar behaviour. The power is nearly flat upto $`\mathrm{}100`$ after which the amplitude falls rapidly. This behaviour is identical to that obtained for primordial gravitational waves. This is expected because, the tensor metric perturbation is sourced by the magnetic field only upto the neutrino-decoupling epoch thereby imprinting an initial nearly scale-invariant spectrum after which the evolution is source-free. The effect of reionization is to reduce the cumulative visibility between recombination ($`z1100`$) and reionization ($`z15`$) epochs. This is why the signal is suppressed for the reionized model. Approximate analytic expressions to primary $`C_\mathrm{}T^T`$ were derived in (mack ). However these give the correct qualitative behaviour $`C_\mathrm{}T^T\mathrm{}^{0.2}`$ only for $`\mathrm{}\stackrel{<}{}100`$. This is because in their analytic results, the lower limit for the time-integral is taken to be zero in Eq. (18) whereas the correct lower limit is $`\eta _{\mathrm{rec}}`$ since the cumulative visibility is zero for $`\eta \stackrel{<}{}\eta _{\mathrm{rec}}`$. We have not neglected this lower limit in our numerical calculation and hence we obtain the damping behaviour for $`\mathrm{}\stackrel{>}{}100`$; Lewis lewis also obtains this damping behaviour. Our results are in reasonable agreement with the results of lewis in the entire range of $`\mathrm{}`$; these results also agree to within factors with the results of mack for $`\mathrm{}\stackrel{<}{}75`$ when the different convention we use for defining $`B_0`$ is taken into account. Our results are quantitatively accurate to better than $`10\%`$ for the lower multipoles $`\mathrm{}\stackrel{<}{}75`$ but begin to differ appreciably from the results of numerical studies for larger $`\mathrm{}`$ or in the damping regime ng . This is because we have not treated the transition regime from radiation-dominated to matter-dominated for the gravitational wave evolution accurately. As described in Appendix B, we have assumed instantaneous transition. This however does not affect the qualitative description of modes whose wavelength is greater than the transition-width $`k\stackrel{<}{}\eta _{\mathrm{eq}}^1`$ which in turn corresponds to multipoles $`\mathrm{}\stackrel{<}{}300`$.
### IV.2 Polarization anisotropies from tensor modes
The line-of-sight solution for the E and B-mode polarization is given as:
$$\frac{\mathrm{\Theta }_\mathrm{}B^T(k,\eta _0)}{2l+1}=\sqrt{6}_0^{\eta _0}๐\eta \dot{\tau }\mathrm{exp}(\tau )P^T\beta _l^T[k(\eta _0\eta )]$$
(19)
$$\frac{\mathrm{\Theta }_\mathrm{}E^T(k,\eta _0)}{2l+1}=\sqrt{6}_0^{\eta _0}๐\eta \dot{\tau }\mathrm{exp}(\tau )P^Tฯต_l^T[k(\eta _0\eta )]$$
(20)
Here, $`\beta _l^T`$ and $`ฯต_l^T`$ are the tensor polarization radial functions as given in Ref. hu4 . The tensor polarization source $`P^T(\eta )`$ in this case will contribute significantly to the above integral only close to the reionization epoch. There are two contributions to the polarization at $`\eta _{\mathrm{reion}}`$: one due to the quadrupole generated at the reionization surface and the other due to the free-streaming primary quadrupole. However, as in the case of vector perturbations, the free-streaming primary quadrupole will give the dominant contribution. We thus have,
$$P^T(k,\eta )=\frac{1}{10}\mathrm{\Theta }_{T2}^T(k,\eta )=\frac{1}{2}_{\eta _{\mathrm{rec}}}^\eta ๐\eta \dot{h}j_2^{(22)}[k(\eta _0\eta )]\mathrm{exp}(\tau )$$
(21)
To simplify the calculations we make the following approximation. Since the visibility function is strongly peaked at $`\eta _{\mathrm{reion}}`$, we take $`P^T`$ outside the integral by evaluating it at the visibility peak $`\eta _{\mathrm{reion}}`$. We have verified that this approximation works extremely well for the lower multipoles where the power is significant. We thus get the following expressions for the polarization angular power spectra:
$$C_\mathrm{}B^T=\frac{6}{\pi }๐kk^2\mathrm{\Pi }_T^2(k)\left[P^T(\eta _{\mathrm{reion}})\right]^2\left(_{\eta _{\mathrm{rec}}}^{\eta _0}๐\eta \dot{\tau }e^\tau \beta _l^T[k(\eta _0\eta )]\right)^2$$
(22)
$$C_\mathrm{}E^T=\frac{6}{\pi }๐kk^2\mathrm{\Pi }_T^2(k)\left[P^T(\eta _{\mathrm{reion}})\right]^2\left(_{\eta _{\mathrm{rec}}}^{\eta _0}๐\eta \dot{\tau }e^\tau ฯต_l^T[k(\eta _0\eta )]\right)^2$$
(23)
As seen in the above expressions, the polarization power spectrum is modulated by the visibility function itself instead of the cumulative visibility in the case of temperature power spectrum. As a result, both $`E`$ as well as $`B`$ mode anisotropies peak close to the multipole corresponding to the horizon scale at reionization. Physically this can be understood as follows: the modes which are super-horizon at reionization experience negligible integrated Sachs-Wolfe effect before $`\eta _{\mathrm{reion}}`$ and hence very small polarization is generated for such modes. Maximum polarization is generated for modes that just enter the horizon at $`\eta _{\mathrm{reion}}`$. For sub-horizon modes, the amplitude of the gravitational wave falls and then sets itself into oscillations which is reflected as a drop in power for higher multipoles.
The polarization power spectra are shown in Figure 4. As seen in the figure, the E-mode power peaks at $`\mathrm{}8`$ whereas the B-mode power peaks at $`\mathrm{}7`$. The corresponding signal strengths at the peaks are $`0.2\mu K`$ in both cases. As expected, the E-mode power is marginally greater than the B-mode power mainly because of the slightly different behaviour of the radial projection factors hu4 . The primary anisotropies for both the polarization modes is sub-dominant on these scales. This enhancement in the net (primary+secondary) signal was also seen in the numerical calculations of Ref. lewis .
We also show the primary CMBR polarization anisotropies from tensor modes in Figure 4. For computing these anisotropies we use the tight-coupling quadrupole, $`P^T(\eta )`$ as derived in Appendix B (Eq. (45)). The primary power spectra are also computed from Eqs. (22) and (23) with lower limit of the time integral replaced by zero. Our results are in agreement with the numerical results of lewis when we take into account the fact that we use different value of $`\eta _{}/\eta _{\mathrm{in}}`$ (Eq. (43)): we use $`\eta _{}/\eta _{\mathrm{in}}=10^{18}`$, which gives the epoch of generation of the tangled magnetic field close to inflationary epoch. While presenting numerical results, Lewis lewis uses $`\eta _{}/\eta _{\mathrm{in}}=10^6`$, which puts the epoch of generation of magnetic field close to the epoch of electro-weak phase transition. Therefore our signal is roughly an order of magnitude larger than the results of lewis .
In Figure 5 we show the expected TE cross-correlation from tensor modes, computed using Eqs. (17), (20), and (21), including the effect of reionization. The effect of reionization is seen as the peak in the TE cross-correlation for $`\mathrm{}\stackrel{<}{}10`$. The signal is dominated by the primary signal for large multipoles. Note that the TE cross-correlation is positive in the entire range $`\mathrm{}\stackrel{<}{}150`$ as was also pointed out by Ref. mack for the primary tensor TE cross-correlation (for details of sign of TE cross-correlation for various modes see hu4 ). In the next section we compare this signal with the WMAP observation of TE cross-correlation.
## V Secondary CMBR anisotropies from scalar modes
In addition to the vortical component of the velocity field, the tangled magnetic fields also generate compressional velocity fields which seed density perturbations. These density perturbations have interesting consequences for the formation of structures in the universe (wasserman , kim1 , subramanian2 , sethi , gopal , sethi1 ). The compressional velocity field also give rise to secondary anisotropies during the epoch of reionization. We compute this anisotropy here. The line-of-sight solution to the temperature anisotropies from these velocity perturbations is:
$$\frac{\mathrm{\Theta }_{\mathrm{}}^S(k,\eta _0)}{2\mathrm{}+1}=_0^{\eta _0}๐\eta e^\tau \dot{\tau }v_\mathrm{b}^\mathrm{s}(k,\eta )j_{\mathrm{}}^{(10)}[k(\eta _0\eta )]$$
(24)
Here $`v_\mathrm{b}^\mathrm{s}`$ is the line-of-sight component of the compressional velocity field. $`j^{(10)}`$ is defined in Ref. hu4 . The growing mode of compressional velocity can be expressed as (wasserman , gopal ):
$$v_\mathrm{b}^\mathrm{s}(๐ค,\eta )=\frac{\eta }{4\pi \rho _{m0}}\widehat{k}.(\mathrm{F}.\mathrm{T}.[\stackrel{~}{๐}(๐ฑ)\mathrm{x}(\mathrm{x}\stackrel{~}{๐}(๐ฑ))])v_\mathrm{b}^0(๐ค)\eta $$
(25)
Here $`\rho _{m0}`$ is the matter density (baryons and the cold dark matter) at the present epoch. The compressional velocity field, unlike the vortical mode, has a growing mode. Also unlike the vortical mode (Eq. (12)), the compressional mode of baryonic velocity couples to the dark matter (gopal , sethi1 ).
In Figure 6 we show the angular power spectrum of the secondary temperature anisotropies generated by the compressional velocity mode. The signal has a peak at roughly the angular scale that corresponds to the width of the visibility function during reionization (for detailed discussion see e.g. dodelson ). The amplitude of this secondary anisotropy is several orders of magnitude smaller than the observed temperature anisotropies and it is unlikely that this signal could be detected.
## VI Detectability
It follows from Figure 1 to 5, that the most important signal at small multipoles arises from tensor polarization anisotropies. In particular, the yet-undetected $`B`$-mode signal holds the promise of unravelling the presence of primordial magnetic fields, as also noted by other authors (e.g. lewis ) In Figure 4, we show the expected errors on the detection of polarization signal from the future CMBR mission, Planck surveyor. The expected 1$`\sigma `$ error, valid for $`\mathrm{}\stackrel{<}{}100`$, is (e.g. zaldarriaga1 , prunet ):
$$\mathrm{\Delta }C_{\mathrm{}}=\left(\frac{2}{(2\mathrm{}+1)f_{\mathrm{sky}}}\right)\left(C_{\mathrm{}}+w^1\right)$$
(26)
For Planck surveyor, $`f_{\mathrm{sky}}1`$ and $`w1.7\times 10^{16}`$ for one-year integration. In Figure 4, we use the primordial tensor $`B`$-mode signal for calculating the expected 1$`\sigma `$ error from Eq. (26). Figure 4 shows that the signal from magnetic fields with strength $`\stackrel{>}{}3\times 10^9\mathrm{G}`$ is detectable by this future mission. However, it is likely that, except for the $`B`$ mode signal, the magnetic field signal will be buried in a larger signal. However, owing to the non-Gaussianity of the magnetic field signal it might still be possible to extract this component of the signal (e.g. lewis ).
In Figure 5 we show the TE cross-correlation signal from tensor modes along with the expected signal from primordial scalar modes with $`\tau _{\mathrm{reion}}=0.17`$, which is in good agreement with the WMAP data of TE cross-correlation kogut . It could be asked if the TE cross-correlation observed by WMAP for $`\mathrm{}\stackrel{<}{}100`$ could be explained as the tensor signal. From Figure 5 it is seen that the tensor signal at small multipoles is roughly a factor of $`5`$ smaller than the scalar signal. And therefore, as the power spectrum from tangled magnetic fields $`B_0^4`$, much of the enhancement observed in the TE cross-correlation for $`\mathrm{}\stackrel{<}{}10`$ could be explainable in terms of the tensor signal from primordial magnetic field for $`B_04.5\times 10^9\mathrm{G}`$. We quantify this notion by computing the $`\chi ^2`$ for $`\mathrm{}15`$ for both the best fit model from WMAP and the tensor model with $`B_04.5\times 10^9\mathrm{G}`$ against the detected WMAP signal <sup>1</sup><sup>1</sup>1for details of WMAP data products http://map.gsfc.nasa.gov kogut ; the $`\chi ^2`$ per degree of freedom in the two cases is $`1.7`$ and $`1.8`$, respectively. Therefore the enhancement can entirely be interpreted in terms of the secondary signal from primordial magnetic fields.
A more realistic possibility is that both primordial scalar and tensor modes gave comparable contribution to the observed signal. As the strength of both these signals for $`\mathrm{}\stackrel{<}{}15`$ is roughly $`\tau _{\mathrm{reion}}^2`$ (for details of secondary scalar signal see e.g. zaldarriaga ), and assuming that there is roughly equal contribution from both, the inferred value of $`\tau _{\mathrm{reion}}`$ from the analysis of the signal could be smaller by a factor of $`\sqrt{2}`$. To quantify this statement, we did a $`\chi ^2`$ test to estimate $`\tau _{\mathrm{reion}}`$ by adding the tensor signal with $`B_04.5\times 10^9\mathrm{G}`$ and the primordial scalar signal with the best-fit cosmological parameters from WMAP. From this analysis we obtain $`\tau _{\mathrm{reion}}0.11\pm 0.02`$ (1$`\sigma `$) with $`\sigma `$ determined by $`\delta \chi ^2=1`$. A possible test of this hypothesis is non-gaussianity of the signal at small multipoles, as the magnetic-field-sourced tensor signal is not Gaussian.
The tensor signal (primary plus secondary) could be appreciable for $`\mathrm{}\stackrel{<}{}100`$. In the range $`15\stackrel{<}{}\mathrm{}\stackrel{<}{}100`$, the tensor and primordial scalar signals are nearly independent of the value of $`\tau _{\mathrm{reion}}`$. While the primordial scalar TE signal anti-correlates for $`\mathrm{}\stackrel{>}{}40`$, the tensor signal shows positive cross-correlation in the range $`\mathrm{}\stackrel{<}{}100`$, as seen in Figure 5. The present WMAP data shows tentative detection of TE anti-correlation for $`\mathrm{}\stackrel{<}{}100`$ peiris . From $`\chi ^2`$ analysis in the range $`15\stackrel{<}{}\mathrm{}\stackrel{<}{}100`$, we notice that the tensor signal alone is a poor fit to the data ($`\chi ^2`$ per degree of freedom of $`2.1`$ as opposed to a value of $`1.6`$ for the primordial scalar model). However a sum of these two signals with $`B_04.5\times 10^9\mathrm{G}`$ is a reasonable fit, as it is dominated by the primordial scalar signal.
It should be noted that for $`B_04.5\times 10^9\mathrm{G}`$, the tensor temperature signal is comparable to the primordial scalar signal (Figure 3). WMAP analysis obtained an upper limit of $`0.7`$ on the ratio of tensor to scalar signal (spergel ). While this limit is rather weak, a more detailed analysis of the temperature signal including the effect of tensor mode signal sourced by primordial magnetic fields might give independent constraints on the strength of primordial magnetic fields.
In our $`\chi ^2`$ analysis we use only the diagonal components of the Fisher matrix. However, owing to incomplete sky, the signal is correlated, especially for small multipoles, across neighboring multipoles. However, a more comprehensive analysis taking into this correlation is likely to yield similar conclusions for the reasons stated above.
Our conclusions are not too sensitive to the value of small scale cut-off $`k_{\mathrm{max}}`$ or the scale of the filter $`k_c`$ used to define the normalization (Eq. (2)) for magnetic field power spectrum index $`n=2.9`$ we use throughout the paper. For $`k_{\mathrm{max}}=k_c=0.05\mathrm{Mpc}^1`$, the foregoing discussion related to tensor mode anisotropies would be valid for $`B_05\times 10^9\mathrm{G}`$. Therefore, the results for TE cross-correlation from tensor perturbations can be interpreted to put bounds on magnetic fields for only large scales $`k\stackrel{<}{}0.05\mathrm{Mpc}^1`$.
The strongest bound on primordial magnetic fields arises from tensor perturbations in the pre-recombination era caprini . These bounds are weakest for nearly scale invariant ($`n3`$) magnetic fields power spectrum (Eq. (33) of caprini ) and largely motivated the choice of the power spectral index we consider here. For $`n=2.9`$, the bound obtained by caprini is considerably weaker than $`B_04.5\times 10^9\mathrm{G}`$, the values of interest to us in this paper. Vector modes might leave observable signature in the temperature and polarization signal for $`\mathrm{}\stackrel{>}{}2000`$; the current observations give weak bound of $`B_0\stackrel{<}{}8\times 10^9\mathrm{G}`$ lewis . Tangled magnetic-field-sourced primary scalar temperature signal gives even weaker bounds koh . More recently, Chen et al. chen obtained, from WMAP data analysis, a limit of $`\stackrel{<}{}10^8\mathrm{G}`$ on the primordial magnetic field strength for nearly scale invariant spectra we consider here; chen consider vector mode temperature signal in their analysis and study possible non-Gaussianity in the WMAP data. Another strong constraint on large scale tangled magnetic fields comes from Faraday rotation of high redshift radio sources (see e.g widrow ); this constraint is also weaker than the value of magnetic field required to explain the enhancement of the TE cross-correlation signal as seen by WMAP sethi . Therefore, the value of $`B_0`$ required to give appreciable contribution to the TE signal is well within the upper limits on $`B_0`$ from other considerations.
It should be noted that the entire foregoing discussion on the tangled-magnetic-field tensor signal can be mapped to primordial tensor modes. The reason for this assertion is that magnetic fields source tensor modes only prior to the epoch of neutrino decoupling, and the subsequent evolution is source free, which is similar to the primordial tensor modes which are generated only during the inflationary epoch and evolve without sources at subsequent times. Therefore, an analysis similar to ours could be used to put constraints on the relative strength of the tensor to scalar mode contribution (for a fixed scale) and the tensor spectral index of the primordial modes. The main observational difference between such an interpretation and the one give here is that tensor signal sourced by magnetic fields will not obey Gaussian statistics as opposed to the primordial tensor modes.
## VII Summary and conclusions
We have computed the secondary anisotropies from the reionization of the universe in the presence of tangled primordial magnetic fields. Throughout our analysis we use the nearly scale invariant magnetic field power spectrum with $`n=2.9`$. For vector modes, we compute the secondary temperature and $`E`$ and $`B`$ mode polarization auto-correlation signal. For scalar modes, the results for secondary temperature angular power spectrum from compressional velocity modes are presented. For tensor modes, in addition to the secondary temperature and polarization angular power spectra, we compute the TE cross-correlation signal and compare it with the existing WMAP data; we also recompute the primary signal for tensor modes. Whenever possible we compare our results with the results existing in the literature. In particular, Lewis lewis recently computed fully-numerically the vector and tensor primary and secondary temperature and polarization power spectra. We compare our semi-analytic results with this analysis and find good agreement. Seshadri and Subramanian seshadri1 computed the secondary temperature anisotropies from vector modes. Our results are in good agreement with their conclusion. Mack et al. mack computed primary signal from vector and tensor modes using the formalism we adopt in this paper. Our results are in disagreement with their results for $`\mathrm{}\stackrel{>}{}75`$, and we have given reasons for our disagreement in the discussion above. In addition to comparison with existing literature, we also give new results for secondary TE cross-correlation from tensor modes and secondary temperature angular power spectrum from scalar modes.
We discuss below the details of expected signal from each of the perturbation mode:
Vector modes: The secondary temperature and polarization signals from the vector modes is shown are Figure 1 and 2. The secondary temperature signal increases $`\mathrm{}^{2.4}`$ for $`\mathrm{}\stackrel{>}{}50`$ and reaches a value $`0.1(\mu \mathrm{k})^2`$ for $`\mathrm{}10^4`$, in agreement with the analysis of seshadri1 . For small $`\mathrm{}`$ the signal is very small ($`\stackrel{<}{}10^4(\mu \mathrm{k})^2`$) and for large $`\mathrm{}`$ the secondary signal is smaller than the primary signal (e.g. lewis ) and therefore it is unlikely that the signature of reionization could be detected in the vector-mode temperature anisotropies. The polarization signal, shown in Figure 2, is sourced by the free-streaming of dipole at the epoch of recombination. This signal dominates the primary signal for $`\mathrm{}\stackrel{<}{}10`$, but is several orders of magnitude smaller than the expected signal from tensor modes.
Scalar modes: We only compute the secondary temperature anisotropies from compressional velocity modes in this case. As seen in Figure 6, this contribution is several orders of magnitude smaller than the already-detected primary signal and therefore its effects are unlikely to be detectable.
Tensor modes: As seen from Figures 4 and 5, the most interesting CMBR anisotropy signal for $`\mathrm{}\stackrel{<}{}100`$ is from these modes. The secondary $`B`$-mode signal from tensor modes is detectable by future CMBR mission Planck surveyor for $`B_03\times 10^9\mathrm{G}`$ . The tensor TE cross correlation from primordial magnetic fields can explain the observed enhancement of the observed signal for $`\mathrm{}\stackrel{<}{}10`$ by WMAP for $`B_04.5\times 10^9\mathrm{G}`$ if the primordial magnetic fields are generated during the epoch of inflation. Assuming that tensor modes make a significant contribution to the observed enhancement, the bounds on the optical depth to the surface of reionization, $`\tau _{\mathrm{reion}}`$ are weaker by roughly a factor of $`\sqrt{2}`$. This hypothesis can be borne/ruled out by testing the Gaussianity of the signal for $`\mathrm{}\stackrel{<}{}10`$
## Acknowledgment
We would like to thank A. Lewis for prompt reply to our queries and to T. R. Seshadri and K. Subramanian for useful discussion.
## Appendix A
In this section, we briefly discuss the terminology and present the complete expressions for the vector and tensor power spectra $`\mathrm{\Pi }^V(k)`$ and $`\mathrm{\Pi }^T(k)`$ (mack ). The energy momentum tensor for magnetic fields for a single Fourier mode is a convolution of different Fourier modes and is given by:
$$T_{ij}(๐ค)=d^3q\left[\stackrel{~}{B}_i(๐ช)\stackrel{~}{B}_j(๐ค๐ช)\frac{1}{2}\delta _{ij}\stackrel{~}{B}_m(๐ช)\stackrel{~}{B}_m(๐ค๐ช)\right]$$
(27)
The energy-momentum tensor has non-vanishing scalar, vector, and tensor components. The vector and tensor components, in Fourier space, are defined as:
$`\mathrm{\Pi }_i^V=P_{ip}\widehat{k}_qT_{pq}`$ (28)
$`\mathrm{\Pi }_{ij}^T=\left(P_{ip}P_{jq}{\displaystyle \frac{1}{2}}P_{ij}P_{pq}\right)T_{pq}`$ (29)
Here $`P_{ij}=\delta _{ij}\widehat{k}_i\widehat{k}_j`$. The vector and tensor anisotropic stress are then defined as the two-point correlations of the above components as:
$`\mathrm{\Pi }_i^V(๐ค)\mathrm{\Pi }_i^V(๐ค^{^{}})2|\mathrm{\Pi }^{(V)}(k)|^2\delta (๐ค+๐ค^{^{}})`$ (30)
$`\mathrm{\Pi }_{ij}^T(๐ค)\mathrm{\Pi }_{ij}^T(๐ค^{^{}})4|\mathrm{\Pi }^{(T)}(k)|^2\delta (๐ค+๐ค^{^{}})`$ (31)
By evaluating the above correlations as also given in mack , we can arrive at the following approximate expression for the power spectra for $`n<3/2`$.
$`|\mathrm{\Pi }_V(k)|^2={\displaystyle \frac{A^2}{64\pi ^4(n+3)}}k^{2n+3}`$ (32)
$`|\mathrm{\Pi }_T(k)|^2={\displaystyle \frac{2A^2}{64\pi ^4(n+3)}}k^{2n+3}`$ (33)
Here, $`A`$ is the normalization of the magnetic power spectrum given in Eq. (2)
## Appendix B
Gravitational waves correspond to transverse,traceless perturbations to the metric: $`\delta g_{ij}=2a^2(\eta )h_{ij}`$ with $`h_{ii}=\widehat{k_i}h_{ij}=0`$. Since $`h_{ij}`$ is a stochastic variable we can define its power spectrum as:
$$h_{ij}(๐ค,\eta )h_{ij}(๐ค^{^{}},\eta )=4|h(k,\eta )|^2\delta (๐ค+๐ค^{^{}})$$
(34)
The evolution of $`h_{ij}`$ then follows from the tensor Einstein equation (see e.g. hu4 ) ,
$$\ddot{h}+2\frac{\dot{a}}{a}\dot{h}+k^2h=8\pi GS(k,\eta )$$
(35)
The source on the RHS is the tensor anisotropic stress of the plasma which is defined as: $`S(k,\eta )=\mathrm{\Pi }^T(k)/a^2`$ (Eq. (31)). We assume that the primordial magnetic fields are generated by some mechanism at a very early epoch $`\eta _{in}`$. It was recently shown by Lewis lewis that after the neutrino decoupling epoch $`\eta _{}`$ the neutrino start free-streaming and develop significant anisotropic stress which cancel the anisotropic stress of the primordial magnetic fields to the leading order for super-horizon modes, resulting in negligible net anisotropic stress in the plasma. We can thus assume that for $`\eta \eta _{}`$, $`S(๐ค,\eta )=0`$ and for $`\eta \eta _{}`$, $`S(๐ค,\eta )=\mathrm{\Pi }^T(k)/a^2`$ where $`\mathrm{\Pi }^T(k)`$ is the magnetic tensor anisotropic stress as defined in Eq. (31). We now derive the solutions to Eq. (35) in various regimes. The evolution of the scale factor $`a(\eta )`$ is given by the Friedmann equation:
$$\dot{a}^2=H_0^2(\mathrm{\Omega }_ma+\mathrm{\Omega }_\gamma +\mathrm{\Omega }_\nu +\mathrm{\Omega }_\mathrm{\Lambda }a^4)$$
(36)
Here, $`\mathrm{\Omega }_{m,\gamma ,\nu ,\mathrm{\Lambda }}`$ are the fractional densities in matter, radiation,neutrinos and cosmological constant respectively. Approximate solutions in the radiation-dominated and matter-dominated epoch are $`a(\eta )=2\sqrt{\frac{\mathrm{\Omega }_\gamma +\mathrm{\Omega }_\nu }{\mathrm{\Omega }_m}}\frac{\eta }{\eta _0}`$ and $`a(\eta )=\left(\frac{\eta }{\eta _0}\right)^2`$ respectively. Using the above form for the scale-factor we can rewrite Eq. (35) for $`\eta _{in}<\eta <\eta _{}`$ as:
$$\ddot{h}+\frac{2}{\eta }\dot{h}+k^2h=\frac{3R_\gamma \mathrm{\Pi }^T(k)}{\rho _\gamma }\frac{1}{\eta ^2}$$
(37)
Here, $`R_\gamma =\mathrm{\Omega }_\gamma /(\mathrm{\Omega }_\gamma +\mathrm{\Omega }_\nu )0.6`$. $`\rho _\gamma `$ is the CMBR energy density. Eq. (37) can be solved exactly using the Greenโs function technique to give mack :
$$h(k,\eta )=\frac{3R_\gamma \mathrm{\Pi }^T(k)}{\rho _\gamma }_{\eta _{in}}^\eta ๐\eta ^{^{}}\frac{\mathrm{sin}[k(\eta \eta ^{^{}})]}{\eta ^{^{}}}$$
(38)
For super-horizon modes $`k\eta 1`$, the above form can be simplified to give:
$$h(k,\eta )\frac{3R_\gamma \mathrm{\Pi }^T(k)}{\rho _\gamma }_{\eta _{in}}^\eta ๐\eta ^{^{}}\frac{k(\eta \eta ^{^{}})}{k\eta \eta ^{^{}}}=\frac{3R_\gamma \mathrm{\Pi }^T(k)}{\rho _\gamma }\mathrm{ln}\left(\frac{\eta }{\eta _{\mathrm{in}}}\right)$$
(39)
For $`\eta \eta _{}`$, the evolution of $`h`$ is given by the homogeneous solutions in the radiation and matter-dominated regimes:
$`h_{rad}(k,\eta )=A_1j_0(k\eta )`$ (40)
$`h_{mat}(k,\eta )=A_2{\displaystyle \frac{j_1(k\eta )}{k\eta }}`$ (41)
The coefficients $`A_1`$ and $`A_2`$ are determined by matching the super-horizon solutions at the two transitions $`\eta _{}`$ and $`\eta _{eq}`$. We thus get
$$A_2=3A_1=\frac{9R_\gamma \mathrm{\Pi }^T(k)}{\rho _\gamma }\mathrm{ln}\left(\frac{\eta _{}}{\eta _{\mathrm{in}}}\right)$$
(42)
Thus, the full expression for the matter-dominated solution can be written as:
$$\dot{h}_{mat}(\eta ,k)=\frac{9R_\gamma \mathrm{\Pi }^T(k)}{\rho _\gamma }\mathrm{ln}\left(\frac{\eta _{}}{\eta _{\mathrm{in}}}\right)\frac{j_2(k\eta )}{\eta }$$
(43)
This solution is used for solving tensor temperature and polarization primary and secondary anisotropies. Few assumptions have been made in deriving the above expression. Firstly, the transition between radiation dominated to matter-dominated region has been assumed to be instantaneous. This however does not affect the evolution of modes with wave-length greater than the width of transition $`k\eta _{\mathrm{eq}}\stackrel{<}{}1`$. Moreover, only super-horizon solutions have been used to match the solutions for $`h`$ at different transitions. These simplifications however do not affect the results quotes for small multipoles as discussed in the main section.
### Tight-coupling tensor quadrupole
In the tight-coupling regime, $`z\stackrel{>}{}1100`$, to lowest order in mean-free path, we have $`P^T=\dot{h}/(3\dot{\tau })`$ mack . We however use the expression accurate to the second order in mean-free path as is done for the scalar modes in zaldarriaga2 . Using the Boltzmann equation for the evolution of tensor modes we get the following equation for $`P^T(๐ค,\eta )`$ in the tight-coupling limit:
$$\dot{P}+\frac{3}{10}\dot{\tau }P=\frac{\dot{h}}{10}$$
(44)
The lowest order solution to this equation is obtained by neglecting the $`\dot{P}`$ in the equation, which gives, $`P=\dot{h}/(3\dot{\tau })`$. The above equation however can be solved exactly to give:
$$P(\eta )=_0^\eta ๐\eta {}_{}{}^{}\dot{h}e^{\frac{3}{10}[\tau (\eta ^{^{}})\tau (\eta )]}$$
(45)
We use the standard recombination history for computing $`\tau `$.
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# Filter Dimension
## 1 Introduction
Throughout the paper, $`K`$ is a field, a module $`M`$ over an algebra $`A`$ means a left module denoted $`{}_{A}{}^{}M`$, $`=_K`$.
Intuitively, the filter dimension of an algebra or a module measures how โcloseโ standard filtrations of the algebra or the module are. In particular, for a simple algebra it also measures the growth of how โfastโ one can prove that the algebra is simple.
The filter dimension appears naturally when one wants to generalize the Bernsteinโs inequality for the Weyl algebras to the class of simple finitely generated algebras.
The $`n`$โth Weyl algebra $`A_n`$ over the field $`K`$ has $`2n`$ generators $`X_1,\mathrm{},X_n`$, $`_1,\mathrm{},_n`$ that satisfy the defining relations
$$_iX_jX_j_i=\delta {}_{ij}{}^{},\text{the Kronecker delta},X_iX_jX_jX_i=_i_j_j_i=0,$$
for all $`i,j=1,\mathrm{},n`$. When char $`K=0`$ the Weyl algebra $`A_n`$ is a simple Noetherian finitely generated algebra canonically isomorphic to the ring of differential operators $`K[X_1,\mathrm{},X_n,\frac{}{X_1},\mathrm{},\frac{}{X_n}]`$ with polynomial coefficients ($`X_iX_i,_i\frac{_i}{X_i}`$, $`i=1,\mathrm{},n`$).
Let $`\mathrm{K}.\mathrm{dim}`$ and $`\mathrm{GK}`$ be the (left) Krull (in the sense of Rentschler and Gabriel, ) and the Gelfand-Kirillov dimension respectively.
###### Theorem 1.1
(The Bernsteinโs inequality, ) Let $`A_n`$ be the $`n`$โth Weyl algebra over a field of characteristic zero. Then $`\mathrm{GK}(M)n`$ for all nonzero finitely generated $`A_n`$-modules $`M`$.
Let $`A`$ be a simple finitely generated infinite dimensional $`K`$-algebra. Then $`\mathrm{dim}_K(M)=\mathrm{}`$ for all nonzero $`A`$-modules $`M`$ (the algebra $`A`$ is simple, so the $`K`$-linear map $`A\mathrm{Hom}_K(M,M)`$, $`a(mam)`$, is injective, and so $`\mathrm{}=\mathrm{dim}_K(A)\mathrm{dim}_K(\mathrm{Hom}_K(M,M))`$ hence $`\mathrm{dim}_K(M)=\mathrm{}`$). So, the Gelfand-Kirillov dimension (over $`K`$) $`\mathrm{GK}(M)1`$ for all nonzero $`A`$-modules $`M`$.
Definition. $`h_A:=inf\{\mathrm{GK}(M)|M`$ is a nonzero finitely generated $`A`$-module$`\}`$ is called the holonomic number for the algebra $`A`$.
Problem. For a simple finitely generated algebra find its holonomic number.
To find an approximation of the holonomic number for simple finitely generated algebras and to generalize the Bernsteinโs inequality for these algebras was a main motivation for introducing the filter dimension, . In this paper $`d`$ stands for the filter dimension $`\mathrm{fd}`$ or the left filter dimension $`\mathrm{lfd}`$. The following two inequalities are central for the proofs of almost all results in this paper.
* The First Filter Inequality, . Let $`A`$ be a simple finitely generated algebra. Then
$$\mathrm{GK}(M)\frac{\mathrm{GK}(A)}{\mathrm{d}(A)+\mathrm{max}\{\mathrm{d}(A),1\}}$$
for all nonzero finitely generated $`A`$-modules $`M`$.
* The Second Filter Inequality, . Under a certain mild conditions (Theorem 4.2) the (left) Krull dimension of the algebra $`A`$ satisfies the following inequality
$$\mathrm{K}.\mathrm{dim}(A)\mathrm{GK}(A)(1\frac{1}{\mathrm{d}(A)+\mathrm{max}\{\mathrm{d}(A),1\}}).$$
The paper is organized as follows. Both filter dimensions are introduced in Section 2. In Sections 3 and 4 the first and the second filter inequalities are proved respectively. In Section 4 we use both filter inequalities for giving short proofs of some classical results about the rings $`๐(X)`$ of differential operators on smooth irreducible affine algebraic varieties. The (left) filter dimension of $`๐(X)`$ is $`1`$ (Section 5). A concept of multiplicity for the filter dimension and a concept of holonomic module for (simple) finitely generated algebras appear in Section 6. Every holonomic module has finite length (Theorem 6.8). In Section 7 an upper bound is given $`(i)`$ for the Gelfand-Kirillov dimension of commutative subalgebras of simple finitely generated infinite dimensional algebras (Theorem 7.2), and $`(ii)`$ for the transcendence degree of subfields of quotient rings of (certain) simple finitely generated infinite dimensional algebras (Theorems 7.4 and 7.5). In Section 8 a similar upper bound is obtained for the Gelfand-Kirillov dimension of isotropic subalgebras of strongly simple Poisson algebras (Theorem 8.1).
## 2 Filter dimension of algebras and modules
In this section, the filter dimension of algebras and modules will be defined.
The Gelfand-Kirillov dimension. Let $``$ be the set of all functions from the set of natural numbers $`=\{0,1,\mathrm{}\}`$ to itself. For each function $`f`$, the non-negative real number or $`\mathrm{}`$ defined as
$$\gamma (f):=inf\{r|f(i)i^r\mathrm{for}i0\}$$
is called the degree of $`f`$. The function $`f`$ has polynomial growth if $`\gamma (f)<\mathrm{}`$. Let $`f,g,p`$, and $`p(i)=p^{}(i)`$ for $`i0`$ where $`p^{}(t)[t]`$ (a polynomial algebra with coefficients from the field of rational numbers). Then
$`\gamma (f+g)\mathrm{max}\{\gamma (f),\gamma (g)\},`$ $`\gamma (fg)\gamma (f)+\gamma (g),`$
$`\gamma (p)=\mathrm{deg}_t(p^{}(t)),`$ $`\gamma (pg)=\gamma (p)+\gamma (g).`$
Let $`A=Ka_1,\mathrm{},a_s`$ be a finitely generated $`K`$-algebra. The finite dimensional filtration $`F=\{A_i\}`$ associated with algebra generators $`a_1,\mathrm{},a_s`$:
$$A_0:=KA_1:=K+\underset{i=1}{\overset{s}{}}Ka_i\mathrm{}A_i:=A_1^i\mathrm{}$$
is called the standard filtration for the algebra $`A`$. Let $`M=AM_0`$ be a finitely generated $`A`$-module where $`M_0`$ is a finite dimensional generating subspace. The finite dimensional filtration $`\{M_i:=A_iM_0\}`$ is called the standard filtration for the $`A`$-module $`M`$.
Definition. $`\mathrm{GK}(A):=\gamma (i\mathrm{dim}_K(A_i))`$ and $`\mathrm{GK}(M):=\gamma (i\mathrm{dim}_K(M_i))`$ are called the Gelfand-Kirillov dimensions of the algebra $`A`$ and the $`A`$-module $`M`$ respectively.
It is easy to prove that the Gelfand-Kirillov dimension of the algebra (resp. the module) does not depend on the choice of the standard filtration of the algebra (resp. and the choice of the generating subspace of the module).
The return functions and the (left) filter dimension.
Definition . The function $`\nu _{F,M_0}:\{\mathrm{}\}`$,
$$\nu _{F,M_0}(i):=\mathrm{min}\{j\{\mathrm{}\}:A_jM_{i,gen}M_0\mathrm{for}\mathrm{all}M_{i,gen}\}$$
is called the return function of the $`A`$-module $`M`$ associated with the filtration $`F=\{A_i\}`$ of the algebra $`A`$ and the generating subspace $`M_0`$ of the $`A`$-module $`M`$ where $`M_{i,gen}`$ runs through all generating subspaces for the $`A`$-module $`M`$ such that $`M_{i,gen}M_i`$.
Suppose, in addition, that the finitely generated algebra $`A`$ is a simple algebra. The return function $`\nu _F`$ and the left return function $`\lambda _F`$ for the algebra $`A`$ with respect to the standard filtration $`F:=\{A_i\}`$ for the algebra $`A`$ are defined by the rules:
$`\nu _F(i)`$ $`:=`$ $`\mathrm{min}\{j\{\mathrm{}\}|\mathrm{\hspace{0.17em}\hspace{0.17em}1}A_jaA_j\mathrm{for}\mathrm{all}\mathrm{\hspace{0.33em}\hspace{0.17em}0}aA_i\},`$
$`\lambda _F(i)`$ $`:=`$ $`\mathrm{min}\{j\{\mathrm{}\}|\mathrm{\hspace{0.17em}1}AaA_j\mathrm{for}\mathrm{all}\mathrm{\hspace{0.33em}\hspace{0.33em}0}aA_i\},`$
where $`A_jaA_j`$ is the vector subspace of the algebra $`A`$ spanned over the field $`K`$ by the elements $`xay`$ for all $`x,yA_j`$; and $`AaA_j`$ is the left ideal of the algebra $`A`$ generated by the set $`aA_j`$. The next result shows that under a mild restriction the maps $`\nu _F(i)`$ and $`\lambda _F(i)`$ are finite.
Recall that the centre of a simple algebra is a field.
###### Lemma 2.1
Let $`A`$ be a simple finitely generated algebra such that its centre $`Z(A)`$ is an algebraic field extension of $`K`$. Then $`\lambda _F(i)\nu _F(i)<\mathrm{}`$ for all $`i0`$.
Proof. The first inequality is evident.
The centre $`Z=Z(A)`$ of the simple algebra $`A`$ is a field that contains $`K`$. Let $`\{\omega _j|jJ\}`$ be a $`K`$-basis for the $`K`$-vector space $`Z`$. Since $`\mathrm{dim}_K(A_i)<\mathrm{}`$, one can find finitely many $`Z`$-linearly independent elements, say $`a_1,\mathrm{},a_s`$, of $`A_i`$ such that $`A_iZa_1+\mathrm{}+Za_s`$. Next, one can find a finite subset, say $`J^{}`$, of $`J`$ such that $`A_iVa_1+\mathrm{}+Va_s`$ where $`V=_{jJ^{}}K\omega _j`$. The field $`K^{}`$ generated over $`K`$ by the elements $`\omega _j`$, $`jJ^{}`$, is a finite field extension of $`K`$ (i.e. $`\mathrm{dim}_K(K^{})<\mathrm{}`$) since $`Z/K`$ is algebraic, hence $`K^{}A_n`$ for some $`n0`$. Clearly, $`A_iK^{}a_1+\mathrm{}+K^{}a_s`$.
The $`A`$-bimodule $`{}_{A}{}^{}A_{A}^{}`$ is simple with ring of endomorphisms $`\mathrm{End}({}_{A}{}^{}A_{A}^{})Z`$. By the Density Theorem, , 12.2, for each integer $`1js`$, there exist elements of the algebra $`A`$, say $`x_1^j,\mathrm{},x_m^j,y_1^j,\mathrm{},y_m^j`$, $`m=m(j)`$, such that for all $`1ls`$
$$\underset{k=1}{\overset{m}{}}x_k^ja_ly_k^j=\delta _{j,l},\mathrm{the}\mathrm{Kronecker}\mathrm{delta}.$$
Let us fix a natural number, say $`d=d_i`$, such that $`A_d`$ contains all the elements $`x_k^j`$, $`y_k^j`$, and the field $`K^{}`$. We claim that $`\nu _F(i)2d`$. Let $`0aA_i`$. Then $`a=\lambda _1a_1+\mathrm{}+\lambda _sa_s`$ for some $`\lambda _iK^{}`$. There exists $`\lambda _j0`$. Then $`_{k=1}^m\lambda _j^1x_k^ja_jy_k^j=1`$, and $`\lambda _j^1x_k^j,y_k^jA_{2d}`$. This proves the claim and the lemma. $`\mathrm{}`$
Remark. If the field $`K`$ is uncountable then automatically the centre $`Z(A)`$ of a simple finitely generated algebra $`A`$ is algebraic over $`K`$ (since $`A`$ has a countable $`K`$-basis and the rational function field $`K(x)`$ has uncountable basis over $`K`$ since elements $`\frac{1}{x+\lambda }`$, $`\lambda K`$, are $`K`$-linearly independent).
It is easy to see that for a finitely generated algebra $`A`$ any two standard finite dimensional filtrations $`F=\{A_i\}`$ and $`G=\{B_i\}`$ are equivalent, $`(FG),`$ that is, there exist natural numbers $`a,b,c,d`$ such that
$$A_iB_{ai+b}\text{and}B_iA_{ci+d}\text{for}i0.$$
If one of the inclusions holds, say the first, we write $`FG`$.
###### Lemma 2.2
Let $`A`$ be a finitely generated algebra equipped with two standard finite dimensional filtrations $`F=\{A_i\}`$ and $`G=\{B_i\}.`$
1. Let $`M`$ be a finitely generated $`A`$-module. Then $`\gamma (\nu _{F,M_0})=\gamma (\nu _{G,N_0})`$ for any finite dimensional generating subspaces $`M_0`$ and $`N_0`$ of the $`A`$-module $`M`$.
2. If, in addition, $`A`$ is a simple algebra then $`\gamma (\nu _F)=\gamma (\nu _G)`$ and $`\gamma (\lambda _F)=\gamma (\lambda _G)`$.
Proof. 1. The module $`M`$ has two standard finite dimensional filtrations $`\{M_i=A_iM_0\}`$ and $`\{N_i=B_iN_0\}`$. Let $`\nu =\nu _{F,M_0}`$ and $`\mu =\nu _{G,N_0}`$.
Suppose that $`F=G`$. Choose a natural number $`s`$ such that $`M_0N_s`$ and $`N_0M_s`$, so $`N_iM_{i+s}`$ and $`M_iN_{i+s}`$ for all $`i0`$. Let $`N_{i,gen}`$ be any generating subspace for the $`A`$-module $`M`$ such that $`N_{i,gen}N_i`$. Since $`M_0A_{\nu (i+s)}N_{i,gen}`$ for all $`i0`$ and $`N_0A_sM_0,`$ we have $`N_0A_{\nu (i+s)+s}N_{i,gen},`$ hence, $`\mu (i)\nu (i+s)+s`$ and finally $`\gamma (\mu )\gamma (\nu ).`$ By symmetry, the opposite inequality is true and so $`\gamma (\mu )=\gamma (\nu )`$.
Suppose that $`M_0=N_0`$. The algebra $`A`$ is a finitely generated algebra, so all standard finite dimensional filtrations of the algebra $`A`$ are equivalent. In particular, $`FG`$ and so one can choose natural numbers $`a,b,c,d`$ such that
$$A_iB_{ai+b}\text{and}B_iA_{ci+d}\text{for}i0.$$
Then $`N_i=B_iN_0A_{ci+d}M_0=M_{ci+d}`$ for all $`i0`$, hence $`N_0=M_0A_{\nu (ci+d)}N_{i,gen}B_{a\nu (ci+d)+b}N_{i,gen}`$, therefore $`\mu (i)a\nu (ci+d)+b`$ for all $`i0`$, hence $`\gamma (\mu )\gamma (\nu )`$. By symmetry, we get the opposite inequality which implies $`\gamma (\mu )=\gamma (\nu )`$. Now, $`\gamma (\nu _{F,M_0})=\gamma (\nu _{F,N_0})=\gamma (\nu _{G,N_0})`$.
2. The algebra $`A`$ is simple, equivalently, it is a simple (left) $`AA^0`$-module where $`A^0`$ is the opposite algebra to $`A`$. The opposite algebra has the standard filtration $`F^0=\{A_i^0\}`$, opposite to the filtration $`F`$. The tensor product of algebras $`AA^0`$, so-called, the enveloping algebra of $`A`$, has the standard filtration $`FF^0=\{C_n\}`$ which is the tensor product of the standard filtrations $`F`$ and $`F^0`$, that is, $`C_n=\{A_iA_j^0,i+jn\}`$. Let $`\nu _{FF^0,K}`$ be the return function of the $`AA^0`$-module $`A`$ associated with the filtration $`FF^0`$ and the generating subspace $`K`$. Then
$$\nu _F(i)\nu _{FF^0,K}(i)2\nu _F(i)\text{for all}i0,$$
and so
$$\gamma (\nu _F)=\gamma (\nu _{FF^0,K}),$$
(1)
and, by the first statement, we have $`\gamma (\nu _F)=\gamma (\nu _{FF^0,K})=\gamma (\nu _{GG^0,K})=\gamma (\nu _G)`$, as required. Using a similar argument as in the proof of the first statement one can proof that $`\gamma (\lambda _F)=\gamma (\lambda _G)`$. We leave this as an exercise. $`\mathrm{}`$
Definition . $`\mathrm{fd}(M)=\gamma (\nu _{F,M_0})`$ is the filter dimension of the $`A`$-module $`M`$, and $`\mathrm{fd}(A):=\mathrm{fd}({}_{AA^0}{}^{}A)`$ is the filter dimension of the algebra $`A`$. If, in addition, the algebra $`A`$ is simple, then $`\mathrm{fd}(A)=\gamma (\nu _F)`$, and $`\mathrm{lfd}(A):=\gamma (\lambda _F)`$ is called the left filter dimension of the algebra $`A`$.
By the previous lemma the definitions make sense (both filter dimensions do not depend on the choice of the standard filtration $`F`$ for the algebra $`A`$).
By Lemma 2.1, $`\mathrm{lfd}(A)\mathrm{fd}(A)`$.
Question. What is the filter dimension of a polynomial algebra?
## 3 The first filter inequality
In this paper, $`\mathrm{d}(A)`$ means either the filter dimension $`\mathrm{fd}(A)`$ or the left filter dimension $`\mathrm{lfd}(A)`$ of a simple finitely generated algebra $`A`$ (i.e. $`\mathrm{d}=\mathrm{fd},\mathrm{lfd}`$). Both filter dimensions appear naturally when one tries to find a lower bound for the holonomic number (Theorem 3.1) and an upper bound (Theorem 4.2) for the (left and right) Krull dimension (in the sense of Rentschler-Gabriel, ) of simple finitely generated algebras.
The next theorem is a generalization of the Bernsteinโs Inequality (Theorem 1.1) to the class of simple finitely generated algebras.
###### Theorem 3.1
(The First Filter Inequality, ) Let $`A`$ be a simple finitely generated algebra. Then
$$\mathrm{GK}(M)\frac{\mathrm{GK}(A)}{\mathrm{d}(A)+\mathrm{max}\{\mathrm{d}(A),1\}}$$
for all nonzero finitely generated $`A`$-modules $`M`$ where $`\mathrm{d}=\mathrm{fd},\mathrm{lfd}`$.
Proof. Let $`\lambda =\lambda _F`$ be the left return function associated with a standard filtration $`F`$ of the algebra $`A`$ and let $`0aA_i`$. It suffices to prove the inequality for $`\lambda `$ (since $`\mathrm{fd}(A)\mathrm{lfd}(A)`$). It follows from the inclusion
$$AaM_{\lambda (i)}=AaA_{\lambda (i)}M_01M_0=M_0$$
that the linear map
$$A_i\mathrm{Hom}(M_{\lambda (i)},M_{\lambda (i)+i}),a(mam),$$
is injective, so dim $`A_i`$ dim $`M_{\lambda (i)}`$ dim $`M_{\lambda (i)+i}`$. Using the above elementary properties of the degree (see also , 8.1.7), we have
$`\mathrm{GK}(A)`$ $`=`$ $`\gamma (\mathrm{dim}A_i)\gamma (\mathrm{dim}M_{\lambda (i)})+\gamma (\mathrm{dim}M_{\lambda (i)+i})`$
$``$ $`\gamma (\mathrm{dim}M_i)\gamma (\lambda )+\gamma (\mathrm{dim}M_i)\mathrm{max}\{\gamma (\lambda ),1\}`$
$`=`$ $`\mathrm{GK}(M)(\mathrm{lfd}A+\mathrm{max}\{\mathrm{lfd}A,1\})`$
$``$ $`\mathrm{GK}(M)(\mathrm{lfd}A+\mathrm{max}\{\mathrm{lfd}A,1\}).\mathrm{}`$
The result above gives a lower bound for the holonomic number of a simple finitely generated algebra
$$h_A\frac{\mathrm{GK}(A)}{\mathrm{d}(A)+\mathrm{max}\{\mathrm{d}(A),1\}}.$$
###### Theorem 3.2
Let $`A`$ be a finitely generated algebra. Then
$$\mathrm{GK}(M)\mathrm{GK}(A)\mathrm{fd}(M)$$
for any simple $`A`$-module $`M`$.
Proof. Let $`\nu =\nu _{F,Km}`$ be the return function of the module $`M`$ associated with a standard finite dimensional filtration $`F=\{A_i\}`$ of the algebra $`A`$ and a fixed nonzero element $`mM`$. Let $`\pi :MK`$ be a non-zero linear map satisfying $`\pi (m)=1`$. Then, for any $`i0`$ and any $`0uM_i`$: $`1=\pi (m)\pi (A_{\nu (i)}u)`$, and so the linear map
$$M_i\mathrm{Hom}(A_{\nu (i)},K),u(a\pi (au)),$$
is an injective map hence dim $`M_i`$ dim $`A_{\nu (i)}`$ and finally $`\mathrm{GK}(M)\mathrm{GK}(A)\mathrm{fd}(M)`$. $`\mathrm{}`$
###### Corollary 3.3
Let $`A`$ be a simple finitely generated infinite dimensional algebra. Then
$$\mathrm{fd}(A)\frac{1}{2}.$$
Proof. The algebra $`A`$ is a finitely generated infinite dimensional algebra hence $`\mathrm{GK}(A)>0`$. Clearly, $`\mathrm{GK}(AA^0)\mathrm{GK}(A)+\mathrm{GK}(A^0)=2\mathrm{G}\mathrm{K}(A)`$. Applying Theorem 3.2 to the simple $`AA^0`$-module $`M=A`$ we finish the proof:
$$\mathrm{GK}(A)=\mathrm{GK}({}_{AA^0}{}^{}A)\mathrm{GK}(AA^0)\mathrm{fd}({}_{AA^0}{}^{}A)2\mathrm{G}\mathrm{K}(A)\mathrm{fd}(A)$$
hence $`\mathrm{fd}(A)\frac{1}{2}`$. $`\mathrm{}`$
Question. Is $`\mathrm{fd}(A)1`$ for all simple finitely generated infinite dimensional algebras $`A`$?
Question. For which numbers $`d\frac{1}{2}`$ there exists a simple finitely generated infinite dimensional algebra $`A`$ with $`\mathrm{fd}(A)=d`$?
###### Corollary 3.4
Let $`A`$ be a simple finitely generated infinite dimensional algebra. Then
$$\mathrm{fd}(M)\frac{1}{\mathrm{fd}(A)+\mathrm{max}\{\mathrm{fd}(A),1\}}$$
for all simple $`A`$-modules $`M`$.
Proof. Applying Theorem 3.1 and Theorem 3.2, we have the result
$$\mathrm{fd}(M)\frac{\mathrm{GK}(M)}{\mathrm{GK}(A)}\frac{\mathrm{GK}(A)}{\mathrm{GK}(A)(\mathrm{fd}(A)+\mathrm{max}\{\mathrm{fd}(A),1\})}=\frac{1}{\mathrm{fd}(A)+\mathrm{max}\{\mathrm{fd}(A),1\}}.\mathrm{}$$
## 4 Krull, Gelfand-Kirillov and filter dimensions of simple finitely generated algebras
In this section, we prove the second filter inequality (Theorem 4.2) and apply both filter inequalities for giving short proofs of some classical results about the rings of differential operators on a smooth irreducible affine algebraic varieties (Theorems 1.1, 4.4, 4.5, 4.7).
We say that an algebra $`A`$ is (left) finitely partitive (, 8.3.17) if, given any finitely generated $`A`$-module $`M`$, there is an integer $`n=n(M)>0`$ such that for every strictly descending chain of $`A`$-submodules of $`M`$:
$$M=M_0M_1\mathrm{}M_m$$
with $`\mathrm{GK}(M_i/M_{i+1})=\mathrm{GK}(M)`$, one has $`mn`$. McConnell and Robson write in their book , 8.3.17, that โyet no examples are known which fail to have this property.
Recall that $`\mathrm{K}.\mathrm{dim}`$ denotes the (left) Krull dimension in the sense of Rentschler and Gabriel, .
###### Lemma 4.1
Let $`A`$ be a finitely partitive algebra with $`\mathrm{GK}(A)<\mathrm{}`$. Let $`a`$, $`b0`$ and suppose that $`\mathrm{GK}(M)a+b`$ for all finitely generated $`A`$-modules $`M`$ with $`\mathrm{K}.\mathrm{dim}(M)=a`$, and that $`\mathrm{GK}(N)`$ for all finitely generated $`A`$-modules $`N`$ with $`\mathrm{K}.\mathrm{dim}(N)a`$. Then $`\mathrm{GK}(M)\mathrm{K}.\mathrm{dim}(M)+b`$ for all finitely generated $`A`$-modules $`M`$ with $`\mathrm{K}.\mathrm{dim}(M)a`$. In particular, $`\mathrm{GK}(A)\mathrm{K}.\mathrm{dim}(A)+b`$.
Remark. It is assumed that a module $`M`$ with $`\mathrm{K}.\mathrm{dim}(M)=a`$ exists.
Proof. We use induction on $`n=\mathrm{K}.\mathrm{dim}(M)`$. The base of induction, $`n=a`$, is true. Let $`n>a`$. There exists a descending chain of submodules $`M=M_1M_2\mathrm{}`$ with $`\mathrm{K}.\mathrm{dim}(M_i/M_{i+1})=n1`$ for $`i1`$. By induction, $`\mathrm{GK}(M_i/M_{i+1})n1+b`$ for $`i1`$. The algebra $`A`$ is finitely partitive, so there exists $`i`$ such that $`\mathrm{GK}(M)>\mathrm{GK}(M_i/M_{i+1})`$, so $`\mathrm{GK}(M)1\mathrm{GK}(M_i/M_{i+1})n1+b`$, since $`\mathrm{GK}(M)`$, hence $`\mathrm{GK}(M)\mathrm{K}.\mathrm{dim}(M)+b`$. Since $`\mathrm{K}.\mathrm{dim}(A)\mathrm{K}.\mathrm{dim}(M)`$ for all finitely generated $`A`$-modules $`M`$ we have $`\mathrm{GK}(A)\mathrm{K}.\mathrm{dim}(A)+b`$. $`\mathrm{}`$
###### Theorem 4.2
() Let $`A`$ be a simple finitely generated finitely partitive algebra with $`\mathrm{GK}(A)<\mathrm{}`$. Suppose that the Gelfand-Kirillov dimension of every finitely generated $`A`$-module is a natural number. Then
$$\mathrm{K}.\mathrm{dim}(M)\mathrm{GK}(M)\frac{\mathrm{GK}(A)}{\mathrm{d}(A)+\mathrm{max}\{\mathrm{d}(A),1\}}$$
for any nonzero finitely generated $`A`$-module $`M`$. In particular,
$$\mathrm{K}.\mathrm{dim}(A)\mathrm{GK}(A)(1\frac{1}{\mathrm{d}(A)+\mathrm{max}\{\mathrm{d}(A),1\}}).$$
Proof. Applying the lemma above to the family of finitely generated $`A`$-modules of Krull dimension $`0`$, by Theorem 3.1, we can put $`a=0`$ and
$$b=\frac{\mathrm{GK}(A)}{\mathrm{d}(A)+\mathrm{max}\{\mathrm{d}(A),1\}},$$
and the result follows. $`\mathrm{}`$
Let $`K`$ be a field of characteristic zero and $`B`$ be a commutative $`K`$-algebra. The ring of ($`K`$-linear) differential operators $`๐(B)`$ on $`B`$ is defined as $`๐(B)=_{i=0}^{\mathrm{}}๐_i(B)`$ where $`๐_0(B)=\mathrm{End}_R(B)B`$, ($`(xbx)b`$),
$$๐_i(B)=\{u\mathrm{End}_K(B):[u,r]๐_{i1}(B)\mathrm{for}\mathrm{each}rB\}.$$
Note that the $`\{๐_i(B)\}`$ is, so-called, the order filtration for the algebra $`๐(B)`$:
$$๐_0(B)๐_1(B)\mathrm{}๐_i(B)\mathrm{}\mathrm{and}๐_i(B)๐_j(B)๐_{i+j}(B),i,j0.$$
The subalgebra $`\mathrm{\Delta }(B)`$ of $`\mathrm{End}_K(B)`$ generated by $`B\mathrm{End}_R(B)`$ and by the set $`\mathrm{Der}_K(B)`$ of all $`K`$-derivations of $`B`$ is called the derivation ring of $`B`$. The derivation ring $`\mathrm{\Delta }(B)`$ is a subring of $`๐(B)`$.
Let the finitely generated algebra $`B`$ be a regular commutative domain of Krull dimension $`n<\mathrm{}`$. In geometric terms, $`B`$ is the coordinate ring $`๐ช(X)`$ of a smooth irreducible affine algebraic variety $`X`$ of dimension $`n`$. Then
* $`\mathrm{Der}_K(B)`$ is a finitely generated projective $`B`$-module of rank $`n`$;
* $`๐(B)=\mathrm{\Delta }(B)`$;
* $`๐(B)`$ is a simple (left and right) Noetherian domain with $`\mathrm{GK}๐(B)=2n`$ ($`n=\mathrm{GK}(B)=\mathrm{K}.\mathrm{dim}(B))`$;
* $`๐(B)=\mathrm{\Delta }(B)`$ is an almost centralizing extension of $`B`$;
* the associated graded ring $`\mathrm{gr}๐(B)=๐_i(B)/๐_{i1}(B)`$ is a commutative domain;
* the Gelfand-Kirillov dimension of every finitely generated $`๐(B)`$-module is a natural number.
For the proofs of the statements above the reader is referred to , Chapter 15. So, the domain $`๐(B)`$ is a simple finitely generated infinite dimensional Noetherian algebra (, Chapter 15).
Example. Let $`P_n=K[X_1,\mathrm{},X_n]`$ be a polynomial algebra. $`\mathrm{Der}_K(P_n)=_{i=1}^nP_n\frac{}{X_i}`$,
$$๐(P_n)=\mathrm{\Delta }(P_n)=K[X_1,\mathrm{},X_n,\frac{}{X_1},\mathrm{},\frac{}{X_n}]$$
is the ring of differential operators with polynomial coefficients, i.e. the $`n`$โth Weyl algebra $`A_n`$.
In Section 5, we prove the following result.
###### Theorem 4.3
() The filter dimension and the left filter dimension of the ring of differential operators $`๐(B)`$ are both equal to $`1`$.
As an application we compute the Krull dimension of $`๐(B)`$.
###### Theorem 4.4
(, Ch. 15)
$$\mathrm{K}.\mathrm{dim}๐(B)=\frac{\mathrm{GK}(๐(B))}{2}=\mathrm{K}.\mathrm{dim}(B).$$
Proof. The second equality is clear $`(\mathrm{GK}(๐(B))=2\mathrm{G}\mathrm{K}(B)=2\mathrm{K}.\mathrm{dim}(B))`$. It follows from Theorems 4.2 and 4.3 that
$$\mathrm{K}.\mathrm{dim}๐(B)\frac{\mathrm{GK}(๐(B))}{2}=\mathrm{K}.\mathrm{dim}(B).$$
The map $`I๐(B)I`$ from the set of left ideals of $`B`$ to the set of left ideals of $`๐(B)`$ is injective, thus $`\mathrm{K}.\mathrm{dim}(B)\mathrm{K}.\mathrm{dim}๐(B)`$. $`\mathrm{}`$
This result shows that for the ring of differential operators on a smooth irreducible affine algebraic variety the inequality in Theorem 4.2 is the equality.
###### Theorem 4.5
(, 15.4.3) Let $`M`$ be a nonzero finitely generated $`๐(B)`$-module. Then
$$\mathrm{GK}(M)\frac{\mathrm{GK}(๐(B))}{2}=\mathrm{K}.\mathrm{dim}(B).$$
Proof. By Theorems 3.1 and 4.3,
$$\mathrm{GK}(M)\frac{\mathrm{GK}(๐(B))}{1+1}=\frac{2\mathrm{G}\mathrm{K}(B)}{2}=\mathrm{GK}(B)=\mathrm{K}.\mathrm{dim}(B).\mathrm{}$$
So, for the ring of differential operators on a smooth affine algebraic variety the inequality in Theorem 3.1 is in fact the equality.
In general, it is difficult to find the exact value for the filter dimension but for the Weyl algebra $`A_n`$ it is easy and one can find it directly.
###### Theorem 4.6
Both the filter dimension and the left filter dimension of the Weyl algebra $`A_n`$ over a field of characteristic zero are equal to $`1`$.
Proof. Denote by $`a_1,\mathrm{},a_{2n}`$ the canonical generators of the Weyl algebra $`A_n`$ and denote by $`F=\{A_{n,i}\}_{i0}`$ the standard filtration associated with the canonical generators. The associated graded algebra $`\mathrm{gr}A_n:=_{i0}A_{n,i}/A_{n,i1}`$, $`(A_{n,1}=0)`$ is a polynomial algebra in $`2n`$ variables, so
$$\mathrm{GK}(A_n)=\mathrm{GK}(\mathrm{gr}A_n)=2n.$$
For every $`i0`$:
$$\mathrm{ad}a_j:A_{n,i}A_{n,i1},x\mathrm{ad}a_j(x):=a_jxxa_j.$$
The algebra $`A_n`$ is central ($`Z(A_n)=K`$), so
$$\mathrm{ad}a_j(x)=0\mathrm{for}\mathrm{all}j=1,\mathrm{},2nxZ(A_n)=K=A_{n,0}.$$
These two facts imply $`\nu _F(i)i`$ for $`i0`$, and so $`\mathrm{d}(A_n)1`$.
The $`A_n`$-module $`P_n:=K[X_1,\mathrm{},X_n]A_n/(A_n_1+\mathrm{}+A_n_n)`$ has Gelfand-Kirillov dimension $`n`$. By Theorem 3.1 applied to the $`A_n`$-module $`P_n`$, we have
$$2n=\mathrm{GK}(A_n)n(\mathrm{d}(A)+\mathrm{max}\{\mathrm{d}(A),1\}),$$
hence $`\mathrm{d}(A_n)1`$, and so $`\mathrm{d}(A_n)=1`$. $`\mathrm{}`$
Proof of the Bernsteinโs inequality (Theorem 1.1).
Since $`\mathrm{GK}(A_n)=2n`$ and $`\mathrm{d}(A_n)=1`$, Theorem 3.1 gives $`\mathrm{GK}(M)\frac{2n}{2}=n`$. $`\mathrm{}`$
One also gets a short proof of the following result of Rentschler and Gabriel.
###### Theorem 4.7
(). If char$`K=0`$ then the Krull dimension of the Weyl algebra $`A_n`$ is
$$\mathrm{K}.\mathrm{dim}(A_n)=n.$$
Proof. Putting $`\mathrm{GK}(A_n)=2n`$ and $`\mathrm{d}(A_n)=1`$ into the second formula of Theorem 4.2 we have $`\mathrm{K}.\mathrm{dim}(A_n)\frac{2n}{2}=n`$. The polynomial algebra $`P_n=K[X_1,\mathrm{},X_n]`$ is the subalgebra of $`A_n`$ such that $`A_n`$ is a free right $`P_n`$-module. The map $`IA_nI`$ from the set of left ideals of the polynomial algebra $`P_n`$ to the set of left ideals of the Weyl algebra $`A_n`$ is injective, thus $`n=\mathrm{K}.\mathrm{dim}(P_n)\mathrm{K}.\mathrm{dim}(A_n)`$, and so $`\mathrm{K}.\mathrm{dim}(A_n)=n`$. $`\mathrm{}`$
## 5 Filter dimension of the ring of differential operators on a smooth irreducible affine algebraic variety (proof of Theorem 4.3)
Let $`K`$ be a field of characteristic $`0`$ and let the algebra $`B`$ be as in the previous section, i.e. $`B`$ is a finitely generated regular commutative algebra which is a domain. We keep the notations of the previous section. Recall that the derivation ring $`\mathrm{\Delta }=\mathrm{\Delta }(B)`$ coincides with the ring of differential operators $`๐(B)`$ (, 15.5.6) and is a simple finitely generated finitely partitive $`K`$-algebra (, 15.3.8, 15.1.21). We refer the reader to , Chapter 15, for basic definitions. We aim to prove Theorem 4.3.
Let $`\{B_i\}`$ and $`\{\mathrm{\Delta }_i\}`$ be standard finite dimensional filtrations on $`B`$ and $`\mathrm{\Delta }`$ respectively such that $`B_i\mathrm{\Delta }_i`$ for all $`i0`$. Then the enveloping algebra $`\mathrm{\Delta }^e:=\mathrm{\Delta }\mathrm{\Delta }^0`$ can be equipped with the standard finite dimensional filtration $`\{\mathrm{\Delta }_i^e\}`$ which is the tensor product of the filtrations $`\{\mathrm{\Delta }_i\}`$ and $`\{\mathrm{\Delta }_i^0\}`$ of the algebras $`\mathrm{\Delta }`$ and $`\mathrm{\Delta }^0`$ respectively.
Then $`B\mathrm{\Delta }/\mathrm{\Delta }\mathrm{Der}_KB`$ is a simple left $`\mathrm{\Delta }`$-module (, 15.3.8\] with $`\mathrm{GK}(\mathrm{\Delta })`$ $`=2\mathrm{G}\mathrm{K}(B)`$ (, 15.3.2). By Theorem 3.1,
$$\mathrm{d}(\mathrm{\Delta })+\mathrm{max}\{\mathrm{d}(\mathrm{\Delta }),1\}\frac{\mathrm{GK}(\mathrm{\Delta })}{\mathrm{GK}(B)}=\frac{2\mathrm{G}\mathrm{K}(B)}{\mathrm{GK}(B)}=2,$$
hence $`\mathrm{d}(\mathrm{\Delta })1`$. It remains to prove the opposite inequality. For, we recall some properties of $`\mathrm{\Delta }`$ (see , Ch. 15, for details).
Given $`0cB`$, denote by $`B_c`$ the localization of the algebra $`B`$ at the powers of the element $`c`$, then $`\mathrm{\Delta }(B_c)\mathrm{\Delta }(B)_c`$ and the map $`\mathrm{\Delta }(B)\mathrm{\Delta }(B)_c`$, $`dd/1`$, is injective (, 5.1.25). There is a finite subset $`\{c_1,\mathrm{},c_t\}`$ of $`B`$ such that the algebra $`_{i=1}^t\mathrm{\Delta }(B_{c_i})`$ is left and right faithfully flat over its subalgebra $`\mathrm{\Delta }`$,
$$\underset{i=1}{\overset{t}{}}Bc_i=B(\mathrm{see}\mathrm{the}\mathrm{proof}\mathrm{of}\mathrm{\hspace{0.33em}15.2.13},\text{[19]}).$$
(2)
For each $`c=c_i`$, $`\mathrm{Der}_K(B_c)`$ is a free $`B_c`$-module with a basis $`_j=\frac{}{x_j}`$, $`j=1,\mathrm{},n`$ for some $`x_1,\mathrm{},x_nB`$ (, 15.2.13). Note that the choice of the $`x_j`$โth depends on the choice of the $`c_i`$. Then
$$\mathrm{\Delta }(B)_c\mathrm{\Delta }(B_c)=B_c<_1,\mathrm{},_n>K<x_1,\mathrm{},x_n,_1,\mathrm{},_n>$$
Fix $`c=c_i`$. We aim to prove the following statement:
* there exist natural numbers $`a`$, $`b`$, $`\alpha `$, and $`\beta `$ such that for any $`0d\mathrm{\Delta }_k`$ there exists
$$w\mathrm{\Delta }_{ak+b}^e:wd=c^{\alpha k+\beta }.$$
$`()`$
Suppose that we are done. Then one can choose the numbers $`a`$, $`b`$, $`\alpha `$, and $`\beta `$ such that (\*) holds for all $`i=1,\mathrm{},t`$. It follows from (2) that
$$\underset{i=1}{\overset{t}{}}f_ic_i=1\mathrm{for}\mathrm{some}f_iA.$$
Choose $`\nu `$: all $`f_ic_i\mathrm{\Delta }_\nu `$ and set $`N(k)=\alpha k+\beta `$, then
$$1=(\underset{i=1}{\overset{t}{}}f_ic_i)^{tN(k)}=\underset{i=1}{\overset{t}{}}g_ic_i^{N(k)}=\underset{i=1}{\overset{t}{}}g_iw_id=wd,$$
where the $`w_i`$ are from (\*), i.e. $`w_i\mathrm{\Delta }_{ak+b}^e`$, $`w_id=c_i^{N(k)}`$. So, $`w=_{i=1}^tg_iw_i`$ $`\mathrm{\Delta }_{\nu tN(k)+ak+b}^e`$ and so $`\mathrm{d}(\mathrm{\Delta })1`$, as required.
Fix $`c=c_i`$. By (, 15.1.24) $`\mathrm{Der}_K(B_c)\mathrm{Der}_K(B)_c`$ and $`\mathrm{Der}_KB`$ can be seen as a finitely generated $`B`$-submodule of $`\mathrm{Der}_K(B_c)`$ (, 15.1.7).
The algebra $`B`$ contains the polynomial subalgebra $`P=K[x_1,\mathrm{},x_n]`$. The polynomial algebra $`P`$ has the natural filtration $`P=_{i0}P_i`$ by the total degree of the variables. Fix a natural number $`l`$ such that $`P_1B_l`$, then $`P_iB_{li}`$ for all $`i0`$. We denote by $`Q=K(x_1,\mathrm{},x_n)`$ the field of fractions of $`P`$. The field of fractions, say $`L`$, of the algebra $`B`$ has the same transcendence degree $`n`$ as the field of rational functions $`Q`$. The algebra $`B`$ is a finitely generated algebra, hence $`L`$ is a finite field extension of $`Q`$ of dimension, say $`m`$, over $`Q`$. Let $`e_1,\mathrm{},e_mB`$ be a $`Q`$-basis for the vector space $`L`$ over $`Q`$. Note that $`L=QB`$. One can find a natural number $`\beta 1`$ and a nonzero polynomial $`pP_\beta `$ such that
$$\{B_1,e_je_k|j,k=1,\mathrm{},m\}\underset{\alpha =1}{\overset{m}{}}p^1P_\beta e_\alpha .$$
Then $`B_k_{j=1}^mp^{2k}P_{2\beta k}e_j`$ and $`B_ke_i_{j=1}^mp^{3k}P_{3\beta k}e_j`$ for all $`k1`$ and $`i=1,\mathrm{},m`$. Let $`0dB_k`$. The $`m\times m`$ matrix of the bijective $`Q`$-linear map $`LL`$, $`xdx`$, with respect to the basis $`e_1,\mathrm{},e_m`$ has entries from the set $`p^{3k}P_{3\beta k}`$. So, its characteristic polynomial
$$\chi _d(t)=t^m+\alpha _{m1}t^{m1}+\mathrm{}+\alpha _0$$
has coefficients in $`p^{3mk}P_{3m\beta k}`$, and $`\alpha _00`$ as $`xdx`$ is a bijection. Now,
$$P_{6m\beta k}p^{3mk}\alpha _0=p^{3mk}(d^{m1}\alpha _{m1}d^{m2}\mathrm{}\alpha _1)dB_{4m\beta k}P_{3m\beta k}B_{m\beta k(4+3l)}d.$$
(3)
Let $`\delta _1,\mathrm{},\delta _t`$ be a set of generators for the left $`B`$-module $`\mathrm{Der}_K(B)`$. Then
$$_i\underset{j=1}{\overset{t}{}}c^{l_1}B_{l_1}\delta _j\mathrm{for}i=1,\mathrm{},n,$$
for some natural number $`l_1`$. Fix a natural number $`l_2`$ such that $`\delta _j(B_1)B_{l_2}`$ and $`\delta _j(c)B_{l_2}`$ for $`j=1,\mathrm{},t`$. Then
$$^\alpha (B_k)c^{2|\alpha |(l_1+1)}B_{k+|\alpha |(l_1+l_2)}\mathrm{for}\mathrm{all}\alpha ^n,k1,$$
where $`\alpha =(\alpha _1,\mathrm{},\alpha _n)`$, $`|\alpha |=\alpha _1+\mathrm{}+\alpha _n`$, $`^\alpha =_1^{\alpha _1}\mathrm{}_n^{\alpha _n}`$. It follows from (3) that one can find $`\alpha ^n`$ such that $`|\alpha |6m\beta k`$ and
$$1K^{}^\alpha (p^{3mk}\alpha _0)^\alpha (B_{m\beta k(4+3l)}d)c^{12m\beta (l_1+1)k}\mathrm{\Delta }_{2(m\beta k(4+3l)+6m\beta k(l_1+l_2))}^ed$$
where $`K^{}=K\backslash \{0\}`$. Now (\*) follows. $`\mathrm{}`$
In fact we have proved the following corollary.
###### Corollary 5.1
There exist natural numbers $`a`$ and $`b`$ such that for any $`0d\mathrm{\Delta }_k`$ there exists an element $`w\mathrm{\Delta }_{ak+b}^e`$ satisfying $`wd=1`$.
## 6 Multiplicity for the filter dimension, holonomic modules over simple finitely generated algebras
In this section, we introduce a concept of multiplicity for the filter dimension and a concept of holonomic module for (some) finitely generated algebras. We will prove that a holonomic module has finite length (Theorem 6.8). The multiplicity for the filter dimension is a key ingredient in the proof.
First we recall the definition of multiplicity in the commutative situation and then for certain non-commutative algebras (somewhat commutative algebras).
Multiplicity in the commutative situation. Let $`B`$ be a commutative finitely generated $`K`$-algebra with a standard finite dimensional filtration $`F=\{B_i\}`$, and let $`M`$ be a finitely generated $`B`$-module with a finite dimensional generating subspace, say $`M_0`$, and with the standard filtration $`\{M_i=B_iM_0\}`$ attached to it. Then there exists a polynomial $`p(t)=lt^d+\mathrm{}[t]`$ with rational coefficients of degree $`d=\mathrm{GK}(M)`$ such that
$$\mathrm{dim}_K(M_i)=p(i)\mathrm{for}\mathrm{all}i0.$$
The polynomial $`p(t)`$ is called the Hilbert polynomial of the $`B`$-module $`M`$. The Hilbert polynomial does depend on the filtration $`\{M_i\}`$ of the module $`M`$ but its leading coefficient $`l`$ does not. The number $`e(M)=d!l`$ is called the multiplicity of the $`B`$-module $`M`$. It is a natural number which does depend on the filtration $`F`$ of the algebra $`B`$.
In the case when $`M=B`$ is the homogeneous coordinate ring of a projective algebraic variety $`X^m`$ equipped with the natural filtration that comes from the grading of the graded algebra $`B`$, the multiplicity is the degree of $`X`$, the number of points in which $`X`$ meets a general plane of complementary degree in $`^m`$ ($`K`$ is an algebraically closed field).
Somewhat commutative algebras. A $`K`$-algebra $`R`$ is called a somewhat commutative algebra if it has a finite dimensional filtration $`R=_{i0}R_i`$ such that the associated graded algebra $`\mathrm{gr}R:=_{i0}R_i/R_{i1}`$ is a commutative finitely generated $`K`$-algebra where $`R_1=0`$ and $`R_0=K`$. Then the algebra $`R`$ is a Noetherian finitely generated algebra since $`\mathrm{gr}R`$ is so. A finitely generated module over a somewhat commutative algebra has the Gelfand-Kirillov dimension which is a natural number. We refer the reader to the books for the properties of somewhat commutative algebras.
Definition. For a somewhat commutative algebra $`R`$ we define the holonomic number,
$$h_R:=\mathrm{min}\{\mathrm{GK}(M)|M0\mathrm{is}\mathrm{a}\mathrm{finitely}\mathrm{generated}R\mathrm{module}\}.$$
Definition. A finitely generated $`R`$-module $`M`$ is called a holonomic module if $`\mathrm{GK}(M)=h_R`$. In other words, a nonzero finitely generated $`R`$-module is holonomic iff it has the least Gelfand-Kirillov dimension. If $`h_R=0`$ then every holonomic $`R`$-module is finite dimensional and vice versa.
Examples. $`1`$. The holonomic number of the Weyl algebra $`A_n`$ is $`n`$. The polynomial algebra $`K[X_1,\mathrm{},X_n]A_n/_{i=1}^nA_n_i`$ with the natural action of the ring of differential operators $`A_n=K[X_1,\mathrm{},X_n,\frac{}{X_1},\mathrm{},\frac{}{X_n}]`$ is a simple holonomic $`A_n`$-module.
$`2`$. Let $`X`$ be a smooth irreducible affine algebraic variety of dimension $`n`$. The ring of differential operators $`๐(X)`$ is a simple somewhat commutative algebra of Gelfand-Kirillov dimension $`2n`$ with holonomic number $`h_{๐(X)}=n`$. The algebra $`๐ช(X)`$ of regular functions of the variety $`X`$ is a simple $`๐(X)`$-module with respect to the natural action of the algebra $`๐(X)`$. In more detail, $`๐ช(X)๐(X)/๐(X)\mathrm{Der}_K(๐ช(X))`$ where $`\mathrm{Der}_K(๐ช(X))`$ is the $`๐ช(X)`$-module of derivations of the algebra $`๐ช(X)`$.
Let $`R=_{i0}R_i`$ be a somewhat commutative algebra. The associated graded algebra $`\mathrm{gr}R`$ is a commutative affine algebra. Let us choose homogeneous algebra generators of the algebra $`\mathrm{gr}R`$, say $`y_1,\mathrm{},y_s`$, of graded degrees $`1k_1,\mathrm{},k_s`$ respectively (that is $`y_iR_{k_i}/R_{k_i1}`$). A filtration $`\mathrm{\Gamma }=\{\mathrm{\Gamma }_i,i0\}`$ of an $`R`$-module $`M=_{i=0}^{\mathrm{}}\mathrm{\Gamma }_i`$ is called good if the associated graded $`\mathrm{gr}R`$-module $`\mathrm{gr}_\mathrm{\Gamma }M:=_{i0}\mathrm{\Gamma }_i/\mathrm{\Gamma }_{i1}`$ is finitely generated. An $`R`$-module $`M`$ has a good filtration iff it is finitely generated, and if $`\{\mathrm{\Gamma }_i\}`$ and $`\{\mathrm{\Omega }_i\}`$ are two good filtrations of $`M`$, then there exists a natural number $`t`$ such that $`\mathrm{\Gamma }_i\mathrm{\Omega }_{i+t}`$ and $`\mathrm{\Omega }_i\mathrm{\Gamma }_{i+t}`$ for all $`i`$. If an $`R`$-module $`M`$ is finitely generated and $`M_0`$ is a finite dimensional generating subspace of $`M`$, then the standard filtration $`\{\mathrm{\Gamma }_i=R_iM_0\}`$ is good (see for details). The following two Lemmas are well-known by specialists (see their proofs, for example, in , Theorem 3.2 and Proposition 3.3 respectively).
###### Lemma 6.1
Let $`R=_{i0}R_i`$ be a somewhat commutative algebra, $`k=\mathrm{lcm}(k_1,\mathrm{},k_s)`$, and let $`M`$ be a finitely generated $`R`$-module with good filtration $`\mathrm{\Gamma }=\{\mathrm{\Gamma }_i\}`$.
1. There exist $`k`$ polynomials $`\gamma _0,\mathrm{},\gamma _{k1}[t]`$ with coefficients from $`[k^{\mathrm{GK}(M)}\mathrm{GK}(M)!]^1`$ such that
$$\mathrm{dim}\mathrm{\Gamma }_i=\gamma _j(i)\mathrm{for}\mathrm{all}i0\mathrm{and}ji(\mathrm{mod}k).$$
2. The polynomials $`\gamma _j`$ have the same degree $`\mathrm{GK}(M)`$ and the same leading coefficient $`e(M)/\mathrm{GK}(M)!`$ where $`e(M)`$ is called the multiplicity of $`M`$. The multiplicity $`e(M)`$ does not depend on the choice of the good filtration $`\mathrm{\Gamma }`$. $`\mathrm{}`$
Remark. A finitely generated $`R`$-module $`M`$ has $`e(M)=0`$ iff $`\mathrm{dim}_K(M)<\mathrm{}`$.
###### Lemma 6.2
Let $`0NML0`$ be an exact sequence of modules over a somewhat commutative algebra $`R`$. Then $`\mathrm{GK}(M)=\mathrm{max}\{\mathrm{GK}(N),\mathrm{GK}(L)\}`$, and if $`\mathrm{GK}(N)=\mathrm{GK}(M)=\mathrm{GK}(L)`$ then $`e(M)=e(N)+e(L)`$. $`\mathrm{}`$
###### Corollary 6.3
Let the algebra $`R`$ be as in Lemma 6.1 with holonomic number $`h>0`$.
1. Let $`M`$ be a holonomic $`R`$-module with multiplicity $`e(M)`$. The $`R`$-module $`M`$ has finite length $`e(M)k^h`$.
2. Every nonzero submodule or factor module of a holonomic $`R`$-module is a holonomic module.
Proof. This follows directly from Lemma 6.2. $`\mathrm{}`$
Multiplicity. Let $`f`$ be a function from $``$ to $`_+=\{r:r0\}`$, the leading coefficient of $`f`$ is a non-zero limit (if it exists)
$$\mathrm{lc}(f)=lim\frac{f(i)}{i^d}0,i\mathrm{},$$
where $`d=\gamma (f)`$. If $`d`$, we define the multiplicity $`e(f)`$ of $`f`$ by
$$e(f)=d!\mathrm{lc}(f).$$
The factor $`d!`$ ensures that the multiplicity $`e(f)`$ is a positive integer in some important cases. If $`f(t)=a_dt^d+a_{d1}t^{d1}+\mathrm{}+a_0`$ is a polynomial of degree $`d`$ with real coefficients then $`\mathrm{lc}(f)=a_d`$ and $`e(f)=d!a_d`$.
###### Lemma 6.4
Let $`A`$ be a finitely generated algebra equipped with a standard finite dimensional filtration $`F=\{A_i\}`$ and $`M`$ be a finitely generated $`A`$-module with generating finite dimensional subspaces $`M_0`$ and $`N_0`$.
1. If $`\mathrm{lc}(\nu _{F,M_0})`$ exists then so does $`\mathrm{lc}(\nu _{F,N_0})`$, and $`\mathrm{lc}(\nu _{F,M_0})=\mathrm{lc}(\nu _{F,N_0})`$.
2. If $`\mathrm{lc}(\mathrm{dim}A_iM_0)`$ exists then so does $`\mathrm{lc}(\mathrm{dim}A_iN_0)`$, and $`\mathrm{lc}(\mathrm{dim}A_iM_0)=\mathrm{lc}(\mathrm{dim}A_iN_0)`$.
Proof. 1. The module $`M`$ has two filtrations $`\{M_i=A_iM_0\}`$ and $`\{N_i=A_iN_0\}`$. Let $`\nu =\nu _{F,M_0}`$ and $`\mu =\nu _{F,N_0}`$. Choose a natural number $`s`$ such that $`M_0N_s`$ and $`N_0M_s`$, so $`N_iM_{i+s}`$ and $`M_iN_{i+s}`$ for $`i0`$. Since $`M_0A_{\nu (i+s)}N_{i,gen}`$ for each $`i`$ and $`N_0A_sM_0,`$ we have $`N_0A_{\nu (i+s)+s}N_{i,gen},`$ hence, $`\mu (i)\nu (i+s)+s`$. By symmetry, $`\nu (i)\mu (i+s)+s`$, so if $`\mathrm{lc}(\mu )`$ exists then so does $`\mathrm{lc}(\nu )`$ and $`\mathrm{lc}(\mu )=\mathrm{lc}(\nu )`$.
2. Since $`\mathrm{dim}N_i\mathrm{dim}M_{i+s}`$ and $`\mathrm{dim}M_i\mathrm{dim}N_{i+s}`$ for $`i0`$, the statement is clear. $`\mathrm{}`$
Lemma 6.4 shows that the leading coefficients of the functions $`\mathrm{dim}A_iM_0`$ and $`\nu _{F,M_0}`$ (if exist) do not depend on the choice of the generating subspace $`M_0`$. So, denote them by
$$l(M)=l_F(M)\mathrm{and}L(M)=L_F(M)$$
respectively (if they exist). If $`\mathrm{GK}(M)`$ (resp. $`\mathrm{d}(A)`$) is a natural number, then we denote by $`e(M)=e_F(M)`$ (resp. $`E(M)=E_F(M)`$) the multiplicity of the function $`\mathrm{dim}A_iM_0`$ (resp. $`\nu _{F,M_0}`$).
We denote by $`L(A)=L_F(A)`$ the leading coefficient $`L_F({}_{AA^0}{}^{}A)`$ of the return function $`\nu _{FF^0,K}`$ of the $`AA^0`$-module $`A`$.
Holonomic modules. Definition. Let $`A`$ be a finitely generated $`K`$-algebra, and $`h_A`$ be its holonomic number. A nonzero finitely generated $`A`$-module $`M`$ is called a holonomic $`A`$-module if $`\mathrm{GK}(M)=h_A`$. We denote by $`\mathrm{hol}(A)`$ the set of all the holonomic $`A`$-modules.
Since the holonomic number is an infimum it is not clear at the outset that there will be modules which achieve this dimension. Clearly, $`\mathrm{hol}(A)\mathrm{}`$ if the Gelfand-Kirillov dimension of every finitely generated $`A`$-module is a natural number.
A nonzero submodule or a factor module of a holonomic is a holonomic module (since the Gelfand-Kirillov dimension of a submodule or a factor module does not exceed the Gelfand-Kirillov of the module). If, in addition, the finitely generated algebra $`A`$ is left Noetherian and finitely partitive then each holonomic $`A`$-module $`M`$ has finite length and each simple sub-factor of $`M`$ is a holonomic module.
Let us consider algebras $`A`$ having the following properties:
* (S) $`A`$ is a simple finitely generated infinite dimensional algebra.
* (N) There exists a standard finite dimensional filtration $`F=\{A_i\}`$ of the algebra $`A`$ such that the associated graded algebra gr $`A:=`$ $`{}_{i0}{}^{}A_{i}^{}/A_{i1}`$, $`A_1=0`$, is left Noetherian.
* (D) $`\mathrm{GK}(A)<\mathrm{}`$, $`\mathrm{fd}(A)<\mathrm{}`$, both $`l(A)=l_F(A)`$ and $`L(A)=L_F(A)`$ exist.
* (H) For every holonomic $`A`$-module $`M`$ there exists $`l(M)=l_F(M)`$.
In many cases we use the weaker form of the condition (D).
* (Dโ) $`\mathrm{GK}(A)<\mathrm{}`$, $`d=\mathrm{fd}(A)<\mathrm{}`$, there exist $`l(A)=l_F(A)`$ and a positive number $`c>0`$ such that $`\nu (i)ci^d`$ for $`i0`$ where $`\nu `$ is the return function $`\nu _{FF^0,K}`$ of the left $`AA^0`$-module $`A`$.
It follows from (N) that $`A`$ is a left Noetherian algebra.
###### Lemma 6.5
()
1. The Weyl algebra $`A_n`$ over a field of characteristic zero with the standard finite dimensional filtration $`F=\{A_{n,i}\}`$ associated with the canonical generators satisfies the conditions (S), (N), (D), (H). The return function $`\nu _F(i)=i`$ for $`i0`$, and so the leading coefficient of $`\nu _F`$ is $`L_F(A_n)=1`$.
2. $`\nu _{G,K}(i)=i`$ for $`i0`$ and $`L_G(P_n)=1`$ where $`\nu _{G,K}`$ is the return function of the $`A_n`$-module $`P_n=K[X_1,\mathrm{},X_n]=A_n/(A_n_1+\mathrm{}+A_n_n)`$ with the usual filtration $`G=\{P_{n,i}\}`$ of the polynomial algebra.
Proof. $`1`$. The only fact that we need to prove is that $`\nu _F(i)=i`$ for $`i0`$. We keep the notation of Theorem 4.6. In the proof of Theorem 4.6 we have seen that $`\nu _F(i)i`$ for $`i0`$. It remains to prove the reverse inequality.
Each element $`u`$ in $`A_n`$ can be written in a unique way as a finite sum $`u=\lambda _{\alpha \beta }X^\alpha ^\beta `$ where $`\lambda _{\alpha \beta }K`$ and $`X^\alpha `$ denotes the monomial $`X_1^{\alpha _1}\mathrm{}X_n^{\alpha _n}`$ and similarly $`^\beta `$ denotes the monomial $`{}_{1}{}^{\beta _1}\mathrm{}_n^{\beta _n}`$. The element $`u`$ belongs to $`A_{n,m}`$ iff $`|\alpha |+|\beta |m`$, where $`|\alpha |=\alpha {}_{1}{}^{}+\mathrm{}+\alpha _n`$. If $`\alpha K[X_1,\mathrm{},X_n]`$, then
$$_i^m\alpha =\underset{j=0}{\overset{m}{}}\left(\genfrac{}{}{0pt}{}{m}{j}\right)\frac{^j\alpha }{X_i^j}_i^{mj},m.$$
It follows that for any $`vA_{n,i}A_{n,j}^0`$, $`i+j<m`$, the element $`vX_1^m=\lambda _{\alpha \beta }X^\alpha ^\beta `$ has the coefficient $`\lambda _{0,0}=0`$, hence it could not be a non-zero scalar, and so $`\nu (i)i`$ for all $`i0`$. Hence $`\nu (i)=i`$ all $`i0`$ and then $`L_F(A_n)=1`$.
$`2`$. The standard filtration of the $`A_n`$-module $`P_n`$ associated with the generating subspace $`K`$ coincides with the usual filtration of the polynomial algebra $`P_n`$. Since $`_j(P_{n,i})P_{n,i1}`$ for all $`i0`$ and $`j`$, $`\nu _{G,K}(i)i`$ for $`i0`$. Using the same arguments as above we see that for any $`u_{j=0}^{i1}A_{n,j}A_{n,ij1}^0`$ the element $`uX_1^i`$ belongs to the ideal of $`P_n`$ generated by $`X_1`$, hence, $`\nu _{G,K}(i)i`$, and so $`\nu _{G,K}(i)=i`$ for all $`i0`$ and $`L_G(P_n)=1`$. $`\mathrm{}`$
###### Theorem 6.6
() Assume that an algebra $`A`$ satisfies the conditions (S), (H), (D) resp. (Dโ) for some standard finite dimensional filtration $`F=\{A_i\}`$ of $`A`$. Then for every holonomic $`A`$-module $`M`$ its leading coefficient is bounded from below by a nonzero constant:
$$l(M)\sqrt{\frac{l(A)}{(L(A)L^{}(A))^{h_A}}},$$
where
$$L^{}(A)=\{\begin{array}{cc}L(A),\hfill & \text{if }\mathrm{d}(A)>1\text{,}\hfill \\ L(A)+1,\hfill & \text{if }\mathrm{d}(A)=1\text{,}\hfill \\ 1,\hfill & \text{if }\mathrm{d}(A)<1\text{.}\hfill \end{array}$$
resp.
$$l(M)\sqrt{\frac{l(A)}{(c(c+1))^{h_A}}}.$$
Proof. Let $`M_0`$ be a generating finite dimensional subspace of $`M`$ and $`\{M_i=A_iM_0\}`$ be the standard finite dimensional filtration of $`M`$. In the proof of Theorem 3.1 we proved that $`\mathrm{dim}A_i\mathrm{dim}M_{\lambda (i)}\mathrm{dim}M_{\lambda (i)+i}`$ for $`i0`$ where $`\lambda `$ is the left return function of the algebra $`A`$ associated with the filtration $`F`$. Since $`\lambda (i)\nu (i)`$ for $`i0`$ we have $`\mathrm{dim}A_i\mathrm{dim}M_{\nu (i)}\mathrm{dim}M_{\nu (i)+i}`$, hence, if (D) holds then
$$l(A)i^{\mathrm{GK}(A)}+\mathrm{}l^2(M)(L(A)L^{}(A))^{\mathrm{GK}(M)}i^{\mathrm{GK}(M)(\mathrm{fd}(A)+\mathrm{max}\{\mathrm{fd}(A),1\})}+\mathrm{},$$
where three dots denote smaller terms.
If (Dโ) holds then
$$l(A)i^{\mathrm{GK}(A)}+\mathrm{}l^2(M)(c(c+1))^{\mathrm{GK}(M)}i^{\mathrm{GK}(M)(\mathrm{fd}(A)+\mathrm{max}\{\mathrm{fd}(A),1\})}+\mathrm{}.$$
The module $`M`$ is holonomic, i.e. $`\mathrm{GK}(A)=\mathrm{GK}(M)(\mathrm{fd}(A)+\mathrm{max}\{\mathrm{fd}(A),1\})`$. Now, comparing the โleadingโ coefficients in the inequalities above we finish the proof. $`\mathrm{}`$
Let $`A`$ be as in Theorem 6.6. We attach to the algebra $`A`$ two positive numbers $`c_A`$ and $`c_A^{}`$ in the cases (D) and (Dโ) respectively:
$$c_A=\sqrt{\frac{l(A)}{(L(A)L^{}(A))^{h_A}}}\mathrm{and}(c_A^{})=\sqrt{\frac{l(A)}{(c(c+1))^{h_A}}}.$$
###### Corollary 6.7
Assume that an algebra $`A`$ satisfies the conditions (S), (N), (H), (D) or (Dโ). Let $`0NML0`$ be an exact sequence of nonzero finitely generated $`A`$-modules. Then $`M`$ is holonomic if and only if $`N`$ and $`L`$ are holonomic, in that case $`l(M)=l(N)+l(L)`$.
Proof . The algebra $`A`$ is left Noetherian, so the module $`M`$ is finitely generated iff both $`N`$ and $`L`$ are so. The proof of Proposition 3.11 (, p. 295) shows that we can choose finite dimensional generating subspaces $`N_0`$, $`M_0`$, $`L_0`$ of the modules $`N`$, $`M`$, $`L`$ respectively such that the sequences
$$0N_i=A_iN_0M_i=A_iM_0L_i=A_iL_00$$
are exact for all $`i`$, hence, $`\mathrm{dim}M_i=\mathrm{dim}N_i+\mathrm{dim}L_i`$ and the results follow. $`\mathrm{}`$
###### Theorem 6.8
() Suppose that the conditions (S), (N), (H), (D) (resp. (Dโ)) hold. Then each holonomic $`A`$-module $`M`$ has finite length which is less or equal to $`l(M)/c_A`$ (resp. $`l(M)/c_A^{}`$).
Proof. If $`M=M_1M_2\mathrm{}M_mM_{m+1}=0`$ is a chain of distinct submodules, then by Corollary 6.7 and Theorem 6.6
$$l(M)=\underset{i=1}{\overset{m}{}}l(M_i/M_{i+1})mc_A,(\mathrm{resp}.l(M)mc_A^{}),$$
thus $`ml(M)/c_A`$ (resp. $`ml(M)/c_A^{}`$). $`\mathrm{}`$
## 7 Filter dimension and commutative subalgebras of simple finitely generated algebras and their division algebras
In this section, using the first and the second filter inequalities, we obtain $`(i)`$ an upper bound for the Gelfand-Kirillov dimension of commutative subalgebras of simple finitely generated infinite dimensional algebras (Theorem 7.2), and $`(ii)`$ an upper bound for the transcendence degree of subfields of quotient division rings of (certain) simple finitely generated infinite dimensional algebras (Theorems 7.4 and 7.5).
For certain classes of algebras and their division algebras the maximum Gelfand-Kirillov dimension/transcendence degree over the commutative subalgebras/subfields were found in , , , , , , , and .
Recall that
$`\mathrm{the}\mathrm{Gelfand}\mathrm{Kirillov}\mathrm{dimension}\mathrm{GK}(C)`$ $`=`$ $`\mathrm{the}\mathrm{Krull}\mathrm{dimension}\mathrm{K}.\mathrm{dim}(C)`$
$`=`$ $`\mathrm{the}\mathrm{transcendence}\mathrm{degree}\mathrm{tr}.\mathrm{deg}_K(C)`$
for every commutative finitely generated algebra $`C`$ which is a domain.
An upper bound for the Gelfand-Kirillov dimensions of commutative subalgebras of simple finitely generated algebras.
###### Proposition 7.1
Let $`A`$ and $`C`$ be finitely generated algebras such that $`C`$ is a commutative domain with field of fractions $`Q`$, $`B:=CA`$, and $`:=QA`$. Let $`M`$ be a finitely generated $`B`$-module such that $`:=_BM0`$. Then $`\mathrm{GK}({}_{B}{}^{}M)\mathrm{GK}_Q({}_{}{}^{})+\mathrm{GK}(C)`$.
Remark. $`\mathrm{GK}_Q`$ stands for the Gelfand-Kirillov dimension over the field $`Q`$.
Proof. Let us fix standard filtrations $`\{A_i\}`$ and $`\{C_i\}`$ for the algebras $`A`$ and $`C`$ respectively. Let $`h(t)[t]`$ be the Hilbert polynomial for the algebra $`C`$, i.e. $`\mathrm{dim}_K(C_i)=h(i)`$ for $`i0`$. Recall that $`\mathrm{GK}(C)=\mathrm{deg}_t(h(t))`$. The algebra $`B`$ has a standard filtration $`\{B_i\}`$ which is the tensor product of the standard filtrations $`\{C_i\}`$ and $`\{A_i\}`$ of the algebras $`C`$ and $`A`$, i.e. $`B_i:=_{j=0}^iC_jA_{ij}`$. By the assumption, the $`B`$-module $`M`$ is finitely generated, so $`M=BM_0`$ where $`M_0`$ is a finite dimensional generating subspace for $`M`$. Then the $`B`$-module $`M`$ has a standard filtration $`\{M_i:=B_iM_0\}`$. The $`Q`$-algebra $``$ has a standard (finite dimensional over $`Q`$) filtration $`\{_i:=QA_i\}`$, and the $``$-module $``$ has a standard (finite dimensional over $`Q`$) filtration $`\{_i:=_iM_0^{}=QA_iM_0^{}\}`$ where $`M_0^{}`$ is the image of the vector space $`M_0`$ under the $`B`$-module homomorphism $`M`$, $`mm^{}:=1_Bm`$.
For each $`i0`$, one can fix a $`K`$-subspace, say $`L_i`$, of $`A_iM_0^{}`$ such that $`\mathrm{dim}_Q(QA_iM_0^{})=\mathrm{dim}_K(L_i)`$. Now, $`B_{2i}C_iA_i`$ implies $`\mathrm{dim}_K(B_{2i}M_0)\mathrm{dim}_K((C_iA_i)M_0)`$, and $`((C_iA_i)M_0)^{}C_iL_i`$ implies $`\mathrm{dim}_K(((C_iA_i)M_0)^{})\mathrm{dim}_K(C_iL_i)=\mathrm{dim}_K(C_i)\mathrm{dim}_K(L_i)=\mathrm{dim}_K(C_i)\mathrm{dim}_Q(_i)`$. It follows that
$`\mathrm{GK}({}_{B}{}^{}M)`$ $`=`$ $`\gamma (\mathrm{dim}_K(M_i))\gamma (\mathrm{dim}_K(M_{2i}))=\gamma (\mathrm{dim}_K(B_{2i}M_0))\gamma (\mathrm{dim}_K((C_iA_i)M_0))`$
$``$ $`\gamma (\mathrm{dim}_K(((C_iA_i)M_0)^{})\gamma (\mathrm{dim}_K(C_i)\mathrm{dim}_Q(_i))`$
$`=`$ $`\gamma (\mathrm{dim}_K(C_i))+\gamma (\mathrm{dim}_Q(_i))(\mathrm{since}\gamma (\mathrm{dim}_K(C_i))=h(i),\mathrm{for}i0)`$
$`=`$ $`\mathrm{GK}(C)+\mathrm{GK}_Q({}_{}{}^{}).\mathrm{}`$
Recall that $`d=\mathrm{fd},\mathrm{lfd}`$. A $`K`$-algebra $`A`$ is called central if its centre $`Z(A)=K`$.
###### Theorem 7.2
() Let $`A`$ be a central simple finitely generated $`K`$-algebra of Gelfand-Kirillov dimension $`0<n<\mathrm{}`$ (over $`K`$). Let $`C`$ be a commutative subalgebra of $`A`$. Then
$$\mathrm{GK}(C)\mathrm{GK}(A)\left(1\frac{1}{f_A+\mathrm{max}\{f_A,1\}}\right)$$
where $`f_A:=\mathrm{max}\{\mathrm{d}_{Q_m}(Q_mA)|\mathrm{\hspace{0.17em}0}mn\}`$, $`Q_0:=K`$, and $`Q_m:=K(x_1,\mathrm{},x_m)`$ is a rational function field in indeterminates $`x_1,\mathrm{},x_m`$.
Proof. Let $`P_m=K[x_1,\mathrm{},x_m]`$ be a polynomial algebra over the field $`K`$. Then $`Q_m`$ is its field of fractions and $`\mathrm{GK}(P_m)=m`$. Suppose that $`P_m`$ is a subalgebra of $`A`$. Then $`m=\mathrm{GK}(P_m)\mathrm{GK}(A)=n`$. For each $`m0`$, $`Q_mA`$ is a central simple $`Q_m`$-algebra (, 9.6.9) of Gelfand-Kirillov dimension (over $`Q_m`$) $`\mathrm{GK}_{Q_m}(Q_mA)=\mathrm{GK}(A)>0`$, hence
$`\mathrm{GK}(A)`$ $`=`$ $`\mathrm{GK}({}_{A}{}^{}A_{A}^{})\mathrm{GK}({}_{A}{}^{}A_{P_m}^{})=\mathrm{GK}({}_{P_mA}{}^{}A)(P_m\mathrm{is}\mathrm{commutative})`$
$``$ $`\mathrm{GK}_{Q_m}({}_{Q_mA}{}^{}(Q_m_{P_m}A))+\mathrm{GK}(P_m)(\mathrm{Lemma}\text{7.1})`$
$``$ $`{\displaystyle \frac{\mathrm{GK}(A)}{\mathrm{d}_{Q_m}(Q_mA)+\mathrm{max}\{\mathrm{d}_{Q_m}(Q_mA),1\}}}+m(\mathrm{Theorem}\text{3.1}).`$
Hence,
$$m\mathrm{GK}(A)\left(1\frac{1}{\mathrm{d}_{Q_m}(Q_mA)+\mathrm{max}\{\mathrm{d}_{Q_m}(Q_mA),1\}}\right)\mathrm{GK}(A),$$
and so
$$\mathrm{GK}(C)\mathrm{GK}(A)\left(1\frac{1}{f_A+\mathrm{max}\{f_A,1\}}\right).\mathrm{}$$
As a consequence we have a short proof of the following well-known result.
###### Corollary 7.3
Let $`K`$ be an algebraically closed field of characteristic zero, $`X`$ be a smooth irreducible affine algebraic variety of dimension $`n:=\mathrm{dim}(X)>0`$, and $`C`$ be a commutative subalgebra of the ring of differential operators $`๐(X)`$. Then $`\mathrm{GK}(C)n`$.
Proof. The algebra $`๐(X)`$ is central since $`K`$ is an algebraically closed field of characteristic zero , Ch. 15. By Theorem 4.3, $`f_{๐(X)}=1`$, and then, by Theorem 7.2,
$$\mathrm{GK}(C)2n(1\frac{1}{1+1})=n.\mathrm{}$$
Remark. For the ring of differential operators $`๐(X)`$ the upper bound in Theorem 7.2 for the Gelfand-Kirillov dimension of commutative subalgebras of $`๐(X)`$ is an exact upper bound since as we mentioned above the algebra $`๐ช(X)`$ of regular functions on $`X`$ is a commutative subalgebra of $`๐(X)`$ of Gelfand-Kirillov dimension $`n`$.
An upper bound for the transcendence degree of subfields of quotient division algebras of simple finitely generated algebras.
Recall that the transcendence degree $`\mathrm{tr}.\mathrm{deg}_K(L)`$ of a field extension $`L`$ of a field $`K`$ coincides with the Gelfand-Kirillov dimension $`\mathrm{GK}_K(L)`$, and, by the Goldieโs Theorem, a left Noetherian algebra $`A`$ has a quotient algebra $`D=D_A`$ (i.e. $`D=S^1A`$ where $`S`$ is the set of regular elements $`=`$ the set of non-zerodivisors of $`A`$). As a rule, the quotient algebra $`D`$ has infinite Gelfand-Kirillov dimension and is not a finitely generated algebra (eg, the quotient division algebra $`D(X)`$ of the ring of differential operators $`๐(X)`$ on each smooth irreducible affine algebraic variety $`X`$ of dimension $`n>0`$ over a field $`K`$ of characteristic zero contains a noncommutative free subalgebra since $`D(X)D(๐ธ^1)`$ and the first Weyl division algebra $`D(๐ธ^1)`$ has this property ). So, if we want to find an upper bound for the transcendence degree of subfields in the quotient algebra $`D`$ we can not apply Theorem 7.2. Nevertheless, imposing some natural (mild) restrictions on the algebra $`A`$ one can obtain exactly the same upper bound for the transcendence degree of subfields in the quotient algebra $`D_A`$ as the upper bound for the Gelfand-Kirillov dimension of commutative subalgebras in $`A`$.
###### Theorem 7.4
() Let $`A`$ be a simple finitely generated $`K`$-algebra such that $`0<n:=\mathrm{GK}(A)<\mathrm{}`$, all the algebras $`Q_mA`$, $`m0`$, are simple finitely partitive algebras where $`Q_0:=K`$, $`Q_m:=K(x_1,\mathrm{},x_m)`$ is a rational function field and, for each $`m0`$, the Gelfand-Kirillov dimension (over $`Q_m`$) of every finitely generated $`Q_mA`$-module is a natural number. Let $`B=S^1A`$ be the localization of the algebra $`A`$ at a left Ore subset $`S`$ of $`A`$. Let $`L`$ be a (commutative) subfield of the algebra $`B`$ that contains $`K`$. Then
$$\mathrm{tr}.\mathrm{deg}_K(L)\mathrm{GK}(A)\left(1\frac{1}{f_A+\mathrm{max}\{f_A,1\}}\right)$$
where $`f_A:=\mathrm{max}\{\mathrm{d}_{Q_m}(Q_mA)|\mathrm{\hspace{0.17em}0}mn\}`$.
Proof. It follows immediately from a definition of the Gelfand-Kirillov dimension that $`\mathrm{GK}_K^{}(K^{}C)=\mathrm{GK}(C)`$ for any $`K`$-algebra $`C`$ and any field extension $`K^{}`$ of $`K`$. In particular, $`\mathrm{GK}_{Q_m}(Q_mA)=\mathrm{GK}(A)`$ for all $`m0`$. By Theorem 4.2,
$$\mathrm{K}.\mathrm{dim}(Q_mA)\mathrm{GK}(A)\left(1\frac{1}{\mathrm{d}_{Q_m}(Q_mA)+\mathrm{max}\{\mathrm{d}_{Q_m}(Q_mA),1\}}\right).$$
Let $`L`$ be a subfield of the algebra $`B`$ that contains $`K`$. Suppose that $`L`$ contains a rational function field (isomorphic to) $`Q_m`$ for some $`m0`$.
$`m`$ $`=`$ $`\mathrm{tr}.\mathrm{deg}_K(Q_m)\mathrm{K}.\mathrm{dim}(Q_mQ_m)`$
$``$ $`\mathrm{K}.\mathrm{dim}(Q_mB)(\mathrm{by}\text{[19]},\mathrm{\hspace{0.17em}6.5.3}\mathrm{since}Q_mB\mathrm{is}\mathrm{a}\mathrm{free}Q_mQ_m\mathrm{module})`$
$`=`$ $`\mathrm{K}.\mathrm{dim}(Q_mS^1A)=\mathrm{K}.\mathrm{dim}(S^1(Q_mA))`$
$``$ $`\mathrm{K}.\mathrm{dim}(Q_mA)(\mathrm{by}\text{[19]},\mathrm{\hspace{0.17em}6.5.3}.(ii).(b))`$
$``$ $`\mathrm{GK}(A)\left(1{\displaystyle \frac{1}{\mathrm{d}_{Q_m}(Q_mA)+\mathrm{max}\{\mathrm{d}_{Q_m}(Q_mA),1\}}}\right)\mathrm{GK}(A).`$
Hence
$$\mathrm{tr}.\mathrm{deg}_K(L)\mathrm{GK}(A)\left(1\frac{1}{f_A+\mathrm{max}\{f_A,1\}}\right).\mathrm{}$$
Recall that every somewhat commutative algebra $`A`$ is a Noetherian finitely generated finitely partitive algebra of finite Gelfand-Kirillov dimension, the Gelfand-Kirillov dimension of every finitely generated $`A`$-modules is an integer, and (Quillenโs Lemma): the ring $`\mathrm{End}_A(M)`$ is algebraic over $`K`$ (see , Ch. 8 or for details).
###### Theorem 7.5
() Let $`A`$ be a central simple somewhat commutative infinite dimensional $`K`$-algebra and let $`D=D_A`$ be its quotient algebra. Let $`L`$ be a subfield of $`D`$ that contains $`K`$. Then the transcendence degree of the field $`L`$ (over $`K`$)
$$\mathrm{tr}.\mathrm{deg}_K(L)\mathrm{GK}(A)\left(1\frac{1}{f_A+\mathrm{max}\{f_A,1\}}\right)$$
where $`f_A:=\mathrm{max}\{\mathrm{d}_{Q_m}(Q_mA)|\mathrm{\hspace{0.17em}0}m\mathrm{GK}(A)\}`$.
Proof. The algebra $`A`$ is a somewhat commutative algebra, so it has a finite dimensional filtration $`A=_{i0}A_i`$ such that the associated graded algebra is a commutative finitely generated algebra. For each integer $`m0`$, the $`Q_m`$-algebra $`Q_mA=_{i0}Q_mA_i`$ has the finite dimensional filtration (over $`Q_m`$) such that the associated graded algebra $`\mathrm{gr}(Q_mA)=_{i0}Q_mA_i/Q_mA_{i1}Q_m\mathrm{gr}(A)`$ is a commutative finitely generated $`Q_m`$-algebra. So, $`Q_mA`$ is a somewhat commutative $`Q_m`$-algebra.
By the assumption $`\mathrm{dim}_K(A)=\mathrm{}`$, hence $`\mathrm{dim}_K(\mathrm{gr}(A))=\mathrm{}`$ which implies $`\mathrm{GK}(\mathrm{gr}(A))>0`$, and so $`\mathrm{GK}(A)>0`$ (since $`\mathrm{GK}(A)=\mathrm{GK}(\mathrm{gr}(A))`$). The algebra $`A`$ is a central simple $`K`$-algebra, so $`Q_mA`$ is a central simple $`Q_m`$-algebra (, 9.6.9). Now, Theorem 7.5 follows from Theorem 7.4 applied to $`B=D`$. $`\mathrm{}`$
###### Theorem 7.6
Let $`K`$ be an algebraically closed field of characteristic zero, $`๐(X)`$ be the ring of differential operators on a smooth irreducible affine algebraic variety $`X`$ of dimension $`n>0`$, and $`D(X)`$ be the quotient division ring for $`๐(X)`$. Let $`L`$ be a (commutative) subfield of $`D(X)`$ that contains $`K`$. Then $`\mathrm{tr}.\mathrm{deg}_K(L)n`$.
Remark. This inequality is, in fact, an exact upper bound for the transcendence degree of subfields in $`D(X)`$ since the field of fractions $`Q(X)`$ for the algebra $`๐ช(X)`$ is a commutative subfield of the division ring $`D(X)`$ with $`\mathrm{tr}.\mathrm{deg}_K(Q(X))=n`$.
Proof. Since $`Q_m๐_K(๐ช(X))๐_{Q_m}(Q_m๐ช(X))`$ and $`\mathrm{d}(๐(Q_m๐ช(X)))=1`$ for all $`m0`$ we have $`f_{๐(X)}=1`$. Now, Theorem 7.6 follows from Theorem 7.5,
$$\mathrm{tr}.\mathrm{deg}_K(L)2n(1\frac{1}{1+1})=n.\mathrm{}$$
Following for a $`K`$-algebra $`A`$ define the commutative dimension
$$\mathrm{Cdim}(A):=\mathrm{max}\{\mathrm{GK}(C)|C\mathrm{is}\mathrm{a}\mathrm{commutative}\mathrm{subalgebra}\mathrm{of}A\}.$$
The commutative dimension $`\mathrm{Cdim}(A)`$ (if finite) is the largest non-negative integer $`m`$ such that the algebra $`A`$ contains a polynomial algebra in $`m`$ variables (, 1.1, or , 8.2.14). So, $`\mathrm{Cdim}(A)=\{\mathrm{}\}`$. If $`A`$ is a subalgebra of $`B`$ then $`\mathrm{Cdim}(A)\mathrm{Cdim}(B)`$.
###### Corollary 7.7
Let $`X`$ and $`Y`$ be smooth irreducible affine algebraic varieties of dimensions $`n`$ and $`m`$ respectively, let $`D(X)`$ and $`D(Y)`$ be quotient division rings for the rings of differential operators $`๐(X)`$ and $`๐(Y)`$. Then there is no $`K`$-algebra embedding $`D(X)D(Y)`$ if $`n>m`$.
Proof. By Theorem 7.6, $`\mathrm{Cdim}(D(X))=n`$ and $`\mathrm{Cdim}(D(Y))=m`$. Suppose that there is a $`K`$-algebra embedding $`D(X)D(Y)`$. Then $`n=\mathrm{Cdim}(D(X))\mathrm{Cdim}(D(Y))=m`$. $`\mathrm{}`$
For the Weyl algebras $`A_n=๐(๐ธ^n)`$ and $`A_m=๐(๐ธ^m)`$ the result above was proved by Gelfand and Kirillov in . They introduced a new invariant of an algebra $`A`$, so-called the (Gelfand-Kirillov) transcendence degree $`\mathrm{GKtr}.\mathrm{deg}(A)`$, and proved that $`\mathrm{GKtr}.\mathrm{deg}(D_n)=2n`$. Recall that
$$\mathrm{GKtr}.\mathrm{deg}(A):=\underset{V}{sup}\underset{b}{inf}\mathrm{GK}(K[bV])$$
where $`V`$ ranges over the finite dimensional subspaces of $`A`$ and $`b`$ ranges over the regular elements of $`A`$. Another proofs of the corollary based on different ideas were given by A. Joseph and R. Resco , see also , 6.6.19. Josephโs proof is based on the fact that the centralizer of any isomorphic copy of the Weyl algebra $`A_n`$ in its division algebra $`D_n:=D(๐ธ^n)`$ reduces to scalars (, 4.2), Resco proved that $`\mathrm{Cdim}(D_n)=n`$ (, 4.2) using the result of Rentschler and Gabriel that $`\mathrm{K}.\mathrm{dim}(A_n)=n`$ (over an arbitrary field of characteristic zero).
## 8 Filter Dimension and Isotropic Subalgebras of Poisson Algebras
In this section, we apply Theorem 7.2 to obtain an upper bound for the Gelfand-Kirillov dimension of isotropic subalgebras of certain Poisson algebras (Theorem 8.1).
Let $`(P,\{,\})`$ be a Poisson algebra over the field $`K`$. Recall that $`P`$ is an associative commutative $`K`$-algebra which is a Lie algebra with respect to the bracket $`\{,\}`$ for which Leibnizโs rule holds:
$$\{a,xy\}=\{a,x\}y+x\{a,y\}\mathrm{for}\mathrm{all}a,x,yP,$$
which means that the inner derivation $`\mathrm{ad}(a):PP`$, $`x\{a,x\}`$, of the Lie algebra $`P`$ is also a derivation of the associative algebra $`P`$. Therefore, to each Poisson algebra $`P`$ one can attach an associative subalgebra $`A(P)`$ of the ring of differential operators $`๐(P)`$ with coefficients from the algebra $`P`$ which is generated by $`P`$ and $`\mathrm{ad}(P):=\{\mathrm{ad}(a)|aP\}`$. If $`P`$ is a finitely generated algebra then so is the algebra $`A(P)`$ with $`\mathrm{GK}(A(P))\mathrm{GK}(๐(P))<\mathrm{}`$.
Example. Let $`P_{2n}=K[x_1,\mathrm{},x_{2n}]`$ be the Poisson polynomial algebra over a field $`K`$ of characteristic zero equipped with the Poisson bracket
$$\{f,g\}=\underset{i=1}{\overset{n}{}}(\frac{f}{x_i}\frac{g}{x_{n+i}}\frac{f}{x_{n+i}}\frac{g}{x_i}).$$
The algebra $`A(P_{2n})`$ is generated by the elements
$$x_1,\mathrm{},x_{2n},\mathrm{ad}(x_i)=\frac{}{x_{n+i}},\mathrm{ad}(x_{n+i})=\frac{}{x_i},i=1,\mathrm{},n.$$
So, the algebra $`A(P_{2n})`$ is canonically isomorphic to the Weyl algebra $`A_{2n}`$.
Definition. We say that a Poisson algebra $`P`$ is a strongly simple Poisson algebra if
1. $`P`$ is a finitely generated (associative) algebra which is a domain,
2. the algebra $`A(P)`$ is central simple, and
3. for each set of algebraically independent elements $`a_1,\mathrm{},a_m`$ of the algebra $`P`$ such that $`\{a_i,a_j\}=0`$ for all $`i,j=1,\mathrm{},m`$ the (commuting) elements $`a_1,\mathrm{},a_m,\mathrm{ad}(a_1),\mathrm{},`$ $`\mathrm{ad}(a_m)`$ of the algebra $`A(P)`$ are algebraically independent.
###### Theorem 8.1
() Let $`P`$ be a strongly simple Poisson algebra, and $`C`$ be an isotropic subalgebra of $`P`$, i.e. $`\{C,C\}=0`$. Then
$$\mathrm{GK}(C)\frac{\mathrm{GK}(A(P))}{2}\left(1\frac{1}{f_{A(P)}+\mathrm{max}\{f_{A(P)},1\}}\right)$$
where $`f_{A(P)}:=\mathrm{max}\{\mathrm{d}_{Q_m}(Q_mA(P))|\mathrm{\hspace{0.17em}0}m\mathrm{GK}(A(P))\}`$.
Proof. By the assumption the finitely generated algebra $`P`$ is a domain, hence the finitely generated algebra $`A(P)`$ is a domain (as a subalgebra of the domain $`๐(Q(P))`$, the ring of differential operators with coefficients from the field of fractions $`Q(P)`$ for the algebra $`P`$). It suffices to prove the inequality for isotropic subalgebras of the Poisson algebra $`P`$ that are polynomial algebras. So, let $`C`$ be an isotropic polynomial subalgebra of $`P`$ in $`m`$ variables, say $`a_1,\mathrm{},a_m`$. By the assumption, the commuting elements $`a_1,\mathrm{},a_m,\mathrm{ad}(a_1),\mathrm{},\mathrm{ad}(a_m)`$ of the algebra $`A(P)`$ are algebraically independent. So, the Gelfand-Kirillov dimension of the subalgebra $`C^{}`$ of $`A(P)`$ generated by these elements is equal to $`2m`$. By Theorem 7.2,
$$2\mathrm{G}\mathrm{K}(C)=2m=\mathrm{GK}(C^{})\mathrm{GK}(A(P))\left(1\frac{1}{f_{A(P)}+\mathrm{max}\{f_{A(P)},1\}}\right),$$
and this proves the inequality. $`\mathrm{}`$
###### Corollary 8.2
1. The Poisson polynomial algebra $`P_{2n}=K[x_1,\mathrm{},x_{2n}]`$ (with the Poisson bracket) over a field $`K`$ of characteristic zero is a strongly simple Poisson algebra, the algebra $`A(P_{2n})`$ is canonically isomorphic to the Weyl algebra $`A_{2n}`$.
2. The Gelfand-Kirillov dimension of every isotropic subalgebra of the polynomial Poisson algebra $`P_{2n}`$ is $`n`$.
Proof. $`1`$. The third condition in the definition of strongly simple Poisson algebra is the only statement we have to prove. So, let $`a_1,\mathrm{},a_m`$ be algebraically independent elements of the algebra $`P_{2n}`$ such that $`\{a_i,a_j\}=0`$ for all $`i,j=1,\mathrm{},m`$. One can find polynomials, say $`a_{m+1},\mathrm{},a_{2n}`$, in $`P_{2n}`$ such that the elements $`a_1,\mathrm{},a_{2n}`$ are algebraically independent, hence the determinant $`d`$ of the Jacobian matrix $`J:=(\frac{a_i}{x_j})`$ is a nonzero polynomial. Let $`X=(\{x_i,x_j\})`$ and $`Y=(\{a_i,a_j\})`$ be, so-called, the Poisson matrices associated with the elements $`\{x_i\}`$ and $`\{a_i\}`$. It follows from $`Y=J^TXJ`$ that $`det(Y)=d^2det(X)0`$ since $`det(X)0`$. The derivations
$$\delta _i:=d^1det\left(\begin{array}{ccccccc}\{a_1,a_1\}& \mathrm{}& \{a_1,a_{i1}\}& \{a_1,\}& \{a_1,a_{i+1}\}& \mathrm{}& \{a_1,a_{2n}\}\\ \{a_2,a_1\}& \mathrm{}& \{a_2,a_{i1}\}& \{a_2,\}& \{a_2,a_{i+1}\}& \mathrm{}& \{a_2,a_{2n}\}\\ & & & \mathrm{}& & & \\ \{a_{2n},a_1\}& \mathrm{}& \{a_{2n},a_{i1}\}& \{a_{2n},\}& \{a_1,a_{i+1}\}& \mathrm{}& \{a_{2n},a_{2n}\}\end{array}\right),$$
$`i=1,\mathrm{},2n`$, of the rational function field $`Q_{2n}=K(x_1,\mathrm{},x_{2n})`$ satisfy the following properties: $`\delta _i(a_j)=\delta _{i,j}`$, the Kronecker delta. For each $`i`$ and $`j`$, the kernel of the derivation $`\mathrm{\Delta }_{ij}:=\delta _i\delta _j\delta _j\delta _i\mathrm{Der}_K(Q_{2n})`$ contains $`2n`$ algebraically independent elements $`a_1,\mathrm{},a_{2n}`$. Hence $`\mathrm{\Delta }_{ij}=0`$ since the field $`Q_{2n}`$ is algebraic over its subfield $`K(a_1,\mathrm{},a_{2n})`$ and $`\mathrm{char}(K)=0`$. So, the subalgebra, say $`W`$, of the ring of differential operators $`๐(Q_{2n})`$ generated by the elements $`a_1,\mathrm{},a_{2n},\delta _1,\mathrm{},\delta _{2n}`$ is isomorphic to the Weyl algebra $`A_{2n}`$, and so $`\mathrm{GK}(W)=\mathrm{GK}(A_{2n})=4n`$.
Let $`U`$ be the $`K`$-subalgebra of $`๐(Q_{2n})`$ generated by the elements $`x_1,\mathrm{},x_{2n},\delta _1,\mathrm{},\delta _{2n}`$, and $`d^1`$. Let $`P^{}`$ be the localization of the polynomial algebra $`P_{2n}`$ at the powers of the element $`d`$. Then $`\delta _1,\mathrm{},\delta _{2n}_{i=1}^{2n}P^{}\mathrm{ad}(a_i)`$ and $`\mathrm{ad}(a_1),\mathrm{},\mathrm{ad}(a_{2n})_{i=1}^{2n}P^{}\delta _i`$, hence the algebra $`U`$ is generated (over $`K`$) by $`P^{}`$ and $`\mathrm{ad}(a_1),\mathrm{},\mathrm{ad}(a_{2n})`$. The algebra $`U`$ can be viewed as a subalgebra of the ring of differential operators $`๐(P^{})`$. Now, the inclusions, $`WU๐(P^{})`$ imply $`4n=\mathrm{GK}(W)\mathrm{GK}(U)\mathrm{GK}(๐(P^{}))=2\mathrm{G}\mathrm{K}(P^{})=4n`$, therefore $`\mathrm{GK}(U)=4n`$. The algebra $`U`$ is a factor algebra of an iterated Ore extension $`V=P^{}[t_1;\mathrm{ad}(a_1)]\mathrm{}[t_{2n};\mathrm{ad}(a_{2n})]`$. Since $`P^{}`$ is a domain, so is the algebra $`V`$. The algebra $`P^{}`$ is a finitely generated algebra of Gelfand-Kirillov dimension $`2n`$, hence $`\mathrm{GK}(V)=\mathrm{GK}(P^{})+2n=4n`$ (by , 8.2.11). Since $`\mathrm{GK}(V)=\mathrm{GK}(U)`$ and any proper factor algebra of $`V`$ has Gelfand-Kirillov dimension strictly less than $`\mathrm{GK}(V)`$ (by , 8.3.5, since $`V`$ is a domain), the algebras $`V`$ and $`U`$ must be isomorphic. Therefore, the (commuting) elements $`a_1,\mathrm{},a_m,\mathrm{ad}(a_1),\mathrm{},`$ $`\mathrm{ad}(a_m)`$ of the algebra $`U`$ (and $`A(P)`$) must be algebraically independent.
$`2`$. Let $`C`$ be an isotropic subalgebra of the Poisson algebra $`P_{2n}`$. Note that $`f_{A(P_{2n})}=f_{A_{2n}}=1`$ and $`\mathrm{GK}(A_{2n})=4n`$. By Theorem 8.1,
$$\mathrm{GK}(C)\frac{4n}{2}(1\frac{1}{1+1})=n.\mathrm{}$$
Remark. This result means that for the Poisson polynomial algebra $`P_{2n}`$ the right hand side of the inequality of Theorem 8.1 is the exact upper bound for the Gelfand-Kirillov dimension of isotropic subalgebras in $`P_{2n}`$ since the polynomial subalgebra $`K[x_1,\mathrm{},x_n]`$ of $`P_{2n}`$ is isotropic.
Department of Pure Mathematics
University of Sheffield
Hicks Building
Sheffield S3 7RH
UK
email: v.bavula@sheffield.ac.uk
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# Particle number renormalization in almost half filled Mott Hubbard superconductors
## I Introduction
Recently, Anderson has underscored the importance of fugacity in wave functions that do not conserve particle number pwa\_tunnel . Following an earlier paper by Laughlin laughlin\_02 , Anderson argued that a fugacity factor should be included in variational wave functions of the form,
$$P|\mathrm{\Psi }_{\mathrm{BCS}}=P\underset{k}{}\left(u_k+v_kc_k^{}c_k^{}\right)|0.$$
(1)
Here, $`P=_i(1n_in_i)`$ is a (Gutzwiller) projection operator which excludes double occupancies at sites $`i`$ (Ref. gutzwiller\_63, ), and $`|\mathrm{\Psi }_{\mathrm{BCS}}`$, a BCS wave function. Projected wave functions of this form were originally proposed to describe the phase diagram of doped Mott Hubbard insulators such as the high temperature superconductors gros\_88 ; RMFT ; atoz ; ys87 . Detailed variational Monte Carlo (VMC) studies have been carried out recently using projected $`d`$-wave BCS states as variational wave functions for the two dimensional Hubbard model, after a suitable canonical transformation paramekanti ; gros87\_almost .
Despite their simple form, projected wave functions exhibit nontrivial properties because the projection operator acts on a quantum many body state. The action of the projection operator (in reducing the allowed states in the Hilbert space) concomitant with the correlations of the quantum state being projected, leads to a variety of physical phenomena variety . Other interesting effects include non-trivial matrix-element renormalization near half-filling FEMG05 and the occurrence of superconductivity near an (antiferromagnetic) Mott insulator afmi .
Approximate analytical calculations with wave functions such as Eq. (1) can be done using a renormalization scheme based on the Gutzwiller approximation. Within this approximation, the effects of projection on the state $`|\mathrm{\Psi }_0`$ are approximated by a classical statistical weight factor multiplying the quantum result vollhardt\_84 ; extensions . Thus, for example,
$$\frac{\mathrm{\Psi }|\widehat{O}|\mathrm{\Psi }}{\mathrm{\Psi }|\mathrm{\Psi }}g\frac{\mathrm{\Psi }_0|\widehat{O}|\mathrm{\Psi }_0}{\mathrm{\Psi }_0|\mathrm{\Psi }_0},$$
(2)
where $`\widehat{O}`$ is any operator, and $`g`$, the so called Gutzwiller factor. For example, the Gutzwiller approximation for the kinetic energy operator $`c_i^{}c_j+c_j^{}c_i`$ and the superexchange interaction between sites $`i`$ and $`j`$, $`\stackrel{}{S}_i\stackrel{}{S}_j`$ yields the Gutzwiller factors,
$$g_t=\frac{1n}{1n/2},g_s=\frac{1}{(1n/2)^2},$$
(3)
where $`n`$ is the density of electrons. In deriving these renormalization factors, one considers the number of terms that contribute to $`\mathrm{\Psi }|\widehat{O}|\mathrm{\Psi }`$ and to $`\mathrm{\Psi }_0|\widehat{O}|\mathrm{\Psi }_0`$ respectively. The ratio of these two contributions is the renormalization factor.
The renormalization factors are functions of the local charge density. This is a well defined quantity when one considers for example, a projected Fermi liquid state,
$$P|\mathrm{\Psi }_{\mathrm{FS}}=P\underset{k<k_F}{}c_k^{}c_k^{}|0.$$
(4)
But suppose instead, we consider wave functions such as the BCS state in Eq. (1), where the particle number fluctuates. In this case, it is not clear what the local charge density in Eq. (3) should be. It may be argued that the correct $`n`$ is set by the average particle number $`\overline{N}`$. But then, projecting a BCS state changes the average particle number; i.e., the average number of electrons in $`|\mathrm{\Psi }_{\mathrm{BCS}}`$ does not equal that in $`P|\mathrm{\Psi }_{\mathrm{BCS}}`$. Clearly, we need a scheme to keep track of this effect.
Note that this problem can be avoided completely, as is done in most variational Monte Carlo (VMC) studies. Here, the particle number is fixed (one works in a canonical ensemble), and Eq. (1) replaced by,
$$P_NP|\mathrm{\Psi }_{\mathrm{BCS}}=P_NP\underset{k}{}\left(u_k+v_kc_k^{}c_k^{}\right)|0.$$
(5)
The operator $`P_N`$ fixes the particle number, and the issue of projection changing the mean particle number does not arise gros\_88 . However, there are also other VMC studies with wave functions that do not have fixed particle number yokoyama\_88 . Moreover, we are often interested in carrying out analytical approximations in the spirit of Eq. (2). Since such manipulations are easier done with BCS wave functions (where the particle number is not fixed), it is desirable to understand the effects of the projection operator on this class of wave functions. In this paper, we present analytical and numerical considerations of this problem. In doing so, we clarify the notion of fugacity introduced by Anderson pwa\_tunnel . We also discuss the relevance of this approach for the Gutzwiller approximation in the grand canonical scheme and the corresponding VMC studies yokoyama\_88 .
## II The fugacity factor
Consider the projected BCS wave function, Eq. (1). It is clear that the projection operator $`P`$ changes the average number, *viz.*,
$$\frac{\mathrm{\Psi }_{\mathrm{BCS}}|\widehat{N}|\mathrm{\Psi }_{\mathrm{BCS}}}{\mathrm{\Psi }_{\mathrm{BCS}}|\mathrm{\Psi }_{\mathrm{BCS}}}\frac{\mathrm{\Psi }_{\mathrm{BCS}}|P\widehat{N}P|\mathrm{\Psi }_{\mathrm{BCS}}}{\mathrm{\Psi }_{\mathrm{BCS}}|P^2|\mathrm{\Psi }_{\mathrm{BCS}}}.$$
The effect of the projection operator can be seen most clearly by examining the particle number distribution in the unprojected and projected Hilbert spaces. Towards this end, let us write the average numbers, $`\overline{N}^{(0)}(\overline{N})`$ in the unprojected (projected) Hilbert space
$`\overline{N}^{(0)}`$ $`={\displaystyle \underset{N}{}}N\rho _N^{(0)},`$ (6)
$`\overline{N}`$ $`={\displaystyle \underset{N}{}}N\rho _N.`$ (7)
Here,
$`\rho _N^{(0)}`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Psi }_{\mathrm{BCS}}|P_N|\mathrm{\Psi }_{\mathrm{BCS}}}{\mathrm{\Psi }_{\mathrm{BCS}}|\mathrm{\Psi }_{\mathrm{BCS}}}},`$
$`\rho _N`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Psi }_{\mathrm{BCS}}|PP_NP|\mathrm{\Psi }_{\mathrm{BCS}}}{\mathrm{\Psi }_{\mathrm{BCS}}|PP|\mathrm{\Psi }_{\mathrm{BCS}}}},`$
are the particle number distributions in the unprojected and projected BCS wave functions respectively; $`P_N`$ is an operator which projects onto terms with particle number $`N`$. The particle number distributions before and after projection may be related by
$`\underset{\rho _N}{\underset{}{{\displaystyle \frac{\mathrm{\Psi }_{\mathrm{BCS}}|PP_NP|\mathrm{\Psi }_{\mathrm{BCS}}}{\mathrm{\Psi }_{\mathrm{BCS}}|PP|\mathrm{\Psi }_{\mathrm{BCS}}}}}}=g_N\underset{\rho _N^{(0)}}{\underset{}{{\displaystyle \frac{\mathrm{\Psi }_{\mathrm{BCS}}|P_N|\mathrm{\Psi }_{\mathrm{BCS}}}{\mathrm{\Psi }_{\mathrm{BCS}}|\mathrm{\Psi }_{\mathrm{BCS}}}}}},`$ (8)
where
$`g_N=\underset{=C(=\mathrm{const})}{\underset{}{{\displaystyle \frac{\mathrm{\Psi }_{\mathrm{BCS}}|\mathrm{\Psi }_{\mathrm{BCS}}}{\mathrm{\Psi }_{\mathrm{BCS}}|PP|\mathrm{\Psi }_{\mathrm{BCS}}}}}}{\displaystyle \frac{\mathrm{\Psi }_{\mathrm{BCS}}|PP_NP|\mathrm{\Psi }_{\mathrm{BCS}}}{\mathrm{\Psi }_{\mathrm{BCS}}|P_N|\mathrm{\Psi }_{\mathrm{BCS}}}}.`$
Eq. (8) constitutes the Gutzwiller approximation for the projection operator $`P_N`$ with the corresponding renormalization factor, $`g_N`$; $`C`$ is an irrelevant constant (the ratio of the normalization of the unprojected and projected wave functions), which does not depend on $`N`$. Following Gutzwiller, we estimate $`g_N`$ by combinatorial means, as being equal to the ratio of the relative sizes of the projected and unprojected Hilbert spaces. Then,
$`g_NC{\displaystyle \frac{\frac{L!}{(LN_{}N_{})!N_{}!N_{}!}}{\frac{L!}{(LN_{})!N_{}!}\frac{L!}{(LN_{})!N_{}!}}},`$ (9)
where $`L`$ is the number of lattice sites and $`N_{}`$ ($`N_{}`$) is the number of up (down)-spins. Since in a BCS wave function, $`N_{}=N_{}=N/2`$, $`N`$ being the total number of particles, the expression for $`g_N`$ can be simplified to
$`g_NC{\displaystyle \frac{((LN/2)!)^2}{L!(LN)!}}.`$ (10)
Hence, if we were to impose the condition that the average particle number before and after projection be identical, a factor $`g_N^1`$ has to be included in Eq. (7). Then, from Eq. (7) and Eq. (8), we obtain the particle number after projection $`\overline{N}_{\mathrm{new}}`$,
$`\overline{N}_{\mathrm{new}}`$ $`{\displaystyle \underset{N}{}}N{\displaystyle \frac{1}{g_N}}\rho _N={\displaystyle \underset{N}{}}N{\displaystyle \frac{g_N\rho _N^{(0)}}{g_N}}`$
$`=\overline{N}^{(0)},`$ (11)
which is the desired result.
Now, let us show how this procedure can be implemented for the wave function $`|\mathrm{\Psi }_{\mathrm{BCS}}`$. Since the BCS wave function is a linear superposition of states with particle number $`\mathrm{},N2,N,N+2,\mathrm{}`$, we consider the effect of projection on two states whose particle numbers differ by $`2`$. Then, the ratio,
$$f^2\frac{g_{N+2}}{g_N}\left(\frac{LN}{LN/2}\right)^2,$$
(12)
in the thermodynamic limit. Eq. 12 shows that the projection operator acts unequally on the $`N`$ and $`N+2`$ particle states; the renormalization of the weight of the $`N+2`$ particle states $`g_{N+2}`$, is $`f^2`$ times the weight of the $`N`$ particle states, $`g_N`$. This effect can be rectified as in Eq. (11) by multiplying every Cooper pair $`c_k^{}c_k^{}`$ by a factor $`\frac{1}{f}`$ in the BCS wave function. It produces the desired result, viz., the projected and unprojected BCS wave functions have the same average particle number.
Alternatively (following Anderson), we can multiply every empty state by the factor $`f`$ and write,
$`|\mathrm{\Psi }_{\mathrm{BCS}}^{(f)}={\displaystyle \underset{k}{}}{\displaystyle \frac{\left(fu_k+v_kc_k^{}c_k^{}\right)}{\sqrt{f^2|u_k|^2+|v_k|^2}}}|0.`$ (13)
Then again by construction, the fugacity factor $`f`$ in Eq. (13) ensures that the projected wave function $`P|\mathrm{\Psi }_{\mathrm{BCS}}^{(f)}`$ and the unprojected wave function $`|\mathrm{\Psi }_{\mathrm{BCS}}`$ have the same particle number. The denominator in Eq. (13) is the new normalization factor.
The following points are in order: (a) the fugacity factor $`f`$ in Eq. (12) depends on the variable particle number $`N`$. However, since the particle number of the BCS wave function is sharply peaked within the range, $`\overline{N}^{(0)}\sqrt{\overline{N}^{(0)}}`$ and $`\overline{N}^{(0)}+\sqrt{\overline{N}^{(0)}}`$, we will assume that the fugacity factor $`f=f(\overline{N}^{(0)})`$ in the thermodynamic limit; (b) in this limit, Eq. (12) reduces to $`f^2=g_t^2`$, where $`g_t`$ is the Gutzwiller factor defined in Eq. (3). Then, Eq. (13) reduces to
$`|\mathrm{\Psi }_{\mathrm{BCS}}^{(f)}={\displaystyle \underset{k}{}}{\displaystyle \frac{\left(g_tu_k+v_kc_k^{}c_k^{}\right)}{\sqrt{g_t^2|u_k|^2+|v_k|^2}}}|0,`$ (14)
which is the wave function proposed by Anderson pwa\_tunnel ; (c) the fugacity factor ensures that projection affects the $`N`$ and $`N+2`$ particle states of the BCS wave function in the same way. In principle, such a factor could depend on the $`k`$-value, but in this paper, we will treat it as a mere combinatorial device; (d) the combinatorial argument fails at half filling when $`L=\overline{N}^{(0)}`$.
## III Particle number renormalization in projected BCS wave functions
In the previous section, we showed that the inclusion of the fugacity factor is necessary for the average particle number in a BCS wave function to be unchanged by projection. Alternatively, one might ask what is the effect of the projection operator on a BCS wave function; viz., if projection changes the mean particle number of a BCS state, how are the particle numbers before and after projection related? In this section, we will use the fugacity factor to answer this question. In particular, we will show how particle density after projection can be determined self consistently by including the fugacity factor in the projected BCS state (Eq. (14)). An alternate derivation of this result is presented in Appendix A, where the effect of the projection operator on particle number fluctuations is calculated.
Consider two BCS states defined by
$`|\mathrm{\Psi }_{\mathrm{BCS}}`$ $`={\displaystyle \underset{k}{}}\left(u_k+v_kc_k^{}c_k^{}\right)|0,`$ (15)
$`|\mathrm{\Psi }_{\mathrm{BCS}}^{(r)}`$ $`={\displaystyle \underset{k}{}}{\displaystyle \frac{\left(u_k+g_tv_kc_k^{}c_k^{}\right)}{\sqrt{|u_k|^2+g_t^2|v_k|^2}}}|0`$
$`={\displaystyle \underset{k}{}}\left(u_k^{(r)}+v_k^{(r)}c_k^{}c_k^{}\right)|0,`$ (16)
where,
$`u_k^{(r)}`$ $`{\displaystyle \frac{u_k}{\sqrt{|u_k|^2+g_t^2|v_k|^2}}},`$ (17)
$`v_k^{(r)}`$ $`{\displaystyle \frac{g_tv_k}{\sqrt{|u_k|^2+g_t^2|v_k|^2}}}.`$ (18)
From Eq. (12), it is clear that the projection operator reduces the ratio of the weights of $`N+2`$ and $`N`$ particle states in a BCS wave function by a factor $`g_t`$. Then, it follows that
$`{\displaystyle \frac{\mathrm{\Psi }_{\mathrm{BCS}}^{(r)}|\widehat{N}|\mathrm{\Psi }_{\mathrm{BCS}}^{(r)}}{\mathrm{\Psi }_{\mathrm{BCS}}^{(r)}|\mathrm{\Psi }_{\mathrm{BCS}}^{(r)}}}{\displaystyle \frac{\mathrm{\Psi }_{\mathrm{BCS}}|P\widehat{N}P|\mathrm{\Psi }_{\mathrm{BCS}}}{\mathrm{\Psi }_{\mathrm{BCS}}|P|\mathrm{\Psi }_{\mathrm{BCS}}}}`$ (19)
when
$`g_t={\displaystyle \frac{L\overline{N}^{(r)}}{L\overline{N}^{(r)}/2}}.`$ (20)
The average particle number $`\overline{N}^{(r)}`$ of the state $`|\mathrm{\Psi }_{\mathrm{BCS}}^{(r)}`$ is given by,
$$\overline{N}^{(r)}=\frac{\mathrm{\Psi }_{\mathrm{BCS}}^{(r)}|\widehat{N}|\mathrm{\Psi }_{\mathrm{BCS}}^{(r)}}{\mathrm{\Psi }_{\mathrm{BCS}}^{(r)}|\mathrm{\Psi }_{\mathrm{BCS}}^{(r)}}=2\underset{k}{}|v_k^{(r)}|^2.$$
(21)
Since the particle numbers of $`|\mathrm{\Psi }_{\mathrm{BCS}}^{(r)}`$ and $`P|\mathrm{\Psi }_{\mathrm{BCS}}`$ are identical, we can use Eq. (18) in Eq. (21) to obtain,
$`\overline{N}^{(r)}\overline{N}_{\mathrm{after}}`$ $`={\displaystyle \frac{\mathrm{\Psi }_{\mathrm{BCS}}|P\widehat{N}P|\mathrm{\Psi }_{\mathrm{BCS}}}{\mathrm{\Psi }_{\mathrm{BCS}}|P|\mathrm{\Psi }_{\mathrm{BCS}}}}`$
$`2{\displaystyle \underset{k}{}}{\displaystyle \frac{g_t^2|v_k|^2}{|u_k|^2+g_t^2|v_k|^2}}.`$ (22)
Note that $`g_t`$ is specified by the particle number after projection, $`N_{\mathrm{after}}(=\overline{N}^{(r)})`$.
Now, since the number of particles in the state $`|\mathrm{\Psi }_{\mathrm{BCS}}`$ before projection is given by
$$\overline{N}_{\mathrm{before}}=2\underset{k}{}|v_k|^2,$$
Eq. (22) provides us with a way to calculate the number of particles in the state $`P|\mathrm{\Psi }_{\mathrm{BCS}}`$ *after* projection, if $`|\mathrm{\Psi }_{\mathrm{BCS}}`$ (i.e., $`u_k`$ and $`v_k`$) is specified before projection. Eq. (22) can be solved self consistently for $`\overline{N}_{\mathrm{after}}`$. We solve Eq. (22) numerically on a square lattice, using the standard BCS expressions for a $`d`$-wave superconductor, $`u_k`$($`v_k`$):
$`v_k^2`$ $`={\displaystyle \frac{1}{2}}\left(1{\displaystyle \frac{\xi _k}{E_k}}\right),`$ (23)
$`u_k^2`$ $`={\displaystyle \frac{1}{2}}\left(1+{\displaystyle \frac{\xi _k}{E_k}}\right),`$ (24)
where,
$`E_k`$ $`=\left(\mathrm{\Delta }_k^2+\xi _k^2\right)^{\frac{1}{2}},`$ (25)
$`\mathrm{\Delta }_k`$ $`=\mathrm{\Delta }_0(\mathrm{cos}(k_x)\mathrm{cos}(k_y)),`$ (26)
$`\xi _k`$ $`=2(\mathrm{cos}(k_x)+\mathrm{cos}(k_y))\mu .`$ (27)
The only free parameters are the chemical potential $`\mu `$ and the variational parameter $`\mathrm{\Delta }_0`$.
By fixing the parameter $`\mathrm{\Delta }_0`$, we determine the particle numbers (before and after projection) for various chemical potentials. The results for particle density ($`n\overline{N}/L`$) are shown in Fig. 1. The results clearly show that the particle density before projection attains its maximal value ($`n_{\mathrm{before}}=2`$), if $`n_{\mathrm{after}}=1`$ (half-filling). This result holds for any finite value of the variational parameter $`\mathrm{\Delta }_0`$. In the opposite limit, viz., low density of electrons, $`n_{\mathrm{before}}`$ converges to the value of $`n_{\mathrm{after}}`$ as expected. The size of the intermediate region depends on the magnitude of the parameter $`\mathrm{\Delta }_0`$, as illustrated by the results in Fig. 1.
The accuracy of Eq. (22) can be checked by comparing our results with those of Yokoyama and Shiba (YS), who performed VMC studies of projected BCS wave functions with fluctuating particle number (but without the fugacity factor) yokoyama\_88 . They determined the particle density of the projected $`d`$-wave state $`P|\mathrm{\Psi }_{\mathrm{BCS}}`$ as a function of the chemical potential $`\mu `$ and the variational parameter $`\mathrm{\Delta }_0`$, within a grand canonical scheme. The unprojected wave function $`|\mathrm{\Psi }_{\mathrm{BCS}}`$ is specified as usual, through Eq. (23)-Eq. (27). Since YS do not include a fugacity factor in their definition of the BCS wave function, projection changes the particle number. So, we use Eq. (22) to determine $`n_{\mathrm{after}}`$ which we compare with their results for particle number.
As seen in Fig. 2, our results are in good qualitative agreement with YS. Discrepancies are mostly due to finite corrections. YS use $`6\times 6`$ and $`8\times 8`$-lattices, while our analytic calculations are for the thermodynamic limitinfiniteS . The results show the singular effect of the projection near the insulating phase (half filling). The chemical potential goes to infinity in this limit.
In Appendix A, we present an alternate derivation of Eq. (22) by a saddle point approximation without using a fugacity-corrected wave function. This approach also allows for the calculation of particle number fluctuations $`\sigma ^2/L`$.
## IV Gutzwiller approximation in the canonical and grand canonical schemes
In this section, we discuss the differences between the Gutzwiller approximation in the canonical and grand canonical schemes. The validity of our statements can be checked by a comparison to nearly exact VMC yokoyama\_88 ; gros\_88 ; paramekanti ; yokoyama\_96 .
Let us first consider the canonical case. Here, we are interested in the expectation value of an operator $`\widehat{O}`$ calculated with a particle number conserving projected wave function $`P_NP|\mathrm{\Psi }_{\mathrm{BCS}}`$. The corresponding Gutzwiller approximation can be understood as follows:
$`{\displaystyle \frac{\mathrm{\Psi }_{\mathrm{BCS}}|PP_N\widehat{O}P_NP|\mathrm{\Psi }_{\mathrm{BCS}}}{\mathrm{\Psi }_{\mathrm{BCS}}|PP_NP|\mathrm{\Psi }_{\mathrm{BCS}}}}`$
$``$ $`g{\displaystyle \frac{\mathrm{\Psi }_{\mathrm{BCS}}|P_N\widehat{O}P_N|\mathrm{\Psi }_{\mathrm{BCS}}}{\mathrm{\Psi }_{\mathrm{BCS}}|P_N|\mathrm{\Psi }_{\mathrm{BCS}}}}`$
$`=`$ $`g{\displaystyle \frac{\mathrm{\Psi }_{\mathrm{BCS}}|\widehat{O}|\mathrm{\Psi }_{\mathrm{BCS}}}{\mathrm{\Psi }_{\mathrm{BCS}}|\mathrm{\Psi }_{\mathrm{BCS}}}},`$ (28)
where $`P_N`$ is the projector on the terms with particle number $`N`$. The Gutzwiller factor $`g`$, corresponds to the operator $`\widehat{O}`$. The first row represents a quantity which can be calculated exactly by fixed particle number VMC gros\_88 ; paramekanti . Since the particle number is fixed, the Gutzwiller approximation can be invoked, leading to the second row. The equality to the third row is guaranteed only if $`N`$ is equal to the average particle number of $`|\mathrm{\Psi }_{\mathrm{BCS}}`$ ($`N=\overline{N}`$). Here, we perform a transformation from a canonical to a grand canonical ensemble, which is valid in the thermodynamic limit.
In the grand canonical scheme, where we calculate the expectation value of $`\widehat{O}`$ with a particle number non-conserving wave function, this scheme must be modified as follows:
$$\frac{\mathrm{\Psi }_{\mathrm{BCS}}^{(f)}|P\widehat{O}P|\mathrm{\Psi }_{\mathrm{BCS}}^{(f)}}{\mathrm{\Psi }_{\mathrm{BCS}}^{(f)}|PP|\mathrm{\Psi }_{\mathrm{BCS}}^{(f)}}g\frac{\mathrm{\Psi }_{\mathrm{BCS}}|\widehat{O}|\mathrm{\Psi }_{\mathrm{BCS}}}{\mathrm{\Psi }_{\mathrm{BCS}}|\mathrm{\Psi }_{\mathrm{BCS}}},$$
(29)
where $`P|\mathrm{\Psi }_{\mathrm{BCS}}^{(f)}`$ is the projected $`d`$-wave state corrected for fugacity, i.e., a fugacity factor is included simultaneously with the projection (see Sec. II). This correction is essential to guarantee the validity of the Gutzwiller approximation; without it, the lhs and rhs of Eq. (29) would correspond to states with different particle numbers.
Comparing Eq. (28) and Eq. (29) we get,
$$\frac{\mathrm{\Psi }_{\mathrm{BCS}}|PP_N\widehat{O}P_NP|\mathrm{\Psi }_{\mathrm{BCS}}}{\mathrm{\Psi }_{\mathrm{BCS}}|PP_NP|\mathrm{\Psi }_{\mathrm{BCS}}}\frac{\mathrm{\Psi }_{\mathrm{BCS}}^{(f)}|P\widehat{O}P|\mathrm{\Psi }_{\mathrm{BCS}}^{(f)}}{\mathrm{\Psi }_{\mathrm{BCS}}^{(f)}|P^2|\mathrm{\Psi }_{\mathrm{BCS}}^{(f)}}.$$
(30)
Eq. (29) and Eq. (30) constitute the main results of this section. Eq. (29) shows that when the Gutzwiller approximation is used for a wave function which does not have a fixed particle number, a fugacity factor must be included along with the projection. Eq. (30) shows that to obtain identical results, one has to use different wave functions in the grand canonical (rhs) and canonical (lhs) schemes. The wave function $`|\mathrm{\Psi }_{\mathrm{BCS}}^{(f)}`$ is a $`d`$-wave state corrected by the fugacity factor, whereas $`|\mathrm{\Psi }_{\mathrm{BCS}}`$ is a pure $`d`$-wave state. Our arguments leading up to Eq. (29) and Eq. (30) can be verified by a comparison with VMC studies. We now proceed to do so.
The expectation values in the canonical and grand canonical schemes can be calculated (nearly exactly) by VMC studies. In Fig. 3, we show VMC results from Gros gros\_88 (fixed particle number VMC, canonical) and from YS yokoyama\_88 (grand canonical VMC). The discrepancy between the two sets of results can be explained readily by Eq. (30). In the case of $`\mathrm{\Delta }_00`$, there is only small room for particle number fluctuations even in the particle non-conserving wavefunction. Then, canonical and grand canonical schemes should give identical results. The VMC calculations in Fig. 3 do not exactly show thit behavior since the grand-canonical scheme becomes inaccurate in this limit yokoyama\_88 . YS consider a pure $`d`$-wave state, i.e., the fugacity factor is not included in their calculations. In their paper, YS argued that the discrepancies between the two results can be removed by introducing an additional variational parameter $`\alpha `$, so that $`a_kv_k/u_k`$ is replaced by $`a_k\alpha v_k/u_k`$ (Eq. 4.1 in Ref. yokoyama\_88, ). We opine that the parameter $`\alpha `$ is directly related to our fugacity factor, i.e., $`\alpha =g_t`$ in the wave function $`|\mathrm{\Psi }_{\mathrm{BCS}}^{(f)}`$. This conclusion is supported by the comparison of VMC data to the corresponding Gutzwiller approximation (see below).
The validity of the approximation in the canonical case (Eq. (28)) is well accepted. It is used for instance, in the renormalized mean field theory (RMFT) of Zhang *et al.*, where all physical quantities are calculated using unprojected wave functions and the corresponding Gutzwiller renormalization factors RMFT . A comparison with VMC studies with fixed particle number shows good agreement RMFT (also illustrated in Fig. 3).
To compare the grand canonical VMC of YS with the Gutzwiller approximation, we need to modify Eq. (29). This is necessary because YS do not include the fugacity factor in their considerations, as pointed out earlier. We modify Eq. (30) by the following procedure:
(i) we start with a $`d`$-wave BCS state $`|\mathrm{\Psi }_{\mathrm{BCS}}`$ for specified values of $`\mathrm{\Delta }_0`$;
(ii) we use Eq. (22) to determine the chemical potential $`\mu `$. This fixes the particle density $`n_{\mathrm{after}}`$ of $`P|\mathrm{\Psi }_{\mathrm{BCS}}`$;
(iii) we remove the fugacity factor to get $`|\mathrm{\Psi }_{\mathrm{BCS}}^{(r)}`$ via Eq. (16). The fugacity factor is determined for $`n_{\mathrm{after}}`$. $`|\mathrm{\Psi }_{\mathrm{BCS}}^{(r)}`$ and $`P|\mathrm{\Psi }_{\mathrm{BCS}}`$ correspond to the same particle density $`n_{\mathrm{after}}`$.
(iv) The expectation values of the wave function $`P|\mathrm{\Psi }_{\mathrm{BCS}}`$ can now be approximated by $`|\mathrm{\Psi }_{\mathrm{BCS}}^{(r)}`$ and Gutzwiller factors, viz.,
$$\frac{\mathrm{\Psi }_{\mathrm{BCS}}|P\widehat{O}P|\mathrm{\Psi }_{\mathrm{BCS}}}{\mathrm{\Psi }_{\mathrm{BCS}}|PP|\mathrm{\Psi }_{\mathrm{BCS}}}g\frac{\mathrm{\Psi }_{\mathrm{BCS}}^{(r)}|\widehat{O}|\mathrm{\Psi }_{\mathrm{BCS}}^{(r)}}{\mathrm{\Psi }_{\mathrm{BCS}}^{(r)}|\mathrm{\Psi }_{\mathrm{BCS}}^{(r)}}.$$
(31)
This Gutzwiller approximation (GA) generalizes Eq. 2 for wave functions that do not conserve particle number. In Appendix B, we discuss this approximation for the different terms in the $`tJ`$ model.
In Fig. 4, we compare the GA of the kinetic energy $`E^{(1)}`$, and the expectation value $`E^{(2)}`$, of the remaining terms in the $`tJ`$ model ($`\widehat{S}_i\widehat{S}_j`$, $`\widehat{n}_i\widehat{n}_j`$, and the 3-site term) to those from the grand canonical VMC. A good agreement between the VMC and Gutzwiller results is seen, which confirms the validity of our grand canonical Gutzwiller approximation (Eq. (29)).
In Fig. 3 and Fig. 4, we also show Gutzwiller approximations for the fixed particle number VMC gros\_88 . Clearly, canonical and grand canonical approaches yield different energies (as do the corresponding VMC studies ). We emphasize this is because of the projection operator $`P`$, which changes the particle number in a grand canonical scheme. For these two methods to yield the same results, a fugacity corrected wave function must be used when working in a grand canonical ensemble. Hence, all previous speculations about the coincidence of these two VMC schemes in the thermodynamic limit have to be reformulated carefully.
## V Summary
In this paper, we considered the effects of Gutzwiller projection on a state which does not have fixed particle number. We showed that it is necessary to include a fugacity factor when invoking the Gutzwiller approximation for such states. The effects of projecting a number non-conserving BCS state were studied by examining the relation between particle number before and after projection. We obtained an analytical expression, Eq. (22), and compared to variational Monte Carlo data (Fig. 2). We discussed the discrepancies in the VMC results for projected BCS wave functions obtained in the canonical and grand canonical schemes, and presented a resolution. In conclusion, we have clarified several subtle properties of the Gutzwiller projection operator $`P`$ acting on a BCS state, and hope that these results lead to a better understanding of the Gutzwiller approximation in the grand canonical scheme.
We thank P. W. Anderson, N. P. Ong, and H. Yokoyama for several discussions. N. F. was supported by the Deutsche Forschungsgemeinschaft. V. N. M. acknowledges partial financial support from The City University of New York, PSC-CUNY Research Award Program.
## Appendix A Saddle point approximation to mean particle number and number fluctuations in projected BCS wave functions
In Sec. III, we used the fugacity factor to derive Eq. (22). Here, we present an alternative approach by a saddle point approximation to discuss the effects of projection on the mean particle number of a BCS state. This approach also describes the particle number fluctuations after projection.
The particle number distribution for an unprojected BCS wave function $`\rho _N^{(0)}`$ can be written as
$$\rho _N^{(0)}=\frac{2}{N_\sigma !}\left(\frac{\mathrm{d}}{\mathrm{d}\lambda }\right)^{N_\sigma }\underset{k}{}\left(|u_k|^2+|v_k|^2\lambda \right)|_{\lambda 0},$$
(32)
where $`N_\sigma =N/2`$ is the number of electron pairs. This relation can be checked by expanding the product in the wave function,
$$|\mathrm{\Psi }_{\mathrm{BCS}}=\underset{k}{}\left(u_k+v_kc_{k,}^{}c_{k,}^{}\right)|0,$$
and considering contribution to $`\mathrm{\Psi }_{\mathrm{BCS}}|\mathrm{\Psi }_{\mathrm{BCS}}`$ from each term. In Sec. III, we showed that the particle number distribution of a projected wave function $`\rho _N`$ is related to the unprojected distribution $`\rho _N^{(0)}`$ by
$$\rho _N=g_N\rho _N^{(0)},$$
(33)
where,
$`g_NC{\displaystyle \frac{((LN/2)!)^2}{L!(LN)!}}.`$
Particle number and number fluctuations of the projected wave function can be derived from Eq. (33) and Eq. (32), upon invoking a saddle point approximation. We define a generating function $`\mathrm{\Xi }_\lambda `$,
$`\mathrm{\Xi }_\lambda `$ $`2{\displaystyle \underset{k}{}}\left(|u_k|^2+|v_k|^2\lambda \right)`$
$`=`$ $`{\displaystyle \underset{N_\sigma =0}{\overset{L}{}}}\lambda ^{N_\sigma }\rho _{2N_\sigma }^{(0)}.`$ (34)
We invert Eq. (34) using a contour integral on the complex $`\lambda `$-plane along a circle around $`\lambda =0`$:
$`\rho _N^{(0)}={\displaystyle \frac{1}{2\pi i}}{\displaystyle \frac{\mathrm{\Xi }_\lambda }{\lambda ^{N_\sigma +1}}๐\lambda }.`$ (35)
Note that, in the integrand, only $`\rho _{2N_\sigma }^{(0)}/\lambda =\rho _N^{(0)}/\lambda `$ gives a finite value. The others powers of $`\lambda `$ vanishes. Multiplying by $`g_N`$ gives
$`\rho _N{\displaystyle \frac{1}{2\pi i}}{\displaystyle g_N\frac{\mathrm{\Xi }_\lambda }{\lambda ^{N_\sigma +1}}๐\lambda }.`$ (36)
Eq. (36) can be written as
$`\rho _N{\displaystyle \frac{1}{2\pi i}}{\displaystyle ๐\lambda e^{f(\lambda ,N)}},`$ (37)
where
$`f(\lambda ,N)=\mathrm{log}\mathrm{\Xi }_\lambda ({\displaystyle \frac{N}{2}}+1)\mathrm{log}\lambda +\mathrm{log}g_N.`$ (38)
Using Stirlingโs formula,
$`\mathrm{log}g_N`$ $`2(L{\displaystyle \frac{N}{2}})\mathrm{log}(L{\displaystyle \frac{N}{2}})(LN)\mathrm{log}(LN)`$
$``$ $`L\mathrm{log}L+\mathrm{log}C`$ (39)
The saddle point ($`\overline{n}=\frac{\overline{N}}{L}`$, $`\overline{\lambda }`$) of Eq. (37) is determined by
$`{\displaystyle \frac{f}{\lambda }}=`$ $`{\displaystyle \frac{\mathrm{log}\mathrm{\Xi }_\lambda }{\lambda }}{\displaystyle \frac{\frac{N}{2}+1}{\lambda }}`$
$`=`$ $`{\displaystyle \frac{}{\lambda }}{\displaystyle \underset{k}{}}\mathrm{log}\left(|u_k|^2+|v_k|^2\lambda \right){\displaystyle \frac{\frac{N}{2}+1}{\lambda }}`$
$`=`$ $`{\displaystyle \underset{k}{}}{\displaystyle \frac{|v_k|^2}{|u_k|^2+|v_k|^2\lambda }}{\displaystyle \frac{\frac{N}{2}+1}{\lambda }}`$
$``$ $`0,`$ (40)
for $`N1`$,
$`\overline{n}`$ $`2{\displaystyle \frac{\frac{\overline{N}}{2}+1}{L}}`$
$`=`$ $`2{\displaystyle \frac{1}{L}}{\displaystyle \underset{k}{}}{\displaystyle \frac{\overline{\lambda }|v_k|^2}{|u_k|^2+\overline{\lambda }|v_k|^2}},`$ (41)
and
$`{\displaystyle \frac{f}{N}}=`$ $`{\displaystyle \frac{1}{2}}\mathrm{log}\lambda \mathrm{log}(L{\displaystyle \frac{N}{2}})+\mathrm{log}(LN)`$
$``$ $`0,`$ (42)
i.e.,
$`\overline{\lambda }=`$ $`\left({\displaystyle \frac{L\overline{N}}{L\frac{\overline{N}}{2}}}\right)^2`$
$`=`$ $`g_t^2.`$ (43)
Eq. (41) and Eq. (43) lead to Eq. (22),
$`n_{\mathrm{after}}=\overline{n}={\displaystyle \frac{\overline{N}}{L}}=2{\displaystyle \frac{1}{L}}{\displaystyle \underset{k}{}}{\displaystyle \frac{g_t^2|v_k|^2}{|u_k|^2+g_t^2|v_k|^2}},`$
for the average particle density of a projected BCS wave function. Without the factor $`g_N`$ this calculation would give the well known result for an unprojected BCS wave function,
$`n_{\mathrm{before}}=\overline{n}`$ $`=2{\displaystyle \frac{1}{L}}{\displaystyle \underset{k}{}}|v_k|^2.`$
To calculate the particle number fluctuations of the projected wave function, we need to expand $`f(\lambda ,N)`$ up to second order in $`N`$ and $`\lambda `$ around the saddle point. Then, integration over $`\lambda `$ in Eq. (37) approximates the particle number distribution $`\rho `$ by a gaussian distribution, yielding an expression for number fluctuations. With
$`f_{\lambda \lambda }(\lambda ,N)`$ $`{\displaystyle \frac{^2f}{\lambda ^2}}`$
$`=`$ $`{\displaystyle \underset{k}{}}{\displaystyle \frac{|v_k|^4}{(|u_k|^2+|v_k|^2\lambda )^2}}+{\displaystyle \frac{\frac{N}{2}+1}{\lambda ^2}},`$
$`f_{\lambda N}(\lambda ,N)`$ $`{\displaystyle \frac{^2f}{\lambda N}}={\displaystyle \frac{1}{2\lambda }},`$
$`f_{NN}(\lambda ,N)`$ $`{\displaystyle \frac{^2f}{N^2}}={\displaystyle \frac{1}{2}}{\displaystyle \frac{1}{L\frac{N}{2}}}{\displaystyle \frac{1}{LN}},`$
the second order expansion can be written as
$`f(\lambda ,N)f(\overline{\lambda },\overline{N})`$ $`f_{\lambda \lambda }(\overline{\lambda },\overline{N}){\displaystyle \frac{(\lambda \overline{\lambda })^2}{2}}`$
$`+`$ $`f_{\lambda N}(\overline{\lambda },\overline{N})(\lambda \overline{\lambda })(N\overline{N})`$
$`+`$ $`f_{NN}(\overline{\lambda },\overline{N}){\displaystyle \frac{(N\overline{N})^2}{2}}.`$ (44)
For this level of the saddle point approximation for $`\lambda `$, the contour around $`\overline{\lambda }`$ for the integral in Eq. (37) must be taken so that $`f_{\lambda \lambda }(\overline{\lambda },\overline{N})(\lambda \overline{\lambda })^2<0`$. Since $`f_{\lambda \lambda }(\overline{\lambda },\overline{N})>0`$ and contribution only near the saddle point is relevant, the path is taken from $`\overline{\lambda }i\mathrm{}`$ to $`\overline{\lambda }+i\mathrm{}`$. By variable transformation $`\lambda =\overline{\lambda }+i\lambda ^{}`$, one can perform a gaussian integral of $`\lambda ^{}`$. Then, we obtain a a gaussian distribution for $`\rho _N`$,
$`\rho _N{\displaystyle \frac{1}{2\pi }}{\displaystyle _{\mathrm{}}^{\mathrm{}}}๐\lambda ^{}e^{f(\overline{\lambda }+i\lambda ^{},N)}`$
$`\mathrm{exp}\left[\left(f_{\overline{N}\overline{N}}{\displaystyle \frac{f_{\overline{\lambda }\overline{N}}^2}{f_{\overline{\lambda }\overline{\lambda }}}}\right){\displaystyle \frac{\left(N\overline{N}\right)^2}{2}}\right]e^{f(\overline{\lambda },\overline{N})}.`$ (45)
The variance of $`\overline{N}`$ (average particle number of a projected BCS wave function) can now be read off from Eq. (45). We get,
$`{\displaystyle \frac{\sigma _{\mathrm{after}}^2}{L}}`$ $`={\displaystyle \frac{1}{L}}\left(f_{\overline{N}\overline{N}}{\displaystyle \frac{f_{\overline{\lambda }\overline{N}}^2}{f_{\overline{\lambda }\overline{\lambda }}}}\right)^1`$
$`=2\left[{\displaystyle \frac{1}{(1\frac{\overline{n}}{2})(1\overline{n})}}+{\displaystyle \frac{1}{\frac{2}{L}_k\frac{\overline{\lambda }|u_k|^2|v_k|^2}{(|u_k|^2+|v_k|^2\overline{\lambda })^2}}}\right]^1,`$ (46)
where we used Eq. (41) in the second term. For $`\overline{\lambda }`$ we must insert $`g_t^2`$. For completeness, we mention that for the unprojected wave function, i.e. $`g_N`$ not included, this approach yields the known result
$`{\displaystyle \frac{\sigma _{\mathrm{before}}^2}{L}}=4{\displaystyle \frac{1}{L}}{\displaystyle \underset{k}{}}|u_k|^2|v_k|^2.`$ (47)
The fluctuations $`\frac{\sigma ^2}{L}`$ are illustrated in Fig. 5 as a function of the particle density after projection $`n_{\mathrm{after}}(=\overline{n})`$ for unprojected ($`|\mathrm{\Psi }_{\mathrm{BCS}}`$) and projected ($`P|\mathrm{\Psi }_{\mathrm{BCS}}`$) BCS $`d`$-wave functions. As expected, the fluctuations vanish at half filling, since projection freezes the charge degrees of freedom entirely.
## Appendix B Gutzwiller approximation for the $`tJ`$ Hamiltonian
We summarize the Gutzwiller approximation for the so-called 3-site terms in the $`tJ`$ model, that are included in the VMC study of Yokoyama and Shiba yokoyama\_88 .
The $`tJ`$ model can be derived from a large $`U`$ expansion of the Hubbard model. The Hamiltonian is valid in the reduced Hilbert space of no double occupied states, and is given by
$$H_{\mathrm{eff}}=T+H_{\mathrm{eff}}^{(2)}$$
(48)
where
$$T=t\underset{i,j,\sigma }{}\left(c_{i,\sigma }^{}c_{j,\sigma }+c_{j,\sigma }^{}c_{i,\sigma }\right)$$
(49)
and
$`H_{\mathrm{eff}}^{(2)}=`$ $`J{\displaystyle \underset{i,j}{}}๐_i๐_j{\displaystyle \frac{J}{4}}{\displaystyle \underset{i,j}{}}n_in_j`$
$``$ $`{\displaystyle \frac{J}{4}}{\displaystyle \underset{i,\tau \tau ^{},\sigma }{}}c_{i+\tau ,\sigma }^{}c_{i,\sigma }^{}c_{i,\sigma }c_{i+\tau ^{},\sigma }`$
$`+`$ $`{\displaystyle \frac{J}{4}}{\displaystyle \underset{i,\tau \tau ^{},\sigma }{}}c_{i+\tau ,\sigma }^{}c_{i,\sigma }^{}c_{i,\sigma }c_{i+\tau ^{},\sigma }.`$ (50)
Here, $`J=4\frac{t^2}{U}`$, $`๐_i`$ are the spin operators on site $`i`$, and $`n_i=n_{i,}+n_{i,_{}}`$ with $`n_{i,\sigma }=c_{i,\sigma }^{}c_{i,\sigma }`$. $`i,j`$ are pairs of n.n sites and $`i+\tau `$ denotes a n.n. site of $`i`$.
We are interested in the energies $`E^{(1)}`$ and $`E^{(2)}`$ calculated in Ref. yokoyama\_88, :
$`E^{(1)}`$ $`={\displaystyle \frac{1}{L}}{\displaystyle \frac{\mathrm{\Psi }_{\mathrm{BCS}}|PTP|\mathrm{\Psi }_{\mathrm{BCS}}}{\mathrm{\Psi }_{\mathrm{BCS}}|PP|\mathrm{\Psi }_{\mathrm{BCS}}}},`$
$`E^{(2)}`$ $`={\displaystyle \frac{1}{L}}{\displaystyle \frac{\mathrm{\Psi }_{\mathrm{BCS}}|PH_{\mathrm{eff}}^{(2)}P|\mathrm{\Psi }_{\mathrm{BCS}}}{\mathrm{\Psi }_{\mathrm{BCS}}|PP|\mathrm{\Psi }_{\mathrm{BCS}}}}.`$ (51)
We invoke the Gutzwiller approximation. The renormalization factors, $`g_t`$ for kinetic energy (Eq. (49)) and $`g_S`$ for spin exchange (first term in Eq. (50)), are given in Eq. (3). The second term of Eq. (50), $`n_in_j`$, is not renormalized. The approximation for the 3-site terms (3rd and 4th term of Eq. (50)) is done as follows ($`|\psi =P|\psi _0`$):
$`{\displaystyle \frac{\mathrm{\Psi }|c_{i+\tau ,}^{}c_{i,}^{}c_{i,}c_{i+\tau ^{},}|\mathrm{\Psi }}{\mathrm{\Psi }|\mathrm{\Psi }}}`$
$`={\displaystyle \frac{\mathrm{\Psi }|c_{i+\tau ,}^{}n_{i,}(1n_{i,})c_{i+\tau ^{},}|\mathrm{\Psi }}{\mathrm{\Psi }|\mathrm{\Psi }}}`$
$`=g_3{\displaystyle \frac{\mathrm{\Psi }_0|c_{i+\tau ,}^{}n_{i,}(1n_{i,})c_{i+\tau ^{},}|\mathrm{\Psi }_0}{\mathrm{\Psi }_0|\mathrm{\Psi }_0}},`$ (52)
$`{\displaystyle \frac{\mathrm{\Psi }|c_{i+\tau ,}^{}c_{i,}^{}c_{i,}c_{i+\tau ^{},}|\mathrm{\Psi }}{\mathrm{\Psi }|\mathrm{\Psi }}}`$
$`=g_3{\displaystyle \frac{\mathrm{\Psi }_0|c_{i+\tau ,}^{}c_{i,}^{}c_{i,}c_{i+\tau ^{},}|\mathrm{\Psi }_0}{\mathrm{\Psi }_0|\mathrm{\Psi }_0}}`$ (53)
The renormalization factor $`g_3`$ is derived by considering the number of terms that contribute to the projected and the unprojected side respectively. The projected side (lhs) contributes only if (i) site $`i+\tau `$ is unoccupied, i.e., probability $`(1n)`$, (ii) site $`i`$ is singly occupied by a $``$-electron, i.e., probability $`n_{}`$, and (iii) site $`i+\tau ^{}`$ is singly occupied by an $``$-electron, i.e., probability $`n_{}`$. On the other hand, the unprojected side (rhs) in Eq. (52)/Eq. (53) contributes only if (i) site $`i+\tau `$ is not occupied by an $``$-electron/$``$-electron, i.e., probability $`(1n_{})`$ / $`(1n_{})`$ , (ii) site $`i`$ is singly occupied by a $``$-electron, i.e., probability $`n_{}(1n_{})`$, and (iii) site $`i+\tau ^{}`$ must have an $``$-electron, i.e., probability $`n_{}`$. These probabilities yield the Gutzwiller factor (ratio of contributions from projected and unprojected states),
$`g_3={\displaystyle \frac{(1n)n_\sigma n_\sigma }{(1n_\sigma )n_\sigma (1n_\sigma )n_\sigma }}={\displaystyle \frac{1n}{(1n_\sigma )^2}},`$ (54)
where we assumed $`n_{}=n_{}=n_\sigma `$.
We can now write down the renormalized $`tJ`$ Hamiltonian $`H_{\mathrm{eff}}^{}`$,
$$H_{\mathrm{eff}}^{}=T^{}+H_{}^{}{}_{\mathrm{eff}}{}^{(2)},$$
(55)
where,
$$T^{}=g_tt\underset{i,j,\sigma }{}\left(c_{i,\sigma }^{}c_{j,\sigma }+c_{j,\sigma }^{}c_{i,\sigma }\right),$$
(56)
$`H_{}^{}{}_{\mathrm{eff}}{}^{(2)}=`$ $`g_SJ{\displaystyle \underset{i,j}{}}๐_i๐_j{\displaystyle \frac{J}{4}}{\displaystyle \underset{i,j}{}}n_in_j`$
$``$ $`g_3{\displaystyle \frac{J}{4}}{\displaystyle \underset{i,\tau \tau ^{},\sigma }{}}c_{i+\tau ,\sigma }^{}n_{i,\sigma }(1n_{i,\sigma })c_{i+\tau ^{},\sigma }`$
$`+`$ $`g_3{\displaystyle \frac{J}{4}}{\displaystyle \underset{i,\tau \tau ^{},\sigma }{}}c_{i+\tau ,\sigma }^{}c_{i,\sigma }^{}c_{i,\sigma }c_{i+\tau ^{},\sigma }.`$ (57)
By using $`T^{}`$ and $`H_{}^{}{}_{\mathrm{eff}}{}^{(2)}`$, Eq. (51) ($`E^{(1)}`$ and $`E^{(2)}`$) can be calculated using unprojected wave functions,
$`E^{(1)}`$ $`={\displaystyle \frac{1}{L}}{\displaystyle \frac{\mathrm{\Psi }_{\mathrm{BCS}}|T^{}|\mathrm{\Psi }_{\mathrm{BCS}}}{\mathrm{\Psi }_{\mathrm{BCS}}|\mathrm{\Psi }_{\mathrm{BCS}}}},`$
$`E^{(2)}`$ $`={\displaystyle \frac{1}{L}}{\displaystyle \frac{\mathrm{\Psi }_{\mathrm{BCS}}|H_{}^{}{}_{\mathrm{eff}}{}^{(2)}|\mathrm{\Psi }_{\mathrm{BCS}}}{\mathrm{\Psi }_{\mathrm{BCS}}|\mathrm{\Psi }_{\mathrm{BCS}}}}.`$ (58)
Evaluating Eq. (58) by Wickโs decomposition for a $`d`$-wave BCS state,
$$\frac{E^{(1)}}{t}=2g_t(\xi _x+\xi _y),$$
$`{\displaystyle \frac{E^{(2)}}{0.25J}}=`$ $`(3g_s1)(\xi _x^2+\xi _y^2)/2`$ (59)
$`(3g_s+1)(|\stackrel{~}{\mathrm{\Delta }}_x|^2+|\stackrel{~}{\mathrm{\Delta }}_y|^2)/2\mathrm{\hspace{0.17em}2}n^2`$
$`g_3n(1n_\sigma )(\xi _{2x}+\xi _{2y}+2\xi _{xy}+2\xi _{x+y})`$
$`g_3(1+n_\sigma )(\xi _x^2+\xi _y^2+4\xi _x\xi _y)`$
$`g_3(2n_\sigma )(|\stackrel{~}{\mathrm{\Delta }}_x|^2+|\stackrel{~}{\mathrm{\Delta }}_y|^2+4\mathrm{}(\stackrel{~}{\mathrm{\Delta }}_x\stackrel{~}{\mathrm{\Delta }}_y^{})),`$
where we defined ($`\tau ,\tau ^{}=x,y`$)
$`\xi _\tau `$ $`={\displaystyle \underset{\sigma }{}}c_{i,\sigma }^{}c_{i+\tau ,\sigma }={\displaystyle \frac{1}{L}}{\displaystyle \underset{k}{}}2\mathrm{cos}(k_\tau )|v_k|^2,`$
$`\xi _{\tau \pm \tau ^{}}`$ $`={\displaystyle \underset{\sigma }{}}c_{i+\tau ,\sigma }^{}c_{i\pm \tau ^{},\sigma }={\displaystyle \frac{1}{L}}{\displaystyle \underset{k}{}}2\mathrm{cos}(k_\tau \pm k_\tau ^{})|v_k|^2,`$
$`\stackrel{~}{\mathrm{\Delta }}_\tau `$ $`=c_{i,}^{}c_{i+\tau ,}^{}c_{i,}^{}c_{i+\tau ,}^{}={\displaystyle \frac{1}{L}}{\displaystyle \underset{k}{}}2\mathrm{cos}(k_\tau )v_ku_k^{}.`$
The last three rows in Eq. (59) correspond to the 3-site terms of the $`tJ`$ model and are renormalized by the Gutzwiller factor $`g_3`$. Here, itโs important to note that the order parameter $`\stackrel{~}{\mathrm{\Delta }}_\tau `$ is related, but not identical, to the previous introduced variational parameter $`\mathrm{\Delta }_0`$.
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# Doubly Special Relativity as a Limit of Gravity
## 1 Introduction
The anniversary of a great idea is usually a good occasion for critical reassessment of its meaning, influence, and future. In theoretical physics, where methodology, instead of hermeneutics, is based on Popperian conjectures and refutations scheme this last issue โ the future โ is, of course, the most important. Thus, in the course of the celebrations of the 100 anniversary of the Theory of Relativity, we are mostly interested in asking the questions: Is Special Relativity still to be regarded as the correct theory describing relativistic phenomena (particles and fields kinematics and dynamics) in flat space-time? Will it survive the next 100 years, and if not, which theory is going to replace it?
One quite often hears the opinion that there is, in fact, no such theory as Special Relativity. What we have to do with is just a very particular, flat space-time limit of General Relativity. And given the fact that few of us doubt that the ultimate theory of gravity should be Quantum Gravity (in the form of Loop Quantum Gravity, or String Theory, or perhaps โ and most likely โ in the disguise we do not really know yet) the question to be posed is: what is the flat space, semiclassical limit of Quantum Gravity?
For a long time it was taken as obvious that such a limit should be just the ordinary Special Relativity. Recent developments, however, put some doubts on this naive conclusion. First, assuming that one or another form of (super)strings theory is indeed the correct theory of Quantum Gravity one may contemplate the idea that in the (super)string vacuum, corresponding to our universe, Lorentz invariance is spontaneously broken. This would certainly lead to some, possibly observable, modifications of Special Relativity . Violation of Lorentz symmetry is also possible in models based on Loop Quantum Gravity . The issue of Lorentz Invariance Violation and its possible observational consequences has been recently reviewed in and .
The core of all the proposals of spontaneous Lorentz symmetry breaking is introduction of some sort of aether. While phenomenological models of this form are certainly interesting, making possible to device precise experimental tests, they are much less appealing theoretically, since they are based on rejection of the most cherished principle of physics, the relativity principle<sup>1</sup><sup>1</sup>1One should note that the presence of aether does not necessarily mean breaking of Lorentz symmetry. It is well possible that the relations between inertial observers still form Lorentz (or Poincarรฉ) group, however, group elements would, in the presence of the aether, depend both on relative velocity of observers and velocity with respect to the aether. Physically, the breaking of Lorentz symmetry means that some physical processes (particle reactions, for example) are possible for some range of velocities relative to the aether and impossible for other velocities. In DSR proposal, which we will describe below, the central postulate is that relativity principle still holds, and thus if some process is observed by one observer it is observed by all other. There is also some specific relation between descriptions of the same process by two observers, which depends only on their relative velocity.. In our view there is no a priori reason, neither theoretical, nor experimental, to contemplate violation of relativity principle. This does not mean that modifications of Special Relativity are not possible, and we will argue below that there are good reasons to believe that if one regards Special Relativity (understood as a kinematical theory of particles and fields on flat semiclassical space-time) as a limit of (quantum) gravity, it is inevitable that some, Planck scale corrections, must be present. In a sense, the presence of this scale reflects the โmemoryโ of flat space-time about its (quantum) gravitational origin.
## 2 Postulates of Doubly Special Relativity
Let us start with the plausible assumption that the flat space-time kinematics of point particles originates in some theory of quantum gravity. Assume further that in taking the flat space-time, semiclassical limit there is no breaking of spacetime symmetries, and in particular that the relativity principle holds (in other words, the diffeomorphism invariance of the theory is not being broken in the process.) After taking this limit the space-time will be, at least locally, the standard Minkowski space-time, however it may well be that the resulting theory still possesses some information about its origin, in the form of the observer independent mass scale, $`\kappa `$, of order of Planck scale. This scale will be still present as a parameter in the transformation laws, relating different inertial observers.
If such a scale indeed is present, it is natural to expect that deviations from the standard Special Relativistic kinematics arise in the processes characterized by the energy scale close to $`\kappa `$, and that these deviations should be rather generic . For example, it may happen that the standard dispersion relation for particles acquires additional terms, to wit
$$E^2=p^2+m^2+\alpha \frac{E^3}{\kappa }+\mathrm{}$$
(1)
with $`\alpha `$ being a dimensionless parameter of order 1.
Naively, modified dispersion relation would immediately imply the existence of preferred frame, since if it holds in one frame it does not hold in any other, related to the original one by the standard Lorentz transformation. However we can demand that along with deforming dispersion relation we deform Lorentz transformations, so that (1) (and their generalizations to be introduced shortly) remain invariant under action of six parameter group of transformations, with generators satisfying Lorentz $`SO(3,1)`$ algebra.
Based on this intuition, Giovanni Amelino-Camelia , formulated a set of postulates that should be satisfied by new theory, replacing Special Relativity, in the regime of ultra-high energies, and he dubbed this theory Doubly Special Relativity.
Doubly Special Relativity is based on the following postulates:
1. The Relativity Principle holds. This means that if two observers describe the same phenomenon, possible differences in their descriptions can only depend on their relative motion (in the case of inertial observers โ their relative velocity). In particular there is no notion of absolute rest, and absolute motion. It follows also that if some process is observed by one observers (for example particle (1) collides with particle (2) producing particle (3)), all observers agree that this process takes place.
2. There exist two observer-independent scales: the velocity scale c, identified with velocity of light, and, the mass scale $`\kappa `$, identified with Planck mass ($`\kappa 10^{19}`$ GeV).
It is worth recalling at this point what is the difference between observer independent scales and other dimensionful quantities one encounters in physical theories. To be more specific, let us consider the question what is the difference, in the framework of the standard Special Relativity, between a coupling constant, like electric charge $`e`$ or Planckโs constant $`\mathrm{}`$, and the velocity of light $`c`$. In the first case, of dimensionful coupling constants, all observers measure their values with the help of identical, low energy experiments, performed in their rest frame. Then the relativity principle<sup>2</sup><sup>2</sup>2In the โpassiveโ form: Identical experiments performed by inertial observers give the same results. guarantees that all the observers will obtain the same numerical values of the constants. These values can be then used in other experiments. In the case of the speed of light the situation is different however. Since this speed is an observer independent scale, we demand also that if two inertial observers measure velocity of the same photon, they will obtain the same result. This is clearly incompatible with Galilean Relativity, but, as we know, has been successfully built into the principles of his Special Relativity. Analogously, in the construction of Doubly Special Relativity we postulate the presence of yet another observer-independent scale, this time of dimension of mass, whose numerical value is presumed to be of order of Planck mass<sup>3</sup><sup>3</sup>3Note, however, that contrary to Special Relativity, in which we know physical objects moving with the velocity of light (massless particles), in the case of DSR we do not know (yet?) physical objects exemplifying the scale $`\kappa `$..
It has soon been realized , (for recent review see ) that these postulates are satisfied in the framework of theories in which Poincarรฉ algebra is replaced by deformed $`\kappa `$-Poincarรฉ algebra , , , <sup>4</sup><sup>4</sup>4 It should be stressed however that DSR is not just the $`\kappa `$-Poincarรฉ algebra โ not only in the sense analogous to the well known fact Special Relativity is not just the Poincarรฉ algebra โ there might be DSR proposals in which this algebra does not play any role. .
One possible realization is provided by the, so-called, bicrossproduct (or MajidโRuegg)basis , in which the commutators between rotation $`M_i`$, boost $`N_i`$, and momentum $`P_\mu =(P_0,P_i)`$ generators are the following
$$[M_i,M_j]=ฯต_{ijk}M_k,[M_i,N_j]=ฯต_{ijk}N_k,$$
$$[N_i,N_j]=ฯต_{ijk}M_k.$$
(2)
and
$$[M_i,P_j]=ฯต_{ijk}P_k,[M_i,P_0]=0$$
$$[N_i,P_j]=\delta _{ij}\left(\frac{1}{2}\left(1e^{2P_0/\kappa }\right)+\frac{๐^2}{2\kappa }\right)\frac{1}{\kappa }P_iP_j,$$
$$[N_i,P_0]=P_i.$$
As one can easily check, the Casimir of the $`\kappa `$-Poincarรฉ algebra (2), (2) reads
$$๐=\kappa ^2\mathrm{cosh}\frac{P_0}{\kappa }\frac{\stackrel{}{P}^2}{2}e^{P_0/\kappa }M^2.$$
(3)
It is easy to check also that the expansion of (3) to the leading order in $`1/\kappa `$ yields (1).
One should note at this point that the bicrossproduct algebra above is not the only possible realization of DSR. For example, in , Magueijo and Smolin proposed and carefully analyzed another DSR proposal, called sometimes DSR2. In DSR2 the Lorentz algebra is still not deformed and there are no deformations in the brackets of rotations and momenta. The boostsโ momenta generators have now the form
$$[N_i,p_j]=i\left(\delta _{ij}p_0\frac{1}{\kappa }p_ip_j\right),$$
(4)
and
$$[N_i,p_0]=i\left(1\frac{p_0}{\kappa }\right)p_i.$$
(5)
It is easy to check that the Casimir for this algebra has the form
$$M^2=\frac{p_0^2\stackrel{}{p}^2}{\left(1\frac{p_0}{\kappa }\right)^2}.$$
(6)
In order to describe kinematics of a particle we must extend the above algebra to the algebra of phase space of the particle. We will not derive here all the results, and the reader could find the derivation with references to the original literature in the review paper . Instead we will just state the main results which will be important below, when we compare DSR with a theory resulting as a flat limit of gravity coupled to point particles.
* Both the bicrossproduct (2)โ(3) and MagueijoโSmolin algebras (4)โ(6) can be understood as examples of larger class of algebras constructed as follows. In the standard Special Relativity four-momentum can be thought of as a point of four dimensional flat manifold of Lorentz Signature โ the flat momentum manifold. Notice that in this case positions, being โtranslations of momentaโ are points of another flat Minkowski space. In transition to DSR, assume instead that the space of momenta is a maximally symmetric manifold of (constant) curvature, $`1/\kappa ^2`$. In the limit when $`\kappa `$ goes to infinity, the curvature goes to zero and we return to Special Relativity, as we should. In given coordinates on this constant curvature momentum space, each point will correspond to some four-momentum. As it is well known, the group of symmetries in this case is the 10-parameter (in 4 dimensions) de Sitter group $`\mathrm{๐ฒ๐ฎ}(4,1)`$. This group possesses a six-parameter subgroup, $`\mathrm{๐ฒ๐ฎ}(3,1)`$, isomorphic with Lorentz group, and one can easily compute what will be infinitesimal action of the group elements on points of the manifold. It turns out that, for example, the bicrossproduct basis corresponds to the standard system of coordinates, used in cosmology. More details can be found in and .
* What about the remaining four parameters of de Sitter group $`\mathrm{๐ฒ๐ฎ}(4,1)`$? It is well known from differential geometry, that while the generators of $`\mathrm{๐ฒ๐ฎ}(3,1)`$ act as โrotationsโ the remaining ones play the role of โtranslationsโ. This means that it is natural to identify them with positions. It turns out to be convenient to arrange the remaining four generators $`x^\mu `$ so as to form the Iwasawa decomposition of the $`\mathrm{๐๐}(4,1)`$ algebra; explicitly their commutators could be brought to the following form
$$[x^i,x^j]=0,[x^0,x^i]=\frac{1}{\kappa }x^i$$
(7)
The noncommutative space-time satisfying (7) is called $`\kappa `$-Minkowski space-time.
* It should be noted that there is a natural Hopf algebra structure associated with an algebra of symmetries of de Sitter space. This algebra turns out to be exactly the quantum $`\kappa `$-Poincarรฉ algebra of , , , . For more details see .
* It should be also noted that one can in principle construct analogous structure starting from the momentum space of the particle being anti-de Sitter space . Explicit models in four dimensions are not known in this case, however they play a role in 2+1 gravity coupled to a particle, as we will discuss in details below.
Given characterization of properties of single particle DSR models let us now turn to the question, what is DSR coming from.
## 3 Constrained BF action for gravity
It is usually considered rather obvious that Special Relativity, regarded as a theory of particle kinematics should emerge somehow from General Relativity coupled to point particles in a limit, in which gravitational interactions are โswitched offโ. It turns out that it is surprisingly difficult to prove this claim in the framework of the standard Einstein formulation of GR. First of all โswitching offโ gravity would presumably mean going to zero with gravitational constant, but this limit is known to be pathological in GR. Moreover it is well known that coupling of GR to point particle is at least problematic, if not impossible.
Our starting point must be therefore some another (but equivalent) form of the gravity action. The convenient form has been derived recently by Freidel and Starodubtsev . The starting point of that paper is the observation theory of gravity can be defined by the action containing two parts: the โvacuumโ one, being a topological field theory with an appropriate gauge group, and the constraints, leading to the emergence of the dynamical degrees of freedom of gravity. Both parts are manifestly diffeomorphism invariant, which opens new perspectives in construction of diffeomorphism-invariant perturbation theory for quantum gravity. From our perspective, however, the most important aspect of this theory would be that it makes it possible both to define a limit, in which local degrees of freedom of gravity are switched off and the point particle coupling. One can say that in this formulation theory of gravity has the well defined โDSR limit.โ
The construction of the FreidelโStarodubtsev theory borrows from earlier works , , and is based on the $`\mathrm{๐ฒ๐ฎ}(4,1)`$ gauge theory. The basic dynamical variables are<sup>5</sup><sup>5</sup>5Below we use BOLD typeface to denote forms (space-time indices suppressed) and SANS SERIF typeface to denote Lie algebra valued fields (group indices suppressed). $`\mathrm{๐๐}(4,1)`$ connection one form $`๐^{IJ}`$, and the $`\mathrm{๐๐}(4,1)`$-valued two-form $`๐^{IJ}`$. The starting point is the action principle
$$S=๐^{IJ}๐
_{IJ}$$
(8)
where $`๐
_{IJ}`$ is the curvature of connection $`๐^{IJ}`$. The equations of motion following from this action
$`๐
_{IJ}=0`$
$`d_๐๐_{IJ}=0`$ (9)
where $`d_๐`$ is the covariant derivative of connection $`๐`$, tell that the connection is flat, while the $`๐^{IJ}`$ field is covariantly constant. The solutions of these equations on locally connected region $`๐ฐ`$ is of the form
$`๐=๐^1d๐๐=๐^1d๐ฟ๐,๐\mathrm{๐ฒ๐ฎ}(4,1),๐ฟ\mathrm{๐๐}(4,1)`$ (10)
The theory is therefore almost trivial, without any local degrees of freedom.
In order to get General Relativity we must break local symmetry of the theory from $`\mathrm{๐ฒ๐ฎ}(4,1)`$ down to the Lorentz group $`\mathrm{๐ฒ๐ฎ}(3,1)`$. To this end we denote by 5 the preferred direction in the algebra space, and add to the action the term which explicitly breaks the $`\mathrm{๐ฒ๐ฎ}(4,1)`$ gauge symmetry, to wit
$$S=๐^{IJ}๐
_{IJ}\frac{\alpha }{2}๐^{IJ}๐^{KL}ฯต_{IJKL5}$$
(11)
Let us now decompose the algebra index $`I=(i,5)`$, $`i=0,\mathrm{},3`$ with $`ฯต^{ijkl}=ฯต^{IJKL5}`$ being an invariant $`\mathrm{๐ฒ๐ฎ}(3,1)`$ tensor. Note that now the first equation in (3) is replaced by
$$๐
_{IJ}=\alpha ๐^{KL}ฯต_{IJKL5}$$
(12)
and is manifestly not $`\mathrm{๐ฒ๐ฎ}(4,1)`$ covariant.
The $`๐`$-field enters the action $`S`$ only algebraically, so we can substitute the solution (12) back to the action (11) to get
$$S=\frac{1}{4\alpha }๐
^{ij}๐
^{kl}ฯต_{ijkl}.$$
(13)
It is convenient at this point to decompose the curvature as follows
$`๐
^{ij}(๐)`$ $`=`$ $`๐^{ij}(๐){\displaystyle \frac{1}{l^2}}๐^i๐^j`$
$`๐
^{i5}(๐)`$ $`=`$ $`{\displaystyle \frac{1}{l}}d_๐๐^i`$ (14)
where $`๐^{ij}=๐^{ij}`$ is the 4-dimensional connection one-form and $`๐^i=e_\mu {}_{}{}^{i}dx^\mu `$ is a frame field which corresponding to the metric $`g_{\mu \nu }=e_\mu ^ie_{i\nu }`$. $`๐^{ij}`$ is the $`\mathrm{๐๐}(3,1)`$ curvature of connection $`๐`$, $`๐^{ij}(๐)=d๐^{ij}+๐_k^i๐^{kj}`$. Notice that for dimensional reasons we had to introduce the scale $`l`$, of dimension of length. Using the equations for the curvature (3), we can rewrite the action in terms of $`\mathrm{๐๐}(3,1)`$ curvature:
$`S`$ $`=`$ $`{\displaystyle \frac{1}{4\alpha }}{\displaystyle (๐^{ij}(๐)\frac{1}{l^2}๐^i๐^i)}(๐^{kl}(๐){\displaystyle \frac{1}{l^2}}๐^k๐^l)ฯต_{ijkl}`$ (15)
$`=`$ $`{\displaystyle \frac{1}{2G}}{\displaystyle (๐^{ij}(๐)๐^k๐^l\frac{\mathrm{\Lambda }}{6}๐^i๐^k๐^l)ฯต_{ijkl}}`$
$`+`$ $`{\displaystyle \frac{1}{4\alpha }}{\displaystyle ๐^{ij}(๐)}๐^{kl}(๐)ฯต_{ijkl}`$
What we get is nothing but the Palatini action of General Relativity with cosmological constant plus additional term whose variation vanishes identically due to Bianchi identity. Note that the Newtonโs constant $`G`$ equals $`\alpha l^2`$, while the cosmological constant $`\mathrm{\Lambda }=3/l^2`$. Thus the coupling constant $`\alpha =G\mathrm{\Lambda }/3`$ is dimensionless and extremely small, which makes it a perfect candidate for a parameter of (both classical and quantum) perturbative expansion. As stressed by Freidel and Starodubtsev , the constrained BF theory is therefore very promising as a starting point for construction of perturbative quantum gravity, where diffeomorphism invariance is manifestly preserved at all steps of perturbative expansion.
To the initial, topological action (8) we can still add the $`\mathrm{๐ฒ๐ฎ}(4,1)`$ cosmological term of the form
$$\frac{\beta }{2}๐^{IJ}๐_{IJ}.$$
(16)
This addition changes the equations of motion:
$`๐
_{IJ}\beta ๐_{IJ}=0`$
$`d_๐๐_{IJ}=0`$ (17)
(the second equation follows in fact from the first and Bianchi identity.)
It can be shown that the action (16) is still topological, i.e., without local degrees of freedom. As before we can add to this action the $`\alpha `$ constraint, in order to obtain the action of General Relativity with the additional term $`\frac{2}{\gamma }๐^{ij}(๐)๐_iwedge๐_j`$ and more more topological terms. The โbare actionโ parameters $`l,\alpha ,\beta `$ are related to the physical ones $`G`$, $`\mathrm{\Lambda }`$, and $`\gamma `$ (Immirzi parameter) as follows $`\mathrm{\Lambda }=3/l^2`$, $`\gamma =\beta /\alpha `$, $`G=\frac{\alpha ^2\beta ^2}{\alpha }l`$ (for more details and discussion of possible physical relevance of $`\gamma `$ parameter see , and also recent paper .)
The convenient basis of $`\mathrm{๐๐}(4,1)`$ algebra is provided by Dirac matrices $`\gamma ^{ij}=\frac{1}{2}[\gamma ^i,\gamma ^j]`$ and $`\gamma ^i\gamma ^5`$. Using the $`\mathrm{๐๐}(4,1)`$ algebra valued fields $`๐ _\mu `$, $`๐ก_{\mu \nu }`$ the constrained BF action for gravity can be rewritten in the following form
$`S`$ $`=`$ $`{\displaystyle d^4xฯต^{\mu \nu \rho \sigma }Tr(๐ก_{\mu \nu }๐ฅ_{\rho \sigma }(๐ ))}`$ (18)
$``$ $`{\displaystyle \frac{\beta }{2}}{\displaystyle d^4xฯต^{\mu \nu \rho \sigma }Tr(๐ก_{\mu \nu }๐ก_{\rho \sigma })}`$
$``$ $`{\displaystyle \frac{\alpha }{2}}{\displaystyle d^4xฯต^{\mu \nu \rho \sigma }Tr(๐ก_{\mu \nu }๐ก_{\rho \sigma }\gamma ^\mathrm{๐ง})}`$
It is quite easy to couple point particles to the constrained BF action. Indeed since $`๐^{IJ}`$ is a one form, it couples naturally to one-dimensional objects โ the particles world-lines.
The general procedure of coupling particles carrying non-abelian charges to Yang-Mills potential has been developed by Balachandran, Marmo, Skagerstam, and Stern (see and references therein.) In the case at hands<sup>6</sup><sup>6</sup>6The results presented below have been obtained in collaboration with L. Freidel and A. Starodubtsev. the gauge group is $`\mathrm{๐ฒ๐ฎ}(4,1)`$. This group acts by conjugation on its algebra, and the orbits can be labelled by two numbers, corresponding to values of two Casimirs, representing mass and spin of the particle, as follows
$$๐ช=\frac{1}{2}m\gamma _1\gamma ^5+\frac{1}{4}s\gamma _2\gamma _3$$
(19)
As the second ingredient we take connection, gauge-transformed by an arbitrary element $`๐`$ of the Lorentz subgroup $`\mathrm{๐ฒ๐ฎ}(3,1)`$ of $`\mathrm{๐ฒ๐ฎ}(4,1)`$
$$๐ _\mu {}_{}{}^{๐}=๐^1๐ _\mu ๐+๐^1_\mu ๐,๐=\mathrm{exp}\left(\frac{1}{4}\alpha ^{ab}\gamma _{ab}\right)$$
(20)
where, as before $`\mathrm{๐๐}(4,1)`$ connection $`๐ `$ decomposes into $`\mathrm{๐๐}(3,1)`$ connection $`\omega `$ and tetrad $`e`$
$$๐ _\mu =\left(\frac{1}{2}e_\mu {}_{}{}^{a}\gamma _{a}^{}\gamma ^5+\frac{1}{4}\omega _\mu {}_{}{}^{ab}\gamma _{ab}^{}\right)$$
(21)
Then the action of the particle with mass $`m`$ and spin $`s`$ coupled to constrained BF gravity is defined to be
$$L(z,h)=Tr\left(\mathrm{๐ช๐ }_\tau {}_{}{}^{๐}(\tau )\right)S=d\tau L,$$
(22)
where $`๐ _\tau {}_{}{}^{๐}๐ _\mu {}_{}{}^{๐}(z(\tau ))\dot{z}^\mu (\tau )`$ is the value of gauge transformed connection (20) on the particle world-line. We see therefore that the dynamics of the particle is described with the help of the charge $`๐ช`$ it carries, and the Lorentz transformation $`๐`$ relating the particle rest frame and the frame of (asymptotic) observer. It can be shown that variation of the action (22) leads to the correct Mathiasson-Papapetrou equations describing the dynamic of spinning particle in the presence of torsion. When the torsion is zero we recover the usual Mathiasson-Papapetrou equation, which in the case of vanishing spin reduces to the usual geodesic equation.
Having defined the coupling of the particle to gravitational field we can address the question as to what would be the effective behavior of the particle in the limiting case, when the local degrees of freedom of gravitational field are being switched off. To answer this question, one should proceed as follows: take the action being the sum of (18) and (22), solve the resulting equations of motion, and then take the limit $`\alpha 0`$. To see which outcomes of this procedure are possible, note that that although in this limit the gravitational field will become flat in the bulk space-time, there might be some nontrivial leftover at the worldline of the particle. This contribution of the gravitational field may lead to deformation of the (otherwise free) particle action (22), leading to DSR like behavior.
Unfortunately, due to the technical difficulties, the programme described above has not been realized in practice yet. What can be done, however, is to turn to a simpler model of gravity coupled to a particle, in 2+1 dimensions. As we will see in the next section the structure of this model is very similar to the four-dimensional case with the parameters $`\alpha ,\beta `$ equal zero, i.e., in the limit we are mostly interested in. Moreover, the three-dimensional case is not purely of academic interest, as the following argument, borrowed from and , clearly shows.
The main idea is to construct an experimental situation that forces a dimensional reduction from the four dimensional to the $`2+1`$ dimensional theory. It is interesting that this can be done in quantum theory, using the uncertainty principle as an essential element of the argument. Let us consider free elementary particles in $`3+1`$ dimensions, whose mass are less than $`G^1=\kappa `$. The motion of the particles will be linear, at least in some classes of coordinates systems, not accelerating with respect to the natural inertial coordinates at infinity. Let us consider the particle as described by an inertial observer who travels perpendicular to the plane of its motion, which we will call the $`z`$ direction. From the point of view of that observer, the particles are in an eigenstate of longitudinal momentum, $`\widehat{P}_z^{total}`$, with some eigenvalue $`P_z`$. Since the particles are in an eigenstate of $`\widehat{P}_z^{total}`$ their wavefunction will be uniform in $`z`$, with wavelength $`L`$ where (note that we assume here that $`L`$ is so large that we can trust the standard uncertainty relation; besides this uncertainty relation is not being modified in some formulations of DSR)
$$L=\frac{1}{P_z^{total}}$$
(23)
At the same time, we assume that the uncertainties in the transverse positions are bounded a scale $`r`$, such that $`r2L`$. Then the wavefunction for the the particles has support on a narrow cylinder of radius $`r`$ which extend uniformly in the $`z`$ direction. Finally, we assume that the state of the gravitational field is semiclassical, so that to a good approximation, within $`๐`$ the semiclassical Einstein equations hold.
Since the wavefunction is uniform in $`z`$, this implies that the gravitational field seen by our observer will have a spacelike Killing field $`k^a=(/z)^a`$.
Thus, if there are no forces other than the gravitational field, the semiclassical particles must be described by an equivalent $`2+1`$ dimensional problem in which the gravitational field is dimensionally reduced along the $`z`$ direction so that the particles, which are the source of the gravitational field, are replaced by punctures.
The dimensional reduction is governed by a length $`d`$, which is the extent in $`z`$ that the system extends. We cannot take $`d<L`$ without violating the uncertainty principle. It is then convenient to take $`d=L`$. Further, since the system consists of the particles, with no intrinsic extent, there is no other scale associated with their extent in the $`z`$ direction. We can then identify $`z=0`$ and $`z=L`$ to make an equivalent toroidal system, and then dimensionally reduce along $`z`$. The relationship between the four dimensional Newtonโs constant $`G^4`$ and the three dimensional Newtonโs constant $`G^3=G`$ is given by
$$G^3=\frac{G^4}{L}=\frac{G^4P_z^{tot}}{\mathrm{}}$$
(24)
Thus, in the analogous $`3`$ dimensional system, which is equivalent to the original system as seen from the point of view of the boosted observer, the Newtonโs constant depends on the longitudinal momentum.
Of course, in general there will be an additional scalar field, corresponding to the dynamical degrees of freedom of the gravitational field. However, since we are interested only in the four-dimensional limit, in which local degrees of freedom of the gravitational field are not present, all these fields will vanish this limit.
Now we note that, if there are no other particles or excited degrees of freedom, the energy of the system can to a good approximation be described by the hamiltonian $`H`$ of the two dimensional dimensionally reduced system. This is described by a boundary integral, which may be taken over any circle that encloses the particle. But it is well known that in $`3d`$ gravity $`H`$ is bounded from above. This may seem strange, but it is easy to see that it has a natural four dimensional interpretation.
The bound is given by
$$M<\frac{1}{4G^3}=\frac{L}{4G^4}$$
(25)
where $`M`$ is the value of the ADM hamiltonian, $`H`$. But this just implies that
$$L>4G^4M=2R_{Sch}$$
(26)
i.e. this has to be true, otherwise the dynamics of the gravitational field in $`3+1`$ dimensions would have collapsed the system to a black hole! Thus, we see that the total bound from above of the energy in $`2+1`$ dimensions is necessary so that one cannot violate the condition in $`3+1`$ dimensions that a system be larger than its Schwarzschild radius.
Note that we also must have
$$M>P_z^{tot}=\frac{\mathrm{}}{L}$$
(27)
Together with (26) this implies $`L>l_{Planck}`$, which is of course necessary if the semiclassical argument we are giving is to hold.
Now, we have put no restriction on any components of momentum or position in the transverse directions. So the system still has symmetries in the transverse directions. Furthermore, the argument extends to any number of particles, so long as their relative momenta are coplanar. Thus, we learn the following.
Let $`^{QG}`$ be the full Hilbert space of the quantum theory of gravity, coupled to some appropriate matter fields, with $`\mathrm{\Lambda }=0`$. Let us consider a subspace of states $`^{weak}`$ which are relevant in the low energy limit in which all energies are small in Planck units. We expect that this will have a symmetry algebra which is related to the Poincarรฉ algebra $`๐ซ^4`$ in $`4`$ dimensions, by some possible small deformations parameterized by $`G^4`$ and $`\mathrm{}`$. Let us call this low energy symmetry group $`๐ซ_G^4`$.
Let us now consider the subspace of $`^{weak}`$ which is described by the system we have just constructed . It contains the particle, and is an eigenstate of $`\widehat{P}_z^{tot}`$ with large $`P_z^{tot}`$ and vanishing longitudinal momentum. Let us call this subspace of Hilbert space $`_{P_z}`$.
The conditions that define this subspace break the generators of the (possibly modified) Poincarรฉ algebra that involve the $`z`$ direction<sup>7</sup><sup>7</sup>7Notice that if we assume that the four-dimensional rotational symmetry is neither broken, nor deformed, we can recover the whole 4d deformed Poincarรฉ algebra from the 3d one.. But they leave unbroken the symmetry in the $`2+1`$ dimensional transverse space. Thus, a subgroup of $`๐ซ_G^{3+1}`$ acts on this space, which we will call $`๐ซ_G^{2+1}๐ซ_G^{3+1}`$.
We have argued that the physics in $`_{P_z}`$ is to good approximation described by an analogue system of a particle in $`2+1`$ gravity. However, as we will see in the next section the symmetry algebra acting there is not the ordinary $`3`$ dimensional Poincarรฉ algebra, but the $`\kappa `$-Poincarรฉ algebra in $`3`$ dimensions, with
$$\kappa ^1=\frac{4G^4P_z^{tot}}{\mathrm{}}$$
(28)
Now we can note the following. Whatever $`๐ซ_G^4`$ is, it must have the following properties:
* It depends on $`G^4`$ and $`\mathrm{}`$, so that itโs action on each subspace $`_{P_z}`$, for each choice of $`P_z`$, is the $`\kappa `$ deformed $`3d`$ Poincarรฉ algebra, with $`\kappa `$ as above.
* It does not satisfy the rule that momenta and energy add, on all states in $``$, since they are not satisfied in these subspaces.
* Therefore, whatever $`๐ซ_G^4`$ is, it is not the classical Poincarรฉ group.
Let us therefore turn to gravity coupled with point particle in 2+1 dimension.
## 4 DSR from 2+1 dimensional gravity
Even if not for the argument given in the preceding section, the 2+1 dimensional gravity coupled with point particles would be a perfect test ground for investigating properties of DSR theories. As it is well known this theory is topological, i.e., does not posses any local degrees of freedom, moreover its action
$$S=d^3xTr\left(๐๐
(๐)\right)$$
resembles very closely the four dimensional action of the constrained BF theory in the DSR limit
$$S=d^4xTr\left(๐๐
(๐)\right)$$
Investigations in 2+1-dimensional gravity have been pioneered by Staruszkiewicz in 1963 , with interest revived by seminal papers by Deser, Jackiw and โt Hooft and Witten . Here we will follow the approach proposed by Matschull and Welling in .
The action for (2+1) gravity reads:
$$S=\frac{1}{16\pi G}_Md^3xฯต^{\mu \nu \rho }Tr(๐พ_\mu ๐ฅ_{\nu \rho })$$
(29)
and the basic fields are dreibein $`๐พ_\mu `$ and antisymmetric spin connection $`\omega _\mu `$, whereas $`G`$ is the gravitational constant, which in (2+1) gravity has a dimension of inverse mass. The Lorentz group in 2+1 dimension $`\mathrm{๐ฒ๐ฎ}(2,1)`$ is isomorphic to $`\mathrm{๐ฒ๐ซ}(2,R)`$ (which we will in the following denote just $`\mathrm{๐ฒ๐ซ}(2)`$) and thus the field strength $`๐ฅ_{\nu \rho }`$ defined as:
$$๐ฅ_{\mu \nu }=_\mu \omega _\nu _\nu \omega _\mu +[\omega _\mu ,\omega _\nu ]$$
(30)
is Lie algebra $`\mathrm{๐๐
}(2)`$-valued. It is convenient to assume that the dreibein $`๐พ_\mu `$ is also $`\mathrm{๐๐
}(2)`$-valued, where this time the algebra is regarded as a vector space, isomorphic to the three dimensional Minkowski space. As a basis of the $`\mathrm{๐๐
}(2)`$ algebra we take three dimensional Dirac matrices in real (Majorana) representation and the trace in (29) is just the matrix trace. The field equation following from (29) are
$$ฯต^{\mu \nu \rho }D_\nu ๐พ_\rho =0,ฯต^{\mu \nu \rho }๐ฅ_{\nu \rho }=0$$
(31)
The first equation implies that connection is torsion free, while the second assures the metric is flat. The general solutions of these equations on a simply connected region $`UM`$ is well known. It consists of the pair of scalar fields $`(g,f)`$, valued in the Lie group $`\mathrm{๐ฒ๐ซ}(2)`$ and Lie algebra $`\mathrm{๐๐
}(2)`$, respectively, such that
$$\omega _\mu =๐^1_\mu ๐,๐พ_\mu =๐^1_\mu ๐ฟ๐$$
(32)
where $`๐ SL(2)`$ and $`๐\mathrm{๐๐
}(2)`$. Note the similarity between this solution and the solution of the BF theory (10).
Introduction of a particle causes the spacetime to assume the shape of a cone with a particle placed at its top (see Figure 1.) The cone is characterized by the mass-dependent deficit angle $`\alpha `$:
$$\alpha =8\pi Gm$$
(33)
where $`m`$ \- particle mass, $`G`$ \- gravitational constant. In what follows we will set $`8\pi G=1`$, so that the allowed range of the mass is $`m[0,\pi )`$.
In polar coordinates ($`t`$, $`r`$, $`\phi `$), the solution of Einstein equations corresponding to a single spinning particle is of the form
$`๐^0`$ $`=`$ $`dt+{\displaystyle \frac{s}{2\pi }}d\varphi `$
$`๐^1`$ $`=`$ $`\left(1{\displaystyle \frac{m}{2\pi }}\right)\mathrm{cos}\varphi drr\mathrm{sin}\varphi d\varphi `$
$`๐^2`$ $`=`$ $`\left(1{\displaystyle \frac{m}{2\pi }}\right)\mathrm{sin}\varphi dr+r\mathrm{cos}\varphi d\varphi `$
$`๐^0`$ $`=`$ $`{\displaystyle \frac{m}{2\pi }}d\varphi ,๐^1=๐^2=0`$ (34)
This solution corresponds to particle described, similarly to (22) as a delta-like singularity with an appropriate Poincarรฉ charge. There is, however another, more convenient way of treating particles proposed by Matschull and Welling , illustrated in Figure 2. As a result we get singularity free, simply connected spacetime, with boundaries. To make cylindrical boundary look like one dimensional worldline of the particle, we take additional assumption that it circumference vanishes, which can be expressed as requirement that the component $`๐_\phi `$ vanishes on the boundary
$$\overline{๐พ}_\phi =0\text{ at }r=0\text{}$$
(35)
where bar marks the value of the field on the boundary.
Since our manifold is simply connected now, a general solutions of Einstein equations in the neighborhood of the boundary is provided by two functions $`(๐ฟ(r,\phi ,t),๐(r,\phi ,t))`$ satisfying (32). On this solution we must impose appropriate boundary conditions, one of which would be (35) at $`r=0`$, and another that guarantees continuity of dreibein and connection along the cut $`r0`$, $`\varphi =0,2\pi `$. For convenience let us denote $`๐ฟ_\pm (r,t)=๐ฟ(r,t;\varphi =0/2\pi )`$ and $`๐_\pm =๐(r,t;\varphi =0/2\pi )`$. Since $`\omega _\mu `$ and $`๐พ_\mu `$ are to be continuous at the boundary, $`๐ฟ_\pm `$ and $`๐_\pm `$ are not independent, and are related by global Poincarรฉ transformation of the form
$$๐_+=๐ด^1๐_{},๐ฟ_+=๐ด^1(๐ฟ_{}๐)๐ด$$
(36)
where $`๐ด\mathrm{๐ฒ๐ซ}(2)`$ and $`๐\mathrm{๐๐
}(2)`$ are constant.
Now the condition (35) along with (32)
$$\overline{๐พ}_\phi =\overline{๐}^1_\phi \overline{๐ฟ}\overline{๐}=0$$
(37)
tells that $`\overline{๐ฟ}=\overline{๐ฟ}(t)`$. This and the fact that boundary represents worldline the particle makes it possible to identify $`\overline{๐ฟ}`$ with the location of the particle in space-time
$$\overline{๐ฟ}(t,\phi )=๐(t)$$
(38)
Moreover, using the condition $`๐ฟ_+(t)=๐_{}(t)`$ and Poincarรฉ transformation, we find that
$$๐=๐(t)\mathrm{๐ด๐}(t)๐ด^1$$
(39)
Taking time derivative of this equation and making use of the fact that $`๐ด`$ and $`๐`$ are constants gives
$$0=\dot{๐}(t)๐ด\dot{๐}๐ด^1$$
(40)
This last equation is satisfied if and only if the group element $`๐ด`$ is of the form
$$๐ด=u\mathrm{๐}+p_a\gamma ^a,p_a\gamma ^a=\frac{1}{m}\dot{๐}$$
(41)
It is natural then to identify $`p_a`$ with components of momentum of the particle. Note however that this momentum has an unusual property, namely it is, geometrically, a point of the three dimensional anti de Sitter space. Indeed, since $`๐ด\mathrm{๐ฒ๐ซ}(2)`$, $`det๐ด=1`$ and thus it follows from (41) that
$$u^2+p_0^2\stackrel{}{p}^2=1$$
(42)
which is just a definition of anti de Sitter space. We see therefore that three dimensional gravity coupled to point particle possesses a fundamental DSR characteristics: the energy-momentum manifold is curved<sup>8</sup><sup>8</sup>8Four dimensional DSR theories with energy-momentum manifolds of the form of anti de Sitter space have not been intensively investigated, though they are known to exist, see . It is not clear if they arise naturally as a limit of 3+1 gravity. It is also not completely clear if de Sitter energy momentum spaces, intensively investigated in the context of DSR in four dimensions, can be obtained in the 2+1 dimensional case..
It can be shown that instead of the standard dispersion relation the particle on shell satisfies the deformed equation
$$p^ap_a\mathrm{sin}^2m=0$$
(43)
Of course, in the limit when $`m1`$ (remember that the mass scale is set equal 1) we recover the standard dispersion relation.
We know from (32) that the gravitational field in the bulk is pure gauge. It follows that when particles are present the only dynamical degrees of freedom may be associated with boundaries. Therefore, if we start with the action for gravity on the manifold with boundaries, like that in Figure 2, and then perform symplectic reduction, as the result we find action defined only on the worldline on the particle. As we will see in the moment this action differs from the free particle one: the presence of gravitational field causes deformation of the particle lagrangian, exactly in the DSR spirit.
The procedure described briefly above has been performed by Matschull and Welling and the resulting action reads
$$L=\frac{1}{2}Tr(๐ด^1\dot{๐ด}๐)\varsigma (\frac{1}{2}Tr(๐ด)\mathrm{cos}m)$$
(44)
where $`\varsigma `$ is the Lagrange multiplier enforcing the mass shell constraint (43). Using the expression (41) and the fact that $`Tr(\gamma ^a\gamma ^b)=2\eta ^{ab}`$ we can rewrite the Lagrangian in the component form as follows
$$L=\left(\sqrt{p^2+1}\eta ^{ab}+ฯต^{abc}p_c\frac{p^ap^b}{\sqrt{p^2+1}}\right)x_b\dot{p}_a\varsigma \left(p^2\mathrm{sin}^2m\right)$$
(45)
It can be shown that, in spite of the complex, nonlinear form of the Lagrangian, the resulting equations of motion are just the standard one, to wit
$$\dot{p}_a=0,\dot{x}^a=\varsigma p^a.$$
(46)
Let us now turn to discussion of the symmetries of the particle action
$$S=๐\tau L$$
(47)
It is clear from the form of the Langrangian (45) that the action is invariant under standard action of Lorentz generators, so that the Lorentz transformations of both position and momentum are the same as in Special Relativity.
To find generators of Lorentz transformations, we should first derive the form of Poisson brackets, resulting from the symplectic potential in (45). These brackets read
$$\{p_a,p_b\}=0$$
(48)
$$\{p_a,x^b\}=\delta _a{}_{}{}^{b}\sqrt{p^2+1}+ฯต_a{}_{}{}^{bc}p_{c}^{}$$
(49)
and
$$\{x^a,x^b\}=2ฯต^{ab}{}_{c}{}^{}x_{}^{c}$$
(50)
Note that this last bracket tells that the positions of the particle do not commute. The exact form of this bracket differs from $`\kappa `$-Minkowski type of non-commutativity (7), but is very closely related to it .
Using these bracket it is not difficult to derive the form of Noether charges $`J_a`$, generating Lorentz transformations through Poisson bracket. These charges have the form
$$J_a=\sqrt{p^2+1}ฯต_{abc}x^bp^c+2x_{[a}p_{b]}p^b$$
(51)
and together with conserved momenta they form the standard Poincarรฉ algebra.
Note however that while the action of Lorentz generators $`J_a`$ on space-time variables $`x^a`$ is purely classical, the action of translations, generated by momenta $`p_a`$ is deformed, as a result of the bracket (49). This means that in spite of the fact that the particle lives just in Minkowski space-time the translational invariance is lost (or being deformed). This reminds somehow the model considered by Wess , in which also the isometry group is not deformed by itself, but by its action on space-time (or, more generally โ phase space) variables.
The form of the particle Lagrangian (44) suggest simple generalization to the case when the energy momentum space is more general than the $`\mathrm{๐ฒ๐ซ}(2)`$ manifold considered by Matschull and Welling. Consider, for example, the case when this space has the form of de Sitter space. It follows from Iwasawa decomposition of $`\mathrm{๐ฒ๐ฎ}(d,1)`$ group (where $`d`$ is dimension of space-time and momentum space) that in this case the relevant group element has the form ,
$$๐ด=\mathrm{exp}(p_0t_0)\mathrm{exp}(p_it_i)$$
(52)
where the generators of the โtranslationalโ part of Lie algebra $`\mathrm{๐๐}(d,1)`$ $`t_0,t_i`$ satisfy the commutational relations reminding the ones of $`\kappa `$-Minkowski space-time
$$[t_0,t_i]=t_i,[t_i,t_j]=0$$
(53)
The kinetic term of the Lagrangian reads in this case
$$L_k=Tr(๐ด^1\dot{๐ด}๐)=\left(x^0p_ix^i\right)\dot{p}_0x^i\dot{p}_i$$
(54)
and is invariant under action of the Lorentz group of the form of (2), appended by an appropriate action on position variables (see , for details.) In order to get the complete lagrangian, one should add to (54) the term $`\varsigma ๐`$, where $`๐`$ is the Casimir (3). It can be then checked that $`\kappa `$-Minkowski type of non-commutativity (7) follows from the Lagrangian (54). It is not clear, however, if this Lagrangian can be obtained from gravity directly. Work on this question is in progress.
## 5 Conclusions
โSeventy-five thousand generations ago, our ancestors set this program in motion,โ the second man said, โand in all that time we will be the first to hear the computer speak.โ
โWe are the ones who will hear,โ said Phouchg, โthe answer to the great question of Life!..โ
โThe Universe!..โ said Loonquawl.
โAnd Everything!..โ
โAlright,โ said Deep Thought. โThe Answer to the Great Questionโฆโ
โYes!..โ
โOf Life, the Universe and Everythingโฆโ said Deep Thought.
โYes!..โ
โIsโฆโ said Deep Thought, and paused.
โYes!..โ
โIsโฆโ
โYes!!!?..โ
โForty-two,โ said Deep Thought, with infinite majesty and calm.<sup>9</sup><sup>9</sup>9Douglas Adams. The Hitch Hikers Guide to Galaxy Fantazy. 1990.
The current status of DSR reminds somehow the Adamsโ โforty-twoโ, the answer to the question, which we do not really know. To be honest, we do not have any proof yet that this answer is correct, though we hope that Pierre Auger Observatory and GLAST satellite will provide us with such a proof. However, as we tried to argue above, there are more and more indications that the right question is โWhat is the semiclassical, flat space limit of quantum gravity?โ It is our hope that it would not require seventy-five thousand generations to convince ourselves that this hypothesis is correct.
## 6 Acknowledgement
These notes are based on the talk presented by one of us (JKG) at the 339th WE Haraeus-Seminar Potsdam โSpecial Relativity 2005: Will it Survive the Next 100 Years?โ For JKG this work is partially supported by the KBN grant 1 P03B 01828.
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# Critical points and supersymmetric vacua, III: String/M models
## 1. Introduction
This is the third in a series of articles \[DSZ1, DSZ2\] (see also \[Ze2\]) by the authors on statistics of critical points of random holomorphic sections and their applications to the vacuum selection problem in string/M theory. We recall that, in these articles, a โvacuumโ in string theory is a Calabi-Yau manifold of complex dimension $`d=3`$ which forms the $`6`$ โsmall dimensionsโ of the $`10`$-dimensional universe, together with a choice of orientifolding and flux. Mathematically, vacua are critical points of a superpotential $`W`$, a holomorphic section of a line bundle $`๐`$ over the configuration space $`๐`$ which will be recalled in ยง1.1. The โvacuum selection problemโ is that there exists no principle at present which selects a unique superpotential, nor a unique critical point of a given superpotential, out of a large ensemble of possible vacua. This motivates the program of studying statistics of vacua, whose basic problems are to count the number of vacua satisfying physically natural constraints and to determine how they are distributed in $`๐`$ (see \[Do, DD1, AD, DGKT, KL, Si\]). In this article, we present the first rigorous results on counting vacua with remainder estimates. In particular, we justify and improve on the approximations made in \[DD1\].
Our previous articles \[DSZ1, DSZ2\] were devoted to the statistics of critical points of Gaussian random holomorphic sections of line bundles over complex manifolds. The principal issue we face in this article is that the physically relevant ensembles of superpotentials are not Gaussian but rather are discrete ensembles of โquantized fluxโ superpotentials which form a set of lattice points in a hyperbolic shell in $`H^3(X,)`$. This hyperbolic shell is defined by the inequality (known as the tadpole constraint)
$$0Q[\phi ]L,$$
(1)
where
$$Q[\phi ]=Q(\phi ,\overline{\phi })=\sqrt{1}_X\phi \overline{\phi }$$
(2)
is the Hodge-Riemann bilinear form. As will be recalled in ยง2.4, $`Q`$ is an indefinite quadratic form, whose โnull coneโ $`\{G:Q[G]=0\}`$ is a real quadric hypersurface which separates $`H^3(X,)`$ into the interior $`\{W:Q[G]>0\}`$ and the exterior where $`Q[G]<0`$. As will be seen below (Propositions 3.1 and 2.1) , only flux superpotentials corresponding to lattice points in $`\{G:Q[G]>0\}`$ contribute vacua, and that is why we consider the shell (1).
Our main results show that as $`L\mathrm{}`$, the statistics of critical points relative to the discrete lattice ensemble is well approximated by the statistics of critical points relative to the continuum ensemble in the shell, which is dual to the Gaussian ensembles of \[DSZ1, DSZ2\] and is therefore well understood. Thus, the vacuum statistics problem in string/M theory is a mixture of two kinds of equidistribution problems:
1. The distribution of radial projections of lattice points onto a quadric hypersurface;
2. The distribution of critical points of a continuous ensemble of random holomorphic sections (related to a Gaussian ensemble) of a negative line bundle, and their interpretation in the special geometry of Calabi-Yau moduli spaces.
The equidistribution problem in (2) is analyzed in detail in \[DSZ1, DD1\], so the main purpose of this paper is to analyze (1) and to combine it with the previous analysis of (2).
At the end of this article in ยง7 and in \[Ze2\], we compare the mathematical results of this article to discussions of vacua in the string theory literature.
### 1.1. Background to the results
To state our results, we will need some notation (see ยง2 for more details). The models we consider in this article are called type IIb flux compactifications \[GVW, GKP\]. We fix a complex $`3`$-dimensional Calabi-Yau manifold $`X`$, i.e. a complex manifold with trivial canonical bundle $`K_X๐ช`$ and with first Betti number $`b_1(X)=0`$. In some of the physics literature, it is also assumed that $`H^{2,0}(X)=0`$, but our results hold without this assumption. For each complex structure $`z`$ on $`X`$, there is a corresponding Hodge decomposition
$$H^3(X,)=H_z^{3,0}(X)H_z^{2,1}(X)H_z^{1,2}(X)H_z^{0,3}(X).$$
(3)
The space $`H_z^{3,0}(X)`$ of $`(3,0)`$-forms relative to $`z`$ is one-dimensional and is spanned by a nowhere vanishing holomorphic volume form $`\mathrm{\Omega }_z.`$ We also put $`b_3=b_3(X)=dimH^3(X,)`$, $`h^{p,q}=h^{p,q}(X)=dim_{}H^{p,q}(X)`$. Thus, $`b_3=2(h^{2,1}+1)`$.
When we speak of vacua of string theory compactified on the Calabi-Yau space $`X`$, we refer to classical vacua of the effective supergravity theory it determines. As discussed in \[St2\], the effective supergravity Lagrangian is derived by โintegrating outโ or neglecting the massive modes (positive eigenvalues) of various operators. The data of effective supergravity consists of $`(๐,,W)`$ where:
1. $`๐`$ is the configuration space;
2. $`๐`$ is a holomorphic line bundle.
3. the superpotential $`W`$ is a holomorphic section of $``$.
In type IIb flux compactifications the configuration space is the moduli space of Calabi-Yau (Ricci flat Kรคhler ) product metrics on $`X\times T^2`$. At this time of writing, the study of vacua in string theory is simplified by replacing the moduli space of Calabi-Yau metrics by the moduli space of complex structures on $`X`$ (see e.g. \[Do, AD\]). In the case where $`h^{2,0}(X)=0`$, this is equivalent to fixing the Kรคhler class $`[\omega ]H^2(X,)`$ of the Calabi-Yau metrics. Hence we define the configuration space to be
$$๐=\times ,$$
(4)
where $``$ is the moduli space of complex structures on $`X`$ and where $`=/SL(2,)`$ is the moduli space of elliptic curves. Throughout this paper we identify $`๐=\times `$ with a fundamental domain $`๐`$ for the modular group $`\mathrm{\Gamma }`$ in the Teichmรผller space $`๐ฏeich(X)\times `$ of complex structures (see ยง2.1). For simplicity of exposition, we refer to restrictions to $`๐`$ of holomorphic objects on $`๐ฏeich(X)\times `$ as holomorphic objects over $`๐`$.
The line bundle $``$ is defined to be the dual line bundle to the Hodge bundle $`H^{3,0}(X)H^{1,0}(T^2)๐`$, where $`T^2=^2/^2`$. We give $`๐`$ the Weil-Petersson Kรคhler form $`\omega _{WP}`$ induced from the Weil-Petersson metric on $``$ (see ยง3.3). To be precise, $``$ is a holomorphic line bundle over $`๐ฏeich(X)\times `$, and $`W`$ is a holomorphic section of $`๐ฏeich(X)\times `$. But as mentioned above, by holomorphic sections $`WH^0(๐,)`$ we mean restrictions to $`๐`$ of holomorphic sections of $`H^0(๐ฏeich(X)\times ,).`$
Type IIb flux compactifications contain two non-zero harmonic $`3`$-forms $`F,HH^3(X,)`$ which are known respectively as the RR (Ramond-Ramond) and NS (Neveu-Schwartz) $`3`$-form field strengths. We combine them into a complex flux $`G=F+iHH^3(X,i)`$. The parameter $`\tau `$ is known as the dilaton-axion and may be viewed as the period of $`\omega _\tau =dx+\tau dy`$ over the one-cycle dual to $`dy`$ in $`T^2`$. Given $`GH^3(X,\sqrt{1}),`$ physicists define the corresponding flux superpotential $`W_G`$ by:
$$W_G(z,\tau )=_X(F+\tau H)\mathrm{\Omega }_z,$$
(5)
where $`\mathrm{\Omega }_zH^{3,0}(X)`$. This is not well-defined as a function on $`๐`$ since $`\mathrm{\Omega }_z`$ and $`\tau `$ depend on a choice of frame. To be more precise, $`GH^3(X,)`$ determines a section $`W_G`$ of the line bundle
$$=(H^{3,0}(X)H^{1,0}(T^2))^{}๐ฏeich(X)\times $$
by making $`G`$ into the following linear functional on $`H_z^{3,0}(X)H_\tau ^{1,0}(T^2):`$
$$W_G(z,\tau ),\mathrm{\Omega }_z\omega _\tau =_{X\times T^2}(FdyHdx)(\mathrm{\Omega }_z\omega _\tau ).$$
(6)
The map $`GW_G`$ defines an injective real (but not complex) linear map which embeds complex integral fluxes
$$H^3(X,\sqrt{1})H^0(๐,)$$
(7)
as a lattice of rank $`2b_3`$ in $`H^0(\times ,)`$ which we call the lattice $`๐ฎ^{}`$ of integral flux superpotentials. The real span
$$๐ฎ=๐ฎ^{}H^0(,)$$
(8)
of $`๐ฎ^{}`$ is also important, and will be referred as the space of flux superpotentials. We emphasize here that $`๐ฎ`$ is not a complex vector space, nor are any of the associated spaces discussed below. We also use the (real-linear) map $`GW_G`$ to regard $`Q`$ as a quadratic form on $`๐ฎ`$, writing
$$Q[W_G]:=Q[G]=\sqrt{1}_XG\overline{G}=2_XFH,G=F+iHH^3(X,).$$
(9)
The bundles $`H_z^{3,0}`$ and $`H_\tau ^{1,0}`$ carry Weil-Petersson Hermitian metrics $`h_{WP}`$ defined by
$$h_{WP}(\mathrm{\Omega }_z,\mathrm{\Omega }_z)=e^{K(z,\overline{z})}=i_X\mathrm{\Omega }_z\overline{\mathrm{\Omega }}_z,$$
(10)
and their associated Chern connections $`_{WP}`$. They induce dual metrics and connections on $``$. We denote the connection simply by $``$.
### 1.2. Statement of the problem
Given a flux superpotential $`W`$, there is an associated potential energy on $`๐`$ defined by
$$V_W(Z)=|W(Z)|^23|W(Z)|^2.$$
(11)
(See \[WB\] for background on $`V`$). By a vacuum we mean a critical point of $`V(Z)`$ on $`๐`$. In this paper, we only study supersymmetric vacua, namely $`Z๐`$ which are connection critical points in the sense that $`_{WP}W(Z)=0.`$ We denote the set of supersymmetric vacua of $`W`$ by
$$Crit(W)=\{Z๐:_{WP}W(Z)=0\}.$$
(12)
Our goal is thus to count and find the distribution law of the supersymmetric vacua
$$\{\text{SUSY vacua}\}=\underset{G๐ฎ^{}:Q\left[G\right]L}{}Crit(W_G)$$
(13)
as $`W_G`$ varies over the lattice $`๐ฎ^{}`$ within the hyperbolic shell (1). To define the distribution law, we introduce the incidence relation
$$=\{(W_G,Z)๐ฎ\times ๐:W_G(Z)=0\}.$$
(14)
We shall view $`๐`$ as a fundamental domain for the modular group $`\mathrm{\Gamma }`$ in Teichmรผller space (cf. ยง2). The incidence variety $``$ is then a real $`2m`$-dimensional subvariety of $`๐\times ๐ฎ`$ with the following diagram of projections:
$$\begin{array}{ccccc}& & ๐\times ๐ฎ& & \\ \rho & \pi & & & \\ ๐& ๐ฎ& & & \end{array}$$
(15)
The fiber $`\pi ^1(W)`$ is the set $`Crit(W)`$ of critical points of $`W`$ in $`๐`$. Since $`๐`$ is regarded as a fundamental domain in Teichmรผller space, the map $`\pi `$ is not surjective: there exist $`W`$ with no critical points in $`๐`$; hence $`\pi (๐)`$ is a domain with boundary in $`๐ฎ`$ (see ยง6.4.1). Critical points can move out of $`๐`$ as $`W`$ varies in $`๐ฎ`$. (There is a similar but more complicated theory of non-supersymmetric vacua \[DD2\].)
The fibers of $`\rho `$ are the subspaces
$$๐ฎ_Z:=\{W๐ฎ:_{WP}W(Z)=0\},$$
(16)
which play a crucial role in this article. They have the remarkable Hodge theoretic identifications,
$$๐ฎ_{z,\tau }H_z^{2,1}(X)H_z^{0,3}(X)(\text{Proposition}\text{3.1}).$$
(17)
It then follows (see Proposition 3.2) that $`\stackrel{\rho }{}๐`$ is a vector bundle (with fiber $`^{b_3/2}`$) over a manifold with boundary. Another key point is that the restrictions of $`Q`$ to the fibers are always positive definite:
$$Q|_{H_z^{2,1}(X)H_z^{0,3}(X)}0(\text{Proposition}\text{2.1}),$$
(18)
i.e. $`๐ฎ_Z`$ lies in the positive cone $`\{Q(\phi ,\overline{\phi })>0\}`$ of the indefinite quadratic (Hodge-Riemann) form (2) (cf. ยง2.4).
We now define the discriminant locus
$$\stackrel{~}{๐}=\{(Z,W):detH^cW(Z)=0\}$$
of points $`(Z,W)`$ such that $`Z`$ is a degenerate critical point of $`W`$, where $`H^cW(Z)`$ is the complex Hessian of $`W`$ at the critical point $`Z`$ as defined in (59)โ(61). Equivalently, $`\stackrel{~}{๐}`$ is the set of critical points of the second projection $`\stackrel{\pi }{}๐ฎ`$ together with the singular points of $``$. Its image $`๐=\pi (\stackrel{~}{๐})`$ under $`\pi `$ is the discriminant variety of superpotentials with degenerate critical points.
For each $`W๐ฎ\{0\}`$, we define its distribution of (non-degenerate) critical points as the measure $`C_W`$ on $`\stackrel{~}{๐}`$ defined by
$$C_W,\psi =\underset{ZCrit(W)}{}\psi (Z,W),$$
(19)
for $`\psi ๐()`$ such that $`\rho (\mathrm{Supp}\psi )`$ is relatively compact in $`๐`$ and $`\mathrm{Supp}\psi `$ is disjoint from $`\stackrel{~}{๐}`$. A more general definition of $`C_W`$ is
$$C_W=|detH^cW(Z)|W^{}\delta _0$$
(20)
which will be discussed in ยง4.2. We make these assumptions on $`\psi `$ so that the sum on the right side is a finite and well-defined sum. Indeed, the pull back is not well-defined (without further work) on $`\stackrel{~}{๐}`$. We will say more about $`\stackrel{~}{๐}`$ after the statement of Theorem 1.4.
The basic sums we study are :
$`๐ฉ_\psi (L)`$ $`=`$ $`{\displaystyle \{C_N,\psi :N๐ฎ^{},Q[N]L\}}`$ (21)
$`=`$ $`{\displaystyle \{\psi (Z,N):(Z,N),N๐ฎ^{},0Q[N]L\}}.`$
For instance, when $`\psi \chi _K`$ is the characteristic function of a compact subset $`K\stackrel{~}{๐}`$, $`N_\psi (L)`$ counts the total number of non-degenerate critical points lying over $`\rho (K)`$ coming from all integral flux superpotentials with $`Q[W]L`$. Physicists are naturally interested in counting the number of vacua with close to the observed values of the cosmological constant and other physical quantities, and hence would study sums relevant to such quantities. For instance, the cosmological constant of the theory defined by a vacuum $`Z`$ is the value $`V(Z)`$ of the potential there (see \[DD1\], ยง3.3). Thus, we may state the main problem of this paper:
###### Problem 1.1.
Find the asymptotics and remainder for $`๐ฉ_\psi (L)`$ as $`L\mathrm{}.`$
As indicated above, this problem is very closely related to the pure lattice point problem of measuring the rate of uniform distribution of radial projections of lattice points onto the surface of a quadric hypersurface. More generally one could consider any smooth strictly convex set $`Q^n`$ ($`n2)`$ with $`0Q^{}`$. Associated to $`Q`$ is the norm $`|X|_Q`$ of $`X^n`$ defined by
$$Q=\{X^n:|X|_Q<1\}.$$
To measure the equidistribution of radial projections of lattice points to $`Q`$, one considers the sums
$$S_f(t)=\underset{k^ntQ\{0\}}{}f\left(\frac{k}{|k|_Q}\right),\text{with }fC^{\mathrm{}}(Q),t>0.$$
(22)
The parallel lattice point problem is then
###### Problem 1.2.
Find the asymptotics and remainder for $`S_f(t)`$ as $`t\mathrm{}.`$
### 1.3. Statement of the results
In Theorem 5.1, we obtain a van der Corput type estimate for the lattice point problem 1.2. For the critical point problem, we first give an elementary formula which is based on a trivial lattice counting estimate (which is useful since it is sometimes sharp), namely where the remainder term is simply a count of the cubes of the lattice which intersect the boundary. We denote by $`\chi _{Q_Z}`$ the characteristic function of the shell $`\{W๐ฎ_Z:0<Q_Z[W]<1\}`$.
###### Proposition 1.3.
Suppose that $`\psi =\chi _K`$ where $`K`$ such that $`(Z,W)K(Z,rW)K`$ for $`r^+`$. Assume further that $`\rho (K)`$ is relatively compact in $`๐`$ and $`\pi (K)`$ is piecewise smooth. Then
$$๐ฉ_\psi (L)=L^{b_3}\left[_๐_{๐ฎ_Z}\psi (Z,W)|detH^cW(Z)|\chi _{Q_Z}(W)๐Wd\mathrm{Vol}_{WP}(Z)+O\left(L^{1/2}\right)\right].$$
Here and in Theorem 1.4 below, $`dW`$ means the multiple of Lebesgue measure on $`๐ฎ_Z`$ which gives the volume form for the positive-definite quadratic form $`Q_Z=Q|_{๐ฎ_Z}`$. We note that the integral converges, since by (18), $`\{Q_Z1\}`$ is an ellipsoid of finite volume.
It would be interesting to know if the remainder estimate is sharp for any domain $`K`$. In the pure lattice point Problem 1.2, the corresponding โtrivial estimateโ is sharp. For instance, consider the domain $`K=S_+^{n1}S^{n1}`$ formed by the northern hemisphere and put $`\psi =\chi _K`$. Then the remainder term
$$\underset{k^n,|k|\sqrt{L}}{}\chi _K\left(\frac{k}{|k|}\right)L^{\frac{n}{2}}_Kf๐A$$
reflects the concentration of projections of lattice points on the boundary $`S_+^{n1}`$, namely a great equatorial sphere. When the equator is defined by $`x_n=0`$, the lattice points projecting over the equator are the lattice points in $`^{n1}^{n1}`$ and the number with $`|k|\sqrt{L}`$ is of size $`L^{\frac{n1}{2}}.`$ Analogously one may ask if there are domains $`K๐`$ along which critical points concentrate to the same maximal degree. Some evidence that the answer is โnoโ will be presented in ยง4.1.
Our main result stated below is a much sharper van der Corput type asymptotic estimate of $`๐ฉ_\psi (L)`$ as $`L\mathrm{}`$ for homogeneous test functions which vanish near the discriminant locus. Here, we say that a function $`\psi ๐()`$ is homogeneous of order $`\alpha `$ if
$$\psi (Z,rW)=r^\alpha \psi (Z,W),(Z,W),r^+.$$
We consider homogeneous functions since they include (smoothed) characteristic functions as well as the cosmological constant (which is homogeneous of degree $`2`$).
###### Theorem 1.4.
Let $`\psi ๐^{\mathrm{}}()`$ be homogeneous of order $`\alpha 0`$ and suppose that $`\rho (\mathrm{Supp}\psi )`$ is a compact subset of $`๐`$ and $`\mathrm{Supp}\psi \stackrel{~}{๐}=\mathrm{}`$. Then
$$๐ฉ_\psi (L)=L^{b_3+\alpha /2}\left[_๐_{๐ฎ_Z}\psi (Z,W)|detH^cW(Z)|\chi _{Q_Z}(W)๐Wd\mathrm{Vol}_{WP}(Z)+O\left(L^{\frac{2b_3}{2b_3+1}}\right)\right].$$
It is reasonable to make the assumption $`\mathrm{Supp}\psi \stackrel{~}{๐}=\mathrm{}`$, because degenerate critical points cannot be physically acceptable vacua in string/M theory. Indeed, the Hessian of $`W`$ at a critical point defines the โfermionic mass matrixโ of the theory, and a degenerate critical point would give rise to massless fermions which are not observed in physics. (See \[WB\] for definitions of the mass matrix.)
Let us note some key features of the geometry of $`\stackrel{~}{๐}`$ which play a role in the assumptions (and proofs) of Proposition 1.3 and Theorem 1.4. First, as observed in \[DSZ1, DSZ2\], its defining equation
$$detH^cW(Z)=det(H^{}H|W|^2I)=0$$
(23)
is real valued; here, $`H`$ is the holomorphic Hessian (see ยง3.2). Hence, $`\stackrel{~}{๐}`$ is a real analytic hypersurface (with boundary). For test functions $`\psi `$ which do not vanish on $`\stackrel{~}{๐}`$, the expression $`C_W,\psi `$ (when well-defined) can jump as one passes from one component of $`๐ฎ๐`$ to another or across the boundary of $`๐`$. It follows from (23) that $`\stackrel{~}{๐}(\{Z\}\times ๐ฎ_Z)`$ is a real conic hypersurface for all $`Z๐`$. Thus $`\stackrel{~}{๐}๐`$ is a bundle of conic hypersurfaces and $`\rho (\stackrel{~}{๐})=๐`$; i.e., every point of moduli space is a degenerate critical point of some superpotential. We further note that $`๐ฎ๐`$ consists of a finite number of connected components, and that $`\pi :\stackrel{~}{๐}\pi (๐ฎ)๐`$ is a finite covering over each connected component of $`\pi (๐ฎ)๐`$.
### 1.4. Special geometry and critical point density
In obtaining reliable order of magnitude results on numbers of vacua in a given string/M model, it is important to estimate the size of the leading coefficient
$$_๐\psi (Z)_{๐ฎ_Z}|detH^cW(Z)|\chi _{Q_Z}(W)๐Wd\mathrm{Vol}_{WP}(Z)$$
and of the remainder. Since little is known about the volume of $`๐`$ at present (cf. \[LS1\]), we concentrate on estimating the integrand
$$๐ฆ^{\mathrm{crit}}(Z):=_{๐ฎ_Z}|detH^cW(Z)|\chi _{Q_Z}๐W$$
(24)
in the $`b_3`$ aspect. It is also important to study the behavior of the $`๐ฆ^{\mathrm{crit}}(Z)`$ as $`Z`$ tends to โinfinityโ in $`๐`$, or to a singular point such as a conifold point (when one exists).
A key feature of $`๐ฆ^{\mathrm{crit}}(Z)`$ is that it is the integral of a homogeneous function of order $`b_3`$ over a space of dimension $`dim_{}๐ฎ_Z=b_3=2(h^{2,1}+1)`$. Among the known Calabi-Yau $`3`$-folds it is common to have $`300<b_3<1000`$, hence the integral is often over a space of large dimension. The $`b_3`$-dependence is sensitive since (e.g.) the ratio of the $`L^{\mathrm{}}`$ norm to the $`L^2`$ norm of a homogeneous function of degree $`b_3`$ in $`b_3`$ variables can be of order $`b_3^{b_3}.`$ It is useful to have alternative formulas for the leading coefficient, and we now present a few. We will use them to suggest conjectures on the order of magnitude of $`๐ฆ^{\mathrm{crit}}(Z)`$ in the $`b_3`$ aspect in ยง7.
First, using the homogeneity of the integrand, we may rewrite the integral in terms of a Gaussian density
$`๐ฆ^{\mathrm{crit}}(Z)`$ $`=`$ $`{\displaystyle \frac{1}{b_3!}}{\displaystyle _{๐ฎ_Z}}|detH^cW(Z)|e^{Q_ZW,W}๐W.`$ (25)
This formula shows that $`๐ฆ^{\mathrm{crit}}`$ is formally analogous to density of critical points of random holomorphic sections relative to a Gaussian measure studied in \[DSZ1\]. For this reason, we call (24) the critical point density. However, the measure $`e^{Q[W]}\chi _{\{0<Q<1\}}(W)dW`$ is of infinite volume, so the analogy should not be taken too literally. The density $`๐ฆ^{\mathrm{crit}}(Z)`$ is well-defined despite the infinite volume of the underlying measure on $`๐ฎ`$ because the fibers $`Q_Z`$ of $`\rho |_Q`$ are of finite volume. Indeed, the conditional measures of $`e^{Q[W]}dW`$ are standard (un-normalized) Gaussian measures $`e^{Q_Z(W)}dW`$.
Next, we rewrite the integrals by the methods in \[DSZ1, DSZ2\]. The first method is to change variables to the Hessian $`H^cW(Z)`$, i.e. to โpush-forwardโ the $`๐ฎ_Z`$ integral under the Hessian map
$$H_Z:๐ฎ_Z\mathrm{Sym}(m,),H_Z(W)=H^cW(Z),$$
(26)
where $`m=dim๐=h^{2,1}+1`$. In \[DSZ1, DSZ2\], we used this change of variables to simplify the formulas for the density of critical points. There, however, the spaces of holomorphic sections of the line bundles $`LM`$ were so large that the image of the Hessian map was the entire space $`\mathrm{Sym}(m,)`$ of complex Hessians of rank equal to the dimension $`m=dimM`$. In the case of type IIb flux compactifications, the dimension of the configuration space $`๐`$ is as large as the dimension of the space $`๐ฎ`$ of sections, and the Hessian map is by no means surjective. Indeed, in Lemma 6.1, we prove that the Hessian map is an isomorphism to a real $`b_3`$-dimensional space $`_Z`$, where $`_Z`$ is spanned (over $``$) by the $`2h^{2,1}`$ Hermitian matrices
$$\xi ^j:=\left(\begin{array}{cc}0& e_j\\ e_j^t& ^j(z)\end{array}\right),\xi ^{h^{2,1}+j}:=\left(\begin{array}{cc}0& ie_j\\ ie_j^t& i^j(z)\end{array}\right),j=1,\mathrm{},h^{2,1}.$$
(27)
Here, $`e_j`$ is the $`j`$-th standard basis element of $`^{h^{2,1}}`$ and $`^j(z)\mathrm{Sym}(h^{2,1},)`$ is the matrix $`\left(_{ik}^{\overline{j}}(z)\right)`$ whose entries define the โYukawa couplingsโ on $``$ (see (46), ยง2.3 or \[St1, CO\]) with respect to normal coordinates at the point $`z`$.
Since $`_Z`$ is not a complex subspace of $`\mathrm{Sym}(m,)`$, we regard $`\mathrm{Sym}(m,)`$ as a real vector space with inner product
$$(A,B)_{}=\mathrm{Re}A,B_{HS}=\mathrm{Re}(\text{Trace}AB^{}).$$
(28)
To state our next result, we let $`\mathrm{\Lambda }_Z`$ be the operator given by the distortion under the Hessian map (see ยง6.2):
$$((\mathrm{\Lambda }_ZI_{})^1H_ZW,H_ZW)_{}=Q[W](W๐ฎ_Z),$$
(29)
where $`Q[W]`$ is given by (9). In terms of the basis $`\{\xi ^a\}_{1a2h^{2,1}}`$,
$$\mathrm{\Lambda }_Z\xi ^a=\underset{b=1}{\overset{2h^{2,1}}{}}\mathrm{\Lambda }_{ab}\xi ^b,\mathrm{\Lambda }_{ab}=(\xi ^a,\xi ^b)_{}.$$
The $`\mathrm{\Lambda }`$ matrix has the block form
$$(\mathrm{\Lambda }_{ab})=\left(\begin{array}{cc}\mathrm{\Lambda }^{}& \mathrm{\Lambda }^{\prime \prime }\\ \mathrm{\Lambda }^{\prime \prime }& \mathrm{\Lambda }^{}\end{array}\right),\mathrm{\Lambda }_{jk}^{}=2\delta _{jk}+\mathrm{Re}\text{Tr}^j^k,\mathrm{\Lambda }_{jk}^{\prime \prime }=\mathrm{Im}\text{Tr}^j^k.$$
(30)
In Proposition 6.2, we show that the $`(1,1)`$ form
$$\omega _\mathrm{\Lambda }:=\frac{i}{2}(\mathrm{\Lambda }_{jk}^{}+i\mathrm{\Lambda }_{jk}^{\prime \prime })dz^jd\overline{z}^k=\frac{i}{2}\left[2\delta _{jk}+\text{Tr}^j(z_0)^k(z_0)\right]dz^jd\overline{z}^k$$
(31)
is the so-called Hodge metric $`(m+3)\omega _{WP}+Ric(\omega _{WP})`$ of the Weil-Petersson metric \[Lu, Wa2\].
By the injectivity of the Hessian map (stated in Lemma 6.1), we can make the change of variables $`W\stackrel{H_Z}{}(H,x)`$ in (24)โ(25) to obtain the following alternate formulas for $`๐ฆ^{\mathrm{crit}}(Z)`$:
$`๐ฆ^{\mathrm{crit}}(Z)`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{det\mathrm{\Lambda }_Z}}}{\displaystyle __Z}\left|detH^{}H|x|^2I\right|\chi _{\mathrm{\Lambda }_Z}(H,x)๐H๐x,`$ (32)
$`=`$ $`{\displaystyle \frac{1}{b_3!\sqrt{det\mathrm{\Lambda }_Z}}}{\displaystyle __Z}\left|detH^{}H|x|^2I\right|e^{(\mathrm{\Lambda }_Z^1H,H)_{}|x|^2}๐H๐x`$
where $`\chi _{\mathrm{\Lambda }_Z}`$ is the characteristic function of the ellipsoid $`\{(\mathrm{\Lambda }_Z^1H,H)_{}+|x|^21\}.`$ These formulas are analogous to Theorem 1 and Corollary 2 of \[DSZ1\], the key difference being that here we integrate over a moving subspace $`_Z`$ of symmetric matrices.
We similarly have the following alternative formulations of Proposition 1.3 and Theorem 1.4:
###### Corollary 1.5.
Let $`\psi =\chi _K`$, where $`K`$ is as in Proposition 1.3, and let $`\stackrel{~}{\psi }(Z,H_ZW)=\psi (Z,W)`$. Then,
$`๐ฉ_\psi \left(L\right)={\displaystyle \frac{L^{b_3}}{b_3!}}\left[{\displaystyle _๐}{\displaystyle \frac{1}{\sqrt{det\mathrm{\Lambda }_Z}}}{\displaystyle __Z}\stackrel{~}{\psi }(Z;H,x)\right|detH^{}H\left|x|^2I\right|e^{(\mathrm{\Lambda }_Z^1H,H)_{}\left|x\right|^2}dHdxd\mathrm{Vol}_{WP}\left(Z\right)`$
$`+O\left(L^{1/2}\right)].`$
###### Corollary 1.6.
Let $`\psi ๐^{\mathrm{}}()`$ be homogeneous of order $`\alpha 0`$ and suppose that $`\rho (\mathrm{Supp}\psi )`$ is a compact subset of $`๐`$ and $`\mathrm{Supp}\psi \stackrel{~}{๐}=\mathrm{}`$. Let $`\stackrel{~}{\psi }(Z,H_ZW)=\psi (Z,W)`$. Then,
$`๐ฉ_\psi (L)`$ $`=`$ $`{\displaystyle \frac{L^{b_3+\alpha /2}}{\mathrm{\Gamma }(b_3+\alpha /2+1)}}[{\displaystyle _๐}{\displaystyle \frac{1}{\sqrt{det\mathrm{\Lambda }_Z}}}{\displaystyle __Z}\stackrel{~}{\psi }(Z;H,x)`$
$`\times |detH^{}H|x|^2I|e^{(\mathrm{\Lambda }_Z^1H,H)_{}|x|^2}dHdxd\mathrm{Vol}_{WP}(Z)+O\left(L^{\frac{2b_3}{2b_3+1}}\right)].`$
It is not obvious how to estimate the dependence of the integral for $`๐ฆ^{\mathrm{crit}}(Z)`$ on the subspace $`_Z`$. There are two natural ways to parameterize this space. One (which is used in \[DD1\]) is to use as a basis of $`_Z`$ the Hessians of a $`Q_Z`$-orthonormal basis of $`๐ฎ_Z`$. A second method is to use the orthonormal basis of eigenmatrices $`\{H_j\}`$ of $`\mathrm{\Lambda }_Z`$ with respect to the inner product (28). We thus put $`\mathrm{\Lambda }_ZH_j(Z)=\mu _j(Z)H_j(Z)`$, and $`H(y,Z)=_jy_jH_j(Z)`$. We also let $`D(\mu )`$ denote the diagonal matrix with entries $`\mu _j`$. Changing variables to $`\mu _j^{1/2}y`$ cancels $`\frac{1}{\sqrt{det\mathrm{\Lambda }_Z}}`$ and we obtain:
###### Corollary 1.7.
We have:
$$๐ฆ^{\mathrm{crit}}(Z)=_{|y|^2+|x|^21}\left|detH(D(\mu )y,Z)^{}H(D(\mu )y,Z)|x|^2I\right|๐y๐x.$$
In ยง7 we will discuss some conjectural bounds on the density of critical point based on the assumption that the subspaces $`_Z`$ are sufficiently random subspaces of $`\mathrm{Sym}(h^{2,1},)`$.
### 1.5. Index density
The absolute value in the expressions for the distribution of critical points $`C_W`$ of a single section (20) and the expected distribution of critical points of a random section (e.g., (32)) make it very difficult to estimate the order of magnitude of the density of critical points. A simplifying โapproximationโ is to drop the absolute value around the determinant. The resulting density is index density for critical points. It was used in \[AD\] and \[DD1\] to give a lower bound for the critical point density.
To be precise, we modify (20) by defining the signed distribution of critical points of $`W`$ as the measure $`C_W`$ on $`\stackrel{~}{๐}`$ given by
$$Ind_W,\psi =\underset{ZCrit(W)}{}\left(\text{sign}detD^2W(Z)\right)\psi (Z,W),$$
(33)
where sign$`a=1,0,1`$ if $`a`$ is positive, 0, or negative, respectively. We then study the sums
$`nd_\psi (L)`$ $`=`$ $`{\displaystyle \{Ind_N,\psi :N๐ฎ^{},Q[N]L\}}.`$ (34)
For instance, if $`\psi (Z,W)=\chi _K(Z)`$ is the characteristic function of a compact set $`K๐`$, then $`nd_\psi (L)`$ is the sum $`_{ZCrit(W)K}\left(\text{sign}detD^2W(Z)\right)`$ over all non-degenerate critical points lying over $`K`$ of all integral flux superpotentials with $`Q[W]L`$.
Simultaneously with Proposition 1.3, we obtain formula (1.5) of Ashok-Douglas \[AD\] with an estimate for the error produced by passing from the sum to the integral (cf. ยง4):
###### Theorem 1.8.
Let $`K`$ be a compact subset of $`๐`$ with piecewise smooth boundary. Then
$$nd_{\chi _K}(L)=\frac{(\pi L)^{b_3}}{b_3!\mathrm{\hspace{0.17em}2}^{b_3/2}}\left[_Kc_m(T^{(1,0)}(๐),\omega _{WP}h_{WP}^{})+O\left(L^{1/2}\right)\right],$$
where $`m=dim๐=b_3/2`$ and $`c_m(T^{(1,0)}(๐),\omega _{WP}h_{WP}^{})=\frac{1}{\pi ^m}det\left(R\omega _{WP}I\right)`$ is the $`m`$-th Chern form of $`T^{(1,0)}(๐)`$ with respect to the Weil-Petersson metric $`\omega _{WP}h_{WP}^{}`$.
Here, $`R=_{ij}R_{\mathrm{}i\overline{j}}^kdz^id\overline{z}^{\overline{j}}`$ is the curvature $`(1,1)`$ form of $`T^{(1,0)}(๐)`$ regarded as an $`m\times m`$ Hermitian-matrix-valued $`2`$-form (with $`m=dim๐`$= $`b_3/2`$) and $`\omega _{WP}I`$ is a scalar 2-form times the $`m\times m`$ identity matrix. The determinant is defined as in Chern-Weil theory. The only additional step in the proof is the evaluation (given in Lemma 6.3) of the analogue of (25) in terms of the curvature form:
$$_{๐ฎ_Z}detH^cW(Z)e^{Q_ZW,W}dW=\left(\frac{\pi }{2}\right)^m\frac{det\left(R\omega _{WP}I\right)}{d\mathrm{Vol}_{WP}}.$$
(35)
Recall that the Chern-Gauss-Bonnet theorem tells us that if $`W`$ is a holomorphic section of a complex line bundle $`LM_m`$ over a compact complex manifold such that $`W`$ has only non-degenerate zeros, then
$$c_m(T^{(1,0)}ML)=IndW:=\underset{p:W(p)=0}{}\text{sign}detH^cW(p).$$
However, the Chern-Gauss-Bonnet theorem does not apply in our setting, and indeed $`IndW`$ is not constant in $`W`$, since $`๐`$ is an incomplete Kรคhler manifold and critical points can occur on the boundary or disappear. There exists a Chern-Gauss-Bonnet theorem for manifolds with boundary which expresses $`IndW`$ as $`c_n(E)`$ plus a boundary correction depending on $`W`$, but the correction term involves integrating a differential form over the boundary and that becomes problematic when the boundary is highly irregular as in the case of $`๐`$. Nevertheless, the theorem shows that asymptotically the average index density equals the Chern-Gauss-Bonnet form.
### 1.6. Relations to prior results in the physics and mathematics literature
We now relate our results to the physics literature on the number of vacua and the complexity of the string theory landscape as well as to the mathematical literature on lattice points. A more detailed discussion of the landscape aspects is given in ยง7.
First, the string/M aspects. Over the last five years or so, many physics articles have been devoted to estimating the number of candidate vacua $`N_{vac}`$ of string/M theory, in particular those which are consistent with the standard model. The candidate vacua are often pictured as valleys in a โstring theory landscapeโ, which is the graph of the effective potential. The number of vacua is often stated as being around $`10^{500}`$. In \[BP\] Bousso-Polchinski related the number of vacua to the number of quantized fluxes $`N`$ satisfying a constraint $`|N|L`$, which implies $`N_{vac}(L)\frac{L^{b_3}}{b_3!}`$ (see also \[AD, Si\]). In the specific type IIb flux compactifications studied in this paper, the constraint is hyperbolic rather than elliptic (as imagined in \[BP\]), and the more precise estimate $`N_{vac}(L)\frac{L^{b_3}}{b_3!}f(b_3)`$ was given in \[AD, DD1\] where $`f(b_3)`$ is the moduli space integral of the Gaussian integral in (32); it will be discussed further in ยง7. There we will also review the heuristics and the mathematics of the landscape in more detail.
What do our results imply about the number of vacua? Since Proposition 1.3 and Theorem 1.4 are asymptotic results as $`L\mathrm{}`$, they are most useful when $`L^{b_3}`$ is very large. But it is difficult to quantify โvery largeโ due to the complexity of the leading coefficient (24), of the remainder and of the volume of $`๐`$. Hence, we cannot make precise estimates on the number of vacua at this time.
However, to bridge our results with estimates in string theory, we make a speculative attempt in ยง7.3 to draw order of magnitude conclusions from Theorem 1.4. We will use the symbol $``$ in an informal sense of โsame order of magnitudeโ (factorial, exponential and so on). There we give a heuristic estimate of $`๐ฆ^{\mathrm{crit}}(Z)\frac{1}{b_3!}(b_3/2)!\mu ^{b_3}`$ for certain $`\mu >0`$. More precisely, we give heuristic upper and lower bounds with different $`\mu `$ which are irrelevant when comparing factorials. To obtain an order of magnitude for $`\frac{f(b_3)}{b_3!}`$ one would need to integrate $`๐ฆ^{\mathrm{crit}}`$ over $`๐`$. At this time, the order of magnitude of the Weil-Petersson volume $`Vol_{WP}(๐)`$ of $`๐`$ is not known, even approximately (Z. Lu). We can however make a plausible estimate for the integral of $`๐ฆ^{\mathrm{crit}}`$ over the region where the norm of $`\mathrm{\Lambda }_Z`$ is bounded by a uniform constant (independent of $`b_3`$). Since $`\mathrm{\Lambda }_Z`$ is essentially the Hodge metric, regions where $`\mathrm{\Lambda }_Z\mu `$ are regions $`K_\mu `$ where the norm of the Ricci curvature of $`\omega _{WP}`$) is bounded above by a uniform constant. It appears likely that the volume of such regions is bounded above by the volume of balls in $`^{b_3/2}`$ of fixed radius (Z. Lu). Since the volume of balls in $`^{b_3/2}`$ decays like $`\frac{1}{(b_3/2)}!`$, we would find that the number of vacua in $`K_\mu `$ would be approximately $`\frac{L^{b_3}}{b_3!}\mu ^{b_3}`$.
Now, in the physical models, $`L`$ is not a free parameter but is determined by $`X`$. In the case when there exists an involution $`g`$ of $`X`$ (an โorientifoldingโ) and a Calabi-Yau $`4`$-fold $`Z`$ which is an elliptic fibration over $`X/g,`$ the โtadpoleโ number is then given by:
$$\text{ tadpole number}:L=\chi (Z)/24.$$
(36)
In many known examples \[KLRY\], one has $`300<b_3<1,000`$ and $`LCb_3`$ where $`1/3C3`$. Hence the number of vacua in $`K_\mu `$ (and possibly in all of $`๐`$) with the tadpole constraint $`LCb_3`$ would have exponential growth $`\frac{(Cb_3)^{b_3}}{b_3!}\mu ^{b_3}`$.
Next we turn to the purely lattice point aspects of the problem. From a mathematical point of view, this article combines statistical algebraic geometry in the sense of \[BSZ1, DSZ1, DSZ2\] with the study of radial projections of lattice points. As far as we know, the radial projection of lattice points problem has not been studied systematically before in mathematics (we thank B. Randol for helping to sort out the historical background on this problem). The much harder problem of equidistribution of lattice points of fixed height $`R`$, i.e. lying on a sphere or hyperboloid of fixed radius $`R`$, has been studied by Yu. Linnik, C. Pommerenke \[Po\], W. Duke and others. But the remainders obtained in this more delicate problem are not as accurate as ours are for the bulk problem of projecting all lattice points of height $`<R`$. Counting projections of lattice points in domains of a hypersurface is equivalent to counting lattice points in certain cones, and there are some additional studies of this by methods of automorphic forms. In certain right circular cones with a flat top, Duke and Imamoglu \[DO\] use Dirichlet series and Shimura lifts to obtain the leading order asymptotics. Radial projections of lattice points additionally bear some resemblance to rational points. Some results and references for this problem are contained in \[DO\]. In \[Ze2\], the general problem of counting radial projections of lattice points in smooth domains of non-degenerate hypersurfaces is studied. In \[NR\], some further results are given on radial projections of lattice points, in particular in the case of hypersurfaces with flat spots or in the case of polyhedra.
Acknowledgement: We would like to thank Zhiqin Lu for many helpful comments regarding the Weil-Petersson and Hodge metrics on the moduli space of a Calabi-Yau 3-fold. In particular, our discussion of the Weil-Petersson volume $`V_{WP}(๐)`$ and estimates of the eigenvalues of $`\mathrm{\Lambda }_Z`$ are based on his remarks.
## 2. Background on Calabi-Yau manifolds and complex geometry
As mentioned in the introduction, the supersymmetric vacua of type IIb flux compactifications on a $`CY_3`$ are critical points of holomorphic sections of the holomorphic line bundle $`๐`$ dual to the Hodge bundle $`H^{4,0}(X\times T^2)`$, where the configuration space $`๐`$ is the moduli space $`\times `$ of product complex structures on $`X\times T^2`$. In this section, we give the geometric background necessary for the analysis of critical points and Hessians of the holomorphic sections $`W_G`$ of (5).
The most significant aspects of Calabi-Yau geometry in the study of critical points of flux superpotentials are the following:
* The space $`๐ฎ_Z`$ of flux superpotentials with $`W_G(Z)=0`$ may be identified with the space $`H_Z^3(X)`$ of fluxes $`G=F+iH`$ with the special Hodge decomposition $`F+\tau HH_z^{2,1}(X)H_z^{0,3}(X).`$ See Proposition 3.1.
* The space $`H_z^{2,1}(X)H_z^{0,3}(X)H^3(X,)`$ is a positive complex Lagrangian subspace. See Proposition 2.1. Hence, $`๐ฎ_Z`$ is endowed with an inner product.
In addition, we review the relation between holomorphic derivatives, covariant derivatives and second fundamental forms for holomorphic frames $`\mathrm{\Omega }_z`$ of the Hodge bundle, and recall the definition of the prepotential. These results are needed for the calculations in Lemmas 108 and 6.1. Much of this material is essentially standard \[CO, St1, DD1\], but it is not always stated precisely in the physics sources. We therefore present it here for the sake of clarity and completeness.
### 2.1. Geometry of Calabi-Yau manifolds
We recall that a Calabi-Yau $`d`$-fold $`M`$ is a complex $`d`$-dimensional manifold with trivial canonical bundle $`K_M`$, i.e. $`c_1(M)=0`$. By the well-known theorem of Yau, there exists a unique Ricci flat Kรคhler metric in each Kรคhler class on $`M`$. In this article, we fix the Kรคhler class, and then the Calabi-Yau metrics correspond to the complex structures on $`M`$ modulo diffeomorphisms. We denote the moduli space of complex structures on $`M`$ by $``$.
As mentioned in the introduction, the Calabi-Yau manifolds of concern in this article are the $`4`$-folds $`M=X\times T^2,`$ where $`T^2=^2/^2`$. The $`T^2`$ factor plays a special role, and the geometric aspects mainly concern $`X`$. We only consider product Calabi-Yau metrics and complex structures on $`M`$. Thus, the configuration space $`๐=\times `$ where $``$ is the moduli space of complex structures on $`X`$ and where $``$ is the moduli space of elliptic curves. We denote a point of $`๐`$ by $`Z=(z,\tau )`$ where $`z`$ denotes a complex structure on $`X`$ and where $`\tau `$ denotes the complex structure on $`T^2`$ corresponding to the elliptic curve $`/\tau .`$
It is often simplest to view the moduli space of complex structures on $`X`$ as the quotient by the mapping class group $`\mathrm{\Gamma }`$ of the Teichmรผller space $`๐ฏeich(X)`$, where
$$๐ฏeich(X)=\{\text{complex structures on }X\}/\mathrm{Diff}_0$$
where $`JJ^{}`$ if there exists a diffeomorphism $`\phi \mathrm{Diff}_0`$ isotopic to the identity satisfying $`\phi ^{}J^{}=J.`$ The mapping class group is the group of connected components of the diffeomorphism group,
$$\mathrm{\Gamma }_X:=\mathrm{Diff}(X)/\mathrm{Diff}_0(X).$$
We shall identify $``$ with a fundamental domain for $`\mathrm{\Gamma }_X`$ in $`๐ฏeich(X)`$, and $``$ with the usual modular domain in $``$.
The mapping class group for a Calabi-Yau $`d`$-fold has a representation on $`H^d(M,)`$ which preserves the intersection form $`Q`$, which is symplectic in odd dimensions, and indefinite symmetric in even dimensions. In odd dimensions, this representation gives a homomorphism $`\phi :\mathrm{\Gamma }_M\text{Sp}(b_d(M),)`$, while in even dimensions it gives a homomorphism to the corresponding orthogonal group. It was proved by D. Sullivan \[Sul\] that if $`d3`$, then $`\phi (\mathrm{\Gamma }_M)`$ is an (arithmetic) subgroup of finite index (in $`\text{Sp}(b_d(M),)`$ if $`d`$ is odd), and the kernel of $`\phi `$ is a finite subgroup.
On any CY manifold $`M`$ of dimension $`d`$, the space $`H_z^{d,0}(M)`$ of holomorphic $`(d,0)`$ forms for a complex structure $`Z`$ is one-dimensional. It depends holomorphically on $`Z`$ and hence defines a complex holomorphic line bundle $`_{}^{}=H^{d,0}`$, which we refer to as the โHodge bundle.โ The Hodge bundle is equipped with the Weil-Petersson (WP) Hermitian metric of (10), which we repeat here:
$$h_{WP}(\mathrm{\Omega }_z,\mathrm{\Omega }_z)=i^{d^2}_X\mathrm{\Omega }\overline{\mathrm{\Omega }}.$$
(37)
For a holomorphic Hermitian line bundle $`(L,h)M`$ and local holomorphic frame $`e_L`$ over an open set $`UM`$, we write
$$|e_L(z)|_h^2=e^{K(z)}.$$
(38)
The connection $`1`$-form in this frame is the $`(1,0)`$ form $`K(z)`$, and the curvature $`(1,1)`$ -form is given by
$$\omega =\frac{i}{2}\mathrm{\Theta }_h=\frac{i}{2}\overline{}K,K=\mathrm{log}|e_L|_h^2.$$
The Hermitian line bundle is said to be positive if $`\omega `$ is a positive $`(1,1)`$ form, in which case $`K`$ is called the Kรคhler potential. The Hermitian line bundle $`(L,h)`$ is negative if $`\omega `$ is a negative $`(1,1)`$ form.
In particular, the curvature of the Weil-Petersson metric on $`H^{d,0}`$ is a positive $`(1,1)`$ form on $``$, and hence it defines a Kรคhler form with potential (with respect to the frame $`\mathrm{\Omega }_z`$)
$$K_{WP}=\mathrm{log}h_{WP}(\mathrm{\Omega }_z,\mathrm{\Omega }_z)=\mathrm{log}i_X\mathrm{\Omega }\overline{\mathrm{\Omega }}.$$
(39)
For instance, consider the Hodge bundle $`H_\tau ^{1,0}`$. It has a standard frame $`dx+\tau dy`$ for which $`K=\mathrm{log}\mathrm{Im}\tau .`$ Here, $`\tau `$ is the standard coordinate on the upper half plane. Then $`K=\frac{1}{\tau \overline{\tau }}d\tau `$ and the Kรคhler form is $`\frac{i}{2(\tau \overline{\tau })^2}d\tau d\overline{\tau }>>0.`$
In the product situation of $`M=X\times T^2`$, $`H_{z,\tau }^{4,0}(X\times T^2)=H_z^{3,0}(X)H_\tau ^{1,0}(T^2)`$. Thus, the line bundle $`H^{4,0}(X\times T^2)H^{3,0}(X)H^{1,0}(T^2)๐`$ is an exterior tensor product and the WP metric is a direct product. We denote an element of $`H_z^{3,0}(X)`$ by $`\mathrm{\Omega }_z`$, and an element of $`H_\tau ^{1,0}(T^2)`$ by $`\omega _\tau `$. We often assume that $`\omega _\tau =dx+\tau dy`$.
### 2.2. Variational derivatives and covariant derivatives
The bundle $`H_z^{3,0}(X)`$ is a holomorphic line bundle. Since $`H_z^{3,0}(X)H^3(X,)`$, one can view a holomorphically varying family $`\mathrm{\Omega }_zH_z^{3,0}(X)`$ as a holomorphic map $`H^3(X,)`$ or as a holomorphic section of $`H_z^{3,0}(X)`$. As a holomorphic vector valued function, $`\mathrm{\Omega }_z`$ may be differentiated in $`z`$. If $`z_1,\mathrm{},z_{h^{2,1}}`$ are local holomorphic coordinates, and if $`\{\frac{}{z_j}\}`$ are the coordinate vector fields, then $`\frac{\mathrm{\Omega }}{z_j}`$ is a well-defined element of $`H^3(X,)`$.
By the Griffiths transversality theorem, see \[GHJ\], \[CO\], (5.4)) or \[Wa1, Wa2\],
$$\frac{\mathrm{\Omega }_z}{z^j}=k_j(z)\mathrm{\Omega }_z+\chi _j,$$
(40)
where $`\chi _jH_z^{2,1}(X)`$ and where $`kC^{\mathrm{}}()`$. Note that although $`\frac{\mathrm{\Omega }_z}{z^j}`$ is holomorphic, neither term on the right hand side is separately holomorphic.
We define a Levi-Civita connection on the bundle $`H^{3,0}`$ by orthogonally projecting the derivatives $`\frac{\mathrm{\Omega }_z}{z^j}`$ onto $`H^{3,0}`$. This defines the Weil-Petersson connection $`_{WP}`$ on $`H^{3,0}`$,
$$_{WP}:C^{\mathrm{}}(,)C^{\mathrm{}}(,T^{}).$$
It follows from (40) that
$$\frac{}{z^j}_X\mathrm{\Omega }_z\overline{\mathrm{\Omega }}_z=k_j_X\mathrm{\Omega }_z\overline{\mathrm{\Omega }}_z,$$
(41)
which by (39) implies
$$k_j=\frac{K}{z^j}.$$
(42)
Hence,
$$_{WP}\mathrm{\Omega }_z=K\mathrm{\Omega }_z=k_jdz_j\mathrm{\Omega }_z$$
is the Chern connection of the Weil-Petersson Hermitian metric.
We also define the forms
$$\{\begin{array}{c}๐_j\mathrm{\Omega }_z=\frac{}{z^j}\mathrm{\Omega }+\frac{K}{z^j}\mathrm{\Omega }\hfill \\ \\ ๐_j๐_k\mathrm{\Omega }_z=(\frac{}{z^j}+\frac{K}{z^j})(\frac{}{z^k}+\frac{K}{z^k})\mathrm{\Omega }_z.\hfill \end{array}$$
(43)
We then have
$$๐_j\mathrm{\Omega }_z=\frac{\mathrm{\Omega }_z}{z^j}k_j\mathrm{\Omega }_z=\chi _jH^{2,1}(X_z).$$
(44)
The operator $`๐_j\mathrm{\Omega }_z`$ is analogous to the second fundamental form $`II(X,Y)=(\stackrel{~}{}_XY)^{}`$ of an embedding, i.e. it is the โnormalโ component of the ambient derivative. It is known that the first variational derivatives span $`H^{2,1}`$ (see e.g. \[Wa1, Wa2\]. (In the physics literature, $`D_\alpha `$ is often described as a connection, and is often identified with $`_{WP}`$, but this is not quite correct as it is applied to $`\mathrm{\Omega }_z`$).
The Weil-Petersson Hermitian metric $`G_{i\overline{j}}dz_id\overline{z}_{\overline{j}}`$ on $``$ is the curvature $`(1,1)`$-form of the Hodge bundle. From (39) and (44), we have:
$$G_{j\overline{k}}=\frac{^2}{z^j\overline{z}^k}K(z,\overline{z})=\frac{_{}๐_j\mathrm{\Omega }_z\overline{๐_k\mathrm{\Omega }_z}}{_{}\mathrm{\Omega }_z\overline{\mathrm{\Omega }}_z}.$$
(45)
### 2.3. Yukawa couplings and special geometry of the moduli space
In formula (32), the density of critical points is expressed as an integral over a space $`_Z`$, where $`_Z`$ is a subspace of the complex symmetric matrices $`\mathrm{Sym}(h^{2,1}+1,)`$ spanned by the special matrices $`\xi ^j`$ given in (27). Their components $`_{ik}^{\overline{j}}(z)`$ are known as Yukawa couplings and defined as follows: A priori, $`๐_k๐_j\mathrm{\Omega }_zH^{2,1}H^{1,2}`$, and moreover its $`H^{2,1}`$ component vanishes (see e.g. \[CO, (5.5)\]). Hence we may define $`_{kj}^{\overline{l}}`$ by
$$๐_k๐_j\mathrm{\Omega }_z=\sqrt{1}e^K_{kj}^{\overline{l}}\overline{๐_l\mathrm{\Omega }}(1j,k,lh^{2,1}).$$
(46)
See also \[St1, (28)\]. It is further shown in \[St1, (37)\] (see also \[AD, (4.8)\], \[LS2, Theorem 3.1\]) that the Riemann tensor of the Weil-Petersson metric on the moduli space $``$ of Calabi-Yau three-folds is related to the Yukawa couplings by
$$R_{i\overline{j}k\overline{\mathrm{}}}=G_{i\overline{j}}G_{k\overline{\mathrm{}}}+G_{i\overline{\mathrm{}}}G_{k\overline{j}}e^{2K}\underset{p,q}{}G^{p\overline{q}}_{ikp}\overline{_{j\mathrm{}q}}.$$
(47)
The Yukawa couplings are related to the periods of $`\mathrm{\Omega }_z`$ and to the so-called prepotential of $``$. We pause to recall the basic relations and to direct the reader to the relevant references.
First, we consider periods. As a basis of $`H_3(X,)`$ we choose the symplectic basis consisting of dually paired Lagrangian subspaces of $`A`$-cycles $`A_a`$ and $`B`$-cycles $`B_a`$. The periods of $`\mathrm{\Omega }_zH_z^{3,0}(X)`$ over the $`A`$-cycles
$$\zeta ^a=_{A_a}\mathrm{\Omega }_z(1ah^{2,1}+1=b_3/2)$$
define holomorphic coordinates on $`_{}^{}=H^{3,0}`$. Alternately, we can view the $`\zeta ^a`$ as โspecialโ projective coordinates on $``$. The periods of $`\mathrm{\Omega }_z`$ over the $`B`$-cycles are then holomorphic functions of the $`\zeta ^a`$. The principal fact is that the image of $`_{}^{}`$ under the period map is a complex Lagrangian submanifold of $`H^3(M,)`$, and thus is determined by a single holomorphic function, the โprepotentialโ $`=(\zeta ^1,\mathrm{},\zeta ^{b_3/2}):_{}^{}`$ such that
$$_{B_a}\mathrm{\Omega }_z=\frac{}{\zeta ^a}.$$
(48)
Furthermore, $``$ is homogeneous of degree $`2`$ in the periods $`\zeta ^a`$,
$$\underset{j=1}{\overset{b_3/2}{}}\zeta ^a\frac{}{\zeta ^a}=2(z),$$
and hence may be viewed as a holomorphic section of $`_{}^2`$.
The local holomorphic $`3`$-form $`\mathrm{\Omega }_z`$ may be expressed in terms of the Poincarรฉ duals of the symplectic basis by:
$$\mathrm{\Omega }_z=\underset{a=1}{\overset{b_3/2}{}}\left(\zeta ^a\widehat{A}_a\frac{}{\zeta ^a}\widehat{B}_a\right).$$
(49)
(See \[CO\], (3.8)). Further, in these coordinates, the Kรคhler potential (39) of the Weil-Petersson metric may be written as
$$K(z,\overline{z})=\mathrm{log}i\left(\underset{a=1}{\overset{b_3/2}{}}\zeta ^a\overline{\frac{}{\zeta ^a}}\overline{\zeta ^a}\frac{}{\zeta ^a}\right).$$
We also have:
$$_{kj}^{\overline{l}}=\underset{r=1}{\overset{h^{2,1}}{}}G^{r\overline{l}}\frac{^3}{z^rz^jz^k}.$$
(50)
See \[CO, (4.5)\] and \[St1, (64)\].
In summary, we reproduce the table from \[CO\]:
$$\begin{array}{cc}\text{Derivatives of the Basis}\hfill & \text{spans}\hfill \\ & \\ \mathrm{\Omega }\hfill & H^{3,0}\hfill \\ & \\ ๐_j\mathrm{\Omega }\hfill & H^{2,1}\hfill \\ & \\ ๐_k๐_j\mathrm{\Omega }=ie^K_{kj}^{\overline{\gamma }}\overline{๐_\gamma \mathrm{\Omega }}\hfill & H^{12}\hfill \\ & \\ ๐_k๐_{\overline{j}}\mathrm{\Omega }=G_{k\overline{j}}\overline{\mathrm{\Omega }}\hfill & H^{03}\hfill \end{array}$$
(51)
#### 2.3.1. $`๐`$ as the moduli space of complex structures on $`X\times T^2`$
Above, we have reviewed the geometry of the moduli space of complex structures on the Calabi-Yau three-fold. Our configuration space $`๐=\times `$ may be viewed as (a component of) the moduli space of complex structures on $`X\times T^2`$. This point of view is used in \[DD1\], but because the $`T^2`$ factor plays a distinguished role we do not emphasize this identification here. Further, formula (47) needs to be modified for the moduli space of complex structures on a Calabi-Yau four-fold. In \[LS2, Theorem 3.1\]), the Riemann tensor of the Weil-Petersson metric on the moduli space of a Calabi-Yau manifold of arbitrary dimension is shown to be
$$R_{i\overline{j}k\overline{\mathrm{}}}=G_{i\overline{j}}G_{k\overline{\mathrm{}}}+G_{i\overline{\mathrm{}}}G_{k\overline{j}}\frac{๐_k๐_i\mathrm{\Omega },\overline{๐_{\mathrm{}}๐_j\mathrm{\Omega }}}{_{}\mathrm{\Omega }\overline{\mathrm{\Omega }}}.$$
(52)
In the case of three-folds, the vectors $`๐_j\mathrm{\Omega }`$ form an orthonormal basis for $`H^{2,1}`$ and one can write the inner product in the form (47).
### 2.4. Hodge-Riemann form and inner products
The Hodge-Riemann bilinear form on $`H^3(X,)`$ is the intersection form $`(\alpha ,\beta )_X\alpha \beta `$. We consider the sesquilinear pairing:
$$(\alpha ,\beta )Q(\alpha ,\overline{\beta })=\sqrt{1}_X\alpha \overline{\beta },\alpha ,\beta H^3(X,).$$
(53)
An important fact is that under the Hodge decomposition (3) for a given complex structure, the Hodge-Riemann form is definite in each summand:
$$(1)^pQ(\alpha ,\overline{\alpha })>0,\alpha H^{p,3p}(X,),$$
(54)
whose sign depends only on the parity of $`p`$. (See \[GH, ยง7\]. Note that our definition of $`Q`$ has the extra sign $`\sqrt{1}`$. The inequality (54) holds only for primitive forms, but in our case all harmonic 3-forms are primitive, since we are assuming that $`H^1(M,)=0`$.) To restate (54):
###### Proposition 2.1.
Let $`dimX=3`$, and let $`b_1(X)=0`$. Then for each $`z`$, the Hodge-Riemann form is positive definite on $`H_z^{2,1}H_z^{0,3}`$ and negative definite on $`H_z^{3,0}H_z^{1,2}.`$
By Griffiths transversality (see (40)), for any local holomorphic frame $`\mathrm{\Omega }_z`$, $`๐_j\mathrm{\Omega }_zH_z^{2,1}`$ and these elements span $`H_z^{2,1}`$. Also, $`\overline{\mathrm{\Omega }}_z`$ spans $`H^{0,3}`$. These forms provide us with an orthonormal basis for $`H_z^{2,1}H_z^{0,3}`$:
###### Proposition 2.2.
If $`\{z_j\}`$ are coordinates at $`z_0`$ such that $`\left\{/z_j|_{z_0}\right\}`$ are orthonormal, and if $`h_{WP}(\mathrm{\Omega }_{z_0},\mathrm{\Omega }_{z_0})=1`$, then the basis $`\{๐_j\mathrm{\Omega }_{z_0},\overline{\mathrm{\Omega }}(z_0)\}`$ is a complex orthonormal basis of $`H_{z_0}^{2,1}H_{z_0}^{0,3}`$ with respect to the Hodge Riemann form $`Q.`$
Remark: Here and below, when we say that a basis of a complex vector space is complex orthonormal we mean that it is a complex basis and is orthonormal for the given inner product. By a real orthonormal basis of the same vector space we mean an orthonormal basis of the underlying real vector space.
###### Proof.
It suffices to show that:
$$\begin{array}{cc}\hfill (\mathrm{i})& Q(๐_j\mathrm{\Omega }_z,\overline{๐_k\mathrm{\Omega }_z})=i_X๐_j\mathrm{\Omega }_z\overline{๐_k\mathrm{\Omega }_z}=G_{j\overline{k}}e^K\hfill \\ & \\ \hfill (\mathrm{ii})& Q(๐_j\mathrm{\Omega }_z,\mathrm{\Omega }_z)=i_X๐_{\overline{j}}\mathrm{\Omega }_z\mathrm{\Omega }_z=0\hfill \\ & \\ \hfill (\mathrm{iii})& Q(\overline{\mathrm{\Omega }}_z,\mathrm{\Omega }_z)=i_X\overline{\mathrm{\Omega }}_z\mathrm{\Omega }_z=h_{WP}(\mathrm{\Omega }_z,\mathrm{\Omega }_z)\hfill \end{array}$$
Equation (i) follows from (45), (ii) is by type considerations, and (iii) follows from (10). โ
Remark: In the language of complex symplectic geometry, Proposition 2.1 says that $`H_z^{2,1}H_z^{0,3}`$ is a positive complex polarization of $`H^3(X,)`$. Let us recall the definitions. The space $`(H^3(X,),Q)`$ of real $`3`$-cycles with its intersection form $`Q(\alpha ,\beta )=i_M\alpha \beta `$ is a real symplectic vector space. After complexifying, we obtain the complex symplectic vector space $`(H^3(X,),Q).`$ In general, if $`(V,\omega )`$ is a real symplectic vector space and if $`(V_{},\omega _{})`$ is its complexification, a complex Lagrangian subspace $`FV_{}`$ is called a polarization. The polarization is called real if $`F=\overline{F}`$ and complex if $`F\overline{F}=\{0\}`$. The polarization $`F`$ is called positive if $`i\omega (v,\overline{w})`$ is positive definite on $`F`$.
In our setting, $`(V,\omega )=(H^3(X,),Q)`$. We observe that for any complex structure $`z`$ on $`X`$ (as a complex manifold), the Hodge decomposition may be written in the form
$$H^3(X,)=F\overline{F},F=H^{2,1}H^{0,3}\overline{F}=H^{3,0}H^{1,2},,$$
where $`F`$ is complex Lagrangian. By Proposition 2.1, this polarization is positive, i.e.
$$Q(v,\overline{v})>0,vF\{0\}.$$
## 3. Critical points of superpotentials
In this section, we assemble some basic facts about critical points and Hessians of flux superpotentials.
### 3.1. Flux superpotentials as holomorphic sections
As discussed in the previous section, $`๐`$ is a negative line bundle. On a compact complex manifold, a negative line bundle has no holomorphic sections. However, $`(๐,\omega _{WP})`$ is a non-compact, incomplete Kรคhler manifold of finite Weil-Petersson volume (see \[LS1\] for the latter statement), and the line bundle $`๐`$ has many holomorphic sections related to the periods of $`X\times T^2`$.
The sections relevant to this article are the flux superpotentials $`W_G`$ of (5)-(6). $`W_G`$ depends on two real fluxes $`F,HH^3(X,)`$, which we combine into a complex integral flux
$$G=F+iHH^3(X,\sqrt{1}).$$
The main reason to form this complex combination is that it relates the tadpole constraint (1) on the pair $`(F,H)`$ to the Hodge-Riemann form (2). However, none of subsequent identifications preserves this complex structure, and the reader may prefer to view $`G`$ as just the pair $`G=(F,H)H^3(X,)H^3(X,)`$. Alternately, we can identifying $`G=F+iHH^3(X,)`$ with the real cohomology class
$$\stackrel{~}{G}:=FdyHdxH^4(X\times T^2,)H^3(X,).$$
We shall consider the (real-linear) embedding
$$๐ฒ:H^3(X,)H^0(๐,),GW_G,$$
where $`W_G`$ is given by formula (6); i.e.,
$$(W_G(z,\tau ),\mathrm{\Omega }_z\omega _\tau )=_{X\times T^2}\stackrel{~}{G}\mathrm{\Omega }_z\omega _\tau .$$
We denote by $`๐ฎ=`$Image$`(๐ฒ)`$ the range of this map, and by
$$๐ฎ^{}=๐ฒ\left(H^3(X,i)\right)$$
the lattice of sections satisfying the integrality condition. The map $`GW_G`$ is not complex linear, so $`๐ฎ`$ is not a complex subspace of $`H^0(\times ,)`$. Rather, it is a real subspace of dimension $`2b_3`$ (over $``$) and $`๐ฎ^{}`$ is a lattice of rank $`2b_3`$ in it. In fact $`๐ฎ^{2b_3}`$ is totally real in $`H^0(๐,)^{2b_3}`$.
We choose local holomorphic frames $`\mathrm{\Omega }_z`$ of the Hodge bundle $`H^{3,0}`$ and $`\omega _\tau =dx+\tau dy`$ of $`H^{1,0}`$ and let $`\mathrm{\Omega }_z^{}\omega _\tau ^{}`$ denote the dual co-frame of $``$. A holomorphic section of $``$ can then be expressed as $`W=f(z,\tau )\mathrm{\Omega }_z^{}\omega _\tau ^{}`$ where $`f๐ช(๐)`$ is a local holomorphic function. If $`W=W_G`$ is a flux superpotential, then the corresponding function $`f_G`$ is given by:
$$f_G(z,\tau )=_{X\times T^2}(FdyHdx)(\mathrm{\Omega }_z\omega _\tau ).$$
(55)
When $`\omega _\tau =dx+\tau dy`$ (on a fundamental domain in Teichmรผller space), we obtain the simpler form:
$$f_G(z,\tau )=_X(F+\tau H)\mathrm{\Omega }_z.$$
(56)
### 3.2. Critical points and Hessians of holomorphic sections
As preparation for critical points of superpotentials, we recall some basic notations and facts concerning critical points and Hessians of holomorphic sections of a general line bundle $`LM`$ (see \[DSZ1\]).
Let $`(L,h)M`$ be a holomorphic Hermitian line bundle, let $`e_L`$ denote a local frame over an open set $`U`$ and write a general holomorphic section as $`s=fe_L`$ with $`f๐ช(U)`$. Recall that the Chern connection $`_h`$ of $`h`$ is given locally as $`(fe_L)=(ffK)e_L`$, where $`K=\mathrm{log}e_L_h^2`$, i.e.
$$s=\underset{j=1}{\overset{m}{}}\left(\frac{f}{z^j}f\frac{K}{z^j}\right)dz^je_L=\underset{j=1}{\overset{m}{}}e^K\frac{}{z^j}\left(e^Kf\right)dz^je_L.$$
(57)
The critical point equation thus reads,
$$\frac{f}{z^j}f\frac{K}{z^j}=0.$$
The Hessian of a holomorphic section $`s`$ of $`(L,h)M`$ at a critical point $`Z_0`$ is the tensor
$$Ds(Z_0)T^{}T^{}L$$
where $`D`$ is a connection on $`T^{}L`$. At a critical point $`Z_0`$, $`Ds(Z_0)`$ is independent of the choice of connection on $`T^{}`$. In a local frame and in local coordinates we have
$$D^{}^{}s(Z_0)=\underset{j,q}{}H_{jq}^{}dz^qdz^je_L,D^{\prime \prime }^{}s(Z_0)=\underset{j,q}{}H_{jq}^{\prime \prime }d\overline{z}^qdz^je_L.$$
(58)
The Hessian $`Ds(Z_0)`$ at a critical point thus determines the complex symmetric matrix $`H^c`$ (which we call the โcomplex Hessianโ):
$$H^c:=\left(\begin{array}{cc}H^{}& H^{\prime \prime }\\ \overline{H^{\prime \prime }}& \overline{H^{}}\end{array}\right)=\left(\begin{array}{cc}H^{}& f(Z_0)\mathrm{\Theta }\\ \overline{f(Z_0)\mathrm{\Theta }}& \overline{H^{}}\end{array}\right),$$
(59)
whose components are given by
$`H_{jq}^{}`$ $`=`$ $`({\displaystyle \frac{}{z^j}}{\displaystyle \frac{K}{z^j}})({\displaystyle \frac{}{z^q}}{\displaystyle \frac{K}{z^q}})f(Z_0),`$ (60)
$`H_{jq}^{\prime \prime }`$ $`=`$ $`f{\displaystyle \frac{^2K}{z^j\overline{z}^q}}|_{Z_0}=f(Z_0)\mathrm{\Theta }_{jq},\mathrm{\Theta }_h(Z_0)={\displaystyle \underset{j,q}{}}\mathrm{\Theta }_{jq}dz^jd\overline{z}^q.`$ (61)
### 3.3. Supersymmetric critical points and the Hodge decomposition
We now specialize to the critical point equations for flux superpotentials $`W_G(z,\tau )`$. An important observation that is now standard in the physics literature is that the complex moduli $`(z,\tau )`$ at which a flux superpotential $`W_G(z,\tau )`$ satisfies $`W_G=0`$ are characterized by the following special Hodge decomposition of $`H^3(X,)`$ at $`z`$ (see \[AD\], (3.5)โ(3.8)).
A local holomorphic frame for the Hodge bundle $`๐`$ is $`e_{}=\mathrm{\Omega }_z^{}\omega _\tau ^{}`$, where $`\mathrm{\Omega }_z^{}`$ is dual to the local frame $`\mathrm{\Omega }_z`$ of the Hodge line bundle $`H^{3,0}`$ and $`\omega _\tau ^{}`$ is dual to the local frame $`\omega _\tau =dx+\tau dy`$ of $`H^{1,0}`$. We let $`K(Z)=K_X(z)+K_{T^2}(\tau )`$ be the Kรคhler potential for the local frame $`\mathrm{\Omega }_z\omega _\tau `$ of the (positive) Hodge bundle $`^{}`$. We then have
$$|e_{}(Z)|_h^2=|\mathrm{\Omega }_z\omega _\tau |_{h_{WP}}^2=e^{K(Z)}=e^{K_X(z)}e^{K_{T^2}(\tau )}.$$
(62)
Hence, the Weil-Petersson Kรคhler potential on $`๐`$ is
$$K(Z)=\mathrm{log}_X\mathrm{\Omega }_z\overline{\mathrm{\Omega }}_z\mathrm{log}(\overline{\tau }\tau ).$$
In particular, the $`\tau `$-covariant derivative on $``$ is given in the local frame $`e_{}`$ by
$$_\tau =\frac{}{\tau }+\frac{1}{\overline{\tau }\tau }.$$
(63)
Hence with $`W_G=f_Ge_{}`$, we have
$`_\tau f_G`$ $`=`$ $`{\displaystyle _X}\left[H+{\displaystyle \frac{1}{\overline{\tau }\tau }}(F+\tau H)\right]\mathrm{\Omega }_z`$ (64)
$`=`$ $`{\displaystyle \frac{1}{\overline{\tau }\tau }}{\displaystyle _X}(F+\overline{\tau }H)\mathrm{\Omega }_z.`$
To compute the $`z`$-derivatives, we see from ยง2.2 and (56)โ(57) that
$`_{z^j}f_G`$ $`=`$ $`\left({\displaystyle \frac{f_G}{z^j}}+{\displaystyle \frac{K}{z^j}}f_G\right)(z,\tau )={\displaystyle _X}(F+\tau H)\left({\displaystyle \frac{\mathrm{\Omega }_z}{z^j}}+{\displaystyle \frac{K}{z^j}}\mathrm{\Omega }_z\right)`$ (65)
$`=`$ $`{\displaystyle _X}(F+\tau H)\chi _j=0,`$
for $`1jh^{2,1}`$. Thus, the supersymmetric critical point equations for the flux superpotential $`W_G`$ read:
$$\{\begin{array}{c}_X(F+\tau H)๐_j\mathrm{\Omega }_z=0(1jh^{2,1})\hfill \\ \\ _X(F+\tau H)\overline{\mathrm{\Omega }}_z=0.\hfill \end{array}$$
(66)
As in (16), we denote by $`๐ฎ_Z`$ $`(Z=(\tau ,z))`$ the space of superpotentials $`W_G`$ with $`W_G(Z)=0`$. Although the equation is complex linear on $`H^0(๐,)`$, $`๐ฎ`$ is not a complex subspace of $`H^0(๐,)`$, so $`๐ฎ_Z`$ is a real but not complex vector space. Put another way, for each $`Z=(z,\tau )`$, the critical point equation determines a real subspace
$$H_Z^3(X,)=๐ฒ^1(๐ฎ_Z)=\{F+iH,F,HH^3(X,),(\text{66})\text{is true}\}.$$
(67)
The critical point equations (66) put $`b_3=2(h^{2,1}+1)`$ independent real linear conditions on $`2b_3`$ real unknowns $`(F,H)`$.
###### Proposition 3.1.
\[AD, DD1\] Let $`G=F+iH`$ with $`F,HH^3(X,)`$, and let $`W_G(z,\tau ),\mathrm{\Omega }_z\omega _\tau =_X(F+\tau H)\mathrm{\Omega }_z`$ be the associated superpotential. If $`_{z,\tau }W_G(z,\tau )=0`$, then $`(F+\tau H)H_z^{2,1}H_z^{0,3}`$. Moreover, the map
$$I_\tau :H^3(X,)H^3(X,),I_\tau (F+iH)=F+\tau H$$
restricts to give real linear isomorphisms
$$I_{z,\tau }:H_{z,\tau }^3H_z^{2,1}(X)H_z^{0,3}(X),$$
of real vector spaces.
###### Proof.
We first prove that $`(F+iH)F+\tau H`$ takes $`H_Z^3H_z^{2,1}H_z^{0,3}`$. Suppose that $`W_G=0`$. Since the $`\chi _j(z)`$ span $`H_z^{2,1}`$, we conclude from the first equation of (66) that $`(F+\tau H)_z^{1,2}=0`$; by the second equation, we also have $`(F+\tau H)_z^{3,0}=0`$. Thus $`F+\tau HH_z^{2,1}H_z^{0,3}.`$
Since $`I_{z,\tau }`$ is injective and since $`dim_{}H_{z,\tau }^3=dim_{}H_z^{2,1}H_z^{0,3}=b_3`$, it is clearly an isomorphism. โ
### 3.4. The map $`(z,\tau )H_{z,\tau }^3`$
As $`(z,\tau )`$ varies over $`๐`$, how do the spaces $`H_{z,\tau }^3`$ move in $`H^3(X,)`$? This question is important in relating the pure lattice point problem in $`H^3(X,)`$ to the vacuum distribution problem in $`๐`$. It depends on the geometry of the diagram
$$\begin{array}{ccccc}& ๐\times H^3(X,)& & & \\ \rho \pi & & & & \\ ๐H^3(X,),& & & & \end{array}$$
(68)
where $`=\{(z,\tau ,F,H):F+iHH_{(z,\tau )}^3(X)\},`$ which is a replica of (15) in which $`๐ฎ`$ is replaced by $`H^3(X,)`$.
To answer this question, we first note that for each $`(z,\tau )๐`$, the real-linear map
$$H_{z,\tau }^3H^3(X,),F+iHH$$
is bijective. Injectivity follows by noting that
$$FH_{z,\tau }^3FH_z^{2,1}H_z^{0,3}F=\overline{F}H_z^{1,2}H_z^{3,0}F=0.$$
Since both spaces have dimension $`b_3`$, bijectivity follows. Thus there is a real linear isomorphism $`\iota _{z,\tau }:H^3(X,)H_{z,\tau }^3`$ of the form
$$\iota _{z,\tau }(H)=F(z,\tau ,H)+iH.$$
To describe $`F(z,\tau ,H)`$, we form the $`z`$-dependent basis
$$\{\mathrm{Re}D_1\mathrm{\Omega }_z,\mathrm{},\mathrm{Re}D_{h^{2,1}}\mathrm{\Omega }_z,\mathrm{Re}\mathrm{\Omega }_z,\mathrm{Im}D_1\mathrm{\Omega }_z,\mathrm{},\mathrm{Im}D_{h^{2,1}}\mathrm{\Omega }_z,\mathrm{Im}\mathrm{\Omega }_z\}$$
(69)
of $`H^3(X,).`$ We then have
$$F(z,\tau ,H)=J_\tau H,$$
(70)
where $`J_\tau `$ is given by the block matrix
$$J_\tau =\left(\begin{array}{cc}\mathrm{Re}\tau I_m& \mathrm{Im}\tau I_m\\ & \\ \mathrm{Im}\tau I_m& \mathrm{Re}\tau I_m\end{array}\right)(m=h^{2,1}+1),$$
(71)
with respect to the basis (69).
This yields the following proposition:
###### Proposition 3.2.
The mapping $`(z,\tau ,H)(z,\tau ,\iota _{z,\tau }(H))`$ gives an isomorphism
$`๐\times H^3(X,)`$.
An important consequence is:
###### Proposition 3.3.
For any open subset $`U๐`$, the cone $`_{(z,\tau )U}H_{(z,\tau )}^3(X)\{0\}`$ is open in $`H^3(X,)\{0\}`$.
###### Proof.
We must show that
$$\pi \left[\{U\times H^3(X,)\}\right]\{0\}$$
is open. By Proposition 3.2, it suffices to show that the image of the map
$$\iota :U\times [H^3(X,)\{0\}]H^3(X,),\iota (z,\tau ,H)=\iota _{z,\tau }(H)=F(z,\tau ,H)+iH,$$
is open. We fix $`(z_0,\tau _0,H_0)`$ and consider the derivative $`D\iota |_{z_0,\tau _0,H_0}`$ on $`T_{z_0,\tau _0}๐\times H^3(X,)`$. since the linear map $`\iota _{z,\tau }`$ is bijective, if we vary $`H`$, we get all of $`H_{z,\tau }^3`$, so the issue is to prove that we obtain the complementary space by taking variations in $`\tau ,z`$.
First, $`H_{z,\tau }^3=I_\tau ^1(H_z^{2,1}H_z^{0,3}).`$ The $`z`$ variations of $`H_z^{2,1}H_z^{0,3}`$ span this space plus $`H_z^{1,2}`$. By (69)โ(71), variations in $`\mathrm{Re}\tau `$, resp. $`\mathrm{Im}\tau `$, produce $`\mathrm{Re}\mathrm{\Omega }_z,\mathrm{Im}\mathrm{\Omega }_z`$ and hence $`H_z^{3,0}=`$span$`(\mathrm{\Omega }_z)`$ is also in the image.
Remark: We could also ask what kind of set is swept out in $`_{zU}H_z^{2,1}H_z^{0,3}`$ as $`z`$ ranges over an open set $`U`$. Since $`dim_{}U=h^{2,1}`$, the image of this map is a real codimension two submanifold.
### 3.5. Inner product on $`๐ฎ_Z`$
In Theorem 1.4, we have expressed $`๐ฉ_\psi (L)`$ in terms of a Gaussian type ensemble of holomorphic sections in $`๐ฎ_Z`$. We now specify the inner product, Gaussian measure and Szegรถ kernel on this space.
###### Proposition 3.4.
The Hodge-Riemann Hermitian inner product on $`H^3(X,)`$ restricts for each $`Z=(z,\tau )`$ to define a complex valued inner product on $`H_Z^3`$ which satisfies $`Q_Z[G]>0`$ for all $`G0`$. Moreover, the map $`I_\tau :H_Z^3H_z^{2,1}H_z^{0,3}`$ satisfies $`Q[I_\tau G]=\mathrm{Im}\tau Q[G].`$
###### Proof.
It follows by Proposition 2.1 that the symmetric bilinear form
$$Q[F+\tau H]=i^3_X(F+\tau H)\overline{(F+\tau H)}=\mathrm{Im}\tau Q[F+iH]$$
(72)
on $`H_{z,\tau }^3(X,)`$ in (67) is positive definite.โ
Recall that we have the real-linear isomorphisms
$`H^3(X,)`$ $`\stackrel{๐ฒ}{}๐ฎH^0(๐,)`$
$`I_\tau `$ $`.`$
$`H^3(X,)`$
where $`I_\tau (F+iH)=F+\tau H`$. Restricting (3.5) to fluxes with a critical point at $`Z=(z,\tau )`$, we have isomorphisms
$`H_Z^3`$ $`\stackrel{๐ฒ}{}๐ฎ_Z`$
$`I_\tau `$ $`.`$
$`H_z^{2,1}H_z^{0,3}`$
We let $`\stackrel{~}{Q}`$ denote the Hermitian inner product on $`H_z^{2,1}H_z^{0,3}`$ transported from $`(H_Z^3,Q)`$ by $`I_\tau `$; i.e.,
$$\stackrel{~}{Q}[C]=Q\left[I_\tau ^1(C)\right],CH_z^{2,1}H_z^{0,3}.$$
(75)
Hence by (72), we have:
$$Q[C]=(\mathrm{Im}\tau )\stackrel{~}{Q}[C].$$
(76)
## 4. Counting critical points: proof of Proposition 1.3
We now prove the first result on counting critical points of flux superpotentials $`W_G`$ where $`G`$ satisfies the tadpole constraint (1). Before starting the proof, we review the geometry of the lattice point problem and the critical point problem.
We wish to count vacua in a region of moduli space as $`G`$ varies over fluxes satisfying the tadpole constraint. Equivalently, we count inequivalent vacua in Teichmรผller space. That is, $`\mathrm{\Gamma }`$ acts on the pairs $`(W,Z)`$ of superpotentials and moduli by
$$\gamma (G,Z)=(\phi (\gamma )G,\gamma Z),$$
Therefore $`\mathrm{\Gamma }`$ acts on the incidence relation (14). We only wish to count critical points modulo the action of $`\mathrm{\Gamma }`$. To do this, there are two choices: we could break the symmetry by fixing a fundamental domain $`๐_\mathrm{\Gamma }๐`$ for $`\mathrm{\Gamma }`$ in $`๐`$, i.e. only count critical points in a fundamental domain. Or we could fix a fundamental domain for $`\phi (\mathrm{\Gamma })`$ in $`H^3(X,)`$ and count all critical points of these special flux superpotentials. When we do not know the group $`\phi (\mathrm{\Gamma })`$ precisely, it seems simpler to take the first option and that is what we do in Proposition 1.3 and Theorem 1.4. We note that the number of critical points of $`W_G`$ in Teichmรผller space equals the number of critical points of the $`\mathrm{\Gamma }`$-orbit of $`W_G`$ in $`๐`$.
The level sets $`Q[G]=C`$ for $`C>0`$ are hyperboloids contained in $`\{G:Q[G]>0\}`$ and thus the tadpole constraint defines a hyperbolic shell in $`\{G:Q[G]>0\}.`$ The critical point equation $`W_G(Z)=0`$ is homogeneous of degree $`1`$ in $`G`$. Hence, summing a homogeneous function over $`G\{G:Q[G]>0\}`$ with $`Q[G]L`$ may be viewed as summing a function on the hyperboloid $`Q[G]=1`$ over the radial projections of the lattice points $`G`$ in the shell $`Q[G]L.`$ The number which project over a compact subset of $`Q[G]=1`$ is finite.
### 4.1. Approximating the sum by an integral
Our main argument in the proof of Proposition 1.3 is the following lemma:
###### Lemma 4.1.
Let $`\psi =\chi _K`$ where $`K`$ is as in Proposition 1.3. Then
$$๐ฉ_\psi (L)=L^{b_3}\left[_๐ฎC_W,\psi \chi _Q(W)๐W+O\left(L^{1/2}\right)\right].$$
###### Proof.
We consider the integer-valued function
$$f(W)=C_W,\psi =\underset{\{Z:W(Z)=0\}}{}\psi (Z,W)=\mathrm{\#}\{Z๐:(Z,W)K\}.$$
We note that the characteristic function of the set $`\{0Q[W]L\}`$ is $`\chi _Q(W/\sqrt{L})`$. Using our symplectic basis to identify $`H^3(X,\sqrt{1})`$ with $`^{2b_3}`$, we have
$$๐ฉ_\psi (L)=\underset{N^{2b_3}}{}f(N)\chi _Q(N/\sqrt{L})=\underset{N^{2b_3}}{}f(N/\sqrt{L})\chi _Q(N/\sqrt{L})=\underset{N^{2b_3}}{}g(N/\sqrt{L}),$$
where
$$g=f\chi _Q.$$
We note that $`f`$ is constant on each connected component of $`๐ฎ[๐\pi (K)]`$. Since the number of these components is finite, $`f`$ is bounded. We let $`S(๐ฎ_Z)=\{N๐ฎ_Z:N=1\}`$, where $`N`$ denotes the norm in $`^{2b_3}`$. Since $`Q_Z`$ is positive definite, the sphere $`S(๐ฎ_Z)`$ is contained in the interior of the cone $`\{W๐ฎ:Q[W]0\}`$. Let
$$A_\psi =\underset{Z\rho (\mathrm{Supp}\psi )}{sup}Q_Z^1<+\mathrm{}.$$
(77)
Then
$$inf\{Q[W]:W\underset{Z\rho (\mathrm{Supp}\psi )}{}S(๐ฎ_Z)\}=1/A_\psi >0.$$
(78)
Now let
$$Q_0:=\{W:Q[W]1,|W|A_\psi \}\mathrm{Supp}g.$$
(79)
Approximating sums by integrals, we have
$$L^{b_3}๐ฉ_\psi (L)=L^{b_3}\underset{N^{2b_3}}{}g(N/\sqrt{L})=_{^{2b_3}}g(W)๐W+\underset{N^{2b_3}}{}E_{N,L},$$
where
$`E_{N,L}`$ $`=`$ $`{\displaystyle _{_{N,L}}}[g(N/\sqrt{L})g(W)]๐W,`$
$`_{N,L}=\{W=(W_1,\mathrm{},W_{2b_3})^{2b_3}:N_j<W_j<N_j+1/\sqrt{L}\}.`$
Let
$$B=Q_0[Q๐\pi (K)].$$
Since $`g`$ is locally constant on $`๐ฎB`$, the error $`E_{N,L}`$ vanishes whenever $`_{N,L}B=\mathrm{}`$. Hence
$$\underset{N^{2b_3}}{}E_{N,L}(supf)L^{b_3}\left[\mathrm{\#}\{N:_{N,L}B\mathrm{}\}\right]=L^{b_3}O\left(\sqrt{L}^{2b_31}\right)=O(L^{1/2}).$$
#### 4.1.1. The index density
By applying precisely the same argument for $`nd_\psi (L),`$ we obtain
###### Lemma 4.2.
Let $`\psi =\chi _K`$ where $`K`$ is as in Proposition 1.3. Then
$$nd_\psi (L)=L^{b_3}\left[_{\{Q[W]1\}}Ind_W,\psi ๐W+O\left(L^{1/2}\right)\right].$$
#### 4.1.2. Non-clustering of critical points
Before concluding the proof of Proposition 1.3, we briefly consider the question of whether there exist real hypersurfaces $`\mathrm{\Gamma }๐`$ with the property that $`\sqrt{L}^{2b_31}`$ critical points of norm $`L`$ cluster within a $`1/L`$ tube around $`\mathrm{\Gamma }`$. A domain in $`๐`$ whose boundary contained a piece of $`\mathrm{\Gamma }`$ would attain the remainder estimate in Proposition 1.3.
Since the number of critical points corresponding to $`GH^3(X,\sqrt{1})`$ is bounded, such clustering of critical points could only occur if a sublattice of rank $`2b_31`$ clustered around the hypersurface
$$\underset{(z,\tau )\mathrm{\Gamma }}{}H_{z,\tau }^3H^3(X,).$$
(80)
There do exist real hypersurfaces in $`H^3(X,)`$ for which such exceptional clustering occurs, namely hyperplanes containing a sublattice of rank $`2b_31`$. We refer to such a hyperplane as a rational hyperplane $`L`$. For instance, any pair of integral cycles $`\gamma _1,\gamma _2`$ defines a rational hyperplane
$$L=L_{\gamma _1,\gamma _2}=\{G=F+iHH^3(X,):\mathrm{}(F+iG):=_{\gamma _1}F+_{\gamma _2}H=0\}.$$
As mentioned in the introduction, projections of the lattice points $`H^3(X,\sqrt{1})`$ to $`Q`$ concentrate to sub-leading order $`\sqrt{L}^{2b_31}`$ around the hypersurface of $`Q`$ obtained by intersecting it with a rational hyperplane.
However, rational hyperplanes never have the form (80). Indeed, under the correspondence $`\rho \pi ^{}`$ defined by the diagram (68), the image of a hyperplane always covers a region and not a hypersurface of $`๐`$. That is,
$$dim(LH_{z,\tau }^3)>1(z,\tau )๐.$$
Indeed, under the identification $`H_{z,\tau }^3H^3(X,),`$ $`L|_{H_{z,\tau }^3}`$ becomes the real linear functional $`L(H)=_{\gamma _1}F(z,\tau ,H)+_{\gamma _2}H`$ on $`H^3(X,)`$. Here, we use that $`F(z,\tau ,H)`$ is linear in $`H`$. Hence, $`dimLH_{z,\tau }^3b_31`$ for any $`(z,\tau )`$.
As will be studied in \[Ze2\], clustering to order $`\sqrt{L}^{2b_31}`$ can only occur if the second fundamental form of (80) is completely degenerate. Hence the fact that rational hyperplanes never have this form is strong evidence that there are no smooth hypersurfaces $`\mathrm{\Gamma }๐`$ for which lattice points cluster to subleading order around (80).
### 4.2. Hessians and density of critical points
The final step in the proof of Proposition 1.3 is to change the order of integration over $`๐`$ and over $`๐ฎ_Z`$:
###### Lemma 4.3.
We have:
$$_{\{Q[W]1\}}C_W,\psi ๐W=_๐_{๐ฎ_Z}\psi (Z,W)|detH^cW(Z)|\chi _{Q_Z}(W)๐Wd\mathrm{Vol}_{WP}(Z).$$
Combining the formulas in Lemmas 4.1 and 4.3, we obtain the formula of Proposition 1.3.
The proof of Lemma 4.3 is in two parts. The first is an elementary exercise in changing variables in an integral, which we accomplish below by relating both sides to pushforwards from the incidence variety in the diagram (15). The second part involves special geometry, and is given in the next section.
We may interpret the integral
$$_{\{Q[W]1\}}C_W,\psi ๐W$$
as an integral over $``$ as follows. Implicitly, it defines a measure $`d\mu _{}`$ so that
$$_{}\psi (Z,W)๐\mu _{}=_{\{Q[W]1\}}C_W,\psi ๐W.$$
(81)
The measure $`d\mu _{}`$ may be expressed in terms of the Leray measure $`d_{}`$ defined by a measure $`d\nu `$ on $`๐ฎ`$ and the โevaluation mapโ
$$\epsilon :(Z,W)๐\times ๐ฎW(Z).$$
The Leray form is the quotient $`d_{}:=\frac{dV_{WP}\times d\nu }{d\epsilon }`$, i.e. the unique form satisfying
$$d_{}\times d\epsilon =dV_{WP}\times d\nu .$$
This measure is often written $`\delta (W(Z))dWdV(Z)`$ in the physics literature.
As suggested by the physics formula, $`d\mu _{}=s(Z)^{}\delta _0.`$ However, this formula is somewhat ambiguous. If we regard $`s`$ as fixed, then it is simply the pullback of $`\delta _0`$ under $`Zs(Z).`$ It is then well-known that
$$s^{}\delta _0=\underset{Z:s(Z)=0}{}\frac{\delta _Z}{|detH^cs(Z)|}.$$
(82)
However, when interchanging the order of integration, we really wish to think of it as a function of $`s`$ for fixed $`Z`$. So we now have a function $`\epsilon _Z(s)=s(Z)`$ which may be viewed as
$$\epsilon _Z:๐ฎ^m^{b_3},$$
where $`m=h^{2,1}+1=\frac{1}{2}b_3`$. So now the zero set $`\epsilon _Z^1(0)`$ is the subspace $`๐ฎ_Z`$ rather than the discrete set $`Crit(s)`$.
To simplify the notation, we now consider the general situation where we have a real $`n`$-dimensional manifold $`M`$ and a space $`๐ฎ`$ of functions $`F:M_n^n`$. In our case, $`F=s`$ and $`M`$ is a coordinate neighborhood in $`๐`$ where $`M`$ has local coordinates $`(x_1,\mathrm{},x_n)`$ and $``$ has a local frame. Suppose that $`0`$ is a regular value of $`F`$, so that $`F`$ is a local diffeomorphism in a neighborhood $`U`$ of any point $`x_0`$ of $`F^1(0)`$. Let $`h=F_{|U}^1`$ in a neighborhood of $`0`$. Then for $`\phi `$ supported in a neighborhood of $`x_0`$, put
$$F^{}\delta _0,\phi =\delta _0,\phi (h(y))|detdh(y)|.$$
Let $`dim_{}๐ฎ=dn`$. In our case, $`d=2b_3>n=b_3`$, so we introduce a supplementary linear map: for a point $`uUM`$, $`๐ฎ_u`$ is the kernel of $`\epsilon _u`$, and we supplement $`\epsilon _u`$ with the projection $`\mathrm{\Pi }_u:๐ฎ๐ฎ_u.`$ Then,
$$(\epsilon _u,\mathrm{\Pi }_u):๐ฎ^n๐ฎ_u$$
is a linear isomorphism. Hence it equals its derivative, so
$$\begin{array}{ccc}\epsilon _u^{}\delta _0,\phi \hfill & =\hfill & \delta _0,\phi ((\epsilon _u,\mathrm{\Pi }_u)^1)|det(\epsilon _u,\mathrm{\Pi }_u)^1|.\hfill \end{array}$$
Now, $`๐ฎ`$ is equipped with an inner product, which induces an inner product on $`^n๐ฎ_u`$. We choose an orthonormal basis $`\{S_1,\mathrm{},S_n\}`$ of $`๐ฎ_u^{},`$ and $`\{S_{n+1},\mathrm{},S_d\}`$ for $`๐ฎ_u`$. Since $`\mathrm{\Pi }_u:๐ฎ_u๐ฎ_u`$ is the identity, $`(\epsilon _u,\mathrm{\Pi }_u)`$ has a block diagonal matrix relative to the bases of $`๐ฎ=๐ฎ_u^{}๐ฎ_u`$ and $`^n๐ฎ_u`$, with the identity in the $`๐ฎ_u`$-$`๐ฎ_u`$ block. Hence, $`det(\epsilon _u,\mathrm{\Pi }_u)=det(\epsilon _u|_๐ฎ^{})`$ where the determinant is with respect to these bases.
The general case of formula (81) states that
$$d\mu _{}=|detDW(u)|\times \frac{\chi _Qdu\times dW}{d\epsilon }.$$
(83)
We then compute the $``$ integral as an iterated integral using the other singular fibration $`\pi `$, i.e. by first integrating over the fibers $`๐ฎ_u`$:
$$_{}\psi (u)๐\mu _{}=_U_{๐ฎ_u}\frac{\psi (u)}{|det(\epsilon _u|_{๐ฎ_u^{}})|}\chi _{Q_u}(W)|detDW(u)|๐W๐u.$$
(84)
Returning to our case where $`F=s`$, (84) becomes
$$_{}\psi (Z)๐\mu _{}=_๐_{๐ฎ_Z}\frac{\psi (Z,W)}{|det(\epsilon _Z|_{๐ฎ_Z^{}})|}|detH^cW(Z)|\chi _{Q_Z}(W)๐Wd\mathrm{Vol}_{WP}(Z).$$
(85)
### 4.3. Completion of the proof of Lemma 4.3
To complete the proof of the lemma, we need to show that $`|det(\epsilon _Z|_{๐ฎ_Z^{}})|=1`$ with respect to normal coordinates and an adapted frame at $`Z_0=(z_0,\tau _0)M`$.
Recalling (3.5)โ(3.5), we write
$$\stackrel{~}{๐ฎ}_Z^{}=I_\tau ๐ฒ^1(๐ฎ_Z^{})=H_z^{3,0}H_z^{1,2}.$$
A complex orthonormal basis for $`\stackrel{~}{๐ฎ}_{Z_0}^{}`$ relative to $`Q`$ is $`\{\overline{\chi }_0,\overline{\chi }_1,\mathrm{},\overline{\chi }_{h^{2,1}}\}`$, where $`\chi _0=\overline{\mathrm{\Omega }}_{z_0}`$. A basis (over $``$) for $`๐ฎ_{Z_0}^{}`$ is
$$\overline{U}_j:=๐ฒI_\tau ^1(\overline{\chi }_j),\overline{V}_j:=๐ฒI_\tau ^1(\sqrt{1}\overline{\chi }_j),0jh^{2,1}.$$
The basis $`\{\overline{U}_j,\overline{V}_j\}`$ is orthogonal with respect to $`Q_{Z_0}`$, but not orthonormal. By (76)
$$Q[\overline{U}_j]=\stackrel{~}{Q}[\overline{\chi }_j]=\frac{1}{\mathrm{Im}\tau }Q[\overline{\chi }_j]=\frac{1}{\mathrm{Im}\tau },Q[\overline{V}_j]=\stackrel{~}{Q}\left[\sqrt{1}\overline{\chi }_j\right]=\frac{1}{\mathrm{Im}\tau }.$$
(86)
To compute $`det(\epsilon _{Z_0}|๐ฎ_{Z_0}^{})`$, we let $`(z_1,\mathrm{},z_{h^{2,1}})`$ be normal coordinates about $`z_0`$, and we let $`_jf`$ be given by
$$_{/z^j}(fe_{})=(_jf)e_{},$$
for $`1jh^{2,1}`$. We find it convenient to use the coordinate $`\tau `$, although it is not normal, and we use the normalized covariant derivative
$$_0:=(\mathrm{Im}\tau )_\tau .$$
(87)
Now we write
$$\overline{U}_j=f_j(z)\mathrm{\Omega }_z^{}\omega _\tau ^{},\overline{V}_j=g_j(z)\mathrm{\Omega }_z^{}\omega _\tau ^{},$$
where the local frame $`\mathrm{\Omega }_z`$ is normal at $`z_0`$, and $`\omega _\tau =dx+\tau dy`$. Note that the Weil-Petersson norm $`|\omega _\tau ^{}|`$ is given by
$$|\omega _\tau ^{}|=|dx+\tau dy|^1=\frac{1}{(\mathrm{Im}\tau )^{1/2}}.$$
(88)
Taking into account (86)โ(88), the $`\mathrm{Im}\tau `$ factors cancel out, and we obtain
$$det(\epsilon _{Z_0}|๐ฎ_{Z_0}^{}))=det\left(\begin{array}{cc}\mathrm{Re}_jf_k& \mathrm{Re}_jg_k\\ \\ \mathrm{Im}_jf_k& \mathrm{Im}_jg_k\end{array}\right)|_{Z_0},\text{for }0j,kh^{2,1}.$$
We now evaluate the entries of the matrix. By Proposition 2.2, we have
$$_kf_j(Z)=_X\overline{๐_j\mathrm{\Omega }_{z_0}}๐_k\mathrm{\Omega }_z,_kg_j(Z)=_Xi\overline{๐_j\mathrm{\Omega }_{z_0}}๐_k\mathrm{\Omega }_z,$$
and hence
$$_jf_k(Z_0)=i\delta _{jk},_jg_k(Z_0)=\delta _{jk},\text{for }j,k1.$$
Also
$$_kf_0=_X\mathrm{\Omega }_{z_0}[๐_k\mathrm{\Omega }_{z_0}(K/z_j)\mathrm{\Omega }_{z_0}]=0,_kg_0=i_kf_0=0\text{for }k1.$$
By (64), we have
$$_0(f_j)=(\mathrm{Im}\tau )_\tau (f_j)=_X๐_j\mathrm{\Omega }_{z_0}\mathrm{\Omega }_{z_0}=0,_0(g_j)=i_X\mathrm{\Omega }_{z_0}\mathrm{\Omega }_{z_0}=0,j1,$$
and
$$_0(f_0)=_X\overline{\mathrm{\Omega }_{z_0}}\mathrm{\Omega }_{z_0}=i,_0(g_0)=_X\overline{i\mathrm{\Omega }_{z_0}}\mathrm{\Omega }_{z_0}=1.$$
Therefore,
$$|det(\epsilon _{Z_0}|๐ฎ_{Z_0}^{}))|=|det\left(\begin{array}{cc}0& I\\ \\ D(1,1,\mathrm{},1)& 0\end{array}\right)|=1.$$
## 5. Proof of Theorem 1.4
In this section we prove Theorem 1.4, which is a combination of an equidistribution theorem for radial projections of lattice points and an equidistribution theorem for critical points.
### 5.1. A local van der Corput Theorem
We first illustrate the method of proof of Theorem 1.4 by providing a van der Corput type asymptotic estimate for the radial distribution of lattice points (Theorem 5.1). The estimate has much in common with the classical van der Corput estimate of Hlawka, Randol and others on lattice points in dilates of smooth convex sets (see for example, \[Ra, Hl\]), and we adapt the proof of the classical estimate to obtain our asymptotic equidistribution theorem.
Let $`Q^n`$ ($`n2)`$ be a bounded, smooth, strictly convex set with $`0Q^{}`$. Let $`|X|_Q`$ denote the norm of $`X^n`$ given by
$$Q=\{X^n:|X|_Q<1\}.$$
(89)
To measure the equidistribution of projections of lattice points, we consider the sums
$$S_f(t)=\underset{k^ntQ\{0\}}{}f\left(\frac{k}{|k|_Q}\right),\text{with }fC^{\mathrm{}}(Q),t>0.$$
We extend $`f`$ to $`^n`$ as a homogeneous function of degree $`0`$, so that $`f(k)=f\left(\frac{k}{|k|_Q}\right)`$. Our purpose is to obtain the following asymptotics of $`S_f(t)`$:
###### Theorem 5.1.
$$S_f(t)=t^n_Qf๐X+O(t^{n\frac{2n}{n+1}}),t\mathrm{}.$$
From this it is simple to obtain asymptotics of $`S_f(t)`$ when $`fC^{\mathrm{}}(Q)`$ is extended as a homogeneous function of any degree $`\alpha `$ to $`^n`$:
###### Corollary 5.2.
Let $`f๐^{\mathrm{}}(^n\{0\})`$ be homogeneous of degree $`\alpha >0`$, and let
$$S_f(t)=\underset{k^ntQ}{}f(k),t>0$$
Then
$$S_f(t)=t^{n+\alpha }_Qf๐X+O(t^{n\frac{2n}{n+1}+\alpha }),t\mathrm{}.$$
#### 5.1.1. Littlewood-Paley
To deal with the singularity of $`f`$ at $`x=0`$ we use a dyadic Littlewood-Paley decomposition in the radial direction. Let $`\eta C_0^{\mathrm{}}`$ with $`\eta (r)=1`$ for $`r1`$ and with $`\eta (r)=0`$ for $`r2.`$ We then define
$$\rho C_0^{\mathrm{}}(),\rho (r)=\eta (r)\eta (2r).$$
Then $`\rho (r)`$ is supported in the shell $`1/2r2`$, hence $`\rho (2^jr)`$ is supported in the shell $`2^{j1}r2^{j+1}`$. We then have:
$$\eta (r)=\underset{j=0}{\overset{\mathrm{}}{}}\rho (2^jr),(r0).$$
Indeed,
$$\underset{j=0}{\overset{J}{}}\rho (2^jr)=\eta (r)\eta (2^Jr)\eta (r)$$
by the assumption that $`\eta C_0^{\mathrm{}}`$.
We then write
$`S_f(t)`$ $`=`$ $`{\displaystyle \underset{k^n}{}}f(k)\chi _{[0,1]}({\displaystyle \frac{|k|_Q}{t}})=S_f^{}(t)+S_f^{\prime \prime }(t),`$ (91)
$`S_f^{}(t)={\displaystyle \underset{k^n}{}}f(k)\eta ({\displaystyle \frac{|k|_Q}{t}}),`$
$`S_f^{\prime \prime }(t)={\displaystyle \underset{k^n}{}}f(k)(\chi _{[0,1]}\eta )({\displaystyle \frac{|k|_Q}{t}})).`$
We can assume without loss of generality that $`f0`$. We begin with the first sum in $`S_f^{}(t)`$:
###### Lemma 5.3.
$$S_f^{}(t)=t^n_^nf(X)\eta (|X|_Q)๐X+O(\mathrm{log}t).$$
###### Proof.
We write the sum as
$$S_f^{}(t)=\underset{j=0}{\overset{\mathrm{}}{}}\underset{k^n}{}f(k)\rho (\frac{2^j|k|_Q}{t}).$$
We further break up the dyadic sum into $`_{j=0}^{J(t)}+_{j=J(t)+1}^{\mathrm{}}`$ with $`J(t)`$ to be determined later. We first consider
$$S_1^{}:=\underset{j=0}{\overset{J(t)}{}}\underset{k^n}{}f(k)\rho (\frac{2^j|k|_Q}{t}).$$
The function $`f(X)\rho (2^j|X|_Q)C_0^{\mathrm{}}(^n)`$ when $`f`$ is homogeneous of degree $`0`$ and smooth on $`Q`$. Hence we may apply the Poisson summation formula to the $`k`$ sum to obtain
$$S_1^{}=\underset{j=0}{\overset{J(t)}{}}\underset{N^n}{}_^ne^{iX,N}f(X)\rho (\frac{2^j|X|_Q}{t})๐X.$$
The terms with $`N=0`$ sum up to
$`t^n{\displaystyle _^n}f(X)\left\{{\displaystyle \underset{j=0}{\overset{J(t)}{}}}\rho (2^j|X|_Q)\right\}๐X`$ $`=`$ $`t^n{\displaystyle _^n}f(X)\{\eta (|X|_Q)\eta (2^{J(t)+1}|X|_Q))\}dX`$
$`=`$ $`t^n{\displaystyle _^n}f(X)\eta (|X|_Q)๐X+O(t^n2^{nJ(t)}),`$
where the last estimate is a consequence of the fact that $`\eta (2^{J(t)+1}|X|_Q)`$ is supported on $`2^{J(t)}Q`$.
To estimate the remaining terms in the sum $`S_1^{}`$, we make the change of variables $`Y=2^jX/t`$ in the integral to obtain
$$2^{nj}t^n_^nf(Y)\rho (|Y|)e^{i2^jtY,N}๐Y.$$
Since the integrand is smooth, this term has the upper bound
$$c\mathrm{\hspace{0.17em}2}^{nj}t^n(1+2^j|N|t)^K,K>0.$$
(Again, we let $`c`$ denote a constant; $`c`$ depends on $`f`$ and $`K`$, but is independent of $`j,t,N`$.) The sum over $`N0`$ is then bounded by
$`ct^n{\displaystyle \underset{jJ(t)}{}}2^{nj}{\displaystyle \underset{N0}{}}(1+2^j|N|t)^K`$ $``$ $`t^n{\displaystyle \underset{jJ(t)}{}}2^{nj}{\displaystyle _0^{\mathrm{}}}(1+2^jrt)^Kr^{n1}๐r`$
$`=`$ $`{\displaystyle \underset{jJ(t)}{}}{\displaystyle _0^{\mathrm{}}}(1+s)^Ks^{n1}๐s=cJ(t).`$
Therefore
$$S_1^{}=t^n_^nf(X)\eta (|X|_Q)๐X+O(t^n2^{nJ(t)})+O(J(t)).$$
Recall that $`S_f^{}(t)=S_1^{}+S_2^{}`$, where
$$S_2^{}=\underset{j=J(t)+1}{\overset{\mathrm{}}{}}\underset{k^n}{}f(\frac{k}{|k|_Q})\rho (\frac{2^j|k|_Q}{t}).$$
Since
$$\underset{j=J(t)+1}{\overset{\mathrm{}}{}}\rho (\frac{2^j|k|_Q}{t})=\eta \left(\frac{2^{J(t)}|k|_Q}{t}\right)\chi _{t2^{J(t)}Q},$$
the remainder $`S_2^{}`$ is bounded by the total number of lattice points in the shell $`|k|_Q2^{J(t)}t`$, hence is of order $`t^n2^{nJ(t)}`$. It follows that
$$S_f^{}(t)=t^n_^nf(X)\eta (|X|_Q)๐X+O(t^n2^{nJ(t)})+O(J(t)).$$
(92)
To balance the terms, we choose $`J(t)=\mathrm{log}_2t`$, and then the last two terms of (92) have the form
$$O(t^nt^n)+O(\mathrm{log}t)=O(\mathrm{log}t).$$
#### 5.1.2. Stationary phase.
Theorem 5.1 is an immediate consequence of Lemma 5.3 and the following assymptotics of the second sum $`S_f^{\prime \prime }(t)`$ from (91):
###### Lemma 5.4.
$$S_f^{\prime \prime }(t)=t^n_^nf(X)(\chi _{[0,1]}\eta )(|X|_Q)๐X+O(t^{n\frac{2n}{n+1}}).$$
###### Proof.
Let
$$g(X)=f(X)(\chi _{[0,1]}\eta )(|X|_Q)$$
and mollify $`g`$ by a radial approximate identify $`\phi _\epsilon `$ to obtain a smooth approximation $`g_\epsilon =g\phi _\epsilon `$. We claim that
$$S_f^{\prime \prime }(t)=\underset{k^n}{}g\left(\frac{k}{t}\right)=\underset{k^n}{}g_\epsilon \left(\frac{k}{t}\right)+O(\epsilon t^n).$$
(93)
To see this, we break the sum into two parts. The first part is over the lattice points $`k`$ with $`k/t`$ in an $`\epsilon `$ tube $`T_\epsilon `$ about $`\{|X|_Q=1\}`$. The number of such $`k`$ is $`O(\epsilon t^n)`$, so this part contributes the stated error. For the remaining sum, the error is
$$\left|\underset{k^ntT_\epsilon }{}\left[g\left(\frac{k}{t}\right)g_\epsilon \left(\frac{k}{t}\right)\right]\right|\underset{k/t\mathrm{Supp}gT_\epsilon }{}\epsilon \underset{|X|_Q>1}{sup}|dg(X)|=O(\epsilon t^n),$$
which verifies (93).
The Poisson summation formula then gives
$$\underset{k^n}{}g_\epsilon (k/t)=t^n\underset{N^n}{}\widehat{g}_\epsilon (2\pi tN)=t^n\widehat{g}(2\pi tN)\widehat{\phi }(2\pi t\epsilon N).$$
The term $`N=0`$ yields
$$t^N_^ng_\epsilon (X)๐X=t^n_^nf(X)(\chi _{[0,1]}\eta )(|X|_Q)๐X+O(\epsilon t^n),$$
where the last inequality is by breaking up the integral into two parts as above.
As for the remainder terms $`N0`$, we now show that
$$\widehat{g}(2\pi tN)=O\left((1+|tN|)^{\frac{(n+1)}{2}}\right).$$
(94)
To verify (94), we write
$`g`$ $`=`$ $`f\rho h=(f\rho )(h\eta _2),\text{with}\eta _2(X)=\eta (\frac{1}{2}|X|_Q),h=\theta \lambda ,`$
$`\lambda (X)=|X|_Q1,\theta (t)=\text{Heaviside function}=\{\begin{array}{cc}0,\hfill & \text{if }t<0,\hfill \\ 1,\hfill & \text{if }t0.\hfill \end{array}`$
Since $`\widehat{g}=\widehat{f\rho }\widehat{h\eta _2}`$ and $`\widehat{f\rho }`$ is rapidly decaying, it suffices to show that $`\widehat{h\eta _2}`$ satisfies (94). (Here, we use the elementary estimate $`\alpha \beta _{(K)}c\alpha _{(K+n+1)}\beta _{(K)}`$, where $`\alpha _{(K)}=sup_{x^n}(1+|x|)^K|\alpha (x)|`$.) Taking partial derivatives,
$$๐_j(h\eta _2)=๐_j\eta _2+(\delta _0\lambda )๐_j\lambda .$$
Since the latter term is given by integration over $`Q`$, which is strictly convex, the standard stationary phase method (see Hรถrmander \[Ho\]) immediately gives $`(\delta _0\lambda )\widehat{}(x)=O(x^{\frac{(n1)}{2}})`$, and hence
$$x_j\widehat{h\eta _2}=[๐_j(h\eta _2)]\widehat{}=O\left(x^{\frac{(n1)}{2}}\right),$$
which implies (94).
Hence the remainder is bounded above by
$$ct^n\underset{N0}{}(1+|tN|)^{\frac{(n+1)}{2}}(1+|\epsilon tN|)^K.$$
(96)
The sum (96) can be replaced by the integral
$`ct^n{\displaystyle _^n}(1+|tN|)^{\frac{(n+1)}{2}}(1+|\epsilon tN|)^K๐N`$ $`=`$ $`ct^n{\displaystyle _0^{\mathrm{}}}(1+tr)^{\frac{(n+1)}{2}}(1+\epsilon tr)^Kr^{n1}๐r`$
$`=`$ $`c\epsilon ^{\frac{1n}{2}}{\displaystyle _0^{\mathrm{}}}(\epsilon +s)^{\frac{(n+1)}{2}}(1+s)^Ks^{n1}๐s`$
$``$ $`c\epsilon ^{\frac{1n}{2}}{\displaystyle _0^{\mathrm{}}}(1+s)^Ks^{\frac{n3}{2}}๐s=c\epsilon ^{\frac{1n}{2}}.`$
Hence
$$S_f^{\prime \prime }(t)=t^n_^nf(X)(\chi _{[0,1]}\eta )(|X|_Q)๐X+O(\epsilon t^n)+O(\epsilon ^{(n1)/2}).$$
To optimize, we choose $`\epsilon `$ so that $`\epsilon t^n=\epsilon ^{(n1)/2}`$, i.e. $`\epsilon =t^{2n/(n+1)}`$, which gives the result. (To be precise, it is the sum of the terms in (96) with $`|N|\sqrt{n}`$ that is bounded by the above integral. But there are only finitely many terms with $`|N|<\sqrt{n}`$, and each of these terms is $`<ct^{n\frac{n+1}{2}}`$, which is better than $`O(t^{n\frac{2n}{n+1}})`$ when $`n2`$.)โ
#### 5.1.3. Van der Corput for homogeneous weights $`f`$. Proof of Corollary 5.2:
This time, we have
$$S_f(t)=\underset{k^ntQ\{0\}}{}|k|_Q^\alpha f\left(\frac{k}{|k|_Q}\right).$$
The set of norms of lattice points $`\{t_j^+:k^n|k|_Q=t_j\}`$ is a countable set without accumulation point. We order the $`t_j`$ so that $`t_jt_{j+1}`$. We then define the monotone increasing step function on $``$
$$\mu (T)=\underset{j:t_jT}{}\left\{\underset{k:|k|_Q=t_j}{}f\left(\frac{k}{|k|_Q}\right)\right\}.$$
It is clear that
$$\mu (T)=S_{f_0}(T),f_0(x)=\frac{f(x)}{|x|_Q}.$$
Hence, by Theorem 5.1,
$$S_{f_0}(t)=t^n_Qf_0๐X+O(t^{n\frac{2n}{n+1}}),t\mathrm{}.$$
(97)
We further have
$$S_f(T)=_0^Tt^\alpha ๐\mu (t).$$
(98)
Indeed,
$$d\mu (t)=\underset{j}{}\left\{\underset{k:|k|_Q=t_j}{}f\left(\frac{k}{|k|_Q}\right)\right\}\delta (t_j),$$
and
$$_0^Tt^\alpha ๐\mu (t)=\underset{j:t_jT}{}\left\{\underset{k:|k|_Q=t_j}{}f\left(\frac{k}{|k|_Q}\right)\right\}t_j^\alpha =S_f(T).$$
Integrating (98) by parts and applying (97), we get
$`S_f(T)`$ $`=`$ $`T^\alpha \mu (T)\alpha {\displaystyle _0^T}t^{\alpha 1}\mu (t)๐t`$
$`=`$ $`T^\alpha \left[T^n{\displaystyle _Q}f_0๐X+O(T^{n\frac{2n}{n+1}})\right]\alpha {\displaystyle _0^T}t^{\alpha 1}\left[t^n{\displaystyle _Q}f_0๐X+O(t^{n\frac{2n}{n+1}})\right]๐t`$
$`=`$ $`T^{n+\alpha }\left[{\displaystyle _Q}f_0๐X\right]{\displaystyle \frac{n}{\alpha +n}}+O(T^{n\frac{2n}{n+1}+\alpha })`$
$`=`$ $`T^{n+\alpha }{\displaystyle _Q}f๐X+O(T^{n\frac{2n}{n+1}+\alpha }).`$
### 5.2. Van der Corput for critical points
We prove Theorem 1.4 by following the arguments of the proofs of Theorem 5.1 and Corollary 5.2 with hardly any changes. We first assume that $`\psi `$ is homogeneous of order 0 in $`๐ฎ`$. We let $`K_\psi =\rho (\mathrm{Supp}\psi )๐`$, a compact set.
To begin, we recall that if $`W`$ has critical points, then $`W`$ is in the โlight coneโ $`Q[W]>0`$. For $`W`$ in the light cone, we write
$$|W|_Q=Q[W]^{\frac{1}{2}},\text{for }Q[W]>0.$$
The main difference between this case and our previous one, is that now the set $`Q`$ given by (89), in addition to not being convex, is not compact. However, since only with those $`W`$ with critical points in the support of $`\psi `$ contribute to the sum, we consider
$$Q_\psi :=Q๐ฎ_\psi ,๐ฎ_\psi =\left(\underset{\tau K_\psi }{}๐ฎ_\tau \right),$$
which is a compact subset of $`๐ฎ`$.
We let $`f(W)=C_W,\psi `$, which is a smooth function supported in $`๐ฎ_\psi `$. Then
$$๐ฉ_\psi (L)=S_f(L)=\underset{k^n\sqrt{L}Q\{0\}}{}f(k),$$
as before. Now we follow the previous proof, with $`t=\sqrt{L}`$. Our first modification is to verify (93), we instead let $`T_\epsilon `$ be the epsilon tube over $`๐ฎ_\psi Q`$. Finally, the estimate $`(\delta _0\lambda )\widehat{}(t)=O(t^{\frac{(n1)}{2}})`$, which was based on the convexity of $`Q`$ in our previous argument, holds in this case, since the phase $`\psi (Y)=LY,N`$ has (two) non-degenerate critical points whenever $`N`$ is in the light cone. Thus we have
$$๐ฉ_\psi (L)=L^{b_3}\left[_{\{Q[W]1\}}C_W,\psi ๐W+O\left(L^{\frac{2b_3}{2b_3+1}}\right)\right].$$
The case $`\alpha =0`$ now follows from Lemma 4.3, and the general case then follows exactly as in the proof of Corollary 5.2. โ
## 6. Special geometry and density of critical points
The aim of this section is to compute the critical point density $`๐ฆ^{\mathrm{crit}}(Z)`$ and verify Corollaries 1.51.6. At the same time, we compute the index density and prove Theorem 1.8. As in \[DSZ1\], we do this by pushing forward the integrand of (25) under the Hessian map. The Hessian map turns out to be an isomorphism, hence the discussion is more elementary than in \[DSZ1\]. To make the change of variables, we first evaluate the image of the Hessian using the special geometry of Calabi-Yau moduli spaces and then check how the Hessian map distorts inner products. Our discussion gives an alternate approach to the formulas in the article \[DD1\], and connects the special critical point density formula in this article with the general ones in \[DSZ1, DSZ2\].
### 6.1. The range of the Hessian map
We now study the complex Hessian map:
$$H^c(Z):W\left(\begin{array}{cc}H^{}& x\mathrm{\Theta }(Z)\\ \overline{x}\overline{\mathrm{\Theta }}(Z)& \overline{H}^{}\end{array}\right).$$
(99)
To describe $`H^c(Z)`$ in local coordinates, we fix a point $`Z_0=(z_0,\tau _0)`$ and choose normal coordinates $`\{z^1,\mathrm{},z^{h^{2,1}}\}`$ at $`z_0`$. We let $`\mathrm{\Omega }`$ be a local normal frame for $`H^{3,0}`$ at $`z_0`$, and we let $`\omega =dx+\tau dy`$. Recall that $`\omega `$ is not a normal frame, since $`|\omega _\tau |=(\mathrm{Im}\tau )^{1/2}`$. We let $`\stackrel{~}{e}_{}=(\mathrm{Im}\tau _0)^{1/2}\mathrm{\Omega }^{}\omega ^{}`$, so that $`|\stackrel{~}{e}_{}(Z_0)|=1`$.
As in ยง3.2, the matrix $`(H_{jk})`$ of the holomorphic Hessian is given by
$$H^{}(Z_0)=\underset{j,q}{}H_{jq}^{}dz^qdz^j\stackrel{~}{e}_{}|_{Z_0},0j,qh^{2,1},$$
(100)
where
$$dz^0|_{Z_0}=\frac{1}{\mathrm{Im}\tau _0}d\tau |_{Z_0}$$
is the unit holomorphic cotangent vector (with respect to the Weil-Petersson, or hyperbolic, metric on $``$) at $`\tau _0`$.
We wish to express formulas (59)โ(60) for the complex Hessian in terms of these coordinates and frames. We write
$$(_jf)e_{}=_{/z^j}(fe_{}),1jh^{2,1},(_0f)e_{}=(\mathrm{Im}\tau _0)_{/\tau }(fe_{}).$$
($`_0`$ is the normalized covariant $`\tau `$-derivative given by (87).) The complex Hessian matrix is given by:
$$H^c(Z_0)=\left(\begin{array}{cc}H^{}(Z_0)& f(Z_0)I\\ \overline{f(Z_0)}I& \overline{H^{}(Z_0)}\end{array}\right),H^{}=\left(_j_qf\right)_{0j,qh^{2,1}}.$$
(101)
Identifying the off-diagonal components with $`f(Z_0)`$, we view the image space as a subspace of $`\mathrm{Sym}(h^{2,1}+1,)`$, so we can write the Hessian map in the form
$$H_{Z_0}:๐ฎ_Z\mathrm{Sym}(h^{2,1}+1,),W(H^{}(Z_0),f(Z_0)).$$
###### Lemma 6.1.
The range of the Hessian map $`H_{Z_0}:๐ฎ_{Z_0}\mathrm{Sym}(h^{2,1}+1,)`$ is of the form $`_{Z_0}`$, where $`_{Z_0}`$ is a real subspace of $`\mathrm{Sym}(h^{2,1}+1,)`$ of real dimension $`2h^{2,1}`$ spanned over $``$ by the matrices
$$\xi ^k=\left(\begin{array}{cc}0& e_k\\ e_k^t& ^k(z)\end{array}\right),\xi ^{h^{2,1}+k}=\left(\begin{array}{cc}0& \sqrt{1}e_k\\ \sqrt{1}e_k^t& \sqrt{1}^k(z)\end{array}\right),1kh^{2,1},$$
given by (27), where $`e_k`$ is the $`k`$-th standard basis element of $`^{h^{2,1}}`$ and $`^k(z)\mathrm{Sym}(h^{2,1},)`$ is the matrix $`\left(_{jq}^{\overline{k}}(z)\right)`$ of (46).
In other words, $`_{Z_0}`$ is the set of matrices of the form
$$\left(\begin{array}{cc}0& (\overline{v}_1,\mathrm{},\overline{v}_{h^{2,1}})\\ (\overline{v}_1,\mathrm{},\overline{v}_{h^{2,1}})^t& _{k=1}^{h^{2,1}}^k(z)v_k\end{array}\right),(v_1,\mathrm{},v_{h^{2,1}})^{h^{2,1}}.$$
(102)
We emphasize that $`_Z\mathrm{Sym}(h^{2,1}+1,)`$ is only a real and not a complex subspace. We also note that $`dim_{}_Z=2h^{2,1}`$ and hence $`dim_{}(_Z)=b_3=dim_{}๐ฎ_Z`$; i.e., $`H_Z`$ is an isomorphism.
Proof of Lemma 6.1: We shall use the notation $`1j,k,lh^{2,1}`$, $`0\alpha ,\beta ,\gamma h^{2,1}`$. By (3.5), we have the (real-linear) isomorphism
$$\stackrel{~}{๐ฒ}_{Z_0}=๐ฒI_\tau ^1:H_{z_0}^{2,1}H_{z_0}^{0,3}\stackrel{}{}๐ฎ_{Z_0}.$$
Recall that $`H_{z_0}^{2,1}H_{z_0}^{0,3}`$ has a complex orthonormal basis $`\{\chi _\alpha \}`$ of the form
$$\chi _j=๐_j\mathrm{\Omega }_{Z_0},1jh^{2,1},\chi _0=\overline{\mathrm{\Omega }}_{Z_0}.$$
By (76), a real orthonormal basis of $`๐ฎ_{Z_0}`$ is
$$U_\alpha :=(\mathrm{Im}\tau )^{1/2}\stackrel{~}{๐ฒ}_{Z_0}(\chi _\alpha ),V_\alpha :=(\mathrm{Im}\tau )^{1/2}\stackrel{~}{๐ฒ}_{Z_0}(\sqrt{1}\chi _\alpha ).$$
We write:
$$U_\alpha =f_\alpha \stackrel{~}{e}_{},V_\alpha =g_\alpha \stackrel{~}{e}_{};$$
equivalently
$$\stackrel{~}{๐ฒ}_{Z_0}(\chi _\alpha )=f_\alpha e_{},\stackrel{~}{๐ฒ}_{Z_0}(\sqrt{1}\chi _\alpha )=g_\alpha e_{}.$$
We must compute the matrices
$$H_{Z_0}^{}(f_\alpha \stackrel{~}{e}_{})=\left(_\beta _\gamma f_\alpha \right)|_{Z_0},H_{Z_0}^{}(g_\alpha \stackrel{~}{e}_{})=\left(_\beta _\gamma g_\alpha \right)|_{Z_0},$$
where $`H_{Z_0}^{}:๐ฎ_{Z_0}\mathrm{Sym}(h^{2,1}+1,)`$ is the holomorphic Hessian map.
We shall show that:
$$\{\begin{array}{cc}\hfill (\mathrm{i})& _0^2f_G(Z_0)=0,GH_{Z_0}^3(X,)(\text{where }W_G=f_Ge_{})\hfill \\ & \text{and thus }_0^2f_\alpha (Z_0)=_0^2g_\alpha (Z_0)=0,\hfill \\ & \\ \hfill (\mathrm{ii})& _j_0f_0(Z_0)=_j_0g_0(Z_0)=0,\hfill \\ & \\ \hfill (\mathrm{iii})& _k_jf_0(Z_0)=_k_jg_0(Z_0)=0,\hfill \\ & \\ \hfill (\mathrm{iv})& _k_0f_j(Z_0)=\sqrt{1}\delta _{jk},_k_0g_j(Z_0)=\delta _{jk},\hfill \\ & \\ \hfill (\mathrm{v})& _k_lf_j(Z_0)=_{kl}^{\overline{j}},_k_lg_j(Z_0)=\sqrt{1}_{kl}^{\overline{j}}.\hfill \end{array}$$
(103)
First,
$$_0f_G(z,\tau )=\frac{|\mathrm{Im}\tau _0|}{\mathrm{Im}\tau }_X(F+\overline{\tau }H)\mathrm{\Omega }_z.$$
(104)
It follows that
$$_0^2f_G(z_0,\tau _0)=\frac{|\mathrm{Im}\tau _0|^2}{\mathrm{Im}\tau }\frac{}{\tau }_X(F+\overline{\tau }H)\mathrm{\Omega }_z=0$$
by the critical point equation $`_0f_G(z_0,\tau _0)=0`$. This proves (i).
Next, differentiating (104) with $`f_G=f_\alpha `$, we get
$$_j_0f_\alpha (Z_0)=\overline{\chi _\alpha }๐_j\mathrm{\Omega }_{Z_0}=\overline{\chi _\alpha }\chi _j=i\delta _{j\alpha },$$
and similarly,
$$_j_0g_\alpha (Z_0)=\overline{i\chi _\alpha }\chi _j=\delta _{j\alpha }.$$
This verifies (ii) and (iv).
Finally, we have by (46),
$$_k_jf_\alpha =\chi _\alpha ๐_k๐_j\mathrm{\Omega }=i\underset{l}{}_{kj}^{\overline{l}}\chi _\alpha \overline{๐_l\mathrm{\Omega }},$$
and hence
$$_k_jf_\alpha (Z_0)=i\underset{l}{}_{kj}^{\overline{l}}\chi _\alpha \overline{\chi }_l=i\underset{l}{}_{kj}^{\overline{l}}\delta _{l\alpha }=\{\begin{array}{cc}i_{kj}^{\overline{\alpha }}\hfill & \text{for }\alpha 1\hfill \\ 0\hfill & \text{for }\alpha =0\hfill \end{array}.$$
We also have $`_k_jg_\alpha (Z_0)=i_k_jf_\alpha (Z_0)`$, verifying (iii) and (v).
Thus, the holomorphic Hessian $`H^{}(Z_0)`$ maps the orthonormal fluxes
$$iU_1,\mathrm{},iU_{h^{2,1}},iV_1,\mathrm{},iV_{h^{2,1}}$$
(105)
to the matrices $`\xi ^1,\mathrm{},\xi ^{2h^{2,1}}`$ given by (27). Furthermore,
$$f_0(Z_0)=1,H^{}(U_0)=0,g_0(Z_0)=i,H^{}(V_0)=0,$$
while
$$f_j(Z_0)=g_j(Z_0)=0.$$
Thus $`H^c(Z_0)`$ maps the orthonormal fluxes (105) to the elements $`\xi ^a0\mathrm{Sym}(h^{2,1}+1,)`$, and maps $`U_0`$ to $`01`$ and $`V_0`$ to $`0i`$. โ
### 6.2. Distortion of inner product under the Hessian map
We recall that the space $`\mathrm{Sym}(h^{2,1}+1,)`$ of complex symmetric matrices, regarded as a real vector space, has the inner product
$$(A,B)_{}=\mathrm{Re}A,B_{HS}=\mathrm{Re}(\text{Trace}AB^{}).$$
(106)
Recalling that $`๐ฎ_Z=\stackrel{~}{๐ฒ}_Z(H_z^{2,1}H_z^{0,3})`$, we consider its codimension 1 subspace
$$๐ฎ_Z^{}=\stackrel{~}{๐ฒ}_Z(H_z^{2,1}).$$
By the proof of Lemma 6.1, the holomorphic Hessian map
$$H_Z:๐ฎ_Z^{}_Z$$
(107)
is bijective, but as a map between inner product spaces, it is not an isometry. The distortion is given by the positive definite operator $`\mathrm{\Lambda }_Z`$. We write
$$\mathrm{\Lambda }_Z\xi ^a=\underset{b=1}{\overset{2h^{2,1}}{}}\mathrm{\Lambda }_{ab}\xi ^b,$$
so that
$$(\xi ^a,\xi ^b)_{}=(\mathrm{\Lambda }_Z^1\mathrm{\Lambda }_Z\xi ^a,\xi ^b)_{}=\underset{c}{}\mathrm{\Lambda }_{ac}(\mathrm{\Lambda }_Z^1\xi ^c,\xi ^b)_{}=\underset{c}{}\mathrm{\Lambda }_{ac}\delta _{cb}=\mathrm{\Lambda }_{ab}.$$
Tracing through the definitions, we obtain that $`(\mathrm{\Lambda }_{ab})`$ is the matrix
$$\left(\begin{array}{cc}\mathrm{\Lambda }^{}& \mathrm{\Lambda }^{\prime \prime }\\ \mathrm{\Lambda }^{\prime \prime }& \mathrm{\Lambda }^{}\end{array}\right),\mathrm{\Lambda }_{jk}^{}=2\delta _{jk}+\mathrm{Re}\text{Tr}^j^k,\mathrm{\Lambda }_{jk}^{\prime \prime }=\mathrm{Im}\text{Tr}^j^k.$$
(108)
of Hilbert-Schmidt inner products of the matrices in Lemma 6.1. Hence,
$$\mathrm{\Lambda }_{jk}^{}+\sqrt{1}\mathrm{\Lambda }_{jk}^{\prime \prime }=2\delta _{jk}+\text{Tr}^j^k,$$
(109)
To tie this discussion together with that in \[AD\] and \[DSZ2, ยง2.1\], we note that we can consider $`_Z`$ as a complex vector space by redefining complex multiplication in $`_Z`$:
$$c\left(\begin{array}{cc}0& u\\ u^t& A\end{array}\right)=\left(\begin{array}{cc}0& \overline{c}u\\ \overline{c}u^t& cA\end{array}\right).$$
We then define a Hermitian inner product on $`_Z`$:
$$(\left(\begin{array}{cc}0& u\\ u^t& A\end{array}\right),\overline{\left(\begin{array}{cc}0& v\\ v^t& B\end{array}\right)})=2\overline{u}v+\text{Tr}(AB^{}).$$
We recall from (29) that
$$\mathrm{\Lambda }_Z=\underset{j=1}{\overset{h^{2,1}}{}}\xi ^j\xi ^j,$$
(110)
where the $`\xi ^j`$ are $`(h^{2,1}+1)\times (h^{2,1}+1)`$ matrices. Each term $`\xi ^j\xi ^j`$ in $`\mathrm{\Lambda }_Z`$ may be expressed in matrix form as $`\left(\xi _{ab}^j\overline{\xi }_{cd}^j\right)`$; i.e.,
$$(\mathrm{\Lambda }_ZH)_{kl}=\underset{p,q}{}[\mathrm{\Lambda }_Z]_{kl}^{pq}H_{pq},[\mathrm{\Lambda }_Z]_{kl}^{pq}=\underset{j=1}{\overset{h^{2,1}}{}}\xi _{kl}^j\overline{\xi }_{pq}^j,0k,l,p,qh^{2,1}.$$
(111)
As in \[DSZ2, ยง2.1\], the result may be expressed in terms of the Szegรถ kernel $`\mathrm{\Pi }_Z`$, i.e. the kernel of the orthogonal projection onto $`๐ฎ_Z.`$ By (103) and (110),we have
$$\left[\mathrm{\Lambda }_Z\right]_{kl}^{pq}=_{\zeta _k}_{\zeta _l}_{\overline{\eta }_p}_{\overline{\eta }_q}F_Z(\zeta ,\eta )|_{\zeta =\eta =Z},$$
(112)
where $`F_Z`$ is the local representative of $`\mathrm{\Pi }_Z`$ in a frame (cf. \[DSZ2\]).
In addition, $`\mathrm{\Lambda }_Z`$ determines an operator $`\stackrel{~}{\mathrm{\Lambda }}_Z`$ on the space $`^c`$ of complex matrices of the form
$$H^c:=\left(\begin{array}{cc}H& xI\\ \\ \overline{x}I& \overline{H}\end{array}\right),H\mathrm{Sym}(h^{21},),$$
(113)
defined by
$$\stackrel{~}{\mathrm{\Lambda }}_Z\left(\begin{array}{cc}H& xI\\ \\ \overline{x}I& \overline{H}\end{array}\right)=\left(\begin{array}{cc}\mathrm{\Lambda }_ZH& xI\\ \\ \overline{x}I& \overline{\mathrm{\Lambda }_ZH}\end{array}\right)$$
(114)
We now relate the $`(1,1)`$-form $`\omega _\mathrm{\Lambda }`$ of (31) and the operator $`\mathrm{\Lambda }`$ to the curvature of the Weil-Petersson metric on $`๐`$.
###### Proposition 6.2.
We have:
1. $`[\mathrm{\Lambda }_Z]_{j^{}q^{}}^{jq}=G^{q\overline{p}}R_{j^{}q^{}\overline{p}}^j+\delta _j^{}^j\delta _q^{}^q+\delta _q^{}^j\delta _j^{}^q`$, where $`R`$ is the curvature tensor of the Weil-Petersson metric on $`๐`$;
2. $`\omega _\mathrm{\Lambda }=(m+3)\omega _{WP}+Ric(\omega _{WP})`$ where $`Ric`$ is the Ricci curvature $`(1,1)`$ form of the Weil-Petersson metric of $``$, i.e.
$$Ric(\omega _{WP})=\frac{i}{2}\underset{i\overline{j}}{}Ric_{i\overline{j}}dz^id\overline{z}^j,Ric_{i\overline{j}}:=G^{k\overline{\mathrm{}}}R_{i\overline{j}k\overline{\mathrm{}}}.$$
Thus, $`\omega _\mathrm{\Lambda }`$ is the Hodge metric \[Lu, Wa2\].
###### Proof.
To prove (i), it suffices to combine (111) and (52), raising and lowering indices as appropriate. (In (111), a normal frame at $`Z`$ is assumed.)
For (ii) we note that the $`(1,1)`$-form
$$\omega _\mathrm{\Lambda }=\frac{i}{2}\left[2\delta _{ij}+\text{Tr}^i(Z)^j(Z)\right]dz^id\overline{z}^j$$
(115)
On the other hand, by (47),
$$\begin{array}{ccc}Ric_{i\overline{j}}\hfill & =\hfill & G^{k\overline{\mathrm{}}}\left[G_{i\overline{j}}G_{k\overline{\mathrm{}}}+G_{i\overline{\mathrm{}}}G_{k\overline{j}}\frac{1}{_{}\mathrm{\Omega }\overline{\mathrm{\Omega }}}_{p,q}G^{p\overline{q}}_{ikp}\overline{_{j\mathrm{}q}}\right]\hfill \\ & & \\ & =\hfill & (m+1)G_{i\overline{j}}+Tr^i^j\hfill \end{array}$$
(116)
Remark: To facilitate comparison with \[AD, DSZ1\], we note that our notational conventions are the same as in \[DSZ1\]. In \[AD\], the Szegรถ kernel $`\mathrm{\Pi }_Z`$ is denoted $`G_Z`$. The formulas in \[AD\] (4.8) are the same as (111), resp. Proposition 6.2(1). Also $`F_{ab|\overline{c}\overline{d}}=\mathrm{\Lambda }_{ab}^{pq}G_{p\overline{c}}G_{q\overline{d}}.`$ The coefficients $`F_{a\overline{b}|c\overline{d}}`$ in \[AD\] correspond to the off-diagonal blocks of $`\stackrel{~}{\mathrm{\Lambda }}`$.
### 6.3. Proof of Theorem 1.8
All but one of the ingredients of the proof are precisely the same as in Theorem 1.4. We first define the analogue of (25) and (32) for the signed sum:
$`nd(Z)`$ $`:=`$ $`{\displaystyle _{๐ฎ_Z}}detH^cW(Z)\chi _{Q_Z}dW`$ (117)
$`=`$ $`{\displaystyle \frac{1}{b_3!\sqrt{det\mathrm{\Lambda }_Z}}}{\displaystyle __Z}det\left(H^{}H|x|^2I\right)e^{(\mathrm{\Lambda }_Z^1H,H)_{}|x|^2}dHdx.`$
By Lemma 4.2 and the proof of Lemma 4.3, we conclude that
$$nd_{\chi _K}(L)=L^{b_3}\left[_Knd(Z)d\mathrm{Vol}_{WP}+O(L^{1/2})\right].$$
(118)
To complete the proof of Theorem 1.8, we evaluate the integral in (117):
###### Lemma 6.3.
We have
$$b_3!nd(Z)d\mathrm{Vol}_{WP}=\frac{\pi ^{2m}}{2^m}c_m(T^{(1,0)}(๐),\omega _{WP}h_{WP}^{})=\left(\frac{\pi }{2}\right)^mdet\left(R\omega I\right).$$
###### Proof.
This follows by a supersymmetric formula for the determinant, used in this context in \[AD\] and also in \[BSZ2\]. We briefly review the fermionic formalism referring to \[BGV, BSZ2\] for further details in a similar setting.
Let $`M=\left(M_j^{}^j\right)`$ be an $`n\times n`$ complex matrix. Then,
$$detM=^{B^{2n}}e^{M\eta ,\overline{\eta }}๐\eta ,M\eta ,\overline{\eta }=\underset{j,j^{}}{}\eta _jM_j^{}^j\overline{\eta }_j^{},$$
(119)
where $`\eta _j,\overline{\eta }_j`$ ($`1jn`$) are anti-commuting (or โfermionicโ) variables. The integral $`^B=^{B^{2n}}`$ is the Berezin integral, a notation for the linear functional $`^B:^{}^{2n}`$ defined by
$$^B|_{^t^{2n}}=0\text{for }t<2n,^B\left(_j\overline{\eta }_j\eta _j\right)=1.$$
We now apply this formalism to $`det\left(H^{}H|x|^2I\right)=detH^c`$ where $`H^c`$ is defined as in (113) and refer to the discussion in ยง6.2. The matrix $`H^c`$ is of rank $`b_3`$, and we write
$$detH^c=^{B^{2b_3}}e^{H^c(\eta ,\overline{\eta }),(\theta ,\overline{\theta })}๐\eta ๐\theta ,$$
(120)
where $`\eta =(\eta _1,\mathrm{},\eta _{b_3/2}),\theta =(\theta _1,\mathrm{},\theta _{b_3/2})`$, and
$$H^c(\eta ,\overline{\eta }),(\theta ,\overline{\theta })=\left(H_{jk}\eta _j\theta _k+x\delta _{jk}\eta _j\overline{\theta }_k+\overline{x}\delta _{jk}\overline{\eta }_j\theta _k+\overline{H}_{jk}\overline{\eta }_j\overline{\theta }_k\right).$$
The quadratic form $`(\mathrm{\Lambda }_Z^1H,H)_{}+|x|^2`$ in the exponent of the Gaussian integral may be expressed in the form $`\frac{1}{2}(\stackrel{~}{\mathrm{\Lambda }}_Z^1H^c,H^c)`$, where $`\stackrel{~}{\mathrm{\Lambda }}_Z`$ is the restriction of the operator defined in (114) to $`_Z^c`$. Indeed, both quadratic forms are equivalent to $`Q_Z(W,W)`$ under a linear change of variables ($`WH_Z(W)`$ in the case of $`\mathrm{\Lambda }_Z`$ and $`WH^c(W)`$ in the case of $`\stackrel{~}{\mathrm{\Lambda }}_Z`$).
Then
$$b_3!nd(Z)=\frac{1}{\sqrt{det\stackrel{~}{\mathrm{\Lambda }}_Z}}_{_Z^c}^{B^{2b_3}}e^{H^c(\eta ,\overline{\eta }),(\theta ,\overline{\theta })\stackrel{~}{\mathrm{\Lambda }}_Z^1H^c,H^c}๐H^c๐\eta ๐\theta .$$
(121)
We let
$$\mathrm{\Omega }=(\eta ,\overline{\eta })(\theta ,\overline{\theta })^t=\left(\begin{array}{cc}(\eta _j\theta _k)& (\eta _j\overline{\theta }_k)\\ (\overline{\eta }_j\theta _k)& (\overline{\eta }_j\overline{\theta }_k)\end{array}\right),$$
so that $`H^c(\eta ,\overline{\eta }),(\theta ,\overline{\theta })=(H^c,\mathrm{\Omega })=\text{Tr}H^c\mathrm{\Omega }^t`$. Then the $`dH^c`$ integral in (121) becomes the Fourier transform of the Gaussian function $`e^{\stackrel{~}{\mathrm{\Lambda }}^1H^c,H^c}`$ evaluated at $`i\mathrm{\Omega }`$. Recalling that the Fourier transform of $`e^{Ax,x/2}`$ equals $`(2\pi )^{n/2}(detA)^{1/2}e^{A^1\xi ,\xi /2}`$, we have that the $`dH^c`$ integral equals $`(det\stackrel{~}{\mathrm{\Lambda }})^{\frac{1}{2}}e^{\frac{1}{4}\stackrel{~}{\mathrm{\Lambda }}\mathrm{\Omega },\mathrm{\Omega }}`$. After cancelling $`(det\stackrel{~}{\mathrm{\Lambda }})^{\frac{1}{2}}`$, we obtain
$$b_3!nd(Z)=\pi ^m^{B^{2b_3}}e^{\frac{1}{4}(\stackrel{~}{\mathrm{\Lambda }}\mathrm{\Omega },\mathrm{\Omega })_{}}๐\eta ๐\theta ,$$
(122)
where in normal coordinates, we have (by (114) and Proposition 6.2)
$`(\stackrel{~}{\mathrm{\Lambda }}_Z\mathrm{\Omega },\mathrm{\Omega })_{}`$ $`=`$ $`\text{Trace}\left[\left(\begin{array}{cc}\mathrm{\Lambda }_Z\eta \theta & \eta \overline{\theta }\\ \overline{\eta }\theta & \overline{\mathrm{\Lambda }}_Z\overline{\eta }\overline{\theta }\end{array}\right)\left(\begin{array}{cc}\eta \theta & \eta \overline{\theta }\\ \overline{\eta }\theta & \overline{\eta }\overline{\theta }\end{array}\right)^{}\right]`$
$`=`$ $`{\displaystyle \underset{jqj^{}q^{}}{}}\left(\mathrm{\Lambda }_{j^{}q^{}}^{jq}\eta _j\theta _q\overline{\eta }_j^{}\overline{\theta _q^{}}+\overline{\mathrm{\Lambda }}_{j^{}q^{}}^{jq}\overline{\eta }_j\overline{\theta }_q\eta _j^{}\theta _q^{}\right)+{\displaystyle \underset{jq}{}}\left(\eta _j\overline{\theta }_q\overline{\eta }_j\theta _q+\overline{\eta }_j\theta _q\eta _j\overline{\theta }_q\right)`$
$`=`$ $`2{\displaystyle \underset{jqj^{}q^{}}{}}\left(\mathrm{\Lambda }_{j^{}q^{}}^{jq}\delta _{jj^{}}\delta _{qq^{}}\right)\eta _j\theta _q\overline{\eta }_j^{}\overline{\theta _q^{}}`$
$`=`$ $`2{\displaystyle \underset{jqj^{}q^{}}{}}\left(R_{j\overline{j^{}}q\overline{q^{}}}+\delta _{jq}\delta _{j^{}q^{}}\right)\eta _j\theta _q\overline{\eta }_j^{}\overline{\theta _q^{}}.`$
(Here we used the fact that $`\overline{\mathrm{\Lambda }}_{j^{}q^{}}^{jq}=\mathrm{\Lambda }_{jq}^{j^{}q^{}}`$; see (111).) Thus
$`b_3!nd(Z)`$ $`=`$ $`\pi ^m{\displaystyle ^{B^{2b_3}}}e^{\frac{1}{2}\left(R_{j\overline{j^{}}q\overline{q^{}}}+\delta _{jq}\delta _{j^{}q^{}}\right)\eta _j\overline{\eta }_j^{}\theta _q\overline{\theta }_q^{}}๐\eta ๐\theta `$
$`=`$ $`\left({\displaystyle \frac{\pi }{2}}\right)^m{\displaystyle \frac{det\left(R\omega I\right)}{d\mathrm{Vol}_{WP}}}.`$
Remark: The index density computation in special geometry is closely related to the asymptotics in \[DSZ2, ยง5\] for critical point densities for powers of a positive line bundle $`L`$ on a compact Kรคhler manifold $`M`$. The expansions in ยง5.1 of \[DSZ2\] can be used to show that the (first few) terms in the asymptotic expansion of the index density equal those of the Chern form corresponding to $`c_m(T^{1,0}L^N)`$.
### 6.4. Examples
We describe in this section the critical point distribution for the cases where the dimension $`h^{2,1}(X)`$ of the moduli space is 0 and 1, i.e. when $`dim๐`$ is 1 and 2, respectively.
#### 6.4.1. $`h^{2,1}(X)=0`$
The simplest example is the case where the Calabi-Yau manifold $`X`$ is rigid, i.e. $`=\{pt\}`$. (See \[AD, DD1\] for further details and computer graphics of critical points in this case.) Then only the parameter $`\tau `$ varies. Let $`G=F+iH`$, and consider the flux superpotential $`W_G`$. Its critical point equation is
$$F+\tau HH^{0,3}$$
(since in this case $`H^{2,1}(X,)=0`$). So we write
$$F=A\mathrm{\Omega }+\overline{A\mathrm{\Omega }},H=B\mathrm{\Omega }+\overline{B\mathrm{\Omega }},A=a_1+ia_2,B=b_1+ib_2+\sqrt{1}.$$
Then writing $`W_G=W_{A,B}`$, we have
$$W_{A,B}=0F+\tau HH^{0,3}A+\tau B=0\tau =\frac{A}{B}.$$
Each flux superpotential $`W_{A,B}๐ฎ`$ (with $`A,B`$) has a unique critical point in $``$, which may or may not lie in the fundamental domain $`๐`$. In the notation of (15),
$$\pi (๐ฎ)=\{W_{A,B}:\frac{A}{B}๐\}$$
is a domain with boundary in $`^2`$. Each $`SL(2,)`$-orbit of fluxes (or superpotentials) contains a unique element whose critical point lies in $`๐`$, so $`\pi (S)`$ is a fundamental domain for the action of $`\mathrm{\Gamma }`$ on $`๐ฎ`$.
Thus, counting critical points is equivalent to counting $`SL(2,)`$ orbits of superpotentials satisfying the tadpole constraint. The pair $`(A,B)`$ corresponds to the element $`\left(\begin{array}{cc}a_1\hfill & b_1\hfill \\ a_2\hfill & b_2\hfill \end{array}\right)GL(2,)`$ and the Hodge-Riemann form quadratic form may be identified with the indefinite quadratic form
$$Q[(A,B)]=a_1b_2b_2a_1$$
on $`^4`$. The modular group $`SL(2,)`$ acts by the standard diagonal action on $`(A,B)^2\times ^2`$ preserving $`Q[(A,B)]`$ or equivalently by left multiplication preserving $`det`$. Thus, the set of superpotentials satisfying the tadpole constraint is parametrized by:
$$\{\left(\begin{array}{cc}a_1\hfill & b_1\hfill \\ a_2\hfill & b_2\hfill \end{array}\right)GL(2,):0<det\left(\begin{array}{cc}a_1\hfill & b_1\hfill \\ a_2\hfill & b_2\hfill \end{array}\right)L\},$$
and we want to count the number of $`SL(2,)`$-orbits in this set. Counting the number of $`SL(2,)`$ orbits in $`๐_L`$ is equivalent to determining the average order of the classical divisor function $`\sigma (m)`$, see for instance Hardy-Wright \[HW, Theorem 324\]:
$$๐ฉ^{\mathrm{crit}}(L)=\underset{m=1}{\overset{L}{}}\underset{k|m}{}k=\underset{m=1}{\overset{L}{}}\sigma (m)\frac{\pi ^2}{12}L^2+O(L\mathrm{log}L).$$
(123)
As verified in \[DD1\] (and as follows very simply from Theorem 1.4), the critical points are uniformly distributed relative to the hyperbolic area form.
#### 6.4.2. $`h^{2,1}(X)=1`$
We now illustrate our notation and results with the case where the moduli space of complex structures on $`X`$ is one-dimensional over $``$. (This case is also studied in \[DD1\] from a slightly different point of view.) In this case, there is a single Yukawa coupling $`_{11}^{\overline{1}}(z)`$ defined by $`D_z^2\mathrm{\Omega }_z=_{11}^{\overline{1}}(z)\overline{D_z\mathrm{\Omega }_z}.`$
The space $`๐ฎ_{z,\tau }H^{2,1}H^{0,3}^2`$. The space is spanned as a real vector space by four superpotentials $`U_0,U_1,V_0,V_1`$ corresponding to $`\{\overline{\mathrm{\Omega }_z},๐_z\mathrm{\Omega }_z,i\overline{\mathrm{\Omega }_z},i๐_z\mathrm{\Omega }_z\}`$. By the proof of Lemma 6.1, the holomorphic Hessians of $`U_0`$ and $`V_0`$ at a critical point equal zero, so we only need to consider the holomorphic Hessian map on $`U_1`$ and $`V_1`$. The corresponding space of Hessians is the real $`2`$-dimensional subspace $`_Z`$ of $`\mathrm{Sym}(2,)`$ spanned by
$$\xi ^1=\left(\begin{array}{cc}0& 1\\ 1& F(z)\end{array}\right),\xi ^2=\sqrt{1}\left(\begin{array}{cc}0& 1\\ 1& F(z)\end{array}\right),$$
where we write $`F=_{11}^{\overline{1}}`$. Hence, we may parameterize the space $`_Z`$ of holomorphic Hessians by
$$w=y_1+iy_2H(w)=\left(\begin{array}{cc}0\hfill & w\hfill \\ & \\ w\hfill & F(z)\overline{w}\hfill \end{array}\right).$$
By (25), we have:
$$๐ฆ^{\mathrm{crit}}(Z)=\frac{1}{2!}_{}|det(H(w)^{}H(w)|x|^2I)|e^{|w|^2+|x|^2}๐w๐x.$$
We note that
$$det(H(w)^{}H(w)|x|^2I)=|w|^4+|x|^4(2+|F(z)|^2)|x|^2|w|^2.$$
Hence
$$๐ฆ^{\mathrm{crit}}(Z)=\frac{1}{2!}_{}\left||w|^4+|x|^4(2+|F(z)|^2)|x|^2|w|^2\right|e^{|w|^2+|x|^2}๐w๐x,$$
agreeing with (3.19) of \[DD1\]. There, the integral is evaluated as
$$๐ฆ^{\mathrm{crit}}(Z)=\frac{\pi ^2}{2}\left(2|F|^2+\frac{2|F|^3}{\sqrt{4+|\stackrel{~}{F}|^2}}\right).$$
Remark: In this example, the discriminant variety is given by
$$\stackrel{~}{๐}=\{(Z,xW_0(Z)+wW_1(Z)):|w|^2|x|^2=\pm |wxF(z)^2|\},$$
where $`W_\alpha =U_\alpha +iV_\alpha `$. The matrix $`\mathrm{\Lambda }`$ is given by
$$\mathrm{\Lambda }=\left(\begin{array}{cc}2+|F|^2& 0\\ 0& 2+|F|^2\end{array}\right).$$
## 7. Problems and heuristics on the string theory landscape
In this section, we continue the discussion begun in ยง1.6 on the bearing of our methods and results on the physicistsโ picture of the string theory landscape. We briefly review some of the heuristic estimates in the physics discussions, and then discuss a number of mathematical pitfalls in the heuristics. In ยง7.2, we state some mathematical problems suggested by the heuristics and by rigorous vacuum statistics. In ยง7.3, we give our own (tentative) heuristic estimate of the dependence of the critical point density $`๐ฆ^{\mathrm{crit}}(Z)`$ on the dimension $`b_3/2`$ of $`๐`$.
### 7.1. Complexity of the string theory landscape
As mentioned in ยง1.6, the possible vacua in string/M theory are often represented as valleys in a complex string theory landscape, and the number of valleys is often estimated at $`10^{500}`$.
L. Susskind and others have argued that such a large number of possible vacua should essentially be a consequence of the large number of variables in the potential. A common and general argument to arrive at this number of vacua without specifying any particular string theory model is to reason that the potential energy is a function of roughly $`1000`$ variables. A generic polynomial $`f`$ of degree $`d`$ on $`^m`$ has $`(d1)^m`$ critical points since critical points are solutions of the $`m`$ equations $`\frac{f}{z_j}(w)=0`$ of degree $`d1`$. Thus, the number of critical points would seem to grow at an exponential rate in the number of variables. Such an exponential growth rate of critical points also appears in the physics of spin glasses, where the growth in the number of metastable states (local minima of the Hamiltonian) in terms of the number of variables is often used to measure the complexity of the energy landscape. In special model of random Hamiltonians on domains in $`^N`$, exponential growth of the number of local minima in $`N`$ has recently been proved rigorously \[Fy\].
In the specific models of type IIb flux compactifications on a CY $`3`$-fold $`X`$, the number of variables is $`b_3(X)`$. As mentioned above, for a typical $`CY`$ $`3`$-fold, $`b_3`$ is often around $`300`$ and sometimes as high as $`1000`$ (cf. \[GHJ, CO\]), and therefore the scalar potential $`V_W`$ in (11) is a function of this number of variables. By naive counting of variables one would thus arrive at a figure like $`10^{500}`$ for such models. The more sophisticated estimate $`N_{vac}\frac{L^{b_3}}{b_3!}f(b_3)`$ in flux compactifications (see ยง1.6 for the notation) does not supplant the naive counting argument since the order of magnitude of $`f(b_3)`$ is unknown. We recall that it is the integral over $`๐`$ of the Gaussian integral in (32) (see (126). The Gaussian integral for $`๐ฆ^{\mathrm{crit}}`$ in that line resembles to some extent the integral formula for the expected number of critical points in spin glass theory, which has exponential growth (see e.g. \[Fy\]).
Although the naive counting of variables or the analogy to complexity of energy landscapes bring some insight into vacuum counting, we now point out some pitfalls in estimating numbers of vacua or the coefficient $`f(b_3)`$ in flux compactifications on this basis.
1. The critical point equation (12) is $`C^{\mathrm{}}`$ but not holomorphic, so vacua are critical points of a real system of equations, and it is not obvious how many connection critical points to expect even a polynomial of a given degree to have. This number depends on the connection, and is studied in detail in \[DSZ1, DSZ2\] and in the present paper.
2. A flux superpotential $`W`$ is not a polynomial and it is not clear how to assign it a โdegreeโ which reflects its number of critical points on all of Teichmรผller space, or equivalently, the number of critical points in $`๐`$ corresponding to the $`\mathrm{\Gamma }`$-orbit of $`W`$. Examples (e.g. in ยง6.4.1) show that this number can be relatively small.
3. It seems reasonable to say that the number of fluxes rather than the number of critical points per flux that dominates the number of vacua. In flux compactifications, the landscape should therefore be viewed as the graph of the scalar potential $`V_W(Z)`$ on $`๐\times ๐ฎ`$, i.e. as a function of both variables $`W,Z`$, and the local minima should be viewed as pairs $`(W_G,Z)`$ with $`GH^3(X,\sqrt{1})`$ and with $`ZCrit(W_G).`$
4. However (see the problems below) it is not straightforward
to define โper vacuaโ, since the tadpole constraint is hyperbolic, and the total number of lattice points in the shell $`0<Q[G]<L`$ is infinite.
5. In estimating $`๐ฆ^{\mathrm{crit}}(Z)`$ we are fixing $`Z`$ in the interior of $`๐`$. But there could exist singular points of $`๐`$ at which $`๐ฆ^{\mathrm{crit}}(Z)`$ blows up (see \[DD1\] for discussion of conifold points). It would also be interesting to study $`๐ฆ^{\mathrm{crit}}(Z)`$ as $`Z๐`$.
6. As mentioned in ยง1.6 (see also ยง7.3), there may be a significant difference between the order of magnitude of the density of critical points and of the number of critical points, since $`๐`$ is an incomplete Kรคhler manifold of possibly quite small volume. See \[LS1\] for the current state of the art on the volume. There is no analogue of the small volume of the configuration space in spin glass complexity.
7. The tadpole constraint (1) becomes much more highly constraining as the number $`b_3`$ of variables increases for fixed $`L`$ and is responsible for the factor $`\frac{1}{(b_3)!}`$ in Theorem 1.4. Again, no such feature exists in complexity estimates in spin glasses.
### 7.2. Problems
The issues mentioned above (and the detailed heuristics in ยง7.3) suggest a number of problems. The ultimate goal is:
###### Problem 7.1.
Does string theory contain a vacuum consistent with the standard model, and if so, how many? Find examples of Calabi-Yau manifolds, and any other postulated structures, for which it is certain that such a vacuum exists.
Now testing consistency with the standard model requires elucidating far more structure of a candidate vacuum โ the gauge group, the matter content, and so forth โ than we are considering here. To address this ultimate problem, one would need many more statistical results, along the lines set out in \[Do\]. However one can make arguments (admittedly quite speculative at this point) that the dominant multiplicity in vacuum counting arises from the multiplicity of flux vacua we are discussing here. An important problem in this context is
###### Problem 7.2.
How large does $`L`$ need to be to ensure that there exists a vacuum with
$$|W_G(Z)|^2\lambda _{}$$
(124)
for a specified $`\lambda _{}`$ ? In that case, how many such vacua are there? Find examples of Calabi-Yau manifolds where it is certain that such a vacuum exists.
To solve this problem for type IIb flux compactifications, we would need to sharpen Theorem 1.4 in many ways which lead to the subsequent problems stated below.
The constraint (124) on $`|W_G(Z)|^2`$ is a simple example of โconsistency with the standard model.โ If the real world were (counter-factually) exactly supersymmetric, this would be the constraint that the vacuum should have a cosmological constant $`V_W(Z)=3|W_G(Z)|^2`$ (as in (11)) consistent with the known value. While the physical discussion requires taking supersymmetry breaking into account, as discussed in \[DD2\], vacua can exist in which supersymmetry is broken by effects not taken into account here, making additional contributions to the vacuum energy which lift the exact vacuum energy to be consistent with the known value (essentially, zero). For such a vacuum, the quantity $`3|W_G(Z)|^2`$ would be the mass squared of the gravitino, a quantity which could be constrained by physical observations.
An independent motivation for (124) is that some proposals for stabilizing the moduli we did not discuss, such as that of \[KKLT\], are believed only to work under such a constraint.
In any case, as discussed in \[DD1\] (ยง3.3), one can count such vacua by choosing the test function to be $`\theta (\lambda _{}|W_G(Z)|^2)`$ where $`\theta (x)=1`$ for $`x>0`$ and $`=0`$ for $`x0.`$ This test function is not homogeneous but can be handled by the methods of this paper (loc. cit.).
Theorem 1.4 is asymptotic in $`L`$ and we have also analyzed to some degree the $`b_3`$ dependence. But as mentioned in ยง1.6, $`L`$ depends on the topology of $`X`$. There, we stated that in many examples $`LCb_3`$ with $`1/3C3`$. To bridge one gap between Theorem 1.4 and Problem 7.2, we state:
###### Problem 7.3.
How are the order of magnitudes of $`b_3(X)`$ and $`L`$ of (36) related as $`X`$ varies over topologically distinct Calabi-Yau manifolds?
We have already mentioned the importance of obtaining effective estimates in $`b_3`$ of the coefficient (24) in Theorem 1.4:
###### Problem 7.4.
Obtain an effective estimate of $`๐ฆ^{\mathrm{crit}}(Z)`$ and of its integral over $`๐`$ in $`b_3`$. Also, obtain such an estimate of the remainder.
Among the difficulties with this problem is that $`๐ฆ^{\mathrm{crit}}(Z)`$ depends on special features of the moduli space $`๐`$ which depend on more than just the dimension $`b_3`$ and which may change in an irregular way as the dimension increases. We consider this problem below in ยง7.3.
To gain insight into the size of the leading coefficient (24), one could write the principal term in Theorem 1.4 in the form $`\frac{L^{b_3}}{b_3!}\times f(b_3)`$ that is often used in string theory (cf. ยง1.6), with $`f(b_3)`$ the Gaussian integral in (32). As mentioned above, it is natural to try to separate out the effects of the number of fluxes and the number of vacua per flux, or more precisely:
1. the number of fluxes $`G`$ satisfying the tadpole constraint with a critical point in a compact subset $`๐ฆ๐`$;
2. the number of critical points โper fluxโ, or more precisely per $`\mathrm{\Gamma }`$-orbit of fluxes, in $`๐ฆ`$ (see ยง6.4.1 to clarify this distinction);
3. the total number of critical points in $`๐ฆ`$ of all fluxes satisfying the tadpole constraint.
We can define the first quantity precisely as the sum
$$\mathrm{\Theta }_K(L)=\underset{GH^3(X,i):Q[G]L}{}\theta \left(\underset{Z๐:W_G(Z)=0}{}\chi _K(Z)\right).$$
Thus, the problem we pose is:
###### Problem 7.5.
Determine the asymptotics of $`\mathrm{\Theta }_K(L)`$ as $`L\mathrm{}`$.
The second quantity is the ratio $`๐ฉ_K(L)/\mathrm{\Theta }_K(L).`$ A possibly more tractable way to restate this problem is in terms of the โaverage number of critical pointsโ of a superpotential $`W_G`$ in $`๐ฆ`$. To define โaverageโ we need to introduce a probability measure on $``$ which is compatible with $`\chi _QdW`$. The most natural probability measures seem to be the normalized Gaussian measures $`\gamma _{Z_0}`$ on the spaces $`๐ฎ_{Z_0}`$ defined by the inner product $`Q_{Z_0}`$.Thus, we ask for the average number of critical points of $`W๐ฎ_{Z_0}`$ with respect to $`\gamma _{Z_0}`$. It would be interesting to study the number of critical points in a fixed $`๐ฆ๐`$ or in all of $`๐`$ or indeed in all of Teichmรผller space (which corresponds to counting critical points in $`๐`$ for a $`\mathrm{\Gamma }`$-orbit of fluxes).
We observe that $`W๐ฎ_{Z_0}`$ has a critical point at $`Z`$ if and only if $`W๐ฎ_{Z_0}๐ฎ_Z`$. In the case of flux superpotentials, $`dim๐ฎ_{Z_0}=\frac{1}{2}dim`$ so for generic pairs $`Z,Z_0`$, $`๐ฎ_{Z_0}๐ฎ_Z=\{0\}`$. Thus, $`๐_{Z_0}(\mathrm{\#}Crits(W))`$ will be an integral over the special variety $`\mathrm{\Sigma }_{Z_0}=\{Z:dim๐ฎ_{Z_0}๐ฎ_Z>0\}`$. This variety is obviously stratified by $`h^{2,1}`$ strata $`\mathrm{\Sigma }_d`$ on which the dimension $`d`$ takes the values $`d=1,2,\mathrm{},h^{2,1}`$, and $`๐_{Z_0}(\mathrm{\#}Crits(W))`$ is a sum of integrals over each strata.
###### Problem 7.6.
Determine the asymptotics of $`๐_{Z_0}(\chi _{Q_{Z_0}(G/L)}\mathrm{\#}Crits(W_G))`$
We also recall that in Theorem 1.4 we ignored the effect of the discriminant variety and the boundary of the region of $`๐`$.
###### Problem 7.7.
Estimate the remainder if $`\psi `$ does not vanish near the discriminant variety $`๐`$, or if $`\psi `$ is a characteristic function of a smooth region $`K๐.`$ Investigate the boundary behavior as $`๐ฆ`$ fills out to $`๐`$.
An analogue problem about studying accumulation of lattice points around boundaries of domains on non-degenerate surfaces is studied in \[Ze1\].
### 7.3. Heuristic estimate of the critical point density
We now present a heuristic estimate on the $`b_3`$-dependence of the critical point density (relative to the Weil-Petersson volume form)
$`๐ฆ^{\mathrm{crit}}(Z)`$ $`=`$ $`{\displaystyle \frac{1}{b_3!\sqrt{det\mathrm{\Lambda }_Z}}}{\displaystyle __Z}\left|detH^{}H|x|^2I\right|e^{(\mathrm{\Lambda }_Z^1H,H)_{}|x|^2}๐H๐x`$ (125)
for $`Z`$ in regions of moduli space where the norm of $`\mathrm{\Lambda }_Z`$ satisfies bounds independent of $`b_3`$. We recall (cf. Proposition 6.2) that $`\mathrm{\Lambda }_Z`$ is the Hodge metric, hence we are studying the density of critical points in regions $`K๐`$ where the absolute values of the eigenvalues of the Ricci curvature of the Weil-Petersson metric $`\omega _{WP}`$ are bounded by a uniform constant. In the notation $`N_{vac}(L)\frac{L^{b_3}}{b_3!}f(b_3)`$, we have
$$f(b_3)=_๐\chi _K(Z)\frac{1}{\sqrt{det\mathrm{\Lambda }_Z}}__Z\left|detH^{}H|x|^2I\right|e^{(\mathrm{\Lambda }_Z^1H,H)_{}|x|^2}๐H๐x,$$
(126)
where $`๐ฆ`$ is the region in which we are counting the critical points.
Our heuristic estimate is that the Gaussian integral (i.e. $`b_3!๐ฆ^{\mathrm{crit}}(Z)`$) has growth rate $`(b_3/2)!N_\mu ^{b_3}`$ for $`Z`$ in a region $`K=K_\mu `$ of moduli space where $`\mathrm{\Lambda }_Z\mu `$. Here, $`N_\mu `$ is a constant depending only on $`\mu `$. It follows that $`๐ฆ^{\mathrm{crit}}(Z)`$ would have the decay rate $`b_3^{b_3/2}`$ for $`Z`$ in $`K_\mu `$. We note that this heuristic estimate is consistent with the heuristic estimate given by Ashok-Douglas \[AD\] that $`๐ฆ^{\mathrm{crit}}(Z)`$ should have the same order of magnitude as $`nd(Z)`$ (117). By Proposition 6.3, $`b_3!nd(Z)`$ is a differential form depending polynomially on the curvature. The density of $`b_3!nd(Z)`$ relative to $`dVol_{WP}=\frac{\omega _{WP}^{b_3/2}}{(b_3/2)!}`$ thus has the growth $`(b_3/2)!N_\mu ^{b_3}`$ we predict. We present the new heuristic to give evidence that the absolute value only changes the coefficient and not the order of magnitude in vacuum counting.
Before going into the heuristic estimate, we first discuss the consequences for vacuum counting. As mentioned in the introduction, it has been tentatively conjectured at this time of writing (Z. Lu) that the Weil-Petersson volume of $`K_\mu `$ is bounded above by the volume of a ball of radius $`r(\mu )`$ in $`^{b_3/2}`$ depending only on $`\mu `$, and the latter volume decays like $`\frac{1}{(b_3/2)!}`$. Thus it would appear that $`N_{vac,K_\mu }(L)\frac{(C_1LN_\mu )^{b_3}}{b_3!}`$. We include a constant $`C_1`$ to take into account the dependence on various parameters including $`r(\mu )`$, factors of $`\pi `$ and so on. If we then take the (often) observed value $`LCb_3`$ with $`C[\frac{1}{3},3]`$, then the number of vacua in $`K_\mu `$ satisfying the tadpole constraint would grow at an exponential rate in $`b_3`$.
We now explain the heuristic estimate regarding the order of magnitude of $`๐ฆ^{\mathrm{crit}}(Z)`$ (24): the latter depends on two inputs, the subspace $`_Z`$ (or equivalently the orthogonal projection $`P_Z`$ onto $`_Z`$) and the eigenvalues of $`\mathrm{\Lambda }_Z`$. To obtain upper and lower bounds on $`๐ฆ^{\mathrm{crit}}(Z)`$ we note that
$$2P_Z\mathrm{\Lambda }_Z\mu _{\mathrm{max}}(Z)P_Z,$$
(127)
where $`\mu _{\mathrm{max}}(Z)`$ is the maximum eigenvalue of $`\mathrm{\Lambda }_Z`$. We recall here that $`\mathrm{\Lambda }_Z`$ is the matrix of the Hodge metric (see (30)), and its eigenvalues can be estimated in terms of the Weil-Petersson metric and its curvature (cf. \[Lu\]). In particular, its minimum eigenvalue satisfies $`\mu _{\mathrm{min}}(Z)2`$, and that explains the lower bound $`2P_Z`$ in (127). For most CY $`3`$-folds $`X`$, the Weil-Petersson metric on $`๐`$ is incomplete, and $`\mu _{\mathrm{max}}(Z)\mathrm{}`$ as $`Z`$ tends to the boundary (Z. Lu).
By (127), we have
$$J_{}(\mu ,P_Z)(b_3!)๐ฆ^{\mathrm{crit}}(Z)J_+(\mu ,P_Z),(\mu \mu _{\mathrm{max}}(Z))$$
(128)
where
$`J_+(\mu ,P_Z):`$ $`=`$ $`{\displaystyle \frac{1}{^{b_3/21}}}{\displaystyle __Z}\left|detH^{}H|x|^2I\right|e^{\left(\mu ^1\text{Tr}H^{}H|x|^2\right)}๐H๐x,`$ (129)
and where
$`J_{}(\mu ,P_Z):`$ $`=`$ $`{\displaystyle \frac{1}{\mu ^{(b_3/21)}}}{\displaystyle __Z}\left|detH^{}H|x|^2I\right|e^{\left(2^1\text{Tr}H^{}H|x|^2\right)}๐H๐x,`$ (130)
Thus we obtain upper and lower bounds for the density in regions $`K_\mu ๐`$ for which the absolute values of the eigenvalues of the Hodge metric relative to the Weil-Petersson metric satisfy $`\mu _{\mathrm{max}}(Z)\mu `$. We have bounded the determinant of $`\mathrm{\Lambda }`$ by a power of an extremal eigenvalue, but it could also be identified with the volume density of the Hodge metric. We note that the lower bound tends to zero and the upper bound tends to infinity in $`\pm b_3`$ powers of $`\mu _{\mathrm{max}}(Z)`$ as $`Z๐`$ when the Weil-Petersson metric is incomplete and the norm of the Ricci curvature of $`\omega _{WP}`$ tends to infinity.
We now estimate $`J_\pm (\mu ,P_Z)`$ under the assumption that $`_Z`$ is a โsufficiently randomโ subspace. The subspace $`_Z`$ is a real subspace of dimension $`b_32`$ of $`\mathrm{Sym}(b_3/21,)`$, but by modifying the definition of the complex structure it becomes a complex $`b_3/2`$-dimensional one. Hence, we may view $`Z_Z`$ as a map $`๐Gr(b_3/21,\mathrm{Sym}(b_3/21,))`$ to the complex Grassmannian of $`b_3/21`$ dimensional complex subspaces. Lacking knowledge of the distribution of the image of $`Z_Z`$, we make the assumption that it is random, or more precisely we approximate $`J_\pm (\mu ,P_Z)`$ by the expected value of $`J_\pm (\mu ,P)`$, where $`P`$ is the projection corresponding to a random element $`Gr(b_3/21,\mathrm{Sym}(b_3/21,))`$.
This approximation by the expected value seems to be reasonable because Grassmannians $`Gr(k,N)`$ are examples of Gromov-Milman โLevy familiesโ of Riemannian manifolds for which concentration of measure phenomena hold as $`N\mathrm{}`$ \[GM, Ta\]. Concentration of measure refers to a metric space $`(X,d)`$ with a probability measure $`P`$ and a concentration function $`\alpha (P,t)`$, which is the smallest number such that the measure of a set $`A`$ and the metric tube $`A_t=\{x:d(x,A)<t\}`$ around $`A`$ are related by $`P(A)1/2P(A_t)1\alpha (P,t).`$ If $`f`$ is a Lipschitz function and if $`M_f`$ is a median for $`f`$, we put $`A=\{x:f(x)M_f\}`$, and then $`P(|fM_f|>t)2\alpha (P,\frac{t}{f_{Lip}}).`$ Concentration of measure occurs if $`\alpha (P,t)`$ decays rapidly in $`t`$, and thus $`f`$ is highly concentrated around its median. In a Lรฉvy family $`(X_N,d_N)`$, the functions $`\alpha _N(P,t)`$ decay at ever faster rates depending on $`N`$. For instance on the unit $`N`$-sphere $`S^N`$, the rate is (a universal constant times) $`e^{\frac{(N1)}{2}t^2}`$.
In our setting, the family consists of Grassmannians $`Gr(b_3/21,\mathrm{Sym}(b_3/21,))`$ equipped with the invariant probability measure $`d\nu `$ and with the standard bi-invariant metric. It is pointed out in \[GM\] that $`Gr(k,N)`$ is a Lรฉvy family for fixed $`k`$ (see section (3.3) of \[GM\]), and the same argument should apply to $`k_NN/2`$. Moreover, $`\{U(N)\}`$ with its Haar probability measure and bi-invariant metric is Lรฉvy, and by section (2.1) of \[GM\] its quotients should be. The function $`f`$ is $`J_\pm (\mu ,P)`$ for fixed $`\mu `$. Since we are mainly interested in factorial dependencies, we set $`\mu =1`$ and change the exponent $`2^1`$ to $`1`$ to make the Gaussian measure a probability measure. In general, the result would be modified by a $`\pm b_3`$ power of $`\mu `$. In this heuristic discussion, we will not attempt to determine $`\alpha _N(P,t)`$ or $`M_f`$ but will assume that $`\alpha (P,\frac{t}{f_{Lip}})`$ has rapid decrease in $`t`$ which improves with the dimension. We also note that when $`\alpha (P,t)`$ is small, we can replace the median of $`J_\pm (\mu ,P)`$ (with $`\mu =1`$) by its mean
$$_{Gr(b_3/21,\mathrm{Sym}(b_3/21,)}\left\{_{}\right|det(H^{}H|x|^2I)|e^{TrH^{}H|x|^2}dHdx\}๐\nu ()$$
with a small error (cf. \[Ta\]). This mean equals
$$_{\mathrm{Sym}(b_3/21,)}|det(H^{}H|x|^2I)|e^{TrH^{}H|x|^2}๐H๐x$$
(131)
since both measures are invariant probability measures and are therefore equal. Here we ignore factors of $`(2\pi )`$ (etc.) for the sake of simplicity, since we are primarily interested in the factorially growing quantities. Due to the concentration of measure, the spaces $`_Z`$ would have to be very โrare eventsโ if $`J_\pm (\mu ,P_Z)`$ differed appreciably from its mean. We note that since $`H_Z^3`$ is a complex polarization, $`P_Z`$ has special features that do not hold for random subspaces, but we have no reason to believe that these special features bias $`J(\mu ,P_Z)`$ away from its mean.
We now observe that (131) (with any choice of $`\mu `$) is similar to the integral for the density of critical points for holomorphic sections of $`๐ช(N)^m`$ with $`m=b_3/21`$ with respect to the Fubini-Study connection for a fixed degree $`N`$ \[DSZ2\] (ยง4). There, the $`\mathrm{\Lambda }_Z`$ matrix was (for every $`Z`$) a two-block diagonal matrix with a large scalar block and a $`1\times 1`$ scalar block. When $`\mu =1`$ (131) agrees with that $`๐ช(N)^m`$ density in the case $`N=1`$. As noted in \[DSZ2\], the total number of critical points of a given Morse index appears to grow at a rate $`N^m`$ times a rational quantity in $`m`$ as $`m\mathrm{}`$. This growth rate may also be easily verified for the Euler characteristic $`\chi (T^{1,0}๐ช(N))`$, i.e. the alternating sum over the Morse indices, which is given by
$`\chi (T^{1,0}๐ช(N))`$ $`=`$ $`({\displaystyle \frac{c(๐ช(N1))^{m+1}}{c(๐ช(N))}},[^m])={\displaystyle \frac{(N1)^{m+1}+(1)^m}{N}}.`$
Since the volume of $`^m`$ is $`\frac{1}{m!}`$, this would imply that the density of critical points grows like $`m!`$ with the dimension. On this basis, we would expect that $`J_\pm (\mu ,P_Z)`$ for $`\mu 1`$ grows with the dimension at the rate $`(b_3/2)!N_\mu ^{b_3}`$ for some $`N_\mu >0`$.
We note that the Ashok-Douglas heuristic that the density of critical points should have the same order of magnitude as the index density is indeed correct in the setting of $`๐ช(N)^m`$. Further, the origin of the factorials $`(b_3/2)!`$ is essentially in both the $`๐`$ and $`^m`$ settings.
Thus our heuristics give $`๐ฆ^{\mathrm{crit}}(Z)\frac{(b_3/2)!N_\mu ^{b_3}}{b_3!}`$. If we integrate over $`K_\mu `$ and apply the conjectural volume bound $`\frac{1}{(b_3/2)!)}`$ for $`K_\mu `$, we would get roughly $`\frac{L^{b_3}N_\mu ^{b_3}}{b_3!}`$. Further applying the observed relation $`LCb_3`$ with $`C[1/3,3]`$ gives an exponential growth rate for numbers of vacua in $`K_\mu `$.
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# LAPTH-1101/05, DTP-MSU/05-09, Comment on โWhat does the Letelier-Galโtsov metric describe?โ
## Abstract
We show that the Letelier-Galโtsov (LG) metric describing multiple crossed strings in relative motion lg does solve the Einstein equations, in spite of the discontinuity uncovered recently by Krasnikov krasn , provided the strings are straight and moving with constant velocities.
In a recent note krasn , Krasnikov raised an interesting question about the continuity of the LG metric. He showed, on the example of a (parallel) two-cosmic string LG metric, that the metric component $`g_{ty}`$ (where $`y`$ is the coordinate orthogonal to strings and their relative displacement) is generically discontinuous across the cuts associated with the positions of strings. Arguing that the metric must be continuous, he concluded that the strings should be at rest with respect to each other.
Before proceeding with a more elaborate analysis, we note it would be too strong to require the metric to be everywhere continuous. To define a spacetime, one needs a set of overlapping, diffeomorphic charts. We will show that, provided the string motion is geodesic and the strings do not collide, the Riemann tensor computed from the discontinuous metric vanishes outside the strings. It follows that there is a set of overlapping charts such that the metric is continuous in each.
Let us first briefly reformulate Krasnikovโs argument. For simplicity, we shall consider here only the case of moving parallel strings, which can be reduced to that of moving conical singularities in 2+1 dimensions. The reduced LG metric can be written in the ADM form,
$$ds^2=N^2dt^2+h_{ij}(dx^i+N^idt)(dx^j+N^jdt),$$
(1)
with
$$N=1,h_{\zeta \overline{\zeta }}=\frac{1}{2}\overline{Z}_{,\overline{\zeta }}Z_{,\zeta },N_\zeta =h_{\zeta \overline{\zeta }}N^{\overline{\zeta }}=\frac{1}{2}\overline{Z}_{,t}Z_{,\zeta }$$
(2)
($`\zeta =x+iy`$), where, in the case of two strings,
$$Z(t,\zeta )=_{\zeta _0}^\zeta \psi (t,\xi )๐\xi ,\psi (t,\xi )=(\xi \alpha _1(t))^{\mu _1}(\xi \alpha _2(t))^{\mu _2}$$
(3)
($`\mu _i>1`$). This is defined only in the cut complex plane with two line cuts extending from the two conical singularities to infinity. We choose as the first cut the horizontal half-axis $`[\alpha _1,\alpha _1+\mathrm{}[`$. The second cut, starting in $`\alpha _2`$, will be assumed not to intersect the first cut (in this choice we differ from krasn ). Following krasn , we also choose a gauge such that $`\alpha _1(t)=0`$, and $`\alpha _2(t)`$ and $`\zeta _0`$ real, with $`\alpha _2<\zeta _0<0`$.
We choose $`\zeta `$ to be on the first cut, and wish to evaluate at a given time $`t`$ the discontinuity of
$$N_\zeta =\frac{1}{2}\psi (\zeta )\overline{Z_{,t}(\zeta )}.$$
(4)
First,
$$N_\zeta (\zeta \pm i0)=\frac{1}{2}\psi (\zeta )\mathrm{e}^{i\pi \mu _1(11)}\overline{Z_{,t}(\zeta \pm i0)},$$
(5)
with
$$Z_{,t}(\zeta \pm i0)=\mu _2\dot{\alpha }_2_{\mathrm{\Gamma }_\pm }\frac{\psi (\xi )}{\xi \alpha _2}๐\xi .$$
(6)
Choosing $`\mathrm{\Gamma }_+`$ ($`\mathrm{\Gamma }_{}`$) to go first along an arbitrary path from $`\zeta _0`$ (a fixed point anywhere in the $`\zeta `$ plane) to $`\alpha _1=0`$, then (after a small upper (lower) half-circle around the origin) to follow the upper (lower) bank $`\gamma _+`$ ($`\gamma _{}`$) of the cut from $`0`$ to $`\zeta `$, we obtain
$$N_\zeta (\zeta \pm i0)=\frac{1}{2}\left(\overline{\chi (\zeta )}\psi (\zeta )+V\mathrm{e}^{i\pi \mu _1}\psi (\zeta )\right),$$
(7)
where
$$V(t)=e^{i\pi \mu _1}\dot{Z}(0)$$
(8)
is the velocity of the first string (real in our gauge), and
$$\chi (\zeta )=\dot{\alpha }_2_0^\zeta \frac{\psi (\xi )}{\xi \alpha _2}๐\xi .$$
(9)
Extending (7) to $`y=\mathrm{Im}(\zeta )0`$, we obtain, in the vicinity of the cut,
$`N_\zeta (\zeta )`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left(\overline{\chi (\zeta )}\psi (\zeta )+V\mathrm{e}^{i\pi \mu _1ฯต(y)}\psi (\zeta )\right)`$ (10)
$`=`$ $`{\displaystyle \frac{1}{2}}\left(\overline{\chi (\zeta )}+V\mathrm{cos}(\pi \mu _1)\right)\psi (\zeta ){\displaystyle \frac{i}{2}}\beta \psi (\zeta )ฯต(y),`$
where $`ฯต(y)`$ is the sign function, and $`\beta =V\mathrm{sin}(\pi \mu _1)`$. So, while the lapse $`N`$ and the two-metric $`h_{ij}`$ are by construction single-valued and thus continuous, the shift $`N_\zeta `$ is not continuous across the cut, except in the static case $`V(t)=0`$. A priori this discontinuity will lead to delta and gradient of delta contributions to the Einstein tensor. We shall see that the gradient of delta contribution vanishes identically. However the delta contribution remains, so that generically there must be a matter source along the cut. We shall now show that this delta contribution vanishes, iff the motion is geodesic, $`\dot{\beta }(t)=0`$. The Riemann tensor (completely determined in 2+1 dimensions by the Einstein tensor) then vanishes outside the point sources $`\zeta =\alpha _i(t)`$.
For this purpose we use the ($`1+2`$)-dimensional vacuum Einstein equations written in the ADM form
$``$ $``$ $`h^{1/2}R(h)+h^{1/2}\left(\pi ^{ij}\pi _{ij}\pi ^2\right)=0,`$ (11)
$`^i`$ $``$ $`2\pi _{}^{ij}{}_{;j}{}^{}=0`$ (12)
(constraint equations), and
$`_th_{ij}`$ $`=`$ $`2h^{1/2}\left(\pi _{ij}h_{ij}\pi \right)+N_{i;j}+N_{j;i},`$ (13)
$`_t\pi ^{ij}`$ $`=`$ $`{\displaystyle \frac{1}{2}}h^{1/2}\left(h^{ij}\left(\pi ^{kl}\pi _{kl}\pi ^2\right)4\left(\pi ^{ik}\pi _{}^{j}{}_{k}{}^{}\pi ^{ij}\pi \right)\right)`$ (14)
$`+(\pi ^{ij}N^k)_{;k}N_{}^{i}{}_{;k}{}^{}\pi ^{kj}N_{}^{j}{}_{;k}{}^{}\pi ^{ki}`$
(evolution equations), where $`h_{ij}`$ is the two-dimensional metric, $`N=1`$, $`\pi =\pi _{}^{k}{}_{k}{}^{}`$, and $`\pi ^{ij}`$ are the momenta conjugate to the $`h_{ij}`$, related to the extrinsic curvature by
$$\pi _{ij}=h^{1/2}\left(Kh_{ij}K_{ij}\right).$$
(15)
Using the result (10) and the non-vanishing Christoffel symbols $`\mathrm{\Gamma }_{\zeta \zeta }^\zeta =\psi _{,\zeta }/\psi `$ (together with its complex conjugate), we obtain in terms of the real spatial coordinates $`x,y`$
$$\pi ^{xx}=\beta \psi ^1(x)\delta (y),\pi ^{xy}=0,\pi ^{yy}=0.$$
(16)
The component $`\pi ^{xx}`$ vanishes only in the static case ($`\beta =0`$), reflecting the fact that in this case the two-dimensional metric (2) is flat outside the conical singularities $`\zeta =\alpha _i`$.
Inserting this result into the constraint equations, we find that the Hamiltonian constraint (11) is identically satisfied (using the fact that $`R(h)=0`$ for the metric (2)), while the momentum constraint (12) is satisfied on account of $`_y(|\psi |^2)|_{y=0}=0`$. The evolution equations (13) have already been used to compute the momenta $`\pi ^{ij}`$. There remain only the evolution equations (14), which may be rewritten
$$_t\pi ^{ij}=(\pi ^{ij}N^k)_{,k}N_{}^{i}{}_{,k}{}^{}\pi ^{kj}N_{}^{j}{}_{,k}{}^{}\pi ^{ki}.$$
(17)
Both sides of the ($`yy`$) equation vanish identically. The ($`xy`$) equation is satisfied on account of $`_xN_y|_{y=0}=0`$. There only remains the potentially dangerous ($`xx`$) equation. For its left-hand side we obtain
$$_t\pi ^{xx}=\left(\dot{\beta }+\frac{\beta \mu _2\dot{\alpha }_2}{x\alpha _2}\right)\psi ^1(x)\delta (y).$$
(18)
For the right-hand side, using
$$_xN^x=\frac{\mu _2\dot{\alpha }_2}{x\alpha _2}\psi ^1\psi _{,x}N^x\text{for}y=0,$$
(19)
we find
$$N^x_x\pi ^{xx}\pi ^{xx}_xN^x=\frac{\beta \mu _2\dot{\alpha }_2}{x\alpha _2}\psi ^1(x)\delta (y).$$
(20)
Comparison of (18) and (20) shows that the ($`xx`$) component of the evolution equation (17) is satisfied iff
$$\dot{\beta }=0.$$
(21)
The argument can be generalized to show that time derivatives
$$\dot{Z}[t,z,\zeta =\alpha _i(t,z)]=Z_{,t}(\alpha _i)+Z_{,\zeta }(\alpha _i)\dot{\alpha _i}=\mathrm{const}.$$
(22)
for each string, i.e. the strings must move with constant velocities. The extension to the case of non-parallel strings is straightforward and leads to the conclusion that the derivatives with respect to $`z`$ must satisfy
$$Z^{}[t,z,\zeta =\alpha _i(t,z)]=\mathrm{const}.,$$
(23)
meaning that the strings are straight. Therefore the complex trajectories $`Z[\alpha _i]`$ must be linear functions of $`t`$ and $`z`$, or, in geometric terms, the world-sheets of the strings must be totally geodesic submanifolds. We have shown here that this well-known property of self-gravitating cosmic strings vickers ; UHIM is the necessary and sufficient condition for the LG metric to represent a system of crossed straight cosmic strings moving in otherwise empty spacetime.
D.G. thanks LAPTH (Annecy) for hospitality while this note was being written. His work was also supported in part by the RFBR grant 02-04-16949.
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# A macroscopic persistent current generation by merons in a spin density wave ordered state
## Acknowledgment
The author acknowledges helpful comments from S. Sugano and Y. Takada. He is also grateful to A. Fujimori for useful information.
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# Explicit formulas for Hecke operators on cusp forms, Dedekind symbols and period polynomials
## 1. Introduction
Throughout the paper, we assume that $`w`$ is an even positive integer, and we use the following notation:
$`\mathrm{\Gamma }`$ $`:=SL_2()\text{ (the full modular group)},`$
$`S_{w+2}`$ $`:=\text{the vector space of cusp forms on }\mathrm{\Gamma }\text{ with weight }w+2,`$
$`S_{w+2}^{}`$ $`:=\mathrm{Hom}_{}(S_{w+2},)\text{ (the dual space of }S_{w+2}\text{)},`$
$`d_w`$ $`:=\{\begin{array}{cc}\frac{w+2}{12}1\hfill & \mathrm{if}w0(mod12)\hfill \\ \frac{w+2}{12}\hfill & \mathrm{if}w0(mod12)\hfill \end{array}`$
where $`x`$ denotes the greatest integer not exceeding $`x`$.
One of the main reasons to study cusp forms stems from the fact that Fourier coefficients of the forms are interesting from arithmetic view point. Furthermore, periods of cusp forms (or values of L-series associated with the forms) play important roles in number theory. A large number of papers have discussed periods and period polynomials of cusp forms (). One of the basic results is, perhaps, the result due to Eichler and Shimura ().
Periods of cusp forms in $`S_{w+2}`$ can be regarded as elements of $`S_{w+2}^{}`$. The Eichler-Shimura isomorphism theorem asserts that odd (or even) periods span $`S_{w+2}^{}`$ (). However, periods are not linearly independent. In fact, they satisfy the so-called Eichler-Shimura relations (). This leads us to a natural question: which periods would form a basis of $`S_{w+2}^{}`$. The first goal of this paper is to determine odd periods which constitute a basis of $`S_{w+2}^{}`$ (Theorem 2.2). Passing to the dual space $`S_{w+2}`$, this gives rise to a new basis for $`S_{w+2}`$.
Next we consider three spacesโ$`S_{w+2}`$, the space of even Dedekind symbols of weight $`w`$ with polynomial reciprocity laws, and the space of even period polynomials of degree $`w`$. It is known that these three spaces are isomorphic modulo trivial elements (). Through these isomorphisms, we can construct bases for these spaces explicitly (Theorems 2.3, 2.6 and 2.7), starting with the basis for $`S_{w+2}^{}`$ determined in Theorem 2.2. Furthermore, it is also known that these three spaces are equipped with compatible actions of Hecke operators (). We will subsequently find explicit forms for the actions of Hecke operators on the elements of these three spaces (Theorem 2.8).
As the final goal, we obtain explicit formulas for Hecke operators on $`S_{w+2}`$ (Theorem 2.9), which seem quite different from the celebrated Eichler-Selberg trace formula. We will do this by obtaining matrices which represent the Hecke operators $`T_m(m=1,2,\mathrm{})`$ on $`S_{w+2}`$ as well as their characteristic polynomials by means of Bernoulli numbers $`B_k`$ and divisor functions $`\sigma _k(n)`$.
It might be interesting to compare our result with the Eichler-Selberg trace formula (,\[13, p. 48\]). Their formula gives traces of Hecke operators in terms of class numbers of imaginary quadratic fields. Their method is based on integrating a kernel function for Hecke operator. Our approach is different from their method. In fact, our formulas give matrices representing the Hecke operators in terms of Bernoulli numbers and divisor functions, as well as their characteristic polynomials. Our method is based on representing the Hecke operator on $`S_{w+2}`$ as Hecke operator on the space of Dedekind symbols, and then as Hecke operator on the space of period polynomials. It should be noted that our argument depends on the fact that $`dimS_{w+2}=d_w`$; while the Eichler-Selberg trace formula gives this fact as a consequence.
This paper is organized as follows. In Section 2, we give precise statement of our results. The Sections from 3 to 6 are devoted to the preparation of establishing Theorem 2.2. In Section 7, we will give proofs of Theorems 2.2 and 2.3. In Section 8, we will present proofs of Theorems 2.6 and 2.7. The Sections from 9 to 13 are devoted to studying the Hecke operators on the three spaces, which include a proof of Theorem 2.8. Finally in Section 14, we will prove Theorem 2.9 and append a computer program for obtaining matrices which represent the Hecke operators.
## 2. Statement of results
In this section, we will state our results in more precise form.
Let $`f`$ be an element of $`S_{w+2}`$. We write $`f`$ as a Fourier series
$$f(z)=\underset{l=1}{\overset{\mathrm{}}{}}a_le^{2\pi ilz}.$$
Let $`L(f,s)`$ be the L-series of $`f`$. Namely,
$$L(f,s):=\underset{l=1}{\overset{\mathrm{}}{}}\frac{a_l}{l^s}(\mathrm{}(s)0).$$
Then $`n`$th period of $`f`$, $`r_{w,n}(f)`$, is defined by
$$r_{w,n}(f):=_0^i\mathrm{}f(z)z^n๐z=\frac{n!}{(2\pi i)^{n+1}}L(f,n+1)(n=0,1,\mathrm{},w).$$
Also the period polynomial of $`f`$ in the variable $`X`$ is defined by
$$r(f)(X):=_0^i\mathrm{}f(z)(Xz)^w๐z.$$
It is clear that $`r(f)(X)`$ has the following expression:
$$r(f)(X)=\underset{n=0}{\overset{w}{}}(1)^n\left(\genfrac{}{}{0pt}{}{w}{n}\right)r_{w,wn}(f)X^n.$$
Here and hereafter $`\left(\genfrac{}{}{0pt}{}{w}{n}\right)`$ denotes a binomial coefficient.
Each period $`r_{w,n}`$ can be regarded as a linear map from $`S_{w+2}`$ to $``$, that is,
$$r_{w,n}S_{w+2}^{}.$$
Here we recall the result of Eichler and Shimura ():
###### Theorem 2.1 (Eichler-Shimura).
The maps
$`r_w^+:S_{w+2}`$ $`^{(w+2)/2}`$
$`f`$ $`(r_{w,0}(f),r_{w,2}(f),\mathrm{},r_{w,w}(f))`$
and
$`r_w^{}:S_{w+2}`$ $`^{w/2}`$
$`f`$ $`(r_{w,1}(f),r_{w,3}(f),\mathrm{},r_{w,w1}(f))`$
are both injective.
In other words,
1. the even periods
$$r_{w,0},r_{w,2},\mathrm{},r_{w,w}$$
span the vector space $`S_{w+2}^{}`$;
2. the odd periods
$$r_{w,1},r_{w,3},\mathrm{},r_{w,w1}$$
also span $`S_{w+2}^{}`$.
However, these periods are not linearly independent. In fact, they satisfy the so-called Eichler-Shimura relations (): For $`n=0,1,\mathrm{},w`$, it holds that
(ES1)
$$r_{w,n}(f)+(1)^nr_{w,wn}(f)=0,$$
(ES2)
$$(1)^nr_{w,n}(f)+\underset{\begin{array}{c}0mn\\ m0(mod2)\end{array}}{}\left(\genfrac{}{}{0pt}{}{n}{m}\right)r_{w,wn+m}(f)+\underset{\begin{array}{c}0mwn\\ mn(mod2)\end{array}}{}\left(\genfrac{}{}{0pt}{}{wn}{m}\right)r_{w,m}(f)=0,$$
(ES3)
$$\underset{\begin{array}{c}1mn\\ m1(mod2)\end{array}}{}\left(\genfrac{}{}{0pt}{}{n}{m}\right)r_{w,wn+m}(f)+\underset{\begin{array}{c}0mwn\\ mn(mod2)\end{array}}{}\left(\genfrac{}{}{0pt}{}{wn}{m}\right)r_{w,m}(f)=0.$$
This leads us to a natural question. Which periods are linearly independent? Or more strictly, which periods form a basis for $`S_{w+2}^{}`$? Our first result provides an answer to this question in the case of odd periods.
Throughout this paper we adopt the following notation and convention:
###### Definition 2.1.
1. For an integer $`i`$ such that $`1id_w`$, let $`4i\pm 1`$ stand for $`4i+1`$ or $`4i1`$ according as $`w0(mod4)`$ or $`w2(mod4)`$. Namely
$$4i\pm 1:=\{\begin{array}{cc}4i+1\hfill & \mathrm{if}w0(mod4)\hfill \\ 4i1\hfill & \mathrm{if}w2(mod4).\hfill \end{array}$$
2. For an integer $`n`$ with $`0nw`$, let $`\stackrel{~}{n}`$ stand for $`wn`$:
$$\stackrel{~}{n}:=wn.$$
3. For an integer $`k`$, a divisor function $`\sigma _k`$ is defined by
$$\sigma _k(n):=\underset{\begin{array}{c}ad=n\\ a>0\end{array}}{}a^k,(n^+).$$
We recall the well-known fact (see e.g. \[1, p. 133\]) that
$$dimS_{w+2}=d_w.$$
Now we can state our first result:
###### Theorem 2.2.
$$\{r_{w,4i\pm 1}|i=1,2,\mathrm{},d_w\}$$
form a basis for $`S_{w+2}^{}`$.
In other words, $`\{r_{w,4i\pm 1}|i=1,2,\mathrm{},d_w\}`$ are linearly independent over $``$, and thus other periods are linear combinations of $`\{r_{w,4i\pm 1}|i=1,2,\mathrm{},d_w\}`$. Furthermore, other odd periods are linear combinations not only over $``$, but over $``$ (confer the proof of Theorem 2.2).
Next we will display a basis for $`S_{w+2}`$. For $`f,gS_{w+2}`$, let $`(f,g)`$ denote the Petersson inner product. Then there is a cusp form $`R_{w,n}`$, which is characterized by the formula:
$$r_{w,n}(f)=(R_{w,n},f)\text{ for any }fS_{w+2}.$$
Explicit form of $`R_{w,n}`$, as a Poincarรฉ series, was given ():
$$R_{w,n}(z):=c_{w,n}^1\underset{\begin{array}{c}\left[\begin{array}{cc}a& b\\ c& d\end{array}\right]\mathrm{\Gamma }\end{array}}{}\frac{1}{(az+b)^{n+1}(cz+d)^{\stackrel{~}{n}+1}},c_{w,n}=(1)^{n+1}2\pi i\left(\genfrac{}{}{0pt}{}{w}{n}\right).$$
In general, the Poincarรฉ series $`R_{w,n}`$ is expected to have transcendental Fourier coefficients.
Passing to the dual space, we obtain a basis of $`S_{w+2}`$.
###### Theorem 2.3.
$$\{R_{w,4i\pm 1}|i=1,2,\mathrm{},d_w\}$$
form a basis for $`S_{w+2}`$.
Several bases are known for $`S_{w+2}`$ (). We believe that the above basis is the first one whose even periods can be described explicitly (\[11, Theorem $`1^{}`$\]). From this fact, we can obtain bases for the spaces of even period polynomials, as well as bases for the spaces of even Dedekind symbols with polynomial reciprocity laws. We now recall the relationship between cusp forms, Dedekind symbols and period polynomials.
A complex-valued function $`E`$ on $`^+\times `$ is called a weighted Dedekind symbol of weight $`w`$ if it satisfies the following two conditions (confer to ):
$$E(h,k)=E(h,k+h),E(ch,ck)=c^wE(h,k)$$
for any $`(h,k)^+\times `$ and $`c^+`$.
Moreover, a weighted Dedekind symbol $`E`$ is said to be even (resp. odd) if $`E`$ satisfies
$$E(h,k)=E(h,k)\text{ }(\text{resp. }E(h,k)=E(h,k)).$$
There are two rather trivial Dedekind symbols $`G_w`$ and $`F_w`$ which are defined by
$$G_w(h,k):=\left\{\mathrm{gcd}(h,k)\right\}^w\text{ and }F_w(h,k):=h^w$$
for any $`(h,k)^+\times `$.
A symbol $`E`$ is determined by its reciprocity law
$$E(h,k)E(k,h)=S(h,k)$$
up to addition of scalar multiples of $`G_w`$. Here $`S`$ is a complex-valued function defined on $`^+\times ^+`$.
Next we would like to demonstrate the relationship between cusp forms, weighted Dedekind symbols, and period polynomials. We need the following notation:
$`๐ฒ_w`$ $`:=\{E|E\text{ is a Dedekind symbol of weight }w\},`$
$`๐ฒ_w^{}`$ $`:=\left\{E๐ฒ_w\right|E\text{ is odd }\},`$
$`๐ฒ_w^+`$ $`:=\left\{E๐ฒ_w\right|E\text{ is even }\},`$
$`_w`$ $`:=\{E๐ฒ_w|E(h,k)E(k,h)\text{ is a homogeneous polynomial in }h\text{ and }k`$
$`\text{ }\text{ }\text{ }\text{of degree }w\},`$
$`_w^{}`$ $`:=\left\{E_w\right|E\text{ is odd }\},`$
$`_w^+`$ $`:=\left\{E_w\right|E\text{ is even }\},`$
$`๐ฐ_w`$ $`:=\{g|g\text{ is a homogeneous polynomial in }h\text{ and }k\text{ of degree }w`$
$`\text{ }\text{ }\text{satisfying }g(h+k,k)+g(h,h+k)=g(h,k)\text{ and }g(1,1)=0\}`$
(an element of $`๐ฐ_w`$ is essentially a period polynomial modulo $`h^wk^w`$ ),
$`๐ฐ_w^{}`$ $`:=\left\{g๐ฐ_w\right|\text{ }g\text{ is an odd polynomial, i.e., }g(h,k)=g(h,k)\},`$
$`๐ฐ_w^+`$ $`:=\left\{g๐ฐ_w\right|\text{ }g\text{ is an even polynomial, i.e., }g(h,k)=g(h,k)\}.`$
For a cusp form $`fS_{w+2}`$ and $`(h,k)^+\times `$, we define $`E_f`$ and $`E_f^\pm `$ by
(2.1)
$$E_f(h,k):=_{k/h}^i\mathrm{}f(z)(hzk)^w๐z,E_f^\pm (h,k):=\frac{1}{2}\{E_f(h,k)\pm E_f(h,k)\}.$$
Then we can show $`E_f`$ is a Dedekind symbols of weight $`w`$, and we can define maps
$$\alpha _{w+2}:S_{w+2}๐ฒ_w,\alpha _{w+2}^\pm :S_{w+2}๐ฒ_w^\pm $$
by
$$\alpha _{w+2}(f)=E_f,\alpha _{w+2}^\pm (f)=E_f^\pm .$$
Furthermore, we know that $`E_f`$ and $`E_f^\pm `$ have polynomial reciprocity laws, that is, $`E_f_w`$ and $`E_f^\pm _w^\pm `$. Hence we have the restricted maps
$$\alpha _{w+2}^\pm :S_{w+2}_w^\pm $$
(to ease the notation, we use the same notation $`\alpha _{w+2}^\pm `$ for the restricted maps). Then we have the following:
###### Theorem 2.4 (\[9, Theorem 1.1\]).
The map
$$\alpha _{w+2}^{}:S_{w+2}_w^{}$$
is an isomorphism $`(`$between vector spaces$`)`$, and the map
$$\alpha _{w+2}^+:S_{w+2}_w^+$$
is a monomorphism such that its image $`\alpha _{w+2}^+(S_{w+2})`$ is a subspace of $`_w^+`$ of codimension two, and that $`\alpha _{w+2}^+(S_{w+2})`$, $`F_w`$ and $`G_w`$ span $`_w^+`$.
Next we will see how weighted Dedekind symbols are linked to period polynomials. For a weighted Dedekind symbol $`E`$ and $`(h,k)^+\times ^+`$, let $`\beta _w(E)`$ be defined by
$$\beta _w(E)(h,k)=E(h,k)E(k,h).$$
In the case of Dedekind symbol $`E_f`$ associated with a cusp form $`f`$, $`\beta _w(E_f)`$ has the following expression:
(2.2)
$$\beta _w(E_f)(h,k)=_0^i\mathrm{}f(z)(hzk)^w๐z.$$
Note that the right hand side of (2.2) is nothing but a homogeneous period polynomial of $`f`$.
For $`E_w`$, we know that $`\beta _w(E)๐ฐ_w`$. Thus, we have a homomorphism
$$\beta _w:_w๐ฐ_w.$$
Then we see that $`\beta _w`$ is almost isomorphism in the following sense:
###### Theorem 2.5 (\[9, Theorem 1.2\]).
The homomorphism $`\beta _w:_w๐ฐ_w`$ is an epimorphism with $`\beta _w(_w^\pm )=๐ฐ_w^\pm `$, and $`\mathrm{ker}\beta _w`$ is one dimensional subspace of $`_w`$ spanned by $`G_w`$.
In particular, the restricted map
$$\beta _w^{}:_w^{}๐ฐ_w^{}$$
is an isomorphism, and
$$\beta _w^+:_w^+๐ฐ_w^+$$
is an epimorphism where $`\mathrm{ker}\beta _w^+`$ is one dimensional subspace of $`_w^+`$ spanned by $`G_w`$.
Here we examine the composed maps
$$\beta _w^\pm \alpha _{w+2}^\pm :S_{w+2}_w^\pm ๐ฐ_w^\pm .$$
Since $`\beta _w^\pm (E_f^\pm )(h,k)=\beta _w^\pm \alpha _{w+2}^\pm (f)(h,k)`$ is the homogeneous period polynomial of $`f`$, the composed maps
$$\beta _w^{}\alpha _{w+2}^{}:S_{w+2}_w^{}๐ฐ_w^{}$$
and
$$\beta _w^+\alpha _{w+2}^+:S_{w+2}_w^+๐ฐ_w^+$$
can be identified with the Eichler-Shimura isomorphisms (refer to \[11, p. 200\], \[7, Theorem 7.3\]). In fact, $`\beta _w^{}\alpha _{w+2}^{}`$ is an isomorphism, and $`\beta _w^+\alpha _{w+2}^+`$ is an monomorphism such that the image $`\beta _w^+\alpha _{w+2}^+(S_{w+2})`$ and $`h^wk^w`$ span $`๐ฐ_w^+`$.
These facts may be summarized in the following commutative diagram:
Using these correspondences, we can obtain bases for the spaces of even period polynomials, and bases for the spaces of even Dedekind symbols with polynomial reciprocity laws. We need the following notation (refer to ).
For an integer $`n`$ such that $`0<n<w`$, a polynomial $`S_{w,n}`$, in $`h`$ and $`k`$, is defined by
$`S_{w,n}(h,k):=`$ $`(1)^n{\displaystyle \frac{B_{n+1}(\frac{k}{h})h^w}{n+1}}+{\displaystyle \frac{B_{n+1}(\frac{h}{k})k^w}{n+1}}{\displaystyle \frac{B_{\stackrel{~}{n}+1}(\frac{k}{h})h^w}{\stackrel{~}{n}+1}}(1)^n{\displaystyle \frac{B_{\stackrel{~}{n}+1}(\frac{h}{k})k^w}{\stackrel{~}{n}+1}}`$
$`\text{ }+\{\begin{array}{cc}\frac{w+2}{B_{w+2}}\frac{B_{n+1}}{n+1}\frac{B_{\stackrel{~}{n}+1}}{\stackrel{~}{n}+1}(h^wk^w)\hfill & \mathrm{if}n1(mod2)\hfill \\ 0\hfill & \mathrm{if}n0(mod2).\hfill \end{array}`$
For any $`(h,k)^+\times `$, and for an integer $`n`$ such that $`0<n<w`$, a Dedekind symbol $`E_{w,n}:^+\times `$ is defined by
$`E_{w,n}(h,k):=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{\begin{array}{c}\left[\begin{array}{cc}a& b\\ c& d\end{array}\right]\mathrm{\Gamma }\\ ac0\\ (k/h+b/a)(k/h+d/c)<0\end{array}}{}}\mathrm{sgn}\left({\displaystyle \frac{k}{h}}+{\displaystyle \frac{b}{a}}\right)(ak+bh)^{\stackrel{~}{n}}(ck+dh)^n`$
$`\text{ }+\left\{(1)^n{\displaystyle \frac{\overline{B}_{n+1}(\frac{k}{h})h^w}{n+1}}{\displaystyle \frac{\overline{B}_{\stackrel{~}{n}+1}(\frac{k}{h})h^w}{\stackrel{~}{n}+1}}\right\}`$
$`\text{ }+\{\begin{array}{cc}\frac{w+2}{B_{w+2}}\frac{B_{n+1}}{n+1}\frac{B_{\stackrel{~}{n}+1}}{\stackrel{~}{n}+1}h^w\hfill & \mathrm{if}n1(mod2)\hfill \\ 0\hfill & \mathrm{if}n0(mod2).\hfill \end{array}`$
Here and hereafter, $`B_m(x)`$ (resp. $`B_m`$) denotes the $`m`$th Bernoulli polynomial (resp. number), and $`\overline{B}_m(x)`$ denotes the $`m`$th Bernoulli function. That is, $`B_m(x)`$ is defined by
$$\frac{te^{xt}}{e^t1}=\underset{m=0}{\overset{\mathrm{}}{}}B_m(x)\frac{t^m}{m!},$$
and $`\overline{B}_m(x)`$ is defined as the periodic function which coincides with $`B_m(x)`$ on $`[0,1)`$. Moreover, $`\mathrm{sgn}(x)`$ denotes the sign of $`x`$.
It is shown (\[11, Theorem $`1^{}`$\], ) that, for $`n`$ odd$`,\beta _w^+\alpha _{w+2}^+`$ maps $`c_{w,n}R_{w,n}`$ to $`S_{w,n}`$:
$$\beta _w^+\alpha _{w+2}^+:c_{w,n}R_{w,n}S_{w,n}(n\text{ odd}).$$
Furthermore, we will show in Lemma 8.1 below that, for $`n`$ odd$`,\alpha _{w+2}^+`$ maps $`c_{w,n}R_{w,n}`$ to $`E_{w,n}`$:
$$\alpha _{w+2}^+:c_{w,n}R_{w,n}E_{w,n}(n\text{ odd}).$$
From these facts and Theorems 2.3, 2.4 and 2.5, we obtain the following two theorems:
###### Theorem 2.6.
$$\{S_{w,4i\pm 1}(h,k)|i=1,2,\mathrm{},d_w\}\{h^wk^w\}$$
form a basis for $`๐ฐ_w^+`$.
###### Theorem 2.7.
$$\{E_{w,4i\pm 1}|i=1,2,\mathrm{},d_w\}\{F_w\}\{G_w\}$$
form a basis for $`_w^+`$.
The latter half of this paper is devoted to obtaining matrices which represent the Hecke operators $`T_m(m=1,2,\mathrm{})`$ on $`S_{w+2}`$ as well as their characteristic polynomials. For this purpose, first we discuss Hecke operators on the three spaces in Diagram ES. Manin (see also \[11, p. 202\]) and Zagier proved that there are well-defined Hecke operators (also denoted by $`T_m`$) on the spaces of period polynomials which are compatible with the Eichler-Shimura isomorphism:
(2.3)
$$\begin{array}{ccc}S_{w+2}& \stackrel{\beta _w^\pm \alpha _{w+2}^\pm }{}& ๐ฐ_w^\pm \\ T_m& & T_m& & \\ S_{w+2}& \stackrel{\beta _w^\pm \alpha _{w+2}^\pm }{}& ๐ฐ_w^\pm .\end{array}$$
Furthermore, in , we introduced Hecke operators on the space of Dedekind symbols by the following formula:
$$(T_mE)(h,k):=\underset{\begin{array}{c}ad=m\\ d>0\end{array}}{}\underset{b(\mathrm{mod}d)}{}E(dh,ak+bh).$$
It was proved that Hecke operators on the spaces of Dedekind symbols are compatible with Hecke operators on the spaces of cusp forms (). As a consequence, we have the following commutative diagram:
(2.4)
$$\begin{array}{ccc}S_{w+2}& \stackrel{\alpha _{w+2}^\pm }{}& _w^\pm \\ T_m& & T_m& & \\ S_{w+2}& \stackrel{\alpha _{w+2}^\pm }{}& _w^\pm .\end{array}$$
(To ease the notation, we will use the same notation $`T_m`$ for the Hecke operators on $`S_{w+2}`$, $`_w^\pm `$ and $`๐ฐ_w^\pm `$.)
Now we need the following definitions to describe the actions of Hecke operators on $`R_{w,n}`$, $`E_{w,n}`$ and $`S_{w,n}`$:
###### Definition 2.2.
1. For a positive integer $`m`$,
$$H_m:=\{\left[\begin{array}{cc}a& b\\ c& d\end{array}\right]|adbc=m;a,b,c,d\};$$
2. For positive integers $`m`$ and $`n`$ such that $`0<n<w`$,
$$R_{w,n}^m(z):=m^{w+1}c_{w,n}^1\underset{\begin{array}{c}\left[\begin{array}{cc}a& b\\ c& d\end{array}\right]H_m\end{array}}{}\frac{1}{(az+b)^{n+1}(cz+d)^{\stackrel{~}{n}+1}};$$
3. For positive integers $`m`$ and $`n`$ such that $`0<n<w`$, we define a map $`E_{w,n}^m:^+\times `$ by
$`E_{w,n}^m(h,k):=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{\begin{array}{c}\left[\begin{array}{cc}a& b\\ c& d\end{array}\right]H_m\\ ac0\\ (k/h+b/a)(k/h+d/c)<0\end{array}}{}}\mathrm{sgn}\left({\displaystyle \frac{k}{h}}+{\displaystyle \frac{b}{a}}\right)(ak+bh)^{\stackrel{~}{n}}(ck+dh)^n`$
$`\text{ }+{\displaystyle \underset{\begin{array}{c}ad=m\\ a>0\end{array}}{}}\left\{(1)^nd^{\stackrel{~}{n}}{\displaystyle \frac{\overline{B}_{n+1}(\frac{ak}{h})h^w}{n+1}}d^n{\displaystyle \frac{\overline{B}_{\stackrel{~}{n}+1}(\frac{ak}{h})h^w}{\stackrel{~}{n}+1}}\right\}`$
$`\text{ }+\{\begin{array}{cc}\sigma _{w+1}(m)\frac{w+2}{B_{w+2}}\frac{B_{n+1}}{n+1}\frac{B_{\stackrel{~}{n}+1}}{\stackrel{~}{n}+1}h^w\hfill & \mathrm{if}n1(mod2)\hfill \\ 0\hfill & \mathrm{if}n0(mod2);\hfill \end{array}`$
4. For positive integers $`m`$ and $`n`$ such that $`0<n<w`$, we define a polynomial $`S_{w,n}^m`$ in $`h`$ and $`k`$ by
$`S_{w,n}^m(h,`$ $`k):={\displaystyle \frac{1}{2}}{\displaystyle }_{\begin{array}{c}\left[\begin{array}{cc}a& b\\ c& d\end{array}\right]H_m\\ abcd<0\end{array}}\mathrm{sgn}(ab)(ak+bh)^{\stackrel{~}{n}}(ck+dh)^n`$
$`\text{ }+{\displaystyle \underset{\begin{array}{c}ad=m\\ a>0\end{array}}{}}d^{\stackrel{~}{n}}\left\{(1)^n{\displaystyle \frac{B_{n+1}(\frac{ak}{h})h^w}{n+1}}+{\displaystyle \frac{B_{n+1}(\frac{ah}{k})k^w}{n+1}}\right\}`$
$`\text{ }{\displaystyle \underset{\begin{array}{c}ad=m\\ a>0\end{array}}{}}d^n\left\{{\displaystyle \frac{B_{\stackrel{~}{n}+1}(\frac{ak}{h})h^w}{\stackrel{~}{n}+1}}+(1)^n{\displaystyle \frac{B_{\stackrel{~}{n}+1}(\frac{ah}{k})k^w}{\stackrel{~}{n}+1}}\right\}`$
$`\text{ }+\{\begin{array}{cc}\sigma _{w+1}(m)\frac{w+2}{B_{w+2}}\frac{B_{n+1}}{n+1}\frac{B_{\stackrel{~}{n}+1}}{\stackrel{~}{n}+1}(h^wk^w)\hfill & \mathrm{if}n1(mod2)\hfill \\ 0\hfill & \mathrm{if}n0(mod2).\hfill \end{array}`$
Properties of these functions will be studied in the later sections. It is plain that
$$R_{w,n}^1=R_{w,n},E_{w,n}^1=E_{w,n}\text{and }S_{w,n}^1=S_{w,n}.$$
Now we can formulate the actions of Hecke operators on $`R_{w,n}`$, $`E_{w,n}`$ and $`S_{w,n}`$:
###### Theorem 2.8.
The actions of the Hecke operators $`T_m`$ on $`R_{w,n}`$, $`E_{w,n}`$ and $`S_{w,n}`$ are expressed as follows:
$$T_m(R_{w,n})=R_{w,n}^m,T_m(E_{w,n})=E_{w,n}^m\mathrm{and}T_m(S_{w,n})=S_{w,n}^m.$$
Finally, as an application of Theorems 2.3 and 2.8, we will give explicit formulas for the Hecke operators on the spaces of cusp forms. Let
$$f(h,k)=\underset{\nu =0}{\overset{w}{}}a_\nu h^\nu k^{w\nu }\text{ and }g(h,k)=\underset{\nu =0}{\overset{w}{}}b_\nu h^\nu k^{w\nu }$$
be homogeneous polynomials in $`h`$ and $`k`$ with degree $`w`$. Then their inner product $`f,g`$ is defined by
$$f,g:=\underset{\nu =0}{\overset{w}{}}a_\nu \overline{b}_\nu $$
where $`\overline{b}_\nu `$ denotes the complex conjugate of $`b_\nu `$.
Under this notation we obtain the following result.
###### Theorem 2.9.
1. Let $`m`$ be a positive integer, and let $`๐_m`$ be the matrix representing the Hecke operator
$$T_m:S_{w+2}S_{w+2}$$
with respect to the basis
$$c_{w,4i\pm 1}R_{w,4i\pm 1}(i=1,2,\mathrm{},d_w).$$
Let $`๐_1`$ and $`๐_2`$ be matrices defined by
$$๐_1:=\left[\begin{array}{c}S_{w,4i\pm 1},S_{w,4j\pm 1}\end{array}\right](i,j=1,2,\mathrm{},d_w),$$
$$๐_2:=\left[\begin{array}{c}S_{w,4i\pm 1},S_{w,4j\pm 1}^m\end{array}\right](i,j=1,2,\mathrm{},d_w).$$
Then $`๐_m`$ can be expressed as
$$๐_m=๐_1^1๐_2.$$
2. Let $`n`$ be an odd integer with $`0<n<w`$. Then $`S_{w,n}^m`$ can be expressed explicitly as a polynomial in $`h`$ and $`k`$ by the following formula:
$`S_{w,n}^m(h,k)=`$ $`2{\displaystyle \underset{\begin{array}{c}\nu =0\\ \nu \mathrm{even}\end{array}}{\overset{w}{}}}\{{\displaystyle \underset{\mu =1}{\overset{m1}{}}}{\displaystyle \underset{\lambda =\mathrm{max}(0,\nu \stackrel{~}{n})}{\overset{\mathrm{min}(n,\nu )}{}}}\mu ^\lambda (\mu m)^{n\lambda }\left({\displaystyle \genfrac{}{}{0pt}{}{\stackrel{~}{n}}{\nu \lambda }}\right)\left({\displaystyle \genfrac{}{}{0pt}{}{n}{\lambda }}\right)\times `$
$`\text{ }\text{ }\sigma _{\stackrel{~}{n}\nu }(\mu )\sigma _{\nu n}(m\mu )\}h^\nu k^{w\nu }`$
$`+{\displaystyle \frac{(1)^nm^{\stackrel{~}{n}}}{n+1}}{\displaystyle \underset{\nu =\stackrel{~}{n}1}{\overset{w}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{n+1}{\nu \stackrel{~}{n}+1}}\right)B_{\nu \stackrel{~}{n}+1}\sigma _{n\nu }(m)h^\nu k^{w\nu }`$
$`{\displaystyle \frac{m^n}{\stackrel{~}{n}+1}}{\displaystyle \underset{\nu =n1}{\overset{w}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{\stackrel{~}{n}+1}{\nu n+1}}\right)B_{\nu n+1}\sigma _{\stackrel{~}{n}\nu }(m)h^\nu k^{w\nu }`$
$`+{\displaystyle \frac{m^{\stackrel{~}{n}}}{n+1}}{\displaystyle \underset{\nu =0}{\overset{n+1}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{n+1}{n\nu +1}}\right)B_{n\nu +1}\sigma _{\nu \stackrel{~}{n}}(m)h^\nu k^{w\nu }`$
$`{\displaystyle \frac{(1)^nm^n}{\stackrel{~}{n}+1}}{\displaystyle \underset{\nu =0}{\overset{\stackrel{~}{n}+1}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{\stackrel{~}{n}+1}{\stackrel{~}{n}\nu +1}}\right)B_{\stackrel{~}{n}\nu +1}\sigma _{\nu n}(m)h^\nu k^{w\nu }`$
$`+\sigma _{w+1}(m){\displaystyle \frac{w+2}{B_{w+2}}}{\displaystyle \frac{B_{n+1}}{n+1}}{\displaystyle \frac{B_{\stackrel{~}{n}+1}}{\stackrel{~}{n}+1}}(h^wk^w).`$
In particular, setting $`m=1`$, we have
$`S_{w,n}(h,k)=`$ $`{\displaystyle \frac{(1)^n}{n+1}}{\displaystyle \underset{\nu =\stackrel{~}{n}1}{\overset{w}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{n+1}{\nu \stackrel{~}{n}+1}}\right)B_{\nu \stackrel{~}{n}+1}h^\nu k^{w\nu }`$
$`{\displaystyle \frac{1}{\stackrel{~}{n}+1}}{\displaystyle \underset{\nu =n1}{\overset{w}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{\stackrel{~}{n}+1}{\nu n+1}}\right)B_{\nu n+1}h^\nu k^{w\nu }`$
$`+{\displaystyle \frac{1}{n+1}}{\displaystyle \underset{\nu =0}{\overset{n+1}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{n+1}{n\nu +1}}\right)B_{n\nu +1}h^\nu k^{w\nu }`$
$`{\displaystyle \frac{(1)^n}{\stackrel{~}{n}+1}}{\displaystyle \underset{\nu =0}{\overset{\stackrel{~}{n}+1}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{\stackrel{~}{n}+1}{\stackrel{~}{n}\nu +1}}\right)B_{\stackrel{~}{n}\nu +1}h^\nu k^{w\nu }`$
$`+{\displaystyle \frac{w+2}{B_{w+2}}}{\displaystyle \frac{B_{n+1}}{n+1}}{\displaystyle \frac{B_{\stackrel{~}{n}+1}}{\stackrel{~}{n}+1}}(h^wk^w).`$
Consequently, this shows that we can express the matrix $`๐_m`$ representing the Hecke operator $`T_m`$ explicitly in terms of Bernoulli numbers $`B_k`$ and divisor functions $`\sigma _k(n)`$.
In the last section, we append a computer program for obtaining matrices which represent the Hecke operators, and their characteristic polynomials. The program is a straightforward implementation of Theorem 2.9.
## 3. The Eichler-Shimura relations for modified periods
The Sections from 3 to 7 are devoted to the study of odd periods of cusp forms. We also present proofs of Theorems 2.2 and 2.3 in Section 7.
Our strategy for proving Theorems 2.2 is as follows. We obtain an integral matrix which express the Eichler-Shimura relations for odd periods. Then we take the reduction modulo 2 of this matrix. The new matrix has a nice โself-similarโ structure so that we can find linearly independent column vectors. Then we choose odd periods corresponding to the column vectors which turn out to form a basis.
For our purpose, it is more convenient to consider โmodified periodsโ instead of periods themselves.
###### Definition 3.1.
For $`fS_{w+2}`$, we define $`s_{w,n}(f)`$ by
(3.1)
$$s_{w,n}(f):=(1)^n\left(\genfrac{}{}{0pt}{}{w}{n}\right)r_{w,wn}(f),$$
and we call $`s_{w,n}(f)`$ the $`n`$th modified period of $`f`$.
We regard $`s_{w,n}`$ as an element of $`S_{w+2}^{}`$, that is,
$$s_{w,n}S_{w+2}^{}.$$
Then the period polynomial $`r(f)(X)`$ has the following expression:
$$r(f)(X)=\underset{n=0}{\overset{w}{}}(1)^n\left(\genfrac{}{}{0pt}{}{w}{n}\right)r_{w,wn}(f)X^n=\underset{n=0}{\overset{w}{}}s_{w,n}(f)X^n.$$
Using the modified period $`s_{w,n}`$, the Eichler-Shimura relations can be expressed as follows (Kohnen-Zagier \[11, p. 199\]):
(KZ1)
$$s_{w,n}+(1)^ns_{w,wn}=0(0nw),$$
(KZ2)
$$\underset{\begin{array}{c}m=0\\ m\mathrm{even}\end{array}}{\overset{n}{}}\left(\genfrac{}{}{0pt}{}{wm}{wn}\right)s_{w,m}+\underset{\begin{array}{c}m=n+1\\ m\mathrm{even}\end{array}}{\overset{w}{}}\left(\genfrac{}{}{0pt}{}{m}{n}\right)s_{w,m}=0(0nw),$$
(KZ3)
$$\underset{\begin{array}{c}m=0\\ m\mathrm{odd}\end{array}}{\overset{n}{}}\left(\genfrac{}{}{0pt}{}{wm}{wn}\right)s_{w,m}+\underset{\begin{array}{c}m=n+1\\ m\mathrm{odd}\end{array}}{\overset{w}{}}\left(\genfrac{}{}{0pt}{}{m}{n}\right)s_{w,m}=0(0nw).$$
These linear relations are the starting point of our discussions which eventually lead to Theorem 2.2.
Hereafter we adopt the following notation and convention:
1. $$w_2:=\frac{w}{2}\text{ and }w_4:=\frac{w}{4};$$
2. Let $`๐=[x_{ij}]`$ be an $`m\times n`$ matrix. Then, to the end of Section 7, we employ the convention that the index $`i`$ (resp. $`j`$) runs from $`0`$ to $`m1`$ (resp. from $`0`$ to $`n1`$).
We also need the following notation:
###### Definition 3.2.
We define $`t_{w,i}S_{w+2}^{}(i=0,1,\mathrm{},w_2)`$ by
(3.2)
$$t_{w,i}:=\{\begin{array}{cc}0\hfill & \mathrm{if}i=0\hfill \\ s_{w,2i1}\hfill & \mathrm{if}1iw_2.\hfill \end{array}$$
So, $`t_{w,i}(i=1,2,\mathrm{},w_2)`$ are odd modified periods, while $`t_{w,0}`$ is a dummy introduced for technical reasons.
Our first task is to express the relations (KZ1) and (KZ3) for odd modified period in matrix forms. For this, we first introduce the following matrices.
###### Definition 3.3.
1. We define a $`(w_2+1)\times 1`$ matrix $`๐ญ`$ by
$$๐ญ:=\left[\begin{array}{c}t_{w,0}\\ t_{w,1}\\ t_{w,2}\\ \mathrm{}\\ t_{w,w_2}\end{array}\right];$$
2. We define a matrix $`๐=[a_{ij}]`$ $`(`$the matrix whose $`i`$th row and $`j`$th column entry is $`a_{ij}`$; $`i=0,1,\mathrm{},w;j=0,1,\mathrm{},w_2`$$`)`$ by
$$a_{ij}:=\{\begin{array}{cc}0\hfill & \mathrm{if}j=0\hfill \\ \left(\genfrac{}{}{0pt}{}{2j1}{i}\right)\hfill & \mathrm{if}j0,2j>i\hfill \\ \left(\genfrac{}{}{0pt}{}{w2j+1}{wi}\right)\hfill & \mathrm{if}j0,2ji;\hfill \end{array}$$
3. We define a matrix $`๐=[b_{ij}](i=0,1,\mathrm{},w_2;j=0,1,\mathrm{},w_2)`$ by
$$b_{ij}:=\{\begin{array}{cc}0\hfill & \mathrm{if}j=0\hfill \\ 1\hfill & \mathrm{if}j0,i=0\hfill \\ \left(\genfrac{}{}{0pt}{}{2j1}{2i}\right)+\left(\genfrac{}{}{0pt}{}{2j1}{2i1}\right)\hfill & \mathrm{if}j0,i0,j>i\hfill \\ \left(\genfrac{}{}{0pt}{}{w2j+1}{w2i}\right)+\left(\genfrac{}{}{0pt}{}{w2j+1}{w2i+1}\right)\hfill & \mathrm{if}j0,i0,ji.\hfill \end{array}$$
Then the Eichler-Shimura relations (KZ3) are expressed as
$$\mathrm{๐๐ญ}=\mathrm{๐}.$$
Here and hereafter the symbol $`\mathrm{๐}`$ stands for a zero matrix.
Furthermore we have:
###### Lemma 3.1.
It holds that
$$\mathrm{๐๐ญ}=\mathrm{๐}.$$
###### Proof.
We note that $`0`$th row of $`๐`$ is equal to $`0`$th row of $`๐`$, and that the $`i`$th row of $`๐`$ is the sum of the $`(2i1)`$th row and $`(2i)`$th row of $`๐`$ for $`1iw_2`$. Thus, from the fact that $`\mathrm{๐๐ญ}=\mathrm{๐}`$, we have that $`\mathrm{๐๐ญ}=\mathrm{๐}`$. โ
Next we need the following definitions:
###### Definition 3.4.
We define a matrix $`๐=[c_{ij}]`$ $`(i=0,1,\mathrm{},w_4;j=0,1,\mathrm{},w_2)`$ as follows, which depends on the congruence conditions:
1. For an integer $`i`$ with $`0i<w_4`$, we define
$$c_{ij}:=\{\begin{array}{cc}0\hfill & \mathrm{if}j=0\hfill \\ 1\hfill & \mathrm{if}j=i+1\hfill \\ 1\hfill & \mathrm{if}j=w_2i\hfill \\ 0\hfill & \mathrm{otherwise};\hfill \end{array}$$
2. For $`w0(mod4)`$ and $`i=w_4`$, we define
$$c_{ij}:=\{\begin{array}{cc}0\hfill & \mathrm{if}0jw_4\hfill \\ 1\hfill & \mathrm{otherwise};\hfill \end{array}$$
3. For $`w2(mod4)`$ and $`i=w_4`$, we define
$$c_{ij}:=\{\begin{array}{cc}0\hfill & \mathrm{if}0jw_4\hfill \\ 1\hfill & \mathrm{if}j=w_4+1\hfill \\ 2\hfill & \mathrm{otherwise}.\hfill \end{array}$$
Here are two examples of the matrix $`๐`$:
1. For $`w=8`$,
$$๐=\left[\begin{array}{ccccc}0& 1& 0& 0& 1\\ 0& 0& 1& 1& 0\\ 0& 0& 0& 1& 1\end{array}\right];$$
2. For $`w=10`$,
$$๐=\left[\begin{array}{cccccc}0& 1& 0& 0& 0& 1\\ 0& 0& 1& 0& 1& 0\\ 0& 0& 0& 1& 2& 2\end{array}\right].$$
From the Eichler-Shimura relations (KZ1) and (KZ3) for $`n=0`$, we obtain:
###### Lemma 3.2.
(3.3)
$$\mathrm{๐๐ญ}=\mathrm{๐}.$$
Finally, we introduce the following matrix $`๐`$, which can be obtained from the matrices $`๐`$ and $`๐`$:
###### Definition 3.5.
$$๐:=\left[\begin{array}{c}๐\\ ๐\end{array}\right].$$
More precisely, for $`๐=[d_{ij}](i=0,1,\mathrm{},w_2+w_4+1;j=0,1,\mathrm{},w_2)`$,
$$d_{ij}:=\{\begin{array}{cc}b_{ij}\hfill & \mathrm{if}0iw_2,0jw_2\hfill \\ c_{(iw_21)j}\hfill & \mathrm{if}w_2+1iw_2+w_4+1,0jw_2.\hfill \end{array}$$
Then, from lemmas 3.1 and 3.2, we have
###### Lemma 3.3.
$$\mathrm{๐๐ญ}=\mathrm{๐}.$$
## 4. The Eichler-Shimura relations modulo 2
In this section, we consider the reductions modulo 2 of the matrices which were introduced in the previous section. By $`/2`$, we denote the set of congruence classes modulo 2. For an integer $`x`$,
$$\overline{x}\text{ or }xmod2$$
denotes the congruence class of $`x`$ modulo $`2`$ so that we have
$$/2=\{\overline{0},\overline{1}\}.$$
Let $`๐`$, $`๐`$ and $`๐`$ denote the reductions modulo 2 of $`๐`$, $`๐`$ and $`๐`$, respectively:
$$๐:=[k_{ij}]=[\overline{b}_{ij}],๐:=[\mathrm{}_{ij}]=[\overline{c}_{ij}]\mathrm{and}๐:=[m_{ij}]=[\overline{d}_{ij}].$$
From Definition 3.5, we know that
$$๐=\left[\begin{array}{c}๐\\ ๐\end{array}\right].$$
Here we recall Lucasโ congruence theorem on binomial coefficients (\[4, p. 271\]). Let $`p`$ be a prime number, and let $`n,k,a,b`$ be nonnegative integers with $`0a,b<p`$. Then it holds that
(4.1)
$$\left(\genfrac{}{}{0pt}{}{np+a}{kp+b}\right)\left(\genfrac{}{}{0pt}{}{n}{k}\right)\left(\genfrac{}{}{0pt}{}{a}{b}\right)(modp).$$
Using this identity we can prove the following lemma:
###### Lemma 4.1.
As for the matrix $`๐=[k_{ij}]`$, it holds that
(4.2)
$$k_{ij}=\{\begin{array}{cc}\left(\genfrac{}{}{0pt}{}{j}{i}\right)mod2\hfill & \mathrm{if}j>i\hfill \\ \overline{0}\hfill & \mathrm{if}ji.\hfill \end{array}$$
###### Proof.
If $`j>i>0`$, we have
$`b_{ij}`$ $`=\left({\displaystyle \genfrac{}{}{0pt}{}{2j1}{2i}}\right)+\left({\displaystyle \genfrac{}{}{0pt}{}{2j1}{2i1}}\right)=\left({\displaystyle \genfrac{}{}{0pt}{}{2(j1)+1}{2i}}\right)+\left({\displaystyle \genfrac{}{}{0pt}{}{2(j1)+1}{2(i1)+1}}\right)`$
$`\left({\displaystyle \genfrac{}{}{0pt}{}{j1}{i}}\right)\left({\displaystyle \genfrac{}{}{0pt}{}{1}{0}}\right)+\left({\displaystyle \genfrac{}{}{0pt}{}{j1}{i1}}\right)\left({\displaystyle \genfrac{}{}{0pt}{}{1}{1}}\right)(mod2)\text{ (by (}\text{4.1}\text{))}`$
$`=\left({\displaystyle \genfrac{}{}{0pt}{}{j1}{i}}\right)+\left({\displaystyle \genfrac{}{}{0pt}{}{j1}{i1}}\right)=\left({\displaystyle \genfrac{}{}{0pt}{}{j}{i}}\right)\text{ (by Pascalโs identity).}`$
This implies that, if $`j>i>0`$, then
$$k_{ij}=\left(\genfrac{}{}{0pt}{}{j}{i}\right)mod2.$$
If $`0<ji`$, we have
$`b_{ij}`$ $`=\left({\displaystyle \genfrac{}{}{0pt}{}{w2j+1}{w2i}}\right)+\left({\displaystyle \genfrac{}{}{0pt}{}{w2j+1}{w2i+1}}\right)`$
$`\left({\displaystyle \genfrac{}{}{0pt}{}{w_2j}{w_2i}}\right)\left({\displaystyle \genfrac{}{}{0pt}{}{1}{0}}\right)+\left({\displaystyle \genfrac{}{}{0pt}{}{w_2j}{w_2i}}\right)\left({\displaystyle \genfrac{}{}{0pt}{}{1}{1}}\right)(mod2)\text{ (by (}\text{4.1}\text{))}`$
$`0(mod2).`$
This shows that, if $`0<ji`$, then
$$k_{ij}=\overline{0}.$$
In the case that $`i=0`$ or $`j=0`$, the identities follow from the definitions of $`b_{ij}`$ and $`k_{ij}`$. This completes the proof. โ
## 5. The Pascal-Sierpinskiโs triangle and the matrix $`๐`$
Here we investigate the matrix $`๐`$. Since $`๐=[k_{ij}]`$ satisfies
$$k_{ij}=\{\begin{array}{cc}\left(\genfrac{}{}{0pt}{}{j}{i}\right)mod2\hfill & \mathrm{if}j>i\hfill \\ \overline{0}\hfill & \mathrm{if}ji,\hfill \end{array}$$
$`๐`$ is an upper triangular matrix, and the upper triangular part of the matrix is nothing but the Pascal-Sierpinskiโs triangle (see e.g. ). In particular it is โself-similarโ. The fact that the Pascal-Sierpinskiโs triangle has โself-similarityโ have been well-known (the reader can refer to, e.g. \[19, pp. 44,53\] for the congruence properties of binomial coefficients, which gives rise to the โself-similarityโ).
Considering these, we inductively define a family of square matrices $`๐_n(n=0,1,\mathrm{})`$, with entries in $`/2`$ as follows:
###### Definition 5.1.
1. $$๐_0:=\left[\begin{array}{c}\overline{1}\end{array}\right];$$
2. For any positive integer $`n`$,
$$๐_n:=\left[\begin{array}{cc}๐_{n1}& ๐_{n1}\\ \mathrm{๐}& ๐_{n1}\end{array}\right].$$
Note that the size of $`๐_n`$ is $`2^n\times 2^n`$.
3. For any positive integer $`n`$, let $`๐_n`$ denote the identity matrix of size $`n\times n`$. Namely
$$๐_n:=[\overline{\delta }_{ij}](i,j=0,1,\mathrm{},n1)$$
where $`\delta _{ij}`$ is the Kroneckerโs delta.
4. For any positive integer $`n`$,
$$๐_n:=๐_n+๐_{2^n}.$$
Note that the size of $`๐_n`$ is again $`2^n\times 2^n`$.
5. Let $`k`$ and $`n`$ be positive integers with $`kn`$, and let $`๐=[x_{ij}](i,j=0,1,\mathrm{},n1)`$ be $`n\times n`$-matrix. Then, by $`๐[[k]]`$, we denote the principal submatrix of $`๐`$ with size $`k\times k`$. Namely
$$๐[[k]]=[x_{ij}](i,j=0,1,\mathrm{},k1).$$
Under the notation we can express $`๐`$ as follows:
###### Lemma 5.1.
Let $`n`$ be an integer such that $`w_2+12^n`$. Then $`๐`$ can be expressed as
$$๐=๐_n[[w_2+1]].$$
###### Proof.
We note that the size of $`๐`$ is $`(w_2+1)\times (w_2+1)`$ while the size of $`๐_n`$ is $`2^n\times 2^n`$. Then the lemma follows from Lemma 4.1 which asserts that
$$k_{ij}=\{\begin{array}{cc}\left(\genfrac{}{}{0pt}{}{j}{i}\right)mod2\hfill & \mathrm{if}j>i\hfill \\ \overline{0}\hfill & \mathrm{if}ji.\hfill \end{array}$$
In the rest of this section, we prove several lemmas which will be used in the proof of Theorem 2.2. The following is obvious:
###### Lemma 5.2.
For any positive integer $`n`$,
$$๐_n=\left[\begin{array}{cc}๐_{n1}& ๐_{n1}+๐_{2^{n1}}\\ \mathrm{๐}& ๐_{n1}\end{array}\right].$$
Here we introduce a special type of operation on matrices.
###### Definition 5.2.
Let $`๐=[x_{ij}],๐^{}=[x_{ij}^{}](i=0,1,\mathrm{},n1;j=0,1,\mathrm{},k1)`$ be two matrices of size $`n\times k`$. Then $`๐^{}`$ is said to be obtained from $`๐`$ by an $`R^+`$-operation if there exist integers $`i_0`$ and $`i_1`$ with $`0i_0<i_1<n`$ such that
$$x_{ij}^{}=\{\begin{array}{cc}x_{ij}\hfill & \mathrm{if}ii_0\hfill \\ x_{i_0j}+x_{i_1j}\hfill & \mathrm{if}i=i_0.\hfill \end{array}$$
In other words, $`๐^{}`$ is obtained from $`๐`$ by adding $`i_1`$th row to $`i_0`$th row with $`i_0<i_1`$ $`(`$this is an elementary row operation$`)`$.
In what follow, the notation
$$๐๐$$
means โ$`๐`$ is obtained from $`๐`$ by a sequence of $`R^+`$-operationsโ. It is clear that $`๐๐`$ if and only if $`๐`$ can be expressed as
$$๐=\left[\begin{array}{cccc}\overline{1}& & \mathrm{}& \\ \overline{0}& \overline{1}& \mathrm{}& \mathrm{}\\ \mathrm{}& \mathrm{}& \mathrm{}& \\ \overline{0}& \mathrm{}& \overline{0}& \overline{1}\end{array}\right]๐.$$
Then we can formulate the following lemma:
###### Lemma 5.3.
For any positive integer $`n`$, it holds that
$$\left[\begin{array}{c}๐_{2^{n+1}}\\ ๐_{n+1}\end{array}\right]\left[\begin{array}{cc}๐_{2^n}& \mathrm{๐}\\ ๐_n& \mathrm{๐}\\ ๐_n& ๐_{2^n}\\ \mathrm{๐}& ๐_n\end{array}\right].$$
###### Proof.
Using Lemma 5.2, we have
$$\left[\begin{array}{cccc}๐_{2^n}& \mathrm{๐}& \mathrm{๐}& \mathrm{๐}\\ \mathrm{๐}& ๐_{2^n}& ๐_{2^n}& ๐_{2^n}\\ \mathrm{๐}& \mathrm{๐}& ๐_{2^n}& ๐_{2^n}\\ \mathrm{๐}& \mathrm{๐}& \mathrm{๐}& ๐_{2^n}\end{array}\right]\left[\begin{array}{c}๐_{2^{n+1}}\\ ๐_{n+1}\end{array}\right]=\left[\begin{array}{cc}๐_{2^n}& \mathrm{๐}\\ ๐_n& \mathrm{๐}\\ ๐_n& ๐_{2^n}\\ \mathrm{๐}& ๐_n\end{array}\right].$$
Now we set
$$๐:=\left[\begin{array}{c}๐_2\\ ๐_1\end{array}\right]=\left[\begin{array}{cc}\overline{1}& \overline{0}\\ \overline{0}& \overline{1}\\ \overline{0}& \overline{1}\\ \overline{0}& \overline{0}\end{array}\right].$$
Then we obtain:
###### Lemma 5.4.
For any nonnegative integer $`n`$, the matrix
$$\left[\begin{array}{c}๐_{2^{n+1}}\\ ๐_{n+1}\end{array}\right]$$
can be transformed to a block matrix of the form
$$\left[\begin{array}{cccc}๐& \mathrm{๐}& \mathrm{}& \mathrm{๐}\\ & ๐& \mathrm{}& \mathrm{}\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{๐}\\ & \mathrm{}& & ๐\end{array}\right]$$
by a sequence of $`R^+`$-operations.
###### Proof.
Repeatedly applying Lemma 5.3, we can transform the matrix
$$\left[\begin{array}{c}๐_{2^{n+1}}\\ ๐_{n+1}\end{array}\right]$$
by $`R^+`$-operations:
$`\left[\begin{array}{c}๐_{2^{n+1}}\\ ๐_{n+1}\end{array}\right]`$ $`\left[\begin{array}{cc}๐_{2^n}& \mathrm{๐}\\ ๐_n& \mathrm{๐}\\ & ๐_{2^n}\\ & ๐_n\end{array}\right]\left[\begin{array}{cccc}๐_{2^{n1}}& \mathrm{๐}& \mathrm{๐}& \mathrm{๐}\\ ๐_{n1}& \mathrm{๐}& \mathrm{๐}& \mathrm{๐}\\ & ๐_{2^{n1}}& \mathrm{๐}& \mathrm{๐}\\ & ๐_{n1}& \mathrm{๐}& \mathrm{๐}\\ & & ๐_{2^{n1}}& \mathrm{๐}\\ & & ๐_{n1}& \mathrm{๐}\\ & & & ๐_{2^{n1}}\\ & & & ๐_{n1}\end{array}\right]`$
$`\mathrm{}\left[\begin{array}{ccccc}๐_2& \mathrm{๐}& \mathrm{}& \mathrm{}& \mathrm{๐}\\ ๐_1& \mathrm{๐}& \mathrm{}& \mathrm{}& \mathrm{๐}\\ & ๐_2& \mathrm{๐}& & \mathrm{}\\ & ๐_1& \mathrm{๐}& \mathrm{}& \mathrm{}\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{๐}\\ \mathrm{}& & \mathrm{}& \mathrm{}& \mathrm{๐}\\ & \mathrm{}& & & ๐_2\\ & \mathrm{}& & & ๐_1\end{array}\right]=\left[\begin{array}{cccc}๐& \mathrm{๐}& \mathrm{}& \mathrm{๐}\\ & ๐& \mathrm{}& \mathrm{}\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{๐}\\ & \mathrm{}& & ๐\end{array}\right].`$
This completes the proof. โ
From Lemma 5.4, we can obtain the following:
###### Lemma 5.5.
Let $`n`$ be a positive integer. Then, by a sequence of $`R^+`$-operations, $`๐_{n+2}`$ can be transformed to a block matrix $`\stackrel{ห}{๐}_{n+2}`$ of the form
(5.1)
$$\stackrel{ห}{๐}_{n+2}=\left[\begin{array}{ccc}๐& \left[\begin{array}{cccc}๐& \mathrm{๐}& \mathrm{}& \mathrm{๐}\\ & ๐& \mathrm{}& \mathrm{}\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{๐}\\ & \mathrm{}& & ๐\end{array}\right]& \left[\begin{array}{cccc}\overline{1}& & \mathrm{}& \\ \overline{0}& \overline{1}& \mathrm{}& \mathrm{}\\ \mathrm{}& \mathrm{}& \mathrm{}& \\ \overline{0}& \mathrm{}& \overline{0}& \overline{1}\end{array}\right]\\ \mathrm{๐}& \mathrm{๐}& ๐_{n+1}\end{array}\right],$$
where $`๐`$ is a $`2^{n+1}\times 2^n`$ matrix. Furthermore we can assume that $`R^+`$-operations used in this transformation are elementary row operations adding $`i`$th rows with $`i<2^{n+1}`$.
###### Proof.
First we see that $`๐_{n+2}`$ have the form
$`๐_{n+2}`$ $`=\left[\begin{array}{cc}๐_{n+1}& ๐_{n+1}+๐_{2^{n+1}}\\ \mathrm{๐}& ๐_{n+1}\end{array}\right]`$
$`=\left[\begin{array}{cc}\left[\begin{array}{cc}๐_n& ๐_n+๐_{2^n}\\ \mathrm{๐}& ๐_n\end{array}\right]& \left[\begin{array}{cccc}\overline{1}& & \mathrm{}& \\ \overline{0}& \overline{1}& \mathrm{}& \mathrm{}\\ \mathrm{}& \mathrm{}& \mathrm{}& \\ \overline{0}& \mathrm{}& \overline{0}& \overline{1}\end{array}\right]\\ \mathrm{๐}& ๐_{n+1}\end{array}\right].`$
Now, applying $`R^+`$-operations, this matrix can be transformed to a matrix of the form
$$\left[\begin{array}{cc}\left[\begin{array}{cc}๐_n& ๐_{2^n}\\ \mathrm{๐}& ๐_n\end{array}\right]& \left[\begin{array}{cccc}\overline{1}& & \mathrm{}& \\ \overline{0}& \overline{1}& \mathrm{}& \mathrm{}\\ \mathrm{}& \mathrm{}& \mathrm{}& \\ \overline{0}& \mathrm{}& \overline{0}& \overline{1}\end{array}\right]\\ \mathrm{๐}& ๐_{n+1}\end{array}\right]=\left[\begin{array}{ccc}\left[\begin{array}{c}๐_n\\ \mathrm{๐}\end{array}\right]& \left[\begin{array}{c}๐_{2^n}\\ ๐_n\end{array}\right]& \left[\begin{array}{cccc}\overline{1}& & \mathrm{}& \\ \overline{0}& \overline{1}& \mathrm{}& \mathrm{}\\ \mathrm{}& \mathrm{}& \mathrm{}& \\ \overline{0}& \mathrm{}& \overline{0}& \overline{1}\end{array}\right]\\ \mathrm{๐}& \mathrm{๐}& ๐_{n+1}\end{array}\right].$$
Furthermore, by Lemma 5.4, this matrix can be transformed to a matrix, say $`๐`$, of the form
$$๐=\left[\begin{array}{ccc}๐& \left[\begin{array}{cccc}๐& \mathrm{๐}& \mathrm{}& \mathrm{๐}\\ & ๐& \mathrm{}& \mathrm{}\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{๐}\\ & \mathrm{}& & ๐\end{array}\right]& \left[\begin{array}{cccc}\overline{1}& & \mathrm{}& \\ \overline{0}& \overline{1}& \mathrm{}& \mathrm{}\\ \mathrm{}& \mathrm{}& \mathrm{}& \\ \overline{0}& \mathrm{}& \overline{0}& \overline{1}\end{array}\right]\\ \mathrm{๐}& \mathrm{๐}& ๐_{n+1}\end{array}\right],$$
where $`๐`$ is a matrix of size $`2^{n+1}\times 2^n`$. We take this matrix $`๐`$ as $`\stackrel{ห}{๐}_{n+2}`$. Here we used the fact that, while $`R^+`$-operations, a matrix of the form
$$\left[\begin{array}{cccc}\overline{1}& & \mathrm{}& \\ \overline{0}& \overline{1}& \mathrm{}& \mathrm{}\\ \mathrm{}& \mathrm{}& \mathrm{}& \\ \overline{0}& \mathrm{}& \overline{0}& \overline{1}\end{array}\right]$$
remains a matrix of this form. We note that, in this transformation, we only used elementary row operations adding $`i`$th rows with $`i<2^{n+1}`$.
This completes the proof. โ
Moreover, from Lemma 5.5, we can obtain the following:
###### Lemma 5.6.
Let $`n`$ be a positive integer. Then, by a sequence of $`R^+`$-operations, $`๐_{n+3}`$ can be transformed to a block matrix $`\stackrel{ห}{\stackrel{ห}{๐}}_{n+3}`$ of the form
(5.2)
$$\stackrel{ห}{\stackrel{ห}{๐}}_{n+3}=\left[\begin{array}{cc}\left[\begin{array}{ccc}๐& \left[\begin{array}{cccc}๐& \mathrm{๐}& \mathrm{}& \mathrm{๐}\\ & ๐& \mathrm{}& \mathrm{}\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{๐}\\ & \mathrm{}& & ๐\end{array}\right]& \left[\begin{array}{cccc}\overline{1}& & \mathrm{}& \\ \overline{0}& \overline{1}& \mathrm{}& \mathrm{}\\ \mathrm{}& \mathrm{}& \mathrm{}& \\ \overline{0}& \mathrm{}& \overline{0}& \overline{1}\end{array}\right]\\ \mathrm{๐}& \mathrm{๐}& ๐_{n+1}\end{array}\right]& ๐\\ \mathrm{๐}& ๐_{n+2}\end{array}\right]$$
where $`๐`$ and $`๐`$ are $`2^{n+1}\times 2^n`$ and $`2^{n+2}\times 2^{n+2}`$ matrices, respectively. Furthermore we can assume that $`R^+`$-operations used in this transformation are elementary row operations adding $`i`$th rows with $`i<2^{n+1}`$.
###### Proof.
From Lemmas 5.2 and 5.5, we have
$`๐_{n+3}`$ $`=\left[\begin{array}{cc}๐_{n+2}& ๐_{n+2}+๐_{2^{n+2}}\\ \mathrm{๐}& ๐_{n+2}\end{array}\right]\left[\begin{array}{cc}\stackrel{ห}{๐}_{n+2}& ๐\\ \mathrm{๐}& ๐_{n+2}\end{array}\right].`$
This completes the proof. โ
Hereafter we assume the following:
###### Assumption 5.7.
(5.3)
$$1<n\mathrm{and}2^n\frac{w_2}{3}<2^{n+1}.$$
Note that, for a given $`w_212`$, there exists uniquely an integer $`n`$ which satisfies (5.3). Under Assumption 5.7, we see obviously that $`w_2+12^{n+3}`$. Hence we know that $`๐`$ can be express as
$$๐=๐_{n+3}[[w_2+1]].$$
Now we set
(5.4)
$$๐:=\stackrel{ห}{\stackrel{ห}{๐}}_{n+3}[[w_2+1]].$$
Then, from Lemma 5.6, we obtain
###### Lemma 5.8.
By a sequence of $`R^+`$-operations, the matrix $`๐`$ can be transformed to the matrix $`๐`$.
###### Proof.
By Lemma 5.6, $`๐_{n+3}`$ can be transformed to the matrix $`\stackrel{ห}{\stackrel{ห}{๐}}_{n+3}`$ by a sequence of $`R^+`$-operations adding $`i`$th rows with $`i<2^{n+1}`$. Furthermore, from Assumption 5.7, it follows that $`2^{n+1}w_2+1`$. Hence we have proved that the matrix $`๐`$ can be transformed to the matrix $`๐`$ by a sequence of $`R^+`$-operations. โ
## 6. Properties of the matrix $`๐`$
In this section, keeping Assumption 5.7, we will study properties of the matrix
$$๐=\stackrel{ห}{\stackrel{ห}{๐}}_{n+3}[[w_2+1]]$$
which is obtained from $`๐`$ by a sequence of $`R^+`$-operations.
We need the following notation:
###### Definition 6.1.
1. For an integer $`i0`$, we define an integer $`\alpha (i)`$ by
$$\alpha (i):=\frac{i}{2}+\{\begin{array}{cc}1\hfill & \mathrm{if}i1(mod4)\hfill \\ 0\hfill & \mathrm{otherwise};\hfill \end{array}$$
2. Let us set
$$w=12k+2a(0<k,0a5).$$
We define
$$\beta _1(w):=2^{n+1}+\{\begin{array}{cc}4k+2\hfill & \mathrm{if}a=0,1\hfill \\ 4k+4\hfill & \mathrm{if}a=2,3,4\hfill \\ 4k+6\hfill & \mathrm{if}a=5;\hfill \end{array}$$
$$\beta _2(w):=2^{n+1}+\{\begin{array}{cc}4k\hfill & \mathrm{if}a=0,1,2\hfill \\ 4k+2\hfill & \mathrm{if}a=3,4,5.\hfill \end{array}$$
Obviously we have that
$$\alpha (i+4)=\alpha (i)+2,\alpha (i)=\frac{i}{2}\mathrm{if}i0(mod2);$$
(6.1)
$$2^n+\alpha (\beta _1(w))=\{\begin{array}{cc}2k+1\hfill & \mathrm{if}a=0,1\hfill \\ 2k+2\hfill & \mathrm{if}a=2,3,4\hfill \\ 2k+3\hfill & \mathrm{if}a=5;\hfill \end{array}$$
(6.2)
$$2^n+\alpha (\beta _2(w))=\{\begin{array}{cc}2k\hfill & \mathrm{if}a=0,1,2\hfill \\ 2k+1\hfill & \mathrm{if}a=3,4,5.\hfill \end{array}$$
Using the notation $`\alpha (i)`$, we can formulate the following lemma:
###### Lemma 6.1.
Let us set
$$๐=[h_{ij}](i,j=0,1,\mathrm{},w_2),$$
and let $`i`$ be an integer such that
$$0i<2^{n+1}.$$
Then it holds that
(6.3) $`h_{ij}`$ $`=\overline{0}`$ $`\mathrm{for}`$ $`2^n+\alpha (i)<j<\mathrm{min}(\mathrm{\hspace{0.17em}2}^{n+1}+i,w_2+1),`$
(6.4) $`h_{ij}`$ $`=\overline{1}`$ $`\mathrm{for}`$ $`j=2^{n+1}+i,j<w_2+1,`$
(6.5) $`h_{ij}`$ $`=\overline{1}`$ $`\mathrm{for}`$ $`j=2^n+\alpha (i),i0(mod2).`$
###### Proof.
Let us set
$$\stackrel{ห}{\stackrel{ห}{๐}}_{n+3}=[\stackrel{ห}{\stackrel{ห}{q}}_{ij}](i,j=0,1,\mathrm{},2^{n+3}1).$$
Then, for $`i`$ such that $`0i<2^{n+1}`$, we can easily read off the following identities
(6.6) $`\stackrel{ห}{\stackrel{ห}{q}}_{ij}`$ $`=\overline{0}`$ $`\mathrm{for}`$ $`2^n+\alpha (i)<j<2^{n+1}+i,`$
(6.7) $`\stackrel{ห}{\stackrel{ห}{q}}_{ij}`$ $`=\overline{1}`$ $`\mathrm{for}`$ $`j=2^{n+1}+i,`$
(6.8) $`\stackrel{ห}{\stackrel{ห}{q}}_{ij}`$ $`=\overline{1}`$ $`\mathrm{for}`$ $`j=2^n+\alpha (i),i3(mod4)`$
from the form
$$\stackrel{ห}{\stackrel{ห}{๐}}_{n+3}=\left[\begin{array}{cc}\left[\begin{array}{ccc}\stackrel{2^n}{\stackrel{}{}}& \stackrel{2^n}{\stackrel{}{}}& \stackrel{2^{n+1}}{\stackrel{}{}}\\ ๐& \left[\begin{array}{cccc}๐& \mathrm{๐}& \mathrm{}& \mathrm{๐}\\ & ๐& \mathrm{}& \mathrm{}\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{๐}\\ & \mathrm{}& & ๐\end{array}\right]& \left[\begin{array}{cccc}\overline{1}& & \mathrm{}& \\ \overline{0}& \overline{1}& \mathrm{}& \mathrm{}\\ \mathrm{}& \mathrm{}& \mathrm{}& \\ \overline{0}& \mathrm{}& \overline{0}& \overline{1}\end{array}\right]\\ \mathrm{๐}& \mathrm{๐}& ๐_{n+1}\end{array}\right]& ๐\\ \mathrm{๐}& ๐_{n+2}\end{array}\right]$$
and
$$๐=\left[\begin{array}{cc}\overline{1}& \overline{0}\\ \overline{0}& \overline{1}\\ \overline{0}& \overline{1}\\ \overline{0}& \overline{0}\end{array}\right]$$
where $`๐`$ (resp. $`๐`$) is a $`2^{n+1}\times 2^n`$ (resp. $`2^{n+2}\times 2^{n+2}`$) matrix.
Next, noting that $`๐=\stackrel{ห}{\stackrel{ห}{๐}}_{n+3}[[w_2+1]]`$, we have the identities (6.3), (6.4) and (6.5) from (6.6), (6.7) and (6.8) respectively. This completes the proof. โ
By Lemma 6.1, we obtain:
###### Lemma 6.2.
Set $`w=12k+2a(0<k,0a5)`$, and let
$$๐=[h_{ij}](i,j=0,1,\mathrm{},w_2).$$
1. Suppose that
$$w<22^{n+2}\mathrm{and}i=\beta _1(w).$$
Then we have
(6.9) $`h_{ij}`$ $`=\overline{0}`$ $`\mathrm{if}`$ $`\{\begin{array}{cc}2k+1<j<4k+2,\hfill & a=0,1\hfill \\ 2k+2<j<4k+4,\hfill & a=2,3,4\hfill \\ 2k+3<j<4k+6,\hfill & a=5,\hfill \end{array}`$
(6.10) $`h_{ij}`$ $`=\overline{1}`$ $`\mathrm{if}`$ $`\{\begin{array}{cc}j=2k+1,\hfill & a=0,1\hfill \\ j=2k+2,\hfill & a=2,3,4\hfill \\ j=2k+3,\hfill & a=5.\hfill \end{array}`$
2. Suppose that
$$22^{n+2}w\mathrm{and}i=\beta _2(w).$$
Then we have
(6.11) $`h_{ij}`$ $`=\overline{0}`$ $`\mathrm{if}`$ $`\{\begin{array}{cc}2k<j<4k,\hfill & a=0,1,2\hfill \\ 2k+1<j<4k+2,\hfill & a=3,4,5,\hfill \end{array}`$
(6.12) $`h_{ij}`$ $`=\overline{1}`$ $`\mathrm{if}`$ $`\{\begin{array}{cc}j=4k,\hfill & a=0,1,2\hfill \\ j=4k+2,\hfill & a=3,4,5.\hfill \end{array}`$
###### Proof.
Under Assumption 5.7, we can easily show the following inequalities (the detail is left to the reader):
(6.13)
$$0\beta _1(w)<2^{n+1}\mathrm{if}w<22^{n+2},$$
(6.14)
$$0\beta _2(w)<2^{n+1}\mathrm{if}22^{n+2}w.$$
Furthermore, when $`i=\beta _1(w)`$ or $`i=\beta _2(w)`$, we can easily check that
$$2^{n+1}+i<w_2+1.$$
Now we are ready to apply Lemma 6.1 to prove the lemma. Noting that $`\beta _1(w)`$ is even, we can apply (6.3) and (6.5) in Lemma 6.1 and (6.1) to obtain (6.9) and (6.10).
We can also apply (6.3) and (6.4) in Lemma 6.1 and (6.2) to obtain (6.11) and (6.12). This completes the proof. โ
In the sequel, we use the following notation; for column vectors $`๐ฏ_1,\mathrm{},๐ฏ_k`$, let
$$\mathrm{Col}(๐ฏ_1,\mathrm{},๐ฏ_k):=\text{the column vector space spanned by }๐ฏ_1,\mathrm{},๐ฏ_k.$$
We also denote the $`i`$th entry of $`๐ฏ_j`$ by $`๐ฏ_j(i)(i=0,1,\mathrm{})`$. Let $`๐ก_j`$, $`๐ค_j`$, $`๐ฅ_j`$ and $`๐ฆ_j`$ denote column vectors of $`๐`$, $`๐`$, $`๐`$ and $`๐`$, that is
$$๐=[๐ก_0\mathrm{}๐ก_{w_2}],๐=[๐ค_0\mathrm{}๐ค_{w_2}],๐=[๐ฅ_0\mathrm{}๐ฅ_{w_2}]\text{ and }๐=[๐ฆ_0\mathrm{}๐ฆ_{w_2}].$$
Then, for $`๐ก_j`$ and $`๐ค_j`$, it holds that $`๐ก_j(i)=h_{ij}`$ and $`๐ค_j(i)=k_{ij}`$, respectively.
Under this notation, we have the following lemma:
###### Lemma 6.3.
$`๐ฆ_1,๐ฆ_2,\mathrm{},๐ฆ_{w_4}`$ are linearly independent.
###### Proof.
By definition of $`๐`$, it is clear that $`๐ฅ_1,๐ฅ_2,\mathrm{},๐ฅ_{w_4}`$ are linearly independent. Furthermore, from the definition of $`๐`$:
$$๐=\left[\begin{array}{c}๐\\ ๐\end{array}\right],$$
we find that $`๐ฆ_1,๐ฆ_2,\mathrm{},๐ฆ_{w_4}`$ are also linearly independent. โ
We need the following lemma to prove Lemma 6.5:
###### Lemma 6.4.
Suppose that
$$w_4<j_0w_2\mathrm{and}๐ฆ_{j_0}\mathrm{Col}(๐ฆ_1,๐ฆ_2,\mathrm{},๐ฆ_{j_01}).$$
Then $`๐ฆ_{w_2+1j_0}+๐ฆ_{j_0}`$ can be expressed as
$$๐ฆ_{w_2+1j_0}+๐ฆ_{j_0}=\underset{j=w_2+2j_0}{\overset{j_01}{}}a_j๐ฆ_j$$
for some $`a_j/2(j=w_2+2j_0,\mathrm{},j_01)`$.
###### Proof.
Since
$$๐ฆ_{j_0}\mathrm{Col}(๐ฆ_1,๐ฆ_2,\mathrm{},๐ฆ_{j_01}),$$
there are $`a_j/2(j=1,2,\mathrm{},j_01)`$ such that
$$๐ฆ_{j_0}=\underset{j=1}{\overset{j_01}{}}a_j๐ฆ_j.$$
Then we have
$$๐ฅ_{j_0}=\underset{j=1}{\overset{j_01}{}}a_j๐ฅ_j.$$
However $`๐`$ has the form
$$๐=\begin{array}{c}\begin{array}{ccc}& \stackrel{w_4=w_2/2}{\stackrel{}{}}& \stackrel{w_4=w_2/2}{\stackrel{}{}}\end{array}\\ \left[\begin{array}{ccccccc}\overline{0}& \overline{1}& \overline{0}& \mathrm{}& \mathrm{}& \overline{0}& \overline{1}\\ \overline{0}& \overline{0}& \mathrm{}& \overline{0}& \overline{0}& \text{.}\text{.}\text{.}& \overline{0}\\ \overline{0}& \overline{0}& \overline{0}& \overline{1}& \overline{1}& \overline{0}& \overline{0}\\ \overline{0}& \overline{0}& \mathrm{}& \overline{0}& \overline{1}& \mathrm{}& \overline{1}\end{array}\right]\end{array}$$
or
$$๐=\begin{array}{c}\begin{array}{cccc}& \stackrel{w_4=(w_21)/2}{\stackrel{}{}}& & \stackrel{w_4=(w_21)/2}{\stackrel{}{}}\end{array}\\ \left[\begin{array}{cccccccc}\overline{0}& \overline{1}& \overline{0}& \mathrm{}& \overline{0}& \mathrm{}& \overline{0}& \overline{1}\\ \overline{0}& \overline{0}& \mathrm{}& \overline{0}& \overline{0}& \overline{0}& \text{.}\text{.}\text{.}& \overline{0}\\ \overline{0}& \overline{0}& \overline{0}& \overline{1}& \overline{0}& \overline{1}& \overline{0}& \overline{0}\\ \overline{0}& \overline{0}& \overline{0}& \overline{0}& \overline{1}& \overline{0}& \overline{0}& \overline{0}\end{array}\right]\end{array}$$
according as $`w0(mod4)`$ or $`w2(mod4)`$.
Thus we see
$$a_j=\overline{0}\mathrm{for}j=1,2,\mathrm{},w_2j_0$$
and
$$a_j=\overline{1}\mathrm{for}j=w_2+1j_0.$$
Hence we have
$$๐ฆ_{w_2+1j_0}+๐ฆ_{j_0}=\underset{j=w_2+2j_0}{\overset{j_01}{}}a_j๐ฆ_j.$$
Now we are ready to prove the following lemma:
###### Lemma 6.5.
Let
$$w=12k+2a(0<k,0a5),$$
and set
$$j_0=\{\begin{array}{cc}4k\hfill & \mathrm{if}a=0\hfill \\ 4k+1\hfill & \mathrm{if}a=1,2\hfill \\ 4k+2\hfill & \mathrm{if}a=3\hfill \\ 4k+3\hfill & \mathrm{if}a=4,5\hfill \end{array}\mathrm{when}w<22^{n+2},$$
$$j_0=\{\begin{array}{cc}4k\hfill & \mathrm{if}a=0,1,2\hfill \\ 4k+2\hfill & \mathrm{if}a=3,4,5\hfill \end{array}\mathrm{when}22^{n+2}w.$$
Then, for $`j_1`$ with $`w_4<j_1j_0`$, we have
$$๐ฆ_{j_1}\mathrm{Col}(๐ฆ_1,๐ฆ_2,\mathrm{},๐ฆ_{j_11}).$$
###### Proof.
First we will prove the case that $`j_1=j_0`$ by reduction to absurdity. So we suppose that
$$๐ฆ_{j_0}\mathrm{Col}(๐ฆ_1,๐ฆ_2,\mathrm{},๐ฆ_{j_01}).$$
Then, from Lemma 6.4, we have
$$๐ฆ_{w_2+1j_0}+๐ฆ_{j_0}=\underset{j=w_2+2j_0}{\overset{j_01}{}}a_j๐ฆ_j$$
for $`a_j/2(j=w_2+2j_0,\mathrm{},j_01)`$. Since
$$๐=\left[\begin{array}{c}๐\\ ๐\end{array}\right],$$
we also have
$$๐ค_{w_2+1j_0}+๐ค_{j_0}=\underset{j=w_2+2j_0}{\overset{j_01}{}}a_j๐ค_j.$$
Furthermore, since $`๐`$ is obtained from $`๐`$ by row operations, we have
(6.15)
$$๐ก_{w_2+1j_0}+๐ก_{j_0}=\underset{j=w_2+2j_0}{\overset{j_01}{}}a_j๐ก_j.$$
We compute that
$`w_2+1j_0`$ $`=\{\begin{array}{cc}2k+1\hfill & \mathrm{if}a=0,1\hfill \\ 2k+2\hfill & \mathrm{if}a=2,3,4\hfill \\ 2k+3\hfill & \mathrm{if}a=5\hfill \end{array}\mathrm{when}w<22^{n+2},`$
and that
$`w_2+1j_0`$ $`=\{\begin{array}{cc}2k+1\hfill & \mathrm{if}a=0\hfill \\ 2k+2\hfill & \mathrm{if}a=1,3\hfill \\ 2k+3\hfill & \mathrm{if}a=2,4\hfill \\ 2k+4\hfill & \mathrm{if}a=5\hfill \end{array}\mathrm{when}22^{n+2}w.`$
Now we suppose that $`w<22^{n+2}`$, and we take $`i=\beta _1(w)`$. Then, from the identities (6.9) and (6.10), we have
$$๐ก_j(i)=\overline{0}(j=w_2+2j_0,\mathrm{},j_0)\mathrm{and}๐ก_{w_2+1j_0}(i)=\overline{1}.$$
This contradicts (6.15).
Next we suppose that $`22^{n+2}w`$, and we take $`i=\beta _2(w)`$. Then, from the identities (6.11) and (6.12), we have that
$$๐ก_j(i)=\overline{0}(j=w_2+1j_0,\mathrm{},j_01)\mathrm{and}๐ก_{j_0}(i)=\overline{1}.$$
This also contradicts (6.15).
Thus we have proved
$$๐ฆ_{j_0}\mathrm{Col}(๐ฆ_1,๐ฆ_2,\mathrm{},๐ฆ_{j_01}).$$
To prove
$$๐ฆ_{j_1}\mathrm{Col}(๐ฆ_1,๐ฆ_2,\mathrm{},๐ฆ_{j_11})$$
for $`j_1`$ with $`w_4<j_1<j_0`$, we take
$$i=\beta _1(w)+2(j_0j_1)\mathrm{if}w<22^{n+2}$$
and
$$i=\beta _2(w)(j_0j_1)\mathrm{if}22^{n+2}w.$$
Then we apply Lemma 6.1. The subsequent argument is similar to that of the case $`j_1=j_0`$. The detail will be left to the reader. โ
Combining Lemmas 6.3 and 6.5, we have
###### Lemma 6.6.
Let $`j_0`$ be as in Lemma 6.5. If $`1<j_1j_0`$, then it holds that
$$๐ฆ_{j_1}\mathrm{Col}(๐ฆ_1,๐ฆ_2,\mathrm{},๐ฆ_{j_11}).$$
We need the following two Lemmas 6.7 and 6.8 to prove Lemma 6.9.
###### Lemma 6.7.
If $`1<j_1w_2`$ and $`j_1`$ is odd, then
$$๐ฆ_{j_1}\mathrm{Col}(๐ฆ_1,\mathrm{},๐ฆ_{j_11}).$$
###### Proof.
We know
$$๐ค_j(j_11)=\overline{0}(j=1,\mathrm{},j_11)$$
and
$$๐ค_{j_1}(j_11)=\left(\genfrac{}{}{0pt}{}{j_1}{j_11}\right)mod2=\overline{j_1}=\overline{1}.$$
This implies
$$๐ค_{j_1}\mathrm{Col}(๐ค_1,\mathrm{},๐ค_{j_11}).$$
Thus we have
$$๐ฆ_{j_1}\mathrm{Col}(๐ฆ_1,\mathrm{},๐ฆ_{j_11}).$$
###### Lemma 6.8.
If $`w_20(mod2)`$, then it holds that
$$๐ฆ_{w_2}\mathrm{Col}(๐ฆ_1,\mathrm{},๐ฆ_{w_21}).$$
###### Proof.
We know the second row of $`๐`$ has the form
(6.16)
$$(\overline{0},\overline{0},\overline{0},\overline{1},\overline{0},\overline{1},\overline{0},\mathrm{},\overline{1},\overline{0},\overline{1},\overline{0})$$
since
$$m_{1j}=\{\begin{array}{cc}\left(\genfrac{}{}{0pt}{}{j}{1}\right)mod2\hfill & \mathrm{if}1<j\hfill \\ \overline{0}\hfill & \mathrm{if}j1\hfill \end{array}=\{\begin{array}{cc}\overline{1}\hfill & \mathrm{if}1<j,j\mathrm{odd}\hfill \\ \overline{0}\hfill & \mathrm{otherwise}.\hfill \end{array}$$
We add row vectors of $`๐`$ to the vector (6.16) to obtain a vector
(6.17)
$$(\overline{0},\overline{0},\overline{0},\mathrm{},\overline{0},\overline{0},\overline{1},\overline{1},\mathrm{},\overline{1},\overline{1},\overline{0}).$$
Finally we add the last row vector of $`๐`$
$$(\overline{0},\overline{0},\overline{0},\mathrm{},\overline{0},\overline{0},\overline{1},\overline{1},\mathrm{},\overline{1},\overline{1},\overline{1})$$
to obtain a vector
(6.18)
$$(\overline{0},\overline{0},\overline{0},\mathrm{},\overline{0},\overline{0},\overline{0},\overline{0},\mathrm{},\overline{0},\overline{0},\overline{1}).$$
This shows that the vector (6.18) is a linear combination of the row vectors of $`๐`$. Hence we know that the last column vector $`๐ฆ_{w_2}`$ of $`๐`$ is linearly independent of $`๐ฆ_j(j=1,2,\mathrm{},w_21)`$. This completes the proof. โ
Now we define sets of column vectors in $`๐`$ as follows:
###### Definition 6.2.
Let
$$w=12k+2a(0<k,0a5).$$
1. If $`a=0`$, we set
$$S:=\{๐ฆ_1,๐ฆ_2,\mathrm{},๐ฆ_{4k1},๐ฆ_{4k}\}\{๐ฆ_{4k+1},๐ฆ_{4k+1+2},\mathrm{},๐ฆ_{4k+1+2(k1)}\}\{๐ฆ_{6k}\};$$
2. If $`a=1`$, we set
$$S:=\{๐ฆ_1,๐ฆ_2,\mathrm{},๐ฆ_{4k1},๐ฆ_{4k}\}\{๐ฆ_{4k+1},๐ฆ_{4k+1+2},\mathrm{},๐ฆ_{4k+1+2k}\};$$
3. If $`a=2`$, we set
$$S:=\{๐ฆ_1,๐ฆ_2,\mathrm{},๐ฆ_{4k1},๐ฆ_{4k}\}\{๐ฆ_{4k+1},๐ฆ_{4k+1+2},\mathrm{},๐ฆ_{4k+1+2k}\}\{๐ฆ_{6k+2}\};$$
4. If $`a=3`$, we set
$$S:=\{๐ฆ_1,๐ฆ_2,\mathrm{},๐ฆ_{4k+1},๐ฆ_{4k+2}\}\{๐ฆ_{4k+3},๐ฆ_{4k+3+2},\mathrm{},๐ฆ_{4k+3+2k}\};$$
5. If $`a=4`$, we set
$$S:=\{๐ฆ_1,๐ฆ_2,\mathrm{},๐ฆ_{4k+1},๐ฆ_{4k+2}\}\{๐ฆ_{4k+3},๐ฆ_{4k+3+2},\mathrm{},๐ฆ_{4k+3+2k}\}\{๐ฆ_{6k+4}\};$$
6. If $`a=5`$, we set
$$S:=\{๐ฆ_1,๐ฆ_2,\mathrm{},๐ฆ_{4k+1},๐ฆ_{4k+2}\}\{๐ฆ_{4k+3},๐ฆ_{4k+3+2},\mathrm{},๐ฆ_{4k+3+2(k+1)}\}.$$
Now, concluding the discussion in this section, we have
###### Lemma 6.9.
The vectors in $`S`$ are linearly independent.
###### Proof.
Suppose that $`j>1`$ and $`๐ฆ_jS`$. Then, if $`w_20(mod2)`$, by Lemmas 6.6, 6.7 and 6.8, and if $`w_21(mod2)`$, by Lemmas 6.6 and 6.7, we know
$$๐ฆ_j\mathrm{Col}(๐ฆ_1,๐ฆ_2,\mathrm{},๐ฆ_{j1}).$$
This implies that $`S`$ is a set of linearly independent vectors. โ
## 7. Proofs of Theorems 2.2 and 2.3
In this section, we give proofs of Theorems 2.2 and 2.3. A straightforward computation gives:
###### Lemma 7.1.
The cardinality of $`S`$, say #S, satisfies
$$w_2\mathrm{\#}S=d_w.$$
Now we are ready to prove Theorems 2.2 and 2.3.
###### Proof of Theorem 2.2.
One can easily check that Theorem 2.2 holds for $`w_2`$ with $`w_2<12`$. Therefore we assume that $`12w_2`$ as in Assumption 5.7.
Noting that $`\mathrm{\#}S=w_2d_w`$ by Lemma 7.1, we can express $`S`$ as follows:
1. If $`w0(mod4)`$, then
$$S=\{๐ฆ_j|1jw_2,jw_22,w_24,\mathrm{},w_22d_w\};$$
2. If $`w2(mod4)`$, then
$$S=\{๐ฆ_j|1jw_2,jw_21,w_212,\mathrm{},w_212(d_w1)\}.$$
Now let $`n_w`$ be the integer defined by
$$n_w:=w_2d_w,\text{ that is, }n_w=\mathrm{\#}S,$$
and let
$$j_a(a=0,1,\mathrm{},n_w1)$$
be integers such that
$$1j_0<j_1<\mathrm{}<j_{n_w1}w_2\text{ and }๐ฆ_{j_a}S.$$
Furthermore, let
$$\widehat{j}_b(b=0,1,\mathrm{},d_w1)$$
be integers such that
$$1\widehat{j}_0<\widehat{j}_1<\mathrm{}<\widehat{j}_{d_w1}w_2\text{ and }๐ฆ_{\widehat{j}_b}S.$$
Clearly, we have
$$S=\{๐ฆ_{j_a}|a=0,1,\mathrm{},n_w1\}$$
and
(7.1)
$$\begin{array}{cc}\hfill \{\widehat{j}_0,& \widehat{j}_1,\mathrm{},\widehat{j}_{d_w1}\}\hfill \\ & =\{\begin{array}{cc}\{w_22,w_24,\mathrm{},w_22d_w\}\hfill & \mathrm{if}w0(mod4)\hfill \\ \{w_21,w_212,\mathrm{},w_212(d_w1)\}\hfill & \mathrm{if}w2(mod4).\hfill \end{array}\hfill \end{array}$$
Next we consider matrices $`๐_1`$ and $`๐_1`$ defined by
$`๐_1`$ $`:=[๐_{j_0},๐_{j_1},\mathrm{},๐_{j_{n_w1}}],`$
$`๐_1`$ $`:=[๐ฆ_{j_0},๐ฆ_{j_1},\mathrm{},๐ฆ_{j_{n_w1}}].`$
Note that $`๐_1`$ is the reduction modulo 2 of $`๐_1`$.
Since the column vectors $`๐ฆ_{j_0},๐ฆ_{j_1},\mathrm{},๐ฆ_{j_{n_w1}}`$ are linearly independent, we can choose $`n_w`$ rows, say rows $`i_0,i_1,\mathrm{},i_{n_w1}`$ of $`๐_1`$, so that the matrix $`๐_2`$ defined by
$$๐_2:=[m_{i_cj_a}](c,a=0,1,\mathrm{},n_w1)$$
has non-zero determinant. Then the matrix $`๐_2`$ defined by
$$๐_2:=[d_{i_cj_a}](c,a=0,1,\mathrm{},n_w1)$$
also has non-zero determinant. This is because $`๐_2`$ is the reduction modulo 2 of $`๐_2`$.
Furthermore, we set
$$๐_3:=[d_{i_c\widehat{j}_b}](c=0,1,\mathrm{},n_w1;b=0,1,\mathrm{},d_w1).$$
Then, from Lemma 3.3, we know that
$$๐_2\left[\begin{array}{c}t_{w,j_0}\\ t_{w,j_1}\\ \mathrm{}\\ t_{w,j_{n_w1}}\end{array}\right]+๐_3\left[\begin{array}{c}t_{w,\widehat{j}_0}\\ t_{w,\widehat{j}_1}\\ \mathrm{}\\ t_{w,\widehat{j}_{d_w1}}\end{array}\right]=\mathrm{๐},$$
from which it follows that
$$\left[\begin{array}{c}t_{w,j_0}\\ t_{w,j_1}\\ \mathrm{}\\ t_{w,j_{n_w1}}\end{array}\right]=๐_2^1๐_3\left[\begin{array}{c}t_{w,\widehat{j}_0}\\ t_{w,\widehat{j}_1}\\ \mathrm{}\\ t_{w,\widehat{j}_{d_w1}}\end{array}\right].$$
This implies that each of $`t_{w,j_0},t_{w,j_1},\mathrm{},t_{w,j_{n_w1}}`$ can be expressed as linear combinations of $`t_{w,\widehat{j}_0},t_{w,\widehat{j}_1},\mathrm{},t_{w,\widehat{j}_{d_w1}}`$ over $``$. On the other hand, by Theorem 2.1, we know that $`r_{w,1},r_{w,3},\mathrm{},r_{w,w1}`$ span $`S_{w+2}^{}`$. Thus $`t_{w,1},t_{w,2},\mathrm{},t_{w,w_2}`$ also span $`S_{w+2}^{}`$. These imply that
(7.2)
$$t_{w,\widehat{j}_0},t_{w,\widehat{j}_1},\mathrm{},t_{w,\widehat{j}_{d_w1}}\text{ span }S_{w+2}^{}.$$
From the fact (7.2) together with the identity (7.1), and the identities
$$t_{w,j}=s_{w,2j1}\text{ and }s_{w,n}=(1)^n\left(\genfrac{}{}{0pt}{}{w}{n}\right)r_{w,wn},$$
it follows that,
$`\{r_{w,w(2j1)}`$ $`|j=w_22,w_24,\mathrm{},w_22d_w\}`$
$`=\{r_{w,4i+1}|i=1,2,\mathrm{},d_w\}`$
span $`S_{w+2}^{}`$, when $`w0(mod4)`$, and
$`\{r_{w,w(2j1)}`$ $`|j=w_21,w_23,\mathrm{},w_212(d_w1)\}`$
$`=\{r_{w,4i1}|i=1,2,\mathrm{},d_w\}`$
span $`S_{w+2}^{}`$, when $`w2(mod4)`$.
Since
$$dimS_{w+2}^{}=d_w,$$
we know that
$$\{r_{w,4i\pm 1}|i=1,2,\mathrm{},d_w\}$$
is a basis of $`S_{w+2}^{}`$. This completes the proof of Theorem 2.2. โ
###### Proof of Theorem 2.3.
The cusp form $`R_{w,m}`$ is characterized by the formula:
$$r_{w,m}(f)=(R_{w,m},f)\mathrm{for}\mathrm{any}fS_{w+2}$$
with the Petersson inner product $`(,)`$. Then, from Theorem 2.2, it follows that
$$R_{w,4i\pm 1}(i=1,2,\mathrm{},d_w)$$
is a basis for $`S_{w+2}`$. This completes the proof. โ
## 8. Relationship between $`R_{w,n}`$, $`E_{w,n}`$ and $`S_{w,n}`$
The Sections from 8 to 14 are devoted to the study of Hecke operators. We also present proofs for Theorems 2.6, 2.7, 2.8 and 2.9.
First we show that, for $`n`$ odd, $`\alpha _{w+2}^+`$ maps $`c_{w,n}R_{w,n}`$ to $`E_{w,n}`$, and that $`\beta _w^+`$ maps $`E_{w,n}`$ to $`S_{w,n}`$. From this we will prove Theorems 2.6 and 2.7.
It is known (\[11, Theorem $`1^{}`$\]) that $`R_{w,n}(0<n<w;n\text{ odd})`$ corresponds to $`S_{w,n}`$ by the Eichler-Shimura monomorphism $`\beta _w^+\alpha _{w+2}^+:S_{w+2}๐ฐ_w^+`$. More precisely we know that
(8.1)
$$\beta _w^+\alpha _{w+2}^+(c_{w,n}R_{w,n})=S_{w,n}.$$
It was also shown (\[9, Theorem 5.2\]) that $`E_{w,n}(0<n<w;n\text{ odd})`$ corresponds to $`S_{w,n}`$ by the map $`\beta _w^+:_w^+๐ฐ_w^+`$. Namely,
(8.2)
$$\beta _w^+(E_{w,n})=S_{w,n}.$$
From (8.1) and (8.2), we obtain:
###### Lemma 8.1.
Let $`n`$ be an odd integer with $`0<n<w`$. Then it holds that
(8.3)
$$\alpha _{w+2}^+(c_{w,n}R_{w,n})=E_{w,n}.$$
###### Proof.
Since $`\mathrm{ker}\beta _w^+`$ is spanned by $`G_w`$, from (8.1) and (8.2), we know that
$$\alpha _{w+2}^+(c_{w,n}R_{w,n})E_{w,n}=cG_w$$
for some constant $`c`$.
On the other hand, we have
$`\alpha _{w+2}^+(c_{w,n}R_{w,n})(1,0)`$ $`={\displaystyle _0^{\mathrm{}}}c_{w,n}R_{w,n}(z)z^w๐z\text{ (by (}\text{2.1}\text{))}`$
$`=\beta _w^+\alpha _{w+2}^+(c_{w,n}R_{w,n})(1,0)\text{ (by (}\text{2.2}\text{))}`$
$`=S_{w,n}(1,0)\text{ (by (}\text{8.1}\text{))}`$
$`={\displaystyle \frac{B_{n+1}}{n+1}}{\displaystyle \frac{B_{\stackrel{~}{n}+1}}{\stackrel{~}{n}+1}}+{\displaystyle \frac{w+2}{B_{w+2}}}{\displaystyle \frac{B_{n+1}}{n+1}}{\displaystyle \frac{B_{\stackrel{~}{n}+1}}{\stackrel{~}{n}+1}}`$
$`=E_{w,n}(1,0).`$
Hence we know the constant $`c`$ is zero, from which we arrive at (8.3). โ
Now we are ready to prove Theorems 2.6 and 2.7.
###### Proofs of Theorems 2.6 and 2.7.
The Eichler-Shimura theorem (\[11, pp. 199โ200\]) asserts that the map
$$\beta _w^+\alpha _{w+2}^+:S_{w+2}_w^+๐ฐ_w^+$$
is an monomorphism and that the image $`\beta _w^+\alpha _{w+2}^+(S_{w+2})`$ and $`h^wk^w`$ span $`๐ฐ_w^+`$. Since $`R_{w,4i\pm 1}(i=1,2,\mathrm{},d_w)`$ is a basis of $`S_{w+2}`$ and $`\beta _w^+\alpha _{w+2}^+(c_{w,n}R_{w,n})=S_{w,n}`$, we know that $`S_{w,4i\pm 1}(i=1,2,\mathrm{},d_w)`$ and $`h^wk^w`$ is a basis of $`๐ฐ_w^+`$. This completes the proof of Theorem 2.6.
Next, by Theorem 2.4, we see that the map
$$\alpha _{w+2}^+:S_{w+2}_w^+$$
is an monomorphism and that the image $`\alpha _{w+2}^+(S_{w+2})`$, $`F_w`$ and $`G_w`$ span $`_w^+`$. Furthermore, from Lemma 8.1, we know that
$$\alpha _{w+2}^+(c_{w,4i\pm 1}R_{w,4i\pm 1})=E_{w,4i\pm 1}.$$
Hence we see that $`E_{w,4i\pm 1}(i=1,2,\mathrm{},d_w)`$, $`F_w`$ and $`G_w`$ is a basis of $`๐ฐ_w^+`$. This completes the proof of Theorem 2.7. โ
Next we show that, for $`n`$ even, $`\alpha _{w+2}^{}`$ maps $`c_{w,n}R_{w,n}`$ to $`E_{w,n}`$. For even $`n`$ with $`0<n<w`$, it was shown (\[11, Theorem $`1^{}`$\], \[9, Theorem 5.2\]) that
(8.4)
$$\beta _w^{}\alpha _{w+2}^{}(c_{w,n}R_{w,n})=S_{w,n}\text{ and }\beta _w^{}(E_{w,n})=S_{w,n}.$$
From the identities (8.4) and the fact that $`\beta _w^{}`$ is an isomorphism, we obtain:
###### Lemma 8.2.
Let $`n`$ be an even integer with $`0<n<w`$. Then it holds that
(8.5)
$$\alpha _{w+2}^{}(c_{w,n}R_{w,n})=E_{w,n}.$$
## 9. The action of the Hecke operator on $`R_{w,n}`$
In Definition 2.2, we introduced $`R_{w,n}^m`$ as a generalization of $`R_{w,n}`$, i.e. $`R_{w,n}^1=R_{w,n}`$. Naturally $`R_{w,n}^m`$ inherits most of the properties from $`R_{w,n}`$. Moreover we can understand the action of the Hecke operator on $`R_{w,n}`$ in terms of $`R_{w,n}^m`$. Indeed, we will show that $`T_mR_{w,n}=R_{w,n}^m`$.
For this purpose, we need the following notation. Let $`\gamma `$ be a matrix such that
$$\gamma =\left[\begin{array}{cc}a& b\\ c& d\end{array}\right]\text{ with }adbc>0\text{ },$$
and let
$$f:\{z|\mathrm{}z>0\}$$
be a map. Then, as usual, we define $`f|\gamma `$ by
$$(f|\gamma )(z):=(cz+d)^{w2}f(\frac{az+b}{cz+d}).$$
Let $`M_m`$ be the complete set of right coset representatives of $`H_m`$ modulo $`\mathrm{\Gamma }`$, defined as follows:
$$M_m:=\{\left[\begin{array}{cc}a& b\\ 0& d\end{array}\right]|ad=m,a>0,bmodd\}.$$
Recall that the action of the Hecke operation $`T_m`$ on $`f`$ is given by
$$T_mf(z):=m^{w+1}\underset{\alpha M_m}{}(f|\alpha )(z).$$
Now we are ready to prove the following lemma:
###### Lemma 9.1.
$$T_mR_{w,n}=R_{w,n}^m.$$
###### Proof.
We have
$`T_mR_{w,n}(z)`$ $`=c_{w,n}^1T_m{\displaystyle \underset{\begin{array}{c}\left[\begin{array}{cc}a& b\\ c& d\end{array}\right]\mathrm{\Gamma }\end{array}}{}}{\displaystyle \frac{1}{(az+b)^{n+1}(cz+d)^{\stackrel{~}{n}+1}}}`$
$`=c_{w,n}^1T_m{\displaystyle \underset{\gamma \mathrm{\Gamma }}{}}z^{n1}|\gamma `$
$`=m^{w+1}c_{w,n}^1{\displaystyle \underset{\alpha M_m}{}}{\displaystyle \underset{\gamma \mathrm{\Gamma }}{}}z^{n1}|\gamma |\alpha `$
$`=m^{w+1}c_{w,n}^1{\displaystyle \underset{\gamma ^{}H_m}{}}z^{n1}|\gamma ^{}`$
$`=m^{w+1}c_{w,n}^1{\displaystyle \underset{\begin{array}{c}\left[\begin{array}{cc}a^{}& b^{}\\ c^{}& d^{}\end{array}\right]H_m\end{array}}{}}{\displaystyle \frac{1}{(a^{}z+b^{})^{n+1}(c^{}z+d^{})^{\stackrel{~}{n}+1}}}`$
$`=R_{w,n}^m(z).`$
Here we used the identity
$$H_m=\underset{\alpha M_m}{}\mathrm{\Gamma }\alpha .$$
This completes the proof. โ
## 10. Properties of $`E_{w,n}^m`$
In this section, we study $`E_{w,n}^m`$ which generalizes $`E_{w,n}`$. We will show that $`E_{w,n}^m`$ is a well-defined Dedekind symbols of weight $`w`$. Furthermore we will show that $`E_{w,n}^m`$ is even or odd depending on $`n`$ is odd or even.
First we will show that the sum in Definition 2.2 (3) is finite by the following two lemmas (the proofs are left to the reader).
###### Lemma 10.1.
Let $`h`$ be a positive integer. Suppose that $`\left[\begin{array}{cc}a& b\\ c& d\end{array}\right]H_m`$ satisfies $`ac0`$ and
(10.1)
$$(\frac{k}{h}+\frac{b}{a})(\frac{k}{h}+\frac{d}{c})<0.$$
Then
$$mh|a|,mh|c|.$$
###### Lemma 10.2.
Let $`h`$ be a positive integer. We consider two conditions for $`\left[\begin{array}{cc}a& b\\ c& d\end{array}\right]H_m`$:
(10.2)
$$ac0;$$
(10.3)
$$(\frac{k}{h}+\frac{b}{a})(\frac{k}{h}+\frac{d}{c})<0.$$
Then there are only finite many $`\left[\begin{array}{cc}a& b\\ c& d\end{array}\right]H_m`$ which satisfy (10.2) and (10.3).
By Lemmas 10.1 and 10.2, we know that $`E_{w,n}^m`$ is well-defined since the sum in Definition 2.2 (3) is finite.
Next we will show that $`E_{w,n}^m`$ is a Dedekind symbol in the following lemma:
###### Lemma 10.3.
$`E_{w,n}^m`$ is an even $`(`$resp. odd$`)`$ Dedekind symbol of weight $`w`$ for $`n`$ odd $`(`$resp. even$`)`$.
###### Proof.
The argument used in \[8, Lemma 3.1 and Theorem 1.1\] for $`D_{w,n}`$ is also valid to establish that $`E_{w,n}^m`$ is an odd or even Dedekind symbol according as $`n`$ is even or odd. The reader should refer to for the detail.
It is clear that $`E_{w,n}^m`$ satisfies
$$E_{w,n}^m(ch,ck)=c^wE_{w,n}^m(h,k)\text{ for any }c^+.$$
Hence the weight of $`E_{w,n}^m`$ is $`w`$. This completes the proof. โ
## 11. The action of the Hecke operator on $`E_{w,n}`$
In this section we will study how $`T_m`$ acts on $`E_{w,n}`$. We will prove the following proposition.
###### Proposition 11.1.
Let $`m`$ and $`n`$ be positive integers such that $`0<n<w`$. Then it hold that
$$T_mE_{w,n}=E_{w,n}^m.$$
###### Proof.
Let $`E_{w,n,0}^m`$ and $`E_{w,n,0}`$ be defined by
$`E_{w,n,0}^m:`$ $`={\displaystyle \frac{1}{2}}{\displaystyle \underset{\begin{array}{c}\left[\begin{array}{cc}a& b\\ c& d\end{array}\right]H_m\\ ac0\\ (k/h+b/a)(k/h+d/c)<0\end{array}}{}}\mathrm{sgn}({\displaystyle \frac{k}{h}}+{\displaystyle \frac{b}{a}})(ak+bh)^{\stackrel{~}{n}}(ck+dh)^n,`$
$`E_{w,n,0}:`$ $`={\displaystyle \frac{1}{2}}{\displaystyle \underset{\begin{array}{c}\left[\begin{array}{cc}a& b\\ c& d\end{array}\right]\mathrm{\Gamma }\\ ac0\\ (k/h+b/a)(k/h+d/c)<0\end{array}}{}}\mathrm{sgn}({\displaystyle \frac{k}{h}}+{\displaystyle \frac{b}{a}})(ak+bh)^{\stackrel{~}{n}}(ck+dh)^n.`$
First we will prove that $`T_mE_{w,n,0}=E_{w,n,0}^m`$.
Noting that
$$H_m=\underset{\alpha M_m}{}\mathrm{\Gamma }\alpha =\underset{\begin{array}{c}ad=m\\ d>0\\ b(\mathrm{mod}d)\end{array}}{}\mathrm{\Gamma }\left[\begin{array}{cc}a& b\\ 0& d\end{array}\right],$$
we have
$`T_m`$ $`E_{w,n,0}(h,k)`$
$`={\displaystyle \underset{\begin{array}{c}ad=m\\ d>0\end{array}}{}}{\displaystyle \underset{b(\mathrm{mod}d)}{}}E_{w,n,0}(dh,ak+bh)`$
$`={\displaystyle \frac{1}{2}}{\displaystyle \underset{\begin{array}{c}ad=m\\ d>0\end{array}}{}}{\displaystyle \underset{b(\mathrm{mod}d)}{}}{\displaystyle \underset{\begin{array}{c}\left[\begin{array}{cc}a^{}& b^{}\\ c^{}& d^{}\end{array}\right]\mathrm{\Gamma }\\ a^{}c^{}0\\ \{(ak+bh)/dh+b^{}/a^{}\}\{(ak+bh)/dh+d^{}/c^{}\}<0\end{array}}{}}\mathrm{sgn}({\displaystyle \frac{ak+bh}{dh}}+{\displaystyle \frac{b^{}}{a^{}}})`$
$`\text{ }\text{ }\times \{a^{}(ak+bh)+b^{}dh\}^{\stackrel{~}{n}}\{c^{}(ak+bh)+d^{}dh\}^n`$
$`={\displaystyle \frac{1}{2}}{\displaystyle \underset{\begin{array}{c}ad=m\\ d>0\end{array}}{}}{\displaystyle \underset{b(\mathrm{mod}d)}{}}{\displaystyle \underset{\begin{array}{c}\left[\begin{array}{cc}a^{}& b^{}\\ c^{}& d^{}\end{array}\right]\mathrm{\Gamma }\\ a^{}c^{}0\\ \{k/h+(a^{}b+b^{}d)/a^{}a\}\{k/h+(c^{}b+d^{}d)/c^{}a\}<0\end{array}}{}}\mathrm{sgn}({\displaystyle \frac{k}{h}}+{\displaystyle \frac{a^{}b+b^{}d}{a^{}a}})`$
$`\text{ }\text{ }\times \{a^{}ak+(a^{}b+b^{}d)h\}^{\stackrel{~}{n}}\{c^{}ak+(c^{}b+d^{}d)h\}^n`$
$`={\displaystyle \frac{1}{2}}{\displaystyle \underset{\begin{array}{c}\left[\begin{array}{cc}a^{\prime \prime }& b^{\prime \prime }\\ c^{\prime \prime }& d^{\prime \prime }\end{array}\right]H_m\\ a^{\prime \prime }c^{\prime \prime }0\\ (k/h+b^{\prime \prime }/a^{\prime \prime })(k/h+d^{\prime \prime }/c^{\prime \prime })<0\end{array}}{}}\mathrm{sgn}({\displaystyle \frac{k}{h}}+{\displaystyle \frac{b^{\prime \prime }}{a^{\prime \prime }}})(a^{\prime \prime }k+b^{\prime \prime }h)^{\stackrel{~}{n}}(c^{\prime \prime }k+d^{\prime \prime }h)^n`$
$`=E_{w,n,0}^m(h,k)`$
where we set
$$\{\begin{array}{cc}\hfill a^{\prime \prime }& =a^{}a\hfill \\ \hfill b^{\prime \prime }& =a^{}b+b^{}d\hfill \\ \hfill c^{\prime \prime }& =c^{}a\hfill \\ \hfill d^{\prime \prime }& =c^{}b+d^{}d,\hfill \end{array}\text{ namely, }\left[\begin{array}{cc}a^{\prime \prime }& b^{\prime \prime }\\ c^{\prime \prime }& d^{\prime \prime }\end{array}\right]=\left[\begin{array}{cc}a^{}& b^{}\\ c^{}& d^{}\end{array}\right]\left[\begin{array}{cc}a& b\\ 0& d\end{array}\right].$$
Next, applying the formula
$$\underset{b=0}{\overset{c1}{}}\overline{B}_{n+1}(x+\frac{b}{c})=c^n\overline{B}_{n+1}(cx)(c^+),$$
we obtain the other terms of $`T_mE_{w,n}(h,k)`$ as follows:
$`T_m{\displaystyle \frac{\overline{B}_{n+1}(\frac{k}{h})h^w}{n+1}}`$ $`={\displaystyle \underset{\begin{array}{c}ad=m\\ d>0\end{array}}{}}{\displaystyle \underset{b(\mathrm{mod}d)}{}}{\displaystyle \frac{\overline{B}_{n+1}(\frac{ak+bh}{dh})(dh)^w}{n+1}}`$
$`={\displaystyle \underset{\begin{array}{c}ad=m\\ d>0\end{array}}{}}{\displaystyle \frac{d^n\overline{B}_{n+1}(d\frac{ak}{dh})(dh)^w}{n+1}}`$
$`={\displaystyle \underset{\begin{array}{c}ad=m\\ d>0\end{array}}{}}d^{\stackrel{~}{n}}{\displaystyle \frac{\overline{B}_{n+1}(\frac{ak}{h})h^w}{n+1}},`$
$`T_m{\displaystyle \frac{\overline{B}_{\stackrel{~}{n}+1}(\frac{k}{h})h^w}{\stackrel{~}{n}+1}}`$ $`={\displaystyle \underset{\begin{array}{c}ad=m\\ d>0\end{array}}{}}d^n{\displaystyle \frac{\overline{B}_{\stackrel{~}{n}+1}(\frac{ak}{h})h^w}{\stackrel{~}{n}+1}}`$
and
$`T_mh^w`$ $`={\displaystyle \underset{\begin{array}{c}ad=m\\ d>0\end{array}}{}}{\displaystyle \underset{b(\mathrm{mod}d)}{}}(dh)^w={\displaystyle \underset{\begin{array}{c}ad=m\\ d>0\end{array}}{}}d^{w+1}h^w=\sigma _{w+1}(m)h^w.`$
From these identities we obtain $`T_mE_{w,n}=E_{w,n}^m`$ completing the proof. โ
## 12. Reciprocity law for Dedekind symbols $`E_{w,n}^m`$
In this section we will prove the following reciprocity law for Dedekind symbols $`E_{w,n}^m`$:
###### Proposition 12.1.
Let $`m`$ and $`n`$ be positive integers such that $`0<n<w`$. Then it holds that
$$E_{w,n}^m(h,k)E_{w,n}^m(k,h)=S_{w,n}^m(h,k)$$
for any $`(h,k)^+\times ^+`$.
We need a few lemmas to prove Proposition 12.1. We consider the sum in Definition 2.2 and express this as the sum of three series:
$$\begin{array}{cc}\hfill \frac{1}{2}\underset{\begin{array}{c}\left[\begin{array}{cc}a& b\\ c& d\end{array}\right]H_m\\ ac0\\ (k/h+b/a)(k/h+d/c)<0\end{array}}{}\mathrm{sgn}(\frac{k}{h}+\frac{b}{a})& (ak+bh)^{\stackrel{~}{n}}(ck+dh)^n\hfill \\ & =E_{w,n,1}^m(h,k)+E_{w,n,2}^m(h,k)+E_{w,n,3}^m(h,k)\hfill \end{array}$$
where
$$E_{w,n,1}^m(h,k):=\frac{1}{2}\underset{\begin{array}{c}\left[\begin{array}{cc}a& b\\ c& d\end{array}\right]H_m\\ abcd>0\\ (k/h+b/a)(k/h+d/c)<0\end{array}}{}\mathrm{sgn}(\frac{k}{h}+\frac{b}{a})(ak+bh)^{\stackrel{~}{n}}(ck+dh)^n,$$
$$E_{w,n,2}^m(h,k):=\frac{1}{2}\underset{\begin{array}{c}\left[\begin{array}{cc}a& b\\ c& d\end{array}\right]H_m\\ abcd<0\\ (k/h+b/a)(k/h+d/c)<0\end{array}}{}\mathrm{sgn}(\frac{k}{h}+\frac{b}{a})(ak+bh)^{\stackrel{~}{n}}(ck+dh)^n$$
and
$$E_{w,n,3}^m(h,k):=\frac{1}{2}\underset{\begin{array}{c}\left[\begin{array}{cc}a& b\\ c& d\end{array}\right]H_m\\ ac0,bd=0\\ (k/h+b/a)(k/h+d/c)<0\end{array}}{}\mathrm{sgn}(\frac{k}{h}+\frac{b}{a})(ak+bh)^{\stackrel{~}{n}}(ck+dh)^n.$$
First we show a reciprocal property of $`E_{w,n,1}^m`$.
###### Lemma 12.2.
Suppose that $`(h,k)^+\times ^+`$. Then
$$E_{w,n,1}^m(h,k)=E_{w,n,1}^m(k,h).$$
###### Proof.
From the definition of $`E_{w,n,1}^m(h,k)`$, we have
$$E_{w,n,1}^m(h,k)=\frac{1}{2}\underset{\begin{array}{c}\left[\begin{array}{cc}a& b\\ c& d\end{array}\right]H_m\\ abcd>0\\ (k/h+b/a)(k/h+d/c)<0\end{array}}{}\mathrm{sgn}(\frac{k}{h}+\frac{b}{a})(ak+bh)^{\stackrel{~}{n}}(ck+dh)^n.$$
Replacing $`\left[\begin{array}{cc}a& b\\ c& d\end{array}\right]`$ with $`\left[\begin{array}{cc}b& a\\ d& c\end{array}\right]`$, we can transform this into the following formula:
(12.1)
$$\begin{array}{cc}\hfill \frac{1}{2}& \underset{\begin{array}{c}\left[\begin{array}{cc}a& b\\ c& d\end{array}\right]H_m\\ abcd>0\\ (k/h+b/a)(k/h+d/c)<0\end{array}}{}\mathrm{sgn}(\frac{k}{h}+\frac{b}{a})(ak+bh)^{\stackrel{~}{n}}(ck+dh)^n\hfill \\ & \text{ }=\frac{1}{2}\underset{\begin{array}{c}\left[\begin{array}{cc}a& b\\ c& d\end{array}\right]H_m\\ abcd>0\\ (k/ha/b)(k/hc/d)<0\end{array}}{}\mathrm{sgn}(\frac{k}{h}\frac{a}{b})(bkah)^{\stackrel{~}{n}}(dkch)^n.\hfill \end{array}$$
On the other hand, we have
(12.2)
$$E_{w,n,1}^m(k,h)=\frac{1}{2}\underset{\begin{array}{c}\left[\begin{array}{cc}a& b\\ c& d\end{array}\right]H_m\\ abcd>0\\ (h/k+b/a)(h/k+d/c)<0\end{array}}{}\mathrm{sgn}(\frac{h}{k}+\frac{b}{a})(ah+bk)^{\stackrel{~}{n}}(ch+dk)^n.$$
Now we compare the right hand sides of (12.1) and (12.2). First we see that the condition $`abcd>0`$ implies that $`\mathrm{sgn}(ac)=\mathrm{sgn}(bd)`$ and $`\mathrm{sgn}(ab)=\mathrm{sgn}(cd)`$. Next we see that
$$h^2bd(k/ha/b)(k/hc/d)=(ah+bk)(ch+dk)=k^2ac(h/k+b/a)(h/k+d/c).$$
From these we know that two conditions $`(k/ha/b)(k/hc/d)<0`$ and $`(h/k+b/a)(h/k+d/c)<0`$ are equivalent. Furthermore, under the condition that $`(k/ha/b)(k/hc/d)<0`$, we get $`0<k/h<a/b`$ or $`0<k/h<c/d`$, and then $`ab>0`$ or $`cd>0`$. This together with $`\mathrm{sgn}(ab)=\mathrm{sgn}(cd)`$ implies $`\mathrm{sgn}(ab)>0`$, and then $`\mathrm{sgn}(a)=\mathrm{sgn}(b)`$. Hence we have
$$\mathrm{sgn}(\frac{k}{h}\frac{a}{b})=\mathrm{sgn}(hb)\mathrm{sgn}(bkah)=\mathrm{sgn}(ka)\mathrm{sgn}(bkah)=\mathrm{sgn}(\frac{h}{k}+\frac{b}{a}).$$
In conclusion we see that the right hand sides of (12.1) and (12.2) are equal. This completes the proof. โ
Next we show a reciprocal property of $`E_{w,n,2}^m`$.
###### Lemma 12.3.
Suppose that $`(h,k)^+\times ^+`$. Then
$$E_{w,n,2}^m(h,k)E_{w,n,2}^m(k,h)=\frac{1}{2}\underset{\begin{array}{c}\left[\begin{array}{cc}a& b\\ c& d\end{array}\right]H_m\\ abcd<0\end{array}}{}\mathrm{sgn}(ab)(ak+bh)^{\stackrel{~}{n}}(ck+dh)^n.$$
###### Proof.
From the definition of $`E_{w,n,2}^m(h,k)`$, we have
$`E`$ $`{}_{w,n,2}{}^{m}(h,k)`$
$`={\displaystyle \frac{1}{2}}{\displaystyle \underset{\begin{array}{c}\left[\begin{array}{cc}a& b\\ c& d\end{array}\right]H_m\\ abcd<0\\ (k/h+b/a)(k/h+d/c)<0\end{array}}{}}\mathrm{sgn}({\displaystyle \frac{k}{h}}+{\displaystyle \frac{b}{a}})(ak+bh)^{\stackrel{~}{n}}(ck+dh)^n`$
$`={\displaystyle \frac{1}{2}}{\displaystyle \underset{\begin{array}{c}\left[\begin{array}{cc}a& b\\ c& d\end{array}\right]H_m\\ abcd<0\\ d/c>k/h>b/a\end{array}}{}}(ak+bh)^{\stackrel{~}{n}}(ck+dh)^n{\displaystyle \frac{1}{2}}{\displaystyle \underset{\begin{array}{c}\left[\begin{array}{cc}a& b\\ c& d\end{array}\right]H_m\\ abcd<0\\ d/c<k/h<b/a\end{array}}{}}(ak+bh)^{\stackrel{~}{n}}(ck+dh)^n`$
$`={\displaystyle \frac{1}{2}}{\displaystyle \underset{\begin{array}{c}\left[\begin{array}{cc}a& b\\ c& d\end{array}\right]H_m\\ abcd<0\\ d/c>k/h\end{array}}{}}(ak+bh)^{\stackrel{~}{n}}(ck+dh)^n{\displaystyle \frac{1}{2}}{\displaystyle \underset{\begin{array}{c}\left[\begin{array}{cc}a& b\\ c& d\end{array}\right]H_m\\ abcd<0\\ k/h<b/a\end{array}}{}}(ak+bh)^{\stackrel{~}{n}}(ck+dh)^n`$
$`\text{ }\text{(noting }\mathrm{sgn}(d/c)=\mathrm{sgn}(b/a)\text{)}.`$
On the other hand, replacing $`\left[\begin{array}{cc}a& b\\ c& d\end{array}\right]`$ with $`\left[\begin{array}{cc}b& a\\ d& c\end{array}\right]`$, we have
$``$ $`E_{w,n,2}^m(k,h)`$
$`={\displaystyle \frac{1}{2}}{\displaystyle \underset{\begin{array}{c}\left[\begin{array}{cc}a& b\\ c& d\end{array}\right]H_m\\ abcd<0\\ (h/k+b/a)(h/k+d/c)<0\end{array}}{}}\mathrm{sgn}({\displaystyle \frac{h}{k}}+{\displaystyle \frac{b}{a}})(ah+bk)^{\stackrel{~}{n}}(ch+dk)^n`$
$`={\displaystyle \frac{1}{2}}{\displaystyle \underset{\begin{array}{c}\left[\begin{array}{cc}a& b\\ c& d\end{array}\right]H_m\\ abcd<0\\ (h/ka/b)(h/kc/d)<0\end{array}}{}}\mathrm{sgn}({\displaystyle \frac{h}{k}}+{\displaystyle \frac{a}{b}})(bhak)^{\stackrel{~}{n}}(dhck)^n`$
(replacing $`\left[\begin{array}{cc}a& b\\ c& d\end{array}\right]`$ with $`\left[\begin{array}{cc}b& a\\ d& c\end{array}\right]`$)
$`={\displaystyle \frac{1}{2}}{\displaystyle \underset{\begin{array}{c}\left[\begin{array}{cc}a& b\\ c& d\end{array}\right]H_m\\ abcd<0\\ (h/ka/b)(h/kc/d)<0\end{array}}{}}\mathrm{sgn}({\displaystyle \frac{h}{k}}+{\displaystyle \frac{a}{b}})(ak+bh)^{\stackrel{~}{n}}(ck+dh)^n`$
$`={\displaystyle \frac{1}{2}}{\displaystyle \underset{\begin{array}{c}\left[\begin{array}{cc}a& b\\ c& d\end{array}\right]H_m\\ abcd<0\\ c/d>h/k>a/b\end{array}}{}}(ak+bh)^{\stackrel{~}{n}}(ck+dh)^n{\displaystyle \frac{1}{2}}{\displaystyle \underset{\begin{array}{c}\left[\begin{array}{cc}a& b\\ c& d\end{array}\right]H_m\\ abcd<0\\ c/d<h/k<a/b\end{array}}{}}(ak+bh)^{\stackrel{~}{n}}(ck+dh)^n`$
$`={\displaystyle \frac{1}{2}}{\displaystyle \underset{\begin{array}{c}\left[\begin{array}{cc}a& b\\ c& d\end{array}\right]H_m\\ abcd<0\\ c/d>h/k\end{array}}{}}(ak+bh)^{\stackrel{~}{n}}(ck+dh)^n{\displaystyle \frac{1}{2}}{\displaystyle \underset{\begin{array}{c}\left[\begin{array}{cc}a& b\\ c& d\end{array}\right]H_m\\ abcd<0\\ h/k<a/b\end{array}}{}}(ak+bh)^{\stackrel{~}{n}}(ck+dh)^n`$
$`={\displaystyle \frac{1}{2}}{\displaystyle \underset{\begin{array}{c}\left[\begin{array}{cc}a& b\\ c& d\end{array}\right]H_m\\ abcd<0\\ 0<d/c<k/h\end{array}}{}}(ak+bh)^{\stackrel{~}{n}}(ck+dh)^n{\displaystyle \frac{1}{2}}{\displaystyle \underset{\begin{array}{c}\left[\begin{array}{cc}a& b\\ c& d\end{array}\right]H_m\\ abcd<0\\ k/h>b/a>0\end{array}}{}}(ak+bh)^{\stackrel{~}{n}}(ck+dh)^n.`$
Hence we have
$`E_{w,n,2}^m(h,k)`$ $`E_{w,n,2}^m(k,h)`$
$`={\displaystyle \frac{1}{2}}{\displaystyle \underset{\begin{array}{c}\left[\begin{array}{cc}a& b\\ c& d\end{array}\right]H_m\\ abcd<0\\ d/c>k/h\end{array}}{}}{\displaystyle \frac{1}{2}}{\displaystyle \underset{\begin{array}{c}\left[\begin{array}{cc}a& b\\ c& d\end{array}\right]H_m\\ abcd<0\\ k/h<b/a\end{array}}{}}+{\displaystyle \frac{1}{2}}{\displaystyle \underset{\begin{array}{c}\left[\begin{array}{cc}a& b\\ c& d\end{array}\right]H_m\\ abcd<0\\ 0<d/c<k/h\end{array}}{}}{\displaystyle \frac{1}{2}}{\displaystyle \underset{\begin{array}{c}\left[\begin{array}{cc}a& b\\ c& d\end{array}\right]H_m\\ abcd<0\\ k/h>b/a>0\end{array}}{}}`$
$`={\displaystyle \frac{1}{2}}{\displaystyle \underset{\begin{array}{c}\left[\begin{array}{cc}a& b\\ c& d\end{array}\right]H_m\\ abcd<0\\ \mathrm{sgn}(d/c)>0\end{array}}{}}{\displaystyle \frac{1}{2}}{\displaystyle \underset{\begin{array}{c}\left[\begin{array}{cc}a& b\\ c& d\end{array}\right]H_m\\ abcd<0\\ \mathrm{sgn}(b/a)>0\end{array}}{}}`$
$`={\displaystyle \frac{1}{2}}{\displaystyle \underset{\begin{array}{c}\left[\begin{array}{cc}a& b\\ c& d\end{array}\right]H_m\\ abcd<0\\ \mathrm{sgn}(ab)>0\end{array}}{}}{\displaystyle \frac{1}{2}}{\displaystyle \underset{\begin{array}{c}\left[\begin{array}{cc}a& b\\ c& d\end{array}\right]H_m\\ abcd<0\\ \mathrm{sgn}(ab)<0\end{array}}{}}`$
$`={\displaystyle \frac{1}{2}}{\displaystyle \underset{\begin{array}{c}\left[\begin{array}{cc}a& b\\ c& d\end{array}\right]H_m\\ abcd<0\end{array}}{}}\mathrm{sgn}(ab)(ak+bh)^{\stackrel{~}{n}}(ck+dh)^n.`$
This completes the proof. โ
Finally we look into the term $`E_{w,n,3}^m(h,k)`$. We prove that $`E_{w,n,3}^m(h,k)`$ is expressed in terms of Bernoulli polynomials and Bernoulli functions.
###### Lemma 12.4.
Suppose that $`h`$ and $`k`$ are positive integers. Then
(12.3)
$$\begin{array}{cc}\hfill E_{w,n,3}^m(h,k)& =\underset{\begin{array}{c}ad=m\\ a>0\end{array}}{}d^{\stackrel{~}{n}}\left\{\frac{B_{n+1}(\frac{ah}{k})}{n+1}\frac{\overline{B}_{n+1}(\frac{ah}{k})}{n+1}\right\}k^w\hfill \\ & \text{ }(1)^n\underset{\begin{array}{c}ad=m\\ a>0\end{array}}{}d^n\left\{\frac{B_{\stackrel{~}{n}+1}(\frac{ah}{k})}{\stackrel{~}{n}+1}\frac{\overline{B}_{\stackrel{~}{n}+1}(\frac{ah}{k})}{\stackrel{~}{n}+1}\right\}k^w.\hfill \end{array}$$
###### Proof.
From the definition of $`E_{w,n,3}^m`$, we have
(12.4)
$$\begin{array}{cc}\hfill E_{w,n,3}^m(h,k)=& \frac{1}{2}\underset{\begin{array}{c}\left[\begin{array}{cc}a& b\\ c& d\end{array}\right]H_m\\ ac0,b=0\\ (k/h+b/a)(k/h+d/c)<0\end{array}}{}\mathrm{sgn}(\frac{k}{h}+\frac{b}{a})(ak+bh)^{\stackrel{~}{n}}(ck+dh)^n\hfill \\ & +\frac{1}{2}\underset{\begin{array}{c}\left[\begin{array}{cc}a& b\\ c& d\end{array}\right]H_m\\ ac0,d=0\\ (k/h+b/a)(k/h+d/c)<0\end{array}}{}\mathrm{sgn}(\frac{k}{h}+\frac{b}{a})(ak+bh)^{\stackrel{~}{n}}(ck+dh)^n.\hfill \end{array}$$
We first consider the first sum in the right hand side of (12.4). We will apply the following formula
(12.5)
$$\frac{B_{m+1}(x+1)}{m+1}\frac{B_{m+1}(x)}{m+1}=x^m(m0)$$
to transform the sum into the following formula:
$`{\displaystyle \frac{1}{2}}`$ $`{\displaystyle \underset{\begin{array}{c}\left[\begin{array}{cc}a& b\\ c& d\end{array}\right]H_m\\ ac0,b=0\\ (k/h+b/a)(k/h+d/c)<0\end{array}}{}}\mathrm{sgn}({\displaystyle \frac{k}{h}}+{\displaystyle \frac{b}{a}})(ak+bh)^{\stackrel{~}{n}}(ck+dh)^n`$
$`\text{ }={\displaystyle \underset{\begin{array}{c}\left[\begin{array}{cc}a& b\\ c& d\end{array}\right]H_m\\ ac0,b=0,a>0\\ (k/h+b/a)(k/h+d/c)<0\end{array}}{}}\mathrm{sgn}({\displaystyle \frac{k}{h}}+{\displaystyle \frac{b}{a}})(ak+bh)^{\stackrel{~}{n}}(ck+dh)^n`$
$`\text{ }={\displaystyle \underset{\begin{array}{c}ad=m,a>0\\ c0\\ (k/h+d/c)<0\end{array}}{}}(ak)^{\stackrel{~}{n}}(ck+dh)^n`$
$`\text{ }={\displaystyle \underset{\begin{array}{c}ad=m\\ a>0\end{array}}{}}{\displaystyle \underset{dh/k<c<0}{}}k^wa^{\stackrel{~}{n}}\left(c+{\displaystyle \frac{dh}{k}}\right)^n`$
$`\text{ }={\displaystyle \underset{\begin{array}{c}ad=m\\ a>0\end{array}}{}}{\displaystyle \underset{dh/k>c>0}{}}k^wa^{\stackrel{~}{n}}\left({\displaystyle \frac{dh}{k}}c\right)^n`$
$`\text{ }={\displaystyle \underset{\begin{array}{c}ad=m\\ a>0\end{array}}{}}{\displaystyle \underset{c=1}{\overset{\frac{dh}{k}}{}}}k^wa^{\stackrel{~}{n}}\left({\displaystyle \frac{dh}{k}}c\right)^n`$
$`\text{ }={\displaystyle \underset{\begin{array}{c}ad=m\\ a>0\end{array}}{}}k^wa^{\stackrel{~}{n}}\left\{{\displaystyle \frac{B_{n+1}(\frac{dh}{k})}{n+1}}{\displaystyle \frac{B_{n+1}(\frac{dh}{k}\frac{dh}{k})}{n+1}}\right\}\text{ (applying (}\text{12.5}\text{) repeatedly)}`$
$`\text{ }={\displaystyle \underset{\begin{array}{c}ad=m\\ a>0\end{array}}{}}a^{\stackrel{~}{n}}\left\{{\displaystyle \frac{B_{n+1}(\frac{dh}{k})}{n+1}}{\displaystyle \frac{\overline{B}_{n+1}(\frac{dh}{k})}{n+1}}\right\}k^w`$
$`\text{ }={\displaystyle \underset{\begin{array}{c}ad=m\\ a>0\end{array}}{}}d^{\stackrel{~}{n}}\left\{{\displaystyle \frac{B_{n+1}(\frac{ah}{k})}{n+1}}{\displaystyle \frac{\overline{B}_{n+1}(\frac{ah}{k})}{n+1}}\right\}k^w.`$
Similarly the second sum can be transformed into the following formula:
$`{\displaystyle \frac{1}{2}}{\displaystyle \underset{\begin{array}{c}\left[\begin{array}{cc}a& b\\ c& d\end{array}\right]H_m\\ ac0,d=0\\ (k/h+b/a)(k/h+d/c)<0\end{array}}{}}`$ $`\mathrm{sgn}({\displaystyle \frac{k}{h}}+{\displaystyle \frac{b}{a}})(ak+bh)^{\stackrel{~}{n}}(ck+dh)^n`$
$`=(1)^n{\displaystyle \underset{\begin{array}{c}ad=m\\ a>0\end{array}}{}}d^n\left\{{\displaystyle \frac{B_{\stackrel{~}{n}+1}(\frac{ah}{k})}{\stackrel{~}{n}+1}}{\displaystyle \frac{\overline{B}_{\stackrel{~}{n}+1}(\frac{ah}{k})}{\stackrel{~}{n}+1}}\right\}k^w.`$
Combining these we obtain (12.3). This completes the proof. โ
Similar computation yields the following result:
###### Lemma 12.5.
Suppose that $`h`$ and $`k`$ are positive integers. Then
$`E_{w,n,3}^m(k,h)`$ $`={\displaystyle \underset{\begin{array}{c}ad=m\\ a>0\end{array}}{}}d^n\left\{{\displaystyle \frac{B_{\stackrel{~}{n}+1}(\frac{ak}{h})}{\stackrel{~}{n}+1}}{\displaystyle \frac{\overline{B}_{\stackrel{~}{n}+1}(\frac{ak}{h})}{\stackrel{~}{n}+1}}\right\}h^w`$
$`\text{ }+(1)^{\stackrel{~}{n}+1}{\displaystyle \underset{\begin{array}{c}ad=m\\ a>0\end{array}}{}}d^{\stackrel{~}{n}}\left\{{\displaystyle \frac{B_{n+1}(\frac{ak}{h})}{n+1}}{\displaystyle \frac{\overline{B}_{n+1}(\frac{ak}{h})}{n+1}}\right\}h^w.`$
Now we are ready to prove Proposition 12.1.
###### Proof of Proposition 12.1.
From Lemmas 12.2, 12.3, 12.4 and 12.5, we have
$`E_{w,n}^m`$ $`(h,k)E_{w,n}^m(k,h)`$
$`=`$ $`E_{w,n,1}^m(h,k)E_{w,n,1}^m(k,h)+E_{w,n,2}^m(h,k)E_{w,n,2}^m(k,h)`$
$`+E_{w,n,3}^m(h,k)E_{w,n,3}^m(k,h)`$
$`+{\displaystyle \underset{\begin{array}{c}ad=m\\ a>0\end{array}}{}}\left\{(1)^nd^{\stackrel{~}{n}}{\displaystyle \frac{\overline{B}_{n+1}(\frac{ak}{h})h^w}{n+1}}d^n{\displaystyle \frac{\overline{B}_{\stackrel{~}{n}+1}(\frac{ak}{h})h^w}{\stackrel{~}{n}+1}}\right\}`$
$`+{\displaystyle \underset{\begin{array}{c}ad=m\\ a>0\end{array}}{}}\left\{d^{\stackrel{~}{n}}{\displaystyle \frac{\overline{B}_{n+1}(\frac{ah}{k})k^w}{n+1}}(1)^nd^n{\displaystyle \frac{\overline{B}_{\stackrel{~}{n}+1}(\frac{ah}{k})k^w}{\stackrel{~}{n}+1}}\right\}`$
$`+\{\begin{array}{cc}\sigma _{w+1}(m)\frac{w+2}{B_{w+2}}\frac{B_{n+1}}{n+1}\frac{B_{\stackrel{~}{n}+1}}{\stackrel{~}{n}+1}(h^wk^w)\hfill & \mathrm{if}n1(mod2)\hfill \\ 0\hfill & \mathrm{if}n0(mod2)\hfill \end{array}`$
$`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{\begin{array}{c}\left[\begin{array}{cc}a& b\\ c& d\end{array}\right]H_m\\ abcd<0\end{array}}{}}\mathrm{sgn}(ab)(ak+bh)^{\stackrel{~}{n}}(ck+dh)^n`$
$`+{\displaystyle \underset{\begin{array}{c}ad=m\\ a>0\end{array}}{}}d^{\stackrel{~}{n}}\left\{{\displaystyle \frac{B_{n+1}(\frac{ah}{k})}{n+1}}{\displaystyle \frac{\overline{B}_{n+1}(\frac{ah}{k})}{n+1}}\right\}k^w`$
$`\text{ }(1)^n{\displaystyle \underset{\begin{array}{c}ad=m\\ a>0\end{array}}{}}d^n\left\{{\displaystyle \frac{B_{\stackrel{~}{n}+1}(\frac{ah}{k})}{\stackrel{~}{n}+1}}{\displaystyle \frac{\overline{B}_{\stackrel{~}{n}+1}(\frac{ah}{k})}{\stackrel{~}{n}+1}}\right\}k^w`$
$`{\displaystyle \underset{\begin{array}{c}ad=m\\ a>0\end{array}}{}}d^n\left\{{\displaystyle \frac{B_{\stackrel{~}{n}+1}(\frac{ak}{h})}{\stackrel{~}{n}+1}}{\displaystyle \frac{\overline{B}_{\stackrel{~}{n}+1}(\frac{ak}{h})}{\stackrel{~}{n}+1}}\right\}h^w`$
$`\text{ }(1)^{\stackrel{~}{n}+1}{\displaystyle \underset{\begin{array}{c}ad=m\\ a>0\end{array}}{}}d^{\stackrel{~}{n}}\left\{{\displaystyle \frac{B_{n+1}(\frac{ak}{h})}{n+1}}{\displaystyle \frac{\overline{B}_{n+1}(\frac{ak}{h})}{n+1}}\right\}h^w`$
$`+{\displaystyle \underset{\begin{array}{c}ad=m\\ a>0\end{array}}{}}\left\{(1)^nd^{\stackrel{~}{n}}{\displaystyle \frac{\overline{B}_{n+1}(\frac{ak}{h})h^w}{n+1}}d^n{\displaystyle \frac{\overline{B}_{\stackrel{~}{n}+1}(\frac{ak}{h})h^w}{\stackrel{~}{n}+1}}\right\}`$
$`+{\displaystyle \underset{\begin{array}{c}ad=m\\ a>0\end{array}}{}}\left\{d^{\stackrel{~}{n}}{\displaystyle \frac{\overline{B}_{n+1}(\frac{ah}{k})k^w}{n+1}}(1)^nd^n{\displaystyle \frac{\overline{B}_{\stackrel{~}{n}+1}(\frac{ah}{k})k^w}{\stackrel{~}{n}+1}}\right\}`$
$`+\{\begin{array}{cc}\sigma _{w+1}(m)\frac{w+2}{B_{w+2}}\frac{B_{n+1}}{n+1}\frac{B_{\stackrel{~}{n}+1}}{\stackrel{~}{n}+1}(h^wk^w)\hfill & \mathrm{if}n1(mod2)\hfill \\ 0\hfill & \mathrm{if}n0(mod2)\hfill \end{array}`$
$`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{\begin{array}{c}\left[\begin{array}{cc}a& b\\ c& d\end{array}\right]H_m\\ abcd<0\end{array}}{}}\mathrm{sgn}(ab)(ak+bh)^{\stackrel{~}{n}}(ck+dh)^n`$
$`+{\displaystyle \underset{\begin{array}{c}ad=m\\ a>0\end{array}}{}}\{(1)^nd^{\stackrel{~}{n}}{\displaystyle \frac{B_{n+1}(\frac{ak}{h})h^w}{n+1}}+d^{\stackrel{~}{n}}{\displaystyle \frac{B_{n+1}(\frac{ah}{k})k^w}{n+1}}`$
$`\text{ }\text{ }d^n{\displaystyle \frac{B_{\stackrel{~}{n}+1}(\frac{ak}{h})h^w}{\stackrel{~}{n}+1}}(1)^nd^n{\displaystyle \frac{B_{\stackrel{~}{n}+1}(\frac{ah}{k})k^w}{\stackrel{~}{n}+1}}\}`$
$`+\{\begin{array}{cc}\sigma _{w+1}(m)\frac{w+2}{B_{w+2}}\frac{B_{n+1}}{n+1}\frac{B_{\stackrel{~}{n}+1}}{\stackrel{~}{n}+1}(h^wk^w)\hfill & \mathrm{if}n1(mod2)\hfill \\ 0\hfill & \mathrm{if}n0(mod2)\hfill \end{array}`$
$`=`$ $`S_{w,n}^m(h,k).`$
Thus we obtain the required reciprocity laws completing the proof of Proposition 12.1. โ
## 13. The action of the Hecke operator on $`S_{w,n}`$
Now we arrive at the following proposition which will be the key of obtaining our explicit formulas for Hecke operators:
###### Proposition 13.1.
Let $`m`$ be a positive integer, and let $`n`$ be an integer such that $`0<n<w`$. Then it hold that
$$T_m(S_{w,n})=S_{w,n}^m.$$
###### Proof.
We notice that the Hecke operators on the spaces of cusp forms, the spaces of Dedekind symbols and the spaces of period polynomials are all compatible (the diagrams (2.3) and (2.4)).
First we prove the case that $`n`$ is odd. We have
$`T_m(S_{w,n})`$ $`=T_m(\beta _w^+\alpha _{w+2}^+(c_{w,n}R_{w,n}))\text{(by (}\text{8.1}\text{))}`$
$`=\beta _w^+\alpha _{w+2}^+T_m(c_{w,n}R_{w,n})\text{(from the diagram (}\text{2.3}\text{))}`$
$`=\beta _w^+T_m\alpha _{w+2}^+(c_{w,n}R_{w,n})\text{(from the diagram (}\text{2.4}\text{))}`$
$`=\beta _w^+T_m(E_{w,n})\text{(by Lemma }\text{8.1}\text{)}`$
$`=\beta _w^+(E_{w,n}^m)\text{(by Proposition }\text{11.1}\text{)}`$
$`=S_{w,n}^m\text{(by Proposition }\text{12.1}\text{)}.`$
Substituting Lemma 8.2, $`\alpha _{w+2}^{}`$ and $`\beta _w^{}`$ for Lemma 8.1, $`\alpha _{w+2}^+`$ and $`\beta _w^+`$, we can prove the case that $`n`$ is even. This completes the proof. โ
We are also ready to prove Theorem 2.8.
###### Proof of Theorem 2.8.
It is obvious that Theorem 2.8 follows from Lemma 9.1, Proposition 11.1 and Proposition 13.1. โ
## 14. Explicit formulas for Hecke operators
In this section, we give a proof to Theorem 2.9.
###### Proof of Theorem 2.9.
We suppose that
$$๐_m=[\tau _{ij}](i,j=1,2,\mathrm{},d_w)$$
is the matrix representing the Hecke operator $`T_m:S_{w+2}S_{w+2}`$ with respect to the basis
$$c_{w,4i\pm 1}R_{w,4i\pm 1}(i=1,2,\mathrm{},d_w).$$
Namely, we suppose that
(14.1)
$$T_m(c_{w,4j\pm 1}R_{w,4j\pm 1})=\underset{i=1}{\overset{d_w}{}}\tau _{ij}c_{w,4i\pm 1}R_{w,4i\pm 1}(j=1,2,\mathrm{},d_w).$$
Then we have
$`S_{w,4j\pm 1}^m`$ $`=T_m(S_{w,4j\pm 1})\text{(by Proposition }\text{13.1}\text{)}`$
$`=T_m\beta _w^+\alpha _{w+2}^+(c_{w,4j\pm 1}R_{w,4j\pm 1})\text{(by (}\text{8.1}\text{))}`$
$`=\beta _w^+\alpha _{w+2}^+T_m(c_{w,4j\pm 1}R_{w,4j\pm 1})\text{(from the diagram (}\text{2.3}\text{))}`$
$`={\displaystyle \underset{i=1}{\overset{d_w}{}}}\tau _{ij}\beta _w^+\alpha _{w+2}^+(c_{w,4i\pm 1}R_{w,4i\pm 1})\text{(by (}\text{14.1}\text{))}`$
$`={\displaystyle \underset{i=1}{\overset{d_w}{}}}\tau _{ij}S_{w,4i\pm 1}\text{(by (}\text{8.1}\text{))}.`$
This gives
(14.2)
$$S_{w,4j\pm 1}^m=\underset{i=1}{\overset{d_w}{}}\tau _{ij}S_{w,4i\pm 1}(j=1,2,\mathrm{},d_w).$$
Furthermore, by taking inner product of $`S_{w,4k\pm 1}`$ and each side of (14.2), we have
$`S_{w,4k\pm 1},S_{w,4j\pm 1}^m`$ $`={\displaystyle \underset{i=1}{\overset{d_w}{}}}S_{w,4k\pm 1},S_{w,4i\pm 1}\tau _{ij}(j,k=1,2,\mathrm{},d_w).`$
Namely, we have
$`\left[\begin{array}{c}S_{w,4i\pm 1},S_{w,4j\pm 1}^m\end{array}\right]=\left[\begin{array}{c}S_{w,4i\pm 1},S_{w,4j\pm 1}\end{array}\right]๐_m(i,j=1,2,\mathrm{},d_w).`$
Now the linear independence of $`S_{w,4i\pm 1}(i=1,2,\mathrm{},d_w)`$ guarantees that the matrix
$`\left[\begin{array}{c}S_{w,4i\pm 1},S_{w,4j\pm 1}\end{array}\right](i,j=1,2,\mathrm{},d_w)`$
is non-singular. Thus we have
$$๐_m=\left[\begin{array}{c}S_{w,4i\pm 1},S_{w,4j\pm 1}\end{array}\right]^1\left[\begin{array}{c}S_{w,4i\pm 1},S_{w,4j\pm 1}^m\end{array}\right](i,j=1,2,\mathrm{},d_w)$$
completing the proof of (1).
Next we will prove (2). First we expand the term
$$\frac{1}{2}\underset{\begin{array}{c}\left[\begin{array}{cc}a& b\\ c& d\end{array}\right]H_m\\ abcd<0\end{array}}{}\mathrm{sgn}(ab)(ak+bh)^{\stackrel{~}{n}}(ck+dh)^n$$
as a polynomial in $`h`$ and $`k`$ as follows:
$`{\displaystyle \frac{1}{2}}`$ $`{\displaystyle \underset{\begin{array}{c}\left[\begin{array}{cc}a& b\\ c& d\end{array}\right]H_m\\ abcd<0\end{array}}{}}\mathrm{sgn}(ab)(ak+bh)^{\stackrel{~}{n}}(ck+dh)^n`$
$`={\displaystyle \underset{\begin{array}{c}\left[\begin{array}{cc}a& b\\ c& d\end{array}\right]H_m\\ a>0abcd<0\end{array}}{}}\mathrm{sgn}(ab)(ak+bh)^{\stackrel{~}{n}}(ck+dh)^n`$
$`={\displaystyle \underset{\begin{array}{c}\left[\begin{array}{cc}a& b\\ c& d\end{array}\right]H_m\\ a>0b>0abcd<0\end{array}}{}}(ak+bh)^{\stackrel{~}{n}}(ck+dh)^n{\displaystyle \underset{\begin{array}{c}\left[\begin{array}{cc}a& b\\ c& d\end{array}\right]H_m\\ a>0b<0abcd<0\end{array}}{}}(ak+bh)^{\stackrel{~}{n}}(ck+dh)^n`$
$`={\displaystyle \underset{\begin{array}{c}adbc=m\\ a>0b>0abcd<0\end{array}}{}}(ak+bh)^{\stackrel{~}{n}}(ck+dh)^n{\displaystyle \underset{\begin{array}{c}adbc=m\\ a>0b<0abcd<0\end{array}}{}}(ak+bh)^{\stackrel{~}{n}}(ck+dh)^n`$
$`={\displaystyle \underset{\mu =1}{\overset{m1}{}}}{\displaystyle \underset{\begin{array}{c}ad=\mu \\ a>0\end{array}}{}}{\displaystyle \underset{\begin{array}{c}bc=\mu m\\ b>0\end{array}}{}}(ak+bh)^{\stackrel{~}{n}}(ck+dh)^n{\displaystyle \underset{\mu =1}{\overset{m1}{}}}{\displaystyle \underset{\begin{array}{c}ad=\mu \\ a>0\end{array}}{}}{\displaystyle \underset{\begin{array}{c}bc=\mu m\\ b>0\end{array}}{}}(akbh)^{\stackrel{~}{n}}(ck+dh)^n`$
$`=2{\displaystyle \underset{\mu =1}{\overset{m1}{}}}{\displaystyle \underset{\begin{array}{c}ad=\mu \\ a>0\end{array}}{}}{\displaystyle \underset{\begin{array}{c}bc=\mu m\\ b>0\end{array}}{}}{\displaystyle \underset{\begin{array}{c}0k\stackrel{~}{n}\\ 0\mathrm{}n\\ k+n\mathrm{}\text{ odd}\end{array}}{}}\left({\displaystyle \genfrac{}{}{0pt}{}{\stackrel{~}{n}}{k}}\right)a^{\stackrel{~}{n}k}k^{\stackrel{~}{n}k}b^kh^k\left({\displaystyle \genfrac{}{}{0pt}{}{n}{\mathrm{}}}\right)c^n\mathrm{}k^n\mathrm{}d^{\mathrm{}}h^{\mathrm{}}`$
$`=2{\displaystyle \underset{\mu =1}{\overset{m1}{}}}{\displaystyle \underset{\begin{array}{c}ad=\mu \\ 0<a\end{array}}{}}{\displaystyle \underset{\begin{array}{c}bc=\mu m\\ 0<b\end{array}}{}}{\displaystyle \underset{\begin{array}{c}0k\stackrel{~}{n}\\ 0\mathrm{}n\\ k+\mathrm{}\text{ even}\end{array}}{}}\left({\displaystyle \genfrac{}{}{0pt}{}{\stackrel{~}{n}}{k}}\right)a^{\stackrel{~}{n}k}k^{\stackrel{~}{n}k}b^kh^k\left({\displaystyle \genfrac{}{}{0pt}{}{n}{\mathrm{}}}\right)({\displaystyle \frac{\mu m}{b}})^n\mathrm{}k^n\mathrm{}({\displaystyle \frac{\mu }{a}})^{\mathrm{}}h^{\mathrm{}}`$
$`=2{\displaystyle \underset{\mu =1}{\overset{m1}{}}}{\displaystyle \underset{\begin{array}{c}ad=\mu \\ 0<a\end{array}}{}}{\displaystyle \underset{\begin{array}{c}bc=\mu m\\ 0<b\end{array}}{}}{\displaystyle \underset{\begin{array}{c}0k\stackrel{~}{n}\\ 0\mathrm{}n\\ k+\mathrm{}\text{ even}\end{array}}{}}\mu ^{\mathrm{}}(\mu m)^n\mathrm{}\left({\displaystyle \genfrac{}{}{0pt}{}{\stackrel{~}{n}}{k}}\right)\left({\displaystyle \genfrac{}{}{0pt}{}{n}{\mathrm{}}}\right)a^{\stackrel{~}{n}k\mathrm{}}b^{k+\mathrm{}n}h^{k+\mathrm{}}k^{\stackrel{~}{n}+nk\mathrm{}}`$
$`=2{\displaystyle \underset{\mu =1}{\overset{m1}{}}}{\displaystyle \underset{\begin{array}{c}ad=\mu \\ 0<a\end{array}}{}}{\displaystyle \underset{\begin{array}{c}bc=\mu m\\ 0<b\end{array}}{}}{\displaystyle \underset{\begin{array}{c}\nu =0\\ \nu \text{ even}\end{array}}{\overset{w}{}}}{\displaystyle \underset{\lambda =\mathrm{max}(0,\nu \stackrel{~}{n})}{\overset{\mathrm{min}(n,\nu )}{}}}\mu ^\lambda (\mu m)^{n\lambda }\left({\displaystyle \genfrac{}{}{0pt}{}{\stackrel{~}{n}}{\nu \lambda }}\right)\left({\displaystyle \genfrac{}{}{0pt}{}{n}{\lambda }}\right)a^{\stackrel{~}{n}\nu }b^{\nu n}h^\nu k^{w\nu }`$
(setting $`\nu =k+\mathrm{}`$ and $`\lambda =\mathrm{}`$)
$`=2{\displaystyle \underset{\mu =1}{\overset{m1}{}}}{\displaystyle \underset{\begin{array}{c}\nu =0\\ \nu \text{ even}\end{array}}{\overset{w}{}}}{\displaystyle \underset{\lambda =\mathrm{max}(0,\nu \stackrel{~}{n})}{\overset{\mathrm{min}(n,\nu )}{}}}\mu ^\lambda (\mu m)^{n\lambda }\left({\displaystyle \genfrac{}{}{0pt}{}{\stackrel{~}{n}}{\nu \lambda }}\right)\left({\displaystyle \genfrac{}{}{0pt}{}{n}{\lambda }}\right)\sigma _{\stackrel{~}{n}\nu }(\mu )\sigma _{\nu n}(m\mu )h^\nu k^{w\nu }`$
$`=2{\displaystyle \underset{\begin{array}{c}\nu =0\\ \nu \text{ even}\end{array}}{\overset{w}{}}}{\displaystyle \underset{\mu =1}{\overset{m1}{}}}{\displaystyle \underset{\lambda =\mathrm{max}(0,\nu \stackrel{~}{n})}{\overset{\mathrm{min}(n,\nu )}{}}}\mu ^\lambda (\mu m)^{n\lambda }\left({\displaystyle \genfrac{}{}{0pt}{}{\stackrel{~}{n}}{\nu \lambda }}\right)\left({\displaystyle \genfrac{}{}{0pt}{}{n}{\lambda }}\right)\sigma _{\stackrel{~}{n}\nu }(\mu )\sigma _{\nu n}(m\mu )h^\nu k^{w\nu }.`$
We also calculate other terms:
$`{\displaystyle \underset{\begin{array}{c}ad=m\\ a>0\end{array}}{}}d^{\stackrel{~}{n}}{\displaystyle \frac{B_{n+1}(\frac{ak}{h})h^w}{n+1}}`$ $`={\displaystyle \underset{\begin{array}{c}ad=m\\ a>0\end{array}}{}}{\displaystyle \frac{d^{\stackrel{~}{n}}}{n+1}}{\displaystyle \underset{\mu =0}{\overset{n+1}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{n+1}{\mu }}\right)B_\mu \left({\displaystyle \frac{ak}{h}}\right)^{n+1\mu }h^w`$
$`={\displaystyle \underset{\begin{array}{c}ad=m\\ a>0\end{array}}{}}{\displaystyle \frac{m^{\stackrel{~}{n}}}{n+1}}{\displaystyle \underset{\mu =0}{\overset{n+1}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{n+1}{\mu }}\right)B_\mu a^{n+1\mu \stackrel{~}{n}}h^{wn1+\mu }k^{n+1\mu }`$
$`={\displaystyle \frac{m^{\stackrel{~}{n}}}{n+1}}{\displaystyle \underset{\mu =0}{\overset{n+1}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{n+1}{\mu }}\right)B_\mu \sigma _{n+1\mu \stackrel{~}{n}}(m)h^{wn1+\mu }k^{n+1\mu }`$
$`={\displaystyle \frac{m^{\stackrel{~}{n}}}{n+1}}{\displaystyle \underset{\nu =\stackrel{~}{n}1}{\overset{w}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{n+1}{\nu \stackrel{~}{n}+1}}\right)B_{\nu \stackrel{~}{n}+1}\sigma _{n\nu }(m)h^\nu k^{w\nu },`$
$`{\displaystyle \underset{\begin{array}{c}ad=m\\ a>0\end{array}}{}}d^{\stackrel{~}{n}}{\displaystyle \frac{B_{n+1}(\frac{ah}{k})k^w}{n+1}}`$ $`={\displaystyle \frac{m^{\stackrel{~}{n}}}{n+1}}{\displaystyle \underset{\nu =0}{\overset{n+1}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{n+1}{n\nu +1}}\right)B_{n\nu +1}\sigma _{\nu \stackrel{~}{n}}(m)h^\nu k^{w\nu }.`$
Summing up these identities, we obtain the formula for $`S_{w,n}^m`$ in (2) of Theorem 2.9. This completes the proof. โ
Appendix : Computing matrices representing the Hecke operators and their characteristic polynomials
Finally in this section, we will demonstrate a Mathematica<sup>1</sup><sup>1</sup>1Mathematica is a trademark of Wolfram Research, Inc. program which, for given $`w`$ and $`m`$, yields a matrix representing the Hecke operator $`T_m`$ on $`S_{w+2}`$, and its characteristic polynomial. The program is a straightforward application of Theorem 2.9 using built-in functions BernoulliB and DivisorSigma.
```
w=28; (* set a positive even integer *)
m=7; (* set a positive integer *)
If[Mod[w,12]==0,dw=Quotient[w+2,12]-1,dw=Quotient[w+2,12]];
If[OddQ[w]||(dw<=0)||(m<=0),
Print["w odd or dw<=0 or m<=0"];Exit[]];
swn=Array[a,{2,dw,w+1}];
For[j=1,j<=2,j++,
If[j<=1,em=1,em=m];
For[i=1,i<=dw,i++,
If[Mod[w,4]==0,n=4*i+1,n=4*i-1]; tn=w-n;
For[nu=1,nu<=w+1,nu++,
swn[[j,i,nu]]=0;
lmin=Min[n,nu-1]; lmax=Max[0,nu-1-tn];
If[OddQ[nu],
For[mu=1,mu<=em-1,mu++,
For[lambda=lmax,lambda<=lmin,lambda++,
swn[[j,i,nu]]=
swn[[j,i,nu]]+2*(mu^lambda)*((mu-em)^(n-lambda))*
Binomial[tn,nu-1-lambda]*Binomial[n,lambda]*
DivisorSigma[tn-nu+1,mu]*DivisorSigma[nu-1-n,em-mu];
];
];
];
];
For[nu=tn,nu<=w+1,nu++,
swn[[j,i,nu]]=swn[[j,i,nu]]-(em^tn)*Binomial[n+1,nu-tn]*
BernoulliB[nu-tn]*DivisorSigma[n+1-nu,em]/(n+1);
];
For[nu=n,nu<=w+1,nu++,
swn[[j,i,nu]]=swn[[j,i,nu]]-(em^n)*Binomial[tn+1,nu-n]*
BernoulliB[nu-n]*DivisorSigma[tn+1-nu,em]/(tn+1);
];
For[nu=1,nu<=n+2,nu++,
swn[[j,i,nu]]=swn[[j,i,nu]]+(em^tn)*Binomial[n+1,n-nu+2]*
BernoulliB[n-nu+2]*DivisorSigma[nu-1-tn,em]/(n+1);
];
For[nu=1,nu<=tn+2,nu++,
swn[[j,i,nu]]=swn[[j,i,nu]]+(em^n)*Binomial[tn+1,tn-nu+2]*
BernoulliB[tn-nu+2]*DivisorSigma[nu-1-n,em]/(tn+1);
];
swn[[j,i,w+1]]=
swn[[j,i,w+1]]+DivisorSigma[w+1,em]*((w+2)*BernoulliB[n+1]*
BernoulliB[tn+1])/(BernoulliB[w+2]*(n+1)*(tn+1));
swn[[j,i,1]]=
swn[[j,i,1]]-DivisorSigma[w+1,em]*((w+2)*BernoulliB[n+1]*
BernoulliB[tn+1])/(BernoulliB[w+2]*(n+1)*(tn+1));
];
];
s1=Array[b,{dw,dw}];
s2=Array[c,{dw,dw}];
For[i=1,i<=dw,i++,
For[j=1,j<=dw,j++,s1[[i,j]]=swn[[1,i]].swn[[1,j]];
s2[[i,j]]=swn[[1,i]].swn[[2,j]]]];
twm=Inverse[s1].s2;
Print["representation matrix=",MatrixForm[twm]];
Print["characteristic polynomial=",Det[x*IdentityMatrix[dw]-twm]];
```
For example, the output in the case that $`w=28`$ and $`m=7`$ is as follows:
representation matrix= $`\left(\begin{array}{cc}\frac{597428921326303528}{6439}& \frac{4321468293778944}{6439}\\ \frac{79904984173167605760}{6439}& \frac{577981127961754328}{6439}\end{array}\right)`$,
characteristic polynomial= $`101633401431659687926336+3020312682800x+x^2`$.
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# Galaxy Luminosity Functions to ๐งโผ1: DEEP2 vs. COMBO-17 and Implications for Red Galaxy FormationBased on observations taken at the W. M. Keck Observatory.
## 1. Introduction
A major handicap in lookback studies of galaxy evolution is the inability to follow the evolution of any one galaxy over time. Instead, we see only snapshots of the galaxy population at different epochs, and it is difficult to identify objects at one epoch with their precursors and descendants at different epochs. One of the most important tools to solve this problem is precision counts of galaxies, which can quantify the โflowโ of galaxies in parameter space as masses, morphologies, and stellar populations change. The luminosity function of galaxies was historically the first such tool, but the concept is rapidly being broadened to include counts as a function of mass, internal velocity, color, and other parameters.
A further difficulty is caused by the fact that nearly all functions known to date are unimodal and lack clear features that demarcate one class of galaxies from another. In other words, galaxies tend to populate one big โcloudโ in most parameter spaces rather than separate clumps. This makes interpretation difficult, as sub-counts depend on how boundaries within these clouds are defined, and it is not clear whether (or how) these boundaries should be adjusted to follow galaxy evolution. As a result, we often cannot tell whether a change in the number of galaxies in any particular bin is due to a change in the overall number of galaxies or to the motion of galaxies in and out of that bin from neighboring bins. The latter problem is further exacerbated by the fact that samples are usually size- or brightness-limited and population numbers on the other side of these limits are not known. Finally, uncertain errors can smear counts from one bin to another. All of these problems will be modeled when a full theory of galaxy formation is available to predict how galaxies evolve in in every measured parameter. In the meantime, it is hard to break the population into well motivated sub-populations, and we therefore lack the means to obtain more finely divided knowledge.
Amidst this sea of unimodal functions, one function stands out on account of its uniquely bimodal character, namely, the color function. This is visible in the color-magnitude diagram, where early-type E/S0s populate a narrow red sequence that is separated from bluer, star-forming spirals by a shallow valley (Strateva et al. 2001; Hogg et al. 2003; Balogh et al. 2004; Baldry et al. 2004 and references therein). A similar division extends to at least $`z1`$ (Im et al. 2002; Bell et al. 2004b; Weiner et al. 2005, Willmer et al. 2005) and possibly beyond (Giallongo et al. 2005). A bimodal distribution is also seen in other parameters such as spectral class (Madgwick et al. 2002, 2003) and morphologies, metallicities, and star formation rates (Kauffmann et al. 2003a,b), but color is by far the easiest to measure. Thus, not only does color sort galaxies cleanly into bins, it is also highly relevant to the emergence of the Hubble sequence.
However, to exploit this opportunity requires highly accurate counts, as the expected effects are not large. For example, counts of red galaxies in the COMBO-17 survey were seen to evolve by only a factor of a few since $`z=1`$ (Bell et al. 2004b). Even this small number has vital implications for galaxy formation (see below), but confirming and improving the measurement clearly requires accuracies of order 10-20%. Few previous measurements of distant luminosity functions have attained this accuracy. Red galaxies are especially difficult because of their high clustering, which necessitates large samples over a large number of statistically uncorrelated regions. We are only just now coming to appreciate how formidable the problem of cosmic variance really is (e.g., Somerville et al. 2004)
The present paper addresses these challenges by combining two large surveys, DEEP2 and COMBO-17, to create the largest database yet analyzed of galaxies with $`z>`$ 0.8, containing 39,000 galaxies in total, with 15,600 beyond $`z=0.8`$. Further checks are provided by pilot measurements from the DEEP1 survey. The entire sample is large enough and dispersed enough over the sky that cosmic variance and Poisson fluctuations are reduced to 7-15% per redshift bin. The samples were selected and measured in different waysโDEEP2 redshifts are spectroscopic, while COMBO-17โs are photometricโand thus provide an important check on one another. Finally, color bimodality is used to divide red galaxies from blue galaxies at all epochs. The red luminosity function is rederived and compared to the previous results of Bell et al. (2004b, hereafter B04), while the blue function is presented in Paper I and here for the first time (in a sample of this size) and offers a important foil for considering the behavior of the red function. DEEP2 data and COMBO-17 data are found to agree well in all major respects, and the principal conclusions appear to be robust.
Our most important result is to confirm the recent rise in the number of massive red galaxies at fixed stellar mass found by B04. In contrast, the number density of massive blue galaxies has remained essentially constant since $`z1`$. A second important conclusion is that the characteristic $`B`$-band luminosity $`M_B^{}`$ of both red and blue populations has dimmed by about the same factor: we find $``$ 1.3 mag per unit redshift for both red and blue populations since $`z=1`$. The rise in the number of massive red galaxies implies that most early-type galaxies assumed their final form at relatively late times, below $`z=1`$, where the process can be studied in detail. The late emergence of spheroidal galaxies disagrees with classic high-redshift, monolithic collapse models for spheroid formation but seems to be consistent with large amounts of other data, as reviewed in ยง6.
The remainder of this Introduction reviews previous measurements of luminosity functions. The subject has a venerable history (e.g., Binggeli et al. 1988; Tresse 1999; de Lapparent et al. 2003), with determinations ranging from low to high redshift using field and cluster galaxies. Accurate determinations of local field luminosity functions have finally become available from the 2 Degree Field Galaxy Redshift Survey (2dFGRS, Norberg et al. 2002) and the Sloan Digital Sky Survey (SDSS, Blanton et al. 2003; Bell et al. 2003), providing reliable local benchmarks against which evolution can be measured.
The dependence of the galaxy luminosity function on the internal properties of galaxies has been known since Sandage, Binggeli & Tammann (1985) showed that the shape and magnitude of luminosity functions in the Virgo cluster depend on galaxy morphology and luminosity class. This dependence on internal characteristics is also seen in local field galaxies when morphologies (Marzke, Huchra & Geller 1994; Marzke et al. 1998; Marinoni et al. 1999), colors (Lilly et al. 1995b; Marzke & da Costa 1997; Lin et al. 1999; Blanton et al. 2001) and spectral types (e.g., Heyl et al. 1997; Bromley et al. 1998; Folkes et al. 1999; Cohen 2002; Magdwick et al. 2002; de Lapparent et al. 2003) are considered.
Early studies of the evolution of the galaxy luminosity function to $`z`$ 1 used samples of a few hundred galaxies (e.g., Cowie et al. 1996; Brinchmann et al. 1998; Lin et al. 1999; Cohen 2002; Im et al. 2002; de Lapparent et al. 2003). In a landmark paper using the Canada France Redshift Survey (hereafter CFRS), Lilly et al. (1995b) claimed that the evolution of the luminosity function is coupled to internal properties, being strongly correlated with color and, to a lesser extent, with luminosity. Dividing red from blue galaxies using the median spectral type of the sample, the authors claimed a steepening in faint-end slope for blue galaxies at redshifts beyond $`z>0.5`$, while red galaxies showed little change in either luminosity or number density over the redshift range covered, $`0.05z1`$.
A conclusion that evolution depends on internal properties was also reached by Cowie et al. (1996), based on a sample reaching to $`z1.6`$. \[OII\] fluxes were used to estimate star formation rates and $`K`$-band photometry to estimate stellar masses. They found that most of the luminosity evolution since $`z`$1 is due to blue galaxies with small masses but high star formation rates. More massive galaxies were relatively stable in numbers, particularly in the $`K`$-band, while the $`B`$-band showed modest number evolution. They concluded that the characteristic mass of galaxies undergoing intense star formation decreases over time, which they termed โdownsizing.โ
Cohen (2002) measured galaxies in a region centered on the Hubble Deep Field and Flanking Fields and, in contrast to CFRS, found that the luminosity functions of several different spectral classes of galaxies showed no strong evidence of change in faint-end slope to $`z1`$, and further that the value of this slope is comparable to the local value. Galaxies with spectra dominated by absorption lines at $`z=1`$ were brighter by $``$1.5 magnitudes relative to local ones, while galaxies with strong \[OII\] brightened by $``$0.75 magnitudes at $`z1`$.
Im et al. (2002) measured evolution in the luminosity function of morphologically normal red early-type galaxies by selecting distant galaxies to match local E/S0s in both morphology and color. This survey is related to the present work, as it utilized DEEP1 redshifts, supplemented by photo-zโs. The luminosity function of early-type galaxies showed a brightening of 1.1-1.9 magnitudes in rest-frame $`B`$ from $`z=0`$ to $`z0.8`$, but number density was relatively static over the same epoch. We will have more to say about this work later.
Similar brightening for early-type galaxies was also found by Bernardi et al. (2003) in the Sloan Digital Sky Survey, where a brightening of $``$1.15 magnitudes per unit redshift back in time was derived based on a sample reaching to $`z0.3`$.
In a later study going 2 magnitudes fainter than Im et al., Cross et al. (2004) studied the faint end of the luminosity functions of both red-selected galaxies and morphologically-selected early-type galaxies using ACS images and photometric redshifts. The red-selected luminosity function was found to turn over steeply at faint magnitudes, whereas the morphologically selected sample was flat. The difference was attributed to blue spheroids, which filled in the counts at faint levels in the morphologically-selected sample.
The evolution of the luminosity function as a function of color since $`z`$ 1 was also investigated by Pozzetti et al. (2003), who used the near-infrared-selected K20 survey of Cimatti et al. (2002b) and divided the sample using the color of Sa galaxies. They found a modest rise of at most 30% in the number of red galaxies after $`z=1`$ and concluded that most bright red galaxies were already in place by $`z`$ 1.3.
The distant surveys just cited clearly disagree on many points, including the numbers of galaxies, shapes of luminosity functions, and degree of fading over time. However, these early surveys typically cover only a few tens of square arc minutes and in retrospect are seen to be subject to large cosmic variance (see below). More recent surveys containing several thousands of galaxies are just now beginning to provide more robust measurements of the luminosity function. A large survey by Ilbert et al. (2004, VVDS) using the VIMOS spectrograph on the ESO Very Large Telescope has measured the evolution of the total galaxy luminosity function to $`z2`$ using a sample of $``$11,000 galaxies to $`I_{AB}=24.0`$ with spectroscopic redshifts. The authors find that $`M_B^{}`$ for all galaxies has faded by 1.6 to 2.2 mag from $`z=2`$ to $`z=0.05`$ but that the number density of all galaxies has remained nearly constant. They also suggest a possible steepening in faint-end slope at $`z=1`$.
The DEEP2 Survey (Davis et al. 2003) employs spectroscopic redshifts to measure distances and internal kinematics for $``$40,000 galaxies in four regions of the sky. In order to probe galaxies at redshifts $`z`$ 1, galaxies are pre-selected in three of the regions to have $`z>0.7`$ using $`BRI`$ colors. The luminosity function analysis for the first third of the DEEP2 survey is presented in a companion paper to this one by Willmer et al. (2005, hereafter Paper I), which uses $`UB`$ color bimodality to study how the luminosity functions of blue and red galaxies change with redshift. The results show significantly fewer blue galaxies with time at fixed absolute magnitude, which is well modeled by a fading in $`M_B^{}`$ together with roughly constant number density. Counts for red galaxies, in contrast, show little change at fixed absolute magnitude and, when fitted with Schechter functions, show a similar fading in $`M_B^{}`$ with time but a formally significant rise in number density.
An alternative strategy to create large samples of galaxies is the use of photometric redshifts. In spite of lower precision, photometric redshifts yield a larger number of redshifts per unit telescope time and enable distances to be measured for galaxies that are too faint for spectroscopy. This approach has been pursued using both space-based (e.g., Takeuchi et al. 2000; Poli et al. 2001; Bolzonella et al. 2002) and ground-based data (e.g., Fried et al. 2001; Drory et al. 2003; Wolf et al. 2003, hereafter W03; Chen et al. 2003; Gabasch et al. 2004). Of the photometric redshift surveys, the ones most comparable to the present work are COMBO-17 (W03) and the FORS Deep Field (Gabasch et al. 2004, FDF). The latter used a sample of more than 5,500 galaxies down to $`I_{AB}=26.8`$ to measure the total restrame $`B`$-band luminosity function from $`z0.4`$ to $`z`$ 4. Like DEEP2 and VVDS, they found a constant number of galaxies back in time but did not find a steepening of faint-end slope despite the fact that their sample goes ten times fainter than VVDS.
COMBO-17 (Wolf et al. 2001, W03) contains $``$28,000 galaxies. Aside from the use of photometric redshifts, it is similar to the DEEP2 survey in terms of depth and coverage. The first luminosity-function analysis of COMBO-17 by W03 divided the sample into bins of fixed spectral type that did not evolve with redshift. Some of the evolutionary trends that were discovered may have reflected color evolution between these fixed spectral bins rather than changes in overall numbers. The approach was changed in a follow-up analysis using the same database by Bell et al. (2004b, B04), who used an evolving color cut based on bimodality to study red galaxies only. As noted earlier, this work obtained the important new result that red galaxies were not only brighter in the past (by $`>`$1 mag) but were also fewer in number, by at least a factor of two at $`z=1`$. This claim based on counts was buttressed by a separate argument based on the luminosity density of bright galaxies. This quantity, $`j_B`$, can be measured more accurately than either $`L^{}`$ or number density alone and was found to hold roughly constant since $`z1`$. Since stellar population models predict a fading of red stellar populations by 1โ2 mag between $`z=1`$ and now (see ยง5), the total stellar mass bound up in (bright) red galaxies must be increasing by about the same factor, providing additional evidence for growth in the red galaxy population since that time.
A specific evolutionary scenario proposed by B04 had the majority of present-day massive E/S0s moving onto the red sequence after $`z1`$. The stellar populations of such galaxies would age passively once galaxies were on the red sequence, but individual galaxies would continue to increase their stellar masses via mergers along the sequence, as predicted by the hierarchical model of galaxy formation.
Aside from B04, few works have used color bimodality to measure the luminosity functions of red and blue galaxies separately. One of those that has is by Giallongo et al. (2005), who used a mixture of deep and shallow data in four fields containing 1,434 galaxies. Dividing galaxies both by $`UV`$ color and by star-formation rate, they found that color bimodality persists to $`z2`$ but that the number density of red galaxies drops steeply beyond $`z=1`$. They do not give numbers for $`z=1`$ specifically, which makes quantitative comparison with our results difficult, but our conclusions below agree at least qualitatively with theirs.
The foregoing summary illustrates that information on distant luminosity functions is still fragmentary and often contradictory. In particular, the claim by Bell et al. (2004b) for the emergence of red galaxies after $`z=1`$ has not yet been checked. A major impediment is cosmic varianceโmany of the above samples, especially the early and/or deepest ones, cover only a few tens of square arc min, for which the rms cosmic variance is $``$50% per $`\mathrm{\Delta }z=0.2`$ at $`z=1`$, being even greater for red galaxies (Somerville et al. 2004). Since number-density evolution by factors of two or three is at issue (B04), definitive results cannot be obtained using such small areas, and larger samples are needed. DEEP2 and COMBO-17 fill that need; each sample is large enough on its own to give statistically meaningful results, thus providing significant checks on one another.
The paper is organized as follows: ยง2 presents the data, which include not only COMBO-17 (W03) and DEEP2 (Paper I) but also data from the smaller DEEP1 pilot survey, which are presented here for the first time. ยง3 briefly summarizes methods used to measure the luminosity functions and their evolution, referring the reader to Paper I and W03 for more details. Readers wanting results quickly can skip directly to ยง4, which presents the luminosity functions, computes values of $`M_B^{}`$, number density, and $`j_B`$ from fitting to Schechter functions, and compares the answers to local and distant values from the literature. These are our core results on evolution. A detailed discussion of possible sources of error is presented in ยง5, which can also be skipped on first reading. The implications for galaxy formation, especially of red galaxies, are discussed in ยง6, which draws heavily on the properties of local E/S0s as well as distant ones. We ultimately favor a stepwise, โmixedโ scenario in which (massive) E/S0s are quenched during a final gas-rich merger event, migrate to the red sequence, and then undergo a small number of further, purely-stellar mergers on the red sequence to attain their final masses. This scenario is similar to the one outlined by B04 and seems to have the right mixture of ingredients to explain the boxy-disky โstructure sequenceโ of local ellipticals plus the narrow E/S0 scaling relations and the age-$`Z`$ anti-correlation that underlies them. Speculation on the nature of the quenching mechanism and its possible downsizing over time ends the discussion. A final summary is presented in ยง7.
Throughout this work, a ($`H_0,\mathrm{\Omega }_M,\mathrm{\Omega }_\mathrm{\Lambda })`$ = (70, 0.3, 0.7) cosmology is used. Unless indicated otherwise, all magnitudes and colors are on the Vega system. Necessary conversions to AB magnitudes are given for reference in Table 1 of Paper I. Luminosity functions are specified per unit co-moving volumes.
## 2. Data
The main data analyzed in this paper come from the DEEP2 and COMBO-17 surveys, with supporting data from DEEP1. Detailed background information on these surveys appears in other references, but core information needed to understand the samples and their selection effects is provided below. Basic properties of the surveys (area, number of galaxies, magnitude and redshift limits) are summarized in Table 1.
### 2.1. DEEP2
The DEEP2 survey strategy, data acquisition, and data reduction pipeline are described in Davis et al. (2003), Faber et al. (in prep.), and Newman et al. (in prep.). A detailed summary was provided in Paper I as background to computing the DEEP2 luminosity functions; the functions from Paper I are adopted here without change. DEEP2 catalogues are derived from Canada-France-Hawaii Telescope images taken with the 12K $`\times `$ 8K mosaic camera (Cuillandre et al. 2001) in $`B`$, $`R`$ and $`I`$ in four different regions of the sky. Reduction of the photometric data, object detection, photometric calibration, and construction of the star-galaxy catalogs are described in Coil et al. (2004). $`R`$-band images used to define the galaxy sample have a limiting magnitude for image detection at $`R_{AB}25.5`$. Apparent-magnitude cuts of $`R_{AB}18.5`$ and $`R_{AB}24.1`$ and a surface brightness cut of $`\mathrm{\Sigma }_R26.5`$ were applied (see Paper I). Separation between stars and galaxies is based on magnitude, size, and color, which were used to assign each object a probability of being a galaxy; star-galaxy separation efficiency is discussed in Coil et al. (2004) and Paper I. In Fields 2, 3, and 4, the spectroscopic sample is pre-selected using $`B`$, $`R`$, and $`I`$ to have estimated redshifts greater than 0.7, which approximately doubles the efficiency of the survey for galaxies near $`z1`$. The fourth field, the Extended Groth Strip (EGS), does not have this pre-selection applied but instead has roughly equal numbers of galaxies below and above $`z=0.7`$, which were selected using a well understood algorithm versus redshift. Redshifts were measured spectroscopically using the DEIMOS spectrograph (Faber et al. 2003) on the Keck 2 telescope. Slitlets are placed (nearly) randomly on 50% of all galaxies after pre-selection (70% without pre-selection in EGS), of which 70% yield successful redshifts (80% in EGS), with a catastrophic failure rate of 1%.
The DEEP2 sample used here combines data from the first season of observations in Fields 2, 3, and 4 with 1/3 of the total EGS data, which provides an initial sample at low redshifts. The total number of galaxies is 11,284, with 4,946 (44%) in EGS, 3948 (35%) in Field 4, 2299 (20%) in Field 3, and 91 (1%) in Field 2. Since the photometric redshift cut at $`z`$ 0.7 provides a soft boundary for the selection of galaxies, only EGS is used to probe the lower-redshift realm $`z<0.8`$, while data in all four fields are used for $`z0.8`$. Color-magnitude diagrams illustrating the sample binned by redshift are shown in Paper I.
### 2.2. COMBO-17
The COMBO-17 survey consists of multi-color imaging data in 17 optical filters covering a total of 1 $`\mathrm{}^{}`$ of sky at high galactic latitudes. The filter set contains five broad-band filters ($`UBVRI`$) plus 12 medium-band filters stretching from 400 to 930 nm. All observations were obtained with the Wide Field Imager (WFI, Baade et al., 1998, 1999) at the MPG/ESO 2.2-m telescope on La Silla, Chile. The total exposure time is $``$160 ksec per field, which includes a $``$20 ksec exposure in the $`R`$-band with seeing below 0$`\stackrel{}{\mathrm{.}}`$8 FWHM. The WFI provides a field of view of $`34\mathrm{}\times 33\mathrm{}`$ on a CCD mosaic consisting of eight 2K $`\times `$ 4K CCDs with $`67`$ million pixels providing a scale of $`0\stackrel{}{\mathrm{.}}238`$/pixel. The observations began during the commissioning phase of the WFI in January 1999 and are continuing as the area is extended to cover more fields. The data used here are from three fields covering an area of 0.78 $`\mathrm{}^{}`$, providing a catalogue of $``$200,000 objects found by SExtractor (Bertin & Arnouts 1996) on $`R`$-band images with a 5-$`\sigma `$ point-source limit of $`R26`$.
SEDs created from these 17 passbands were used to classify all objects into stars, galaxies, and QSOs by comparison with template SEDs. Less than 1% of the sources have spectra that are peculiar, not yielding an object class (W03), and star-galaxy separation is highly efficient. A first analysis of the COMBO-17 luminosity function was published by W03, but the galaxy catalogue has changed slightly since then. Basic details of the classification algorithm and choice of templates were given in Wolf, Meisenheimer & Rรถser (2001). In 2003, an improved set of SED templates was introduced after it was found that the accuracy of galaxy redshifts was limited by template mismatch for bright galaxies, for which more subtle SED details could be seen. The new set of galaxy templates contains a grid of synthetic spectra based on the PEGASE code for population synthesis models (Fioc & Rocca-Volmerange 1997), whereas in the past the redshift determination relied only on the observed templates by Kinney et al. (1996). The new redshifts are accurate to within $`\delta z/(1+z)<0.01`$ at $`R<21`$ and to within 0.05 down to $`R<24`$ (Wolf et al. 2004, hereafter W04). The changes from the old redshifts are relatively small (see Figure 4 in W04) and are within the errors of the old estimates, but residual errors were reduced by up to a factor of 3 for galaxies with selected SED shapes. Typical catastrophic failure rates for COMBO-17 galaxies are $`1`$% in the magnitude range used for the LFs, as measured using galaxies in common between W04 and Le Fรจvre et al. (2004) in Chandra Deep Field South. The resulting luminosity functions of galaxies are also unchanged within the errors published in W03 if the same color divisions are used. The luminosity functions of red-sequence galaxies published by B04 were already based on the new redshift catalogue.
In this paper, we recalculate luminosity functions from the COMBO-17 galaxy sample using different color divisions than before. The red-sequence cut is similar to the one in B04, with a small difference: B04 measured the mean color of the red-sequence, which is affected to a small degree by K-correction errors that vary in a non-stochastic way with redshift. They then fitted a smooth evolution to the measured colors to identify the most likely trend. Here we use the measured red-sequence colors in each individual redshift bin to define the color valley, and thus the cut. As stated by B04, the difference this makes to luminosity functions is small. However, we consider it preferable to follow the small but systematic variations in the data for the purpose of splitting the population. The new method is identical to that used for DEEP2 (see Paper I).
In COMBO-17, a galaxy redshift measurement is considered successful when the error expected from the probability distribution is below a threshold of $`\sigma _z/(1+z)0.1`$. The completeness of successful redshifts depends on galaxy rest-frame color, and simulated completeness maps are shown in W03 and W04. The large redshift incompleteness among blue galaxies at the faint end of the COMBO-17 sample would lead to unreliable completeness corrections in the last few LF bins just above the survey limit. These bins have thus been dropped, and all correction factors used in the remaining bins (plotted data points) are below 1.5. For red-sequence galaxies at $`z<1.2`$, which are the special focus of this paper, redshifts are measured successfully for the entire sample at $`R<24`$ used for the luminosity function. This claim can be tested by identifying likely red-sequence galaxies from an apparent color-magnitude diagram like that in Figure 1$`a`$. In COMBO-17, galaxies with failed redshifts occupy solely the region of blue galaxies. This is consistent with (a) the known lower COMBO-17 completeness for redshifts of blue galaxies close to the faint limit, and (b) the zero completeness (by design) of COMBO-17 with respect to $`z>1.4`$ galaxies. Blue galaxies with $`z>`$ 1.4 are expected to be a much larger part of a flux-limited sample with $`R<24`$ than red-sequence galaxies at $`z>1.4`$. Owing to their red colors and faint near-UV fluxes, only red-sequence galaxies of extremely high luminosity could pass the flux limit at high redshift, of which there are evidently very few.
### 2.3. DEEP1
DEEP1 was a pilot survey for DEEP2 that was conducted using the LRIS spectrograph on the Keck telescopes in 1995-1999. Since the only published luminosity function using DEEP1 data treated E/S0s (Im et al. 2002) above $`I814`$ = 22, a more detailed description of DEEP1 data will be presented here. This summary also serves to convey the flavor of our treatments of DEEP2 and COMBO-17. Readers not interested in these details should skip to ยง3.
Background on DEEP1 sample selection and photometry is presented in Vogt et al. (2005), on photometry and bulge-disk decompositions by Simard et al. (2002), and on spectroscopy and redshifts by Weiner et al. (2005). The DEEP1 sample is drawn from objects detected in a set of 28 contiguous Wide Field and Planetary Camera (WFPC2) Hubble Space Telescope pointings at approximately 14<sup>h</sup> 16<sup>m</sup> 30<sup>s</sup> +52<sup>o</sup>1550<sup>โฒโฒ</sup> (J2000.0), (PIs E. Groth, GTO 5090; and J. Westphal GTO 5109). The solid angle covered by the GSS is 127 square arc minutes.
The DEEP1 photometric catalogue was created by Groth et al. (1994) and contains several parameters measured using FOCAS (Tyson & Jarvis 1981). To this catalogue were added total magnitudes and colors measured fitting 2-dimensional bulge+disk models using the GIM2D package (Simard et al. 1999, 2002). For galaxies without GIM2D measurements, magnitudes and colors measured by Ratnatunga et al. (1999, MDS) in the Medium Deep Survey were used, which were shown by Simard et al. (2002) to have comparable quality to GIM2D measurements. When neither GIM2D nor MDS had magnitudes and colors, the FOCAS measurements were transformed into the same system as the GIM2D magnitudes by adding the median offset between GIM2D and FOCAS magnitudes in each color ($`I814`$ = $`I_{FOCAS}`$ โ 0.295 and $`V606I814`$ = $`(VI)_{FOCAS}`$ \+ 0.106) (all HST magnitudes are on the Vega system.)
DEEP1 spectroscopic data were obtained over several observing runs using the Low Resolution Imaging Spectrograph (LRIS, Oke et al. 1995) on the Keck 1 and 2 telescopes. Most galaxies were observed using two different gratings, with a blue side ranging from 4500 ร
to 6500 ร
and a red side from 6000 ร
to 9500 ร
. Total exposure times ranged from about 50 minutes for galaxies observed on one mask up to $``$500 minutes for galaxies placed on several masks. In the analysis below, the DEEP1 sample is restricted to galaxies with 16.5 $``$ $`I814`$ $``$ 23.5 and with redshift quality A or B, as explained in Weiner et al. (2005) (no surface-brightness cut was applied). Galaxies generally have more than one identified spectral feature, and the redshift confidence level is better than 90%. The total number of galaxies in the region that satisfy the apparent magnitude limit is 2,438, of which 621 have good quality redshifts. The typical sampling rate of DEEP1 redshifts is $``$40%, and the typical redshift success rate is $``$70%, for a final overall sampling density of $``$25%.
The apparent color-magnitude diagram of DEEP1 galaxies is shown in Figure 1. The spectroscopic sample was selected from a โpseudo-$`R`$โ band magnitude: \[($`V606+I814`$)/2 $`24.0`$\]; this has been converted to an approximate $`I814`$ magnitude shown as the vertical dotted line in all four panels. Figure 1$`a`$ shows the full sample of galaxies, and Figure 1$`b`$ shows the distribution of galaxies placed in slits. Because of the $`R`$-band selection, shown as the inclined dashed line, there is a dearth of red galaxies fainter than $`I814=23.5`$; brighter than this, the sample of galaxies placed on masks is a good representation of the total galaxy sample. Figure 1$`c`$ presents the distribution of galaxies that were successfully measured; green symbols show galaxies below the adopted high-$`z`$ cutoff (see below), while red symbols show galaxies beyond the cutoff, which were not used. Black crosses show galaxies with redshifts from CFRS (Lilly et al. 1995b) and from Brinchmann et al. (1998), for which DEEP1 measured no spectrum. Figure 1$`d`$ shows the distribution of galaxies with failed redshifts. The few bright cases of failed redshifts resulted from short integrations or spectra of galaxies at mask edges, while the majority of failures are of faint and generally blue galaxies. As in DEEP2 and COMBO-17, most failures are likely beyond the adopted high-redshift cutoff of the survey, here taken to be $`z_h=1.0`$. This is the cutoff for the DEEP1 analysis and is the redshift where O II $`\lambda `$3727 becomes heavily confused with strong OH sky lines in the LRIS data (Weiner et al. 2005).
Figure 2 divides the apparent color-magnitude diagram into magnitude and color bins. For each bin, the histogram of the distribution of galaxies as a function of redshift is shown, where the filled histogram represents successful measures and the open bar at the right the number of failures inside each bin. The number of failed redshifts increases at magnitudes fainter than $`I81422.5`$.
Figure 3 shows the distribution of restframe $`UB`$ vs. redshift. The method used to measure $`UB`$ is described in Weiner et al. (2005), using the procedure described in Paper I but limited to the two observed filters ($`V606`$ and $`I814`$). A bimodal color distribution is clearly seen, as well as large-scale structure fluctuations due to galaxy clustering (vertical stripes). The number of successful redshifts above $`z=1`$ falls drastically owing to OH confusion (see above).
Rest-frame color-magnitude diagrams for different redshift intervals are shown in Figure 4. Similar diagrams for DEEP2 were shown in Paper I and for COMBO-17 in B04. The solid line in each panel represents the limiting absolute magnitude that corresponds to apparent magnitude $`I814=23.5`$ at the far edge of the bin as a function of restframe color. The color dependence was calculated using the K-correction code from Paper I. The changing slope of the line as a function of redshift is caused by the fact that the $`I814`$-band filter used to select the sample coincides with rest $`B`$ at $`z`$ 0.8 but differs from it increasingly as the redshift is either greater or smaller than 0.8. Intrinsically red galaxies are included to fainter absolute magnitudes when observed $`I814`$ is redder than rest $`B`$, while intrinsically blue galaxies are favored when observed $`I814`$ is bluer than rest $`B`$, thus causing the line to swing with redshift.
The upper dashed line in each figure represents the cut used to separate red from blue galaxies. This cut is identical to that used by Paper I for DEEP2 since the restframe colors and magnitudes are on the same system. The equation for the line is
$$UB=0.032(M_B+21.52)+0.4540.25,$$
(1)
which is taken from the van Dokkum et al. (2000) color-magnitude relation for red galaxies in distant clusters, converted to the cosmological model used in this paper, and corrected downward by 0.25 mag in order to pass through the valley between red and blue galaxies (Paper I). Although the colors of red galaxies may evolve with redshift, this effect is not strongly seen in either DEEP1 or DEEP2 colors, and a line with constant zero point independent of redshift is used for all redshift bins.
We conclude this section by comparing the strengths and weaknesses of the two major data sets used in this paper, DEEP2 and COMBO-17. Both data sets go to nearly the same apparent magnitude, $`R24`$, and have comparable numbers of galaxies beyond $`z=0.8`$ (see Table 1). The square root of cosmic variances are shown for each sample in Tables 2-4 by galaxy color and by redshift bin. They are comparable for the two surveys beyond $`z=0.8`$ and range between 10-20% for all redshifts and color classes. When combined, the two surveys have a total (square root) cosmic variance of $``$ 7-15% per redshift bin at $`z1`$. The strengths of DEEP2 are rock-solid redshifts and high completeness for blue galaxies all the way to $`z=1.4`$ owing to the sensitivity to \[O II\] $`\lambda `$3727, which is strong in distant blue galaxies. The strengths of COMBO-17 are higher completeness overall at all redshifts, particularly for distant red galaxies near $`z1`$. This is offset by a tendency to lose redshifts for blue galaxies towards the faint limits of the survey, which has forced us to cut off the COMBO-17 All and Blue luminosity functions at a shallower point than DEEP2 to keep completeness corrections small. The two data sets thus complement each other well at high $`z`$, making a parallel, head-to-head analysis extremely useful.
## 3. Methods
The luminosity function is most frequently expressed using the Schechter (1976) parameterization, which in magnitudes is:
$$\varphi (M)dM=0.4ln\mathrm{\hspace{0.33em}10}\varphi ^{}10^{0.4(M^{}M)(\alpha +1)}\{10^{0.4(M^{}M)}\}dM,$$
(2)
where $`\varphi ^{}`$ is a normalizing constant that is proportional to the total number density of galaxies, and $`\alpha `$ is the slope of the power law that describes the behavior of the faint end of this relation. Changes in these parameters with time quantify how galaxy populations evolve.
The methods used for the DEEP2 (and DEEP1) luminosity functions are described in Paper I. The methods used for COMBO-17 are described in W03 and are very similar. A brief overview of all methods is provided here.
Two statistical estimators have traditionally been used in the calculation of the luminosity function. These are the parametric maximum-likelihood method of Sandage, Tammann & Yahil (1979, STY; also Efstathiou, Ellis & Peterson 1988; and Marzke, Huchra & Geller 1994) and the non-parametric $`1/V_{max}`$ method of Schmidt (1968; also Felten 1976 and Eales 1993). The STY method fits an analytic Schechter function, yielding values of the shape parameters $`L^{}`$ and faint-end slope $`\alpha `$, but not the density normalization $`\varphi ^{}`$, which is estimated using the minimum-variance density estimator of Davis & Huchra (1982). The STY method also does not produce any visual check of the fit. In this paper, a visual check both on shape and normalization for each redshift bin is obtained using $`1/V_{max}`$ since it yields the average number density of galaxies in bins of redshift and absolute magnitude. Formulae used for obtaining the STY parameters, $`1/V_{max}`$ points, and the density normalizations are given in Paper I.
Since the STY method does not yield $`\varphi ^{}`$, it is not suitable for calculating the correlated errors between $`\varphi ^{}`$ and $`M^{}`$. For DEEP, these errors were calculated from the 1-$`\sigma `$ error ellipsoid (Press et al. 1992) that results from fitting a Schechter function to the $`1/V_{max}`$ data points (Paper I). Although the luminosity functions obtained from the STY and $`1/V_{max}`$ methods are not quite identical (see Figure 6 below), the differences are small and errors from the $`1/V_{max}`$ method should also apply to the STY method. For COMBO-17, the errors in parameters were calculated first for $`M_B^{}`$ using STY, and then the $`\varphi ^{}`$ errors were calculated using the field-to-field variations (W03).
Weights are needed for every data set to correct for missing galaxies. The adopted weights need to take into account the fact that 1) objects may be missing from the photometric catalogue, 2) stars may be identified as galaxies and vice versa, 3) not all objects in the photometric catalogue are targeted for redshifts, and 4) not all targets yield successful redshifts. For the DEEP surveys, factors 1) and 2) are small or zero (see discussion in Paper I), and only factors 3) and 4) need to be taken into account. The basic assumption to deal with 3) is that all unobserved galaxies share the same average properties as the observed ones in a given color-magnitude bin. Factor 4) is dealt with by assigning a model redshift distribution to the failed galaxies. We use two such models, as explained in Paper I. The โminimalโ model assumes that all failed galaxies lie entirely beyond the high-redshift cutoff of the survey, which is $`z_h=1.4`$ for DEEP2 and $`z_h=1.0`$ for DEEP1. As discussed in Paper I, this model should provide an adequate description for blue galaxies. The second model to treat failed redshifts is the โaverageโ model, which assumes that failed redshifts have the same distribution as the successful redshifts in the same color-color-magnitude bin; this is a reasonable assumption for red galaxies. The difference in weights between the minimal and average models is usually $``$25% (average weights are higher), and most of the large differences occur for galaxies with extreme colors at faint magnitudes. Luminosity functions calculated with the average model are slightly higher than those using the minimal model, by an amount that averages 10-20%. When the combined All-galaxy sample is considered, we use an โoptimalโ model in which red galaxies are modeled using the average model and blue galaxies are modeled using the minimal model.
Two small alterations to this general scheme were applied. In the case of EGS in DEEP2, a final correction (described in Paper I) was applied to account for the different sampling strategy used in this field, which includes low-redshift galaxies but de-weights them so that they do not dominate the sample. In the case of DEEP1, the weights were modified to use additional size information from HST images, which show that all galaxies with angular half-light radius $`r_{hl}1`$ arc sec (from GIM2D) lie within the legal redshift range $`z1.0`$. Figure 5 illustrates these results for DEEP1 by plotting sampling rates, redshift-success rates, and weights; analogous figures are given for DEEP2 in Paper I.
For the COMBO-17 survey, effect 1) from the list above is small, as only galaxies very close to very bright stars are lost from the object catalogue. Effect (3) is zero, as the photo-z code works on the entire catalogue. Effects (2) and (4) are linked, since in COMBO-17 both object classification and redshift estimation are one single process. In one direction, a few K stars are misidentified as galaxies, but their number is negligible. In the other direction, the misclassification of galaxies as stars is modeled together with redshift incompleteness using simulations as described in Wolf et al. (2001) and W03, which take into account the photometry S/N, SED, and redshift and are calibrated using Monte-Carlo simulations. Weights are calculated as a function of apparent magnitude and color and are close to unity for all red galaxies, for which we calculate the luminosity to the full sample depth, but drop rapidly for blue galaxies towards the survey limit. We have not used any data points that involve corrections by more than a factor of 1.5, and as a result the COMBO-17 luminosity function points for the Blue and All samples do not quite reach luminosities corresponding to apparent magnitude $`R=24`$.
## 4. Analysis
### 4.1. The DEEP and COMBO-17 Luminosity Functions
This section compares the luminosity functions derived from DEEP2, DEEP1, and COMBO-17 with one another and with published data. The DEEP2 functions are best estimates from Paper I that use the optimal missing-redshift model for All galaxies, minimal for Blue galaxies and the average model for Red galaxies (see ยง3). For the DEEP1 sample we use the minimal model for all galaxy colors, while COMBO-17 weights are as described in ยง3 and W03. Galaxies are analyzed all together (the โAllโ sample) and divided into โRedโ and โBlueโ sub-samples using color-magnitude bimodality. The method used to divide blue and red galaxies in DEEP2 and DEEP1 is based on the slanting line that goes through the color valley in the $`UB`$ vs. $`M_B`$ CM diagram (see Equation 1 and Figure 4). The line used for COMBO-17 is similar to the one used by B04 based on $`UV`$ vs. $`M_V`$ except that the smoothly-evolving zero point of the line through the color valley is replaced by a zero point adjusted in each redshift bin to make the line go through the valley at that redshift.
Figure 6 shows the resulting luminosity functions for the All data (top row), Blue data (middle row), and Red data (bottom row). Redshift increases from left to right across a row. Non-parametric 1/$`V_{max}`$ data points are shown for DEEP2 by the solid black squares, for DEEP1 by the grey triangles, and for COMBO-17 by the red circles. For all samples, the calculation of the luminosity function is truncated at the faint end using dashed lines analogous to those in Figure 4, taking the limiting absolute magnitude at each color and in each redshift bin into account; details are given in Paper I. Blue galaxies were further trimmed in COMBO-17 as described in ยง3, to allow for greater redshift incompletness.
The error bars on each DEEP2 and DEEP1 point represent Poisson statistics only. Cosmic variance estimates are shown as the separate error bar at the top left corner of each panel and were estimated using the procedure of Newman & Davis (2002) to account for evolution of the correlation function. The bias factors derived by Coil et al. (2004) for red galaxies ($`b=1.32`$) and blue galaxies ($`b=0.93`$) relative to dark-matter halos are included in these estimates. The values plotted are for DEEP2. To first order, Poisson variance is random from point to point, whereas cosmic variance should mainly move all points in a given bin up and down together. Since these effects are different, they are shown separately. For COMBO-17, the error bars combine the Poisson errors in $`\varphi (M)`$ with the cosmic variance estimated from the field-to-field variations.
Also shown in the top row of Figure 6 are 1/$`V_{max}`$ data points by Ilbert et al. (2004, VVDS), represented by blue diamonds. This sample uses $``$11,000 spectroscopic redshifts from the VVDS survey to $`I_{AB}=24`$ (7,800 redshifts are termed โsecureโ). Finally, the grey dashed lines show Schechter fits to local red and blue SDSS samples at $`z0.05`$ from Bell et al. (2003), who divided galaxies both by color and by concentration, getting similar results. The exact Schechter parameters used are given in Table 5.
The conclusions from Figure 6 are as follows:
All galaxies (top row): Measurements of the All-galaxy luminosity function from all four surveys agree well out to $`z1`$ and down to the apparent magnitude limit of DEEP2 and COMBO-17 ($`R24`$). Below this, VVDS claim to see a steepening in faint-end slope from $`\alpha 1.2`$ at $`z=0.05`$ to $`\alpha 1.5`$ at $`z=1`$. Neither DEEP2 nor COMBO-17 go deep enough to test this, but, as noted above, Gabasch et al. (2004, FDF) go nearly 2.5 magnitudes fainter and do not see it, getting $`\alpha =1.25`$ at all redshifts. Relative to the local Schechter total function, the data in successive redshift bins march to dimmer magnitudes ($`M_B^{}`$) with time but stay roughly constant in number density ($`\varphi ^{}`$). This visual assessment is confirmed by Schechter fits below. In short, for the population as a whole (All sample), galaxies are getting dimmer with time but their number density has remained much the same, since $`z1`$.
Blue galaxies (middle row): The results found above for the All sample are replicated for the Blue sample, as expected since blue galaxies account for the majority of objects at all redshifts. Results here are available only from DEEP1, DEEP2, and COMBO-17 since VVDS do not divide their samples by color. However, these three data sets agree well. Relative to the local blue Schechter function (dashed grey line), $`M_B^{}`$ dims with time while $`\varphi ^{}`$ remains constant, again confirmed by Schechter fits below.
Red galaxies (bottom row): Before considering red galaxies, we review the conclusions of Bell et al. (2004b, B04), which offered the first analysis of evolution in $`\varphi ^{}`$ and $`M_B^{}`$ for red galaxies, based on COMBO-17. The main finding was that $`M_B^{}`$ for red galaxies dims over time by $``$1.5 mag from $`z=1`$ to 0 and that $`\varphi ^{}`$ rises by at least a factor of two. This evidence for evolution from the luminosity function was further bolstered by consideration of the total $`B`$-band luminosity density of red galaxies, $`j_B`$, which is measured with smaller (formal) errors than $`M^{}`$ or $`\varphi ^{}`$ separately. The quantity $`j_B`$ was found to hold nearly constant since $`z=1`$. Since models of stellar evolution for red galaxies predict a rise in $`B`$-band stellar mass-to-light ratio by 1-2 mag since $`z=1`$ (see more on this below), constant $`j_B`$ implies that the total stellar mass contained in red galaxies has at least doubled since $`z=1`$, providing further evidence for significant growth and change in red galaxies over this epoch.
Bell et al.โs finding of recent strong evolution among red galaxies disagrees with the classic scenario for red-galaxy formation in which E/S0 galaxies assembled their mass and formed stars very early and have been passively fading ever since (e.g., Eggen, Lynden-Bell & Sandage 1962; Larson 1975). The monolithic-collapse picture predicts constant $`\varphi ^{}`$ accompanied by equal dimming in both $`M_B^{}`$ and $`j_B`$, but neither of these trends was seen by B04. Checking these conclusions by remeasuring these quantities with both DEEP2 and COMBO-17 was therefore a major goal of the present study.
The bottom row of Figure 6 presents the new data for red galaxies. As before, DEEP2 and COMBO-17 agree well. The most striking impression is the relative lack of evolution in the red luminosity function, especially when compared to the large shift to brighter magnitudes seen in the blue function. What evolution there is is quantified below by Schechter fits, which are shown in Figure 6 as the black lines. These fits indicate a formal dimming of $`M_B^{}`$ over time, accompanied by a rise in number density, $`\varphi ^{}`$. The sense of these shifts is such that the data translate nearly parallel to themselves, leaving the raw counts at a fixed absolute magnitude relatively constant. Since the actual counts are not changing a great deal, to first order, the fitted values of $`M_B^{}`$ and $`\varphi ^{}`$ must depend on slight curvature signals in the data, which could be weak and unreliable. We return to this question below, where the relative constancy of the red counts is considered from various points of view. For now we simply note that both the raw data and the fitted Schechter function parameters from DEEP2 and COMBO-17 agree extremely well, and that the formal values of $`\varphi ^{}`$ from both data sets agree with the rise found by B04.
Another important result in Figure 6 is the marked turnover in the slope of the Red luminosity function at the faint end. This turnover is well established in both DEEP2 and COMBO-17 at intermediate redshifts and is seen by Cross et al. (2004) and by Giallongo et al. (2005) at even higher redshifts. However, DEEP2 and COMBO-17 may disagree with one another in the lowest redshift bin ($`z=0.2`$-0.4), where the number of faint red galaxies continues to turn over according to DEEP2 but flattens according to COMBO-17. This is noteworthy as only the potential discrepancy between DEEP2 and COMBO-17, but the error bars on COBMO-17 are large, reflecting large field-to-field variations. Other data sets have also yielded conflicting values for the nearby red faint-end slope. For example, by identifying early-type galaxies in SDSS using both concentration and color, Bell et al. (2003) found only a modest turnover at the faint end, as in COMBO-17, whereas Madgwick et al. (2002) identified red galaxies spectroscopically in 2dF and found a strong turnover, more like DEEP2 (see Table 2). The question of the red faint-end slope and its possible evolution with redshift is very important for understanding the processes that created red-sequence galaxies. We return to this question in ยง5 below when discussing errors in the Schechter function parameters caused by possible evolution in the red faint-end slope versus redshift.
In passing, we note the grey triangles in Figure 6, which show luminosity functions from DEEP1. These agree rather well with the functions from DEEP2 and COMBO-17 except in bin $`z=0.81.0`$, where total DEEP1/Red is a factor of 1.5 too high. Two โwallsโ due to large-scale structure appear in that redshift bin, one at $`z`$ 0.81 and a larger one at $`z`$ 0.98 (Le Fรจvre et al. 1994; Koo et al. 1996; see also Figure 7). However, the observed fluctuation is not much larger than the expected cosmic variance limits ($``$30%).
### 4.2. Schechter Fits
This section presents the results of fitting Schechter functions to DEEP2 and COMBO-17 using the STY method. Aside from the possible low-redshift flattening of the Red function in COMBO-17, we see no variations in faint-end slopes that are statistically significant in different redshift bins, motivating the use of constant values of $`\alpha `$ obtained from averaging over several bins. (In fact, small changes are expected in the shape of the All galaxy function with redshift because the shapes of the Red and Blue functions differ and their relative numbers are changing with redshift; however, this effect is small.) We decided to average the faint-end slope values found within the range $`z=0.2`$ to 0.6 for COMBO-17 (because of its larger number of galaxies in this redshift range), which yielded $`\alpha =0.5`$ for the Red sample and $`\alpha =1.3`$ for the All and Blue samples (these values were also used for DEEP2 in Paper I). The latter slope agrees well with the value $`\alpha =1.25`$ found for all galaxies by FDF, while the former is close to the average value $`0.59`$ found for distant red galaxies by Giallongo et al. (2005).
Schechter function parameters for both DEEP2 and COMBO-17 are presented in Table 2 for the All sample and in Tables 3 and 4 for the Blue and Red samples. Column (1) shows the central redshift of the bin; column (2) the number of galaxies used in the luminosity function calculation in each redshift bin; column (3) the value of the adopted faint-end slope $`\alpha `$; column (4) the value of $`M_B^{}`$, followed by the upper and lower 68% errors in columns (5) and (6); the mean density $`\varphi ^{}`$ is given in column (7), followed by the 68% errors in columns (8) and (9); the square root of the cosmic variance error is shown in column (10); and the luminosity density (see ยง4.3), followed by the 68% errors in columns (11) and (12). Column (13) indicates the the weighting scheme described in ยง3.3 adopted for the DEEP2 fits. As explained above, we adopted minimal weighting for the DEEP2 Blue sample and average weighting for the Red sample because we think that failed redshifts in the two color classes have different redshift distributions. The All sample combines each of these populations with its preferred weighting scheme (called โoptimalโ in Table 2).
For DEEP2, the 68% errors are Poisson estimates for $`M_B^{}`$ and $`\varphi ^{}`$ and are taken from the $`\mathrm{\Delta }\chi ^2`$ = 1 contour levels in the ($`M_B^{}`$, $`\varphi ^{}`$) plane, computed from the $`1/V_{max}`$ residuals and their errors. Errors for $`j_B`$ are conservatively calculated by adding the fractional Poisson errors for $`M_B^{}`$, $`\varphi ^{}`$, and cosmic variance in quadrature; these latter are overestimates because they neglect correlated errors in $`M_B^{}`$ and $`\varphi ^{}`$, which tend to conserve $`j_B`$. However, since the biggest error term is usually cosmic variance, the overestimate is small. For COMBO-17, the 68% errors in $`M_B^{}`$, $`\varphi ^{}`$, and $`j_B`$ are rms estimates from field-to-field variations, which are particularly large for the redshift bin centered at $`z=`$ 0.9, caused by a big downward fluctuation in CDF-South. Tbe tabulated cosmic variance errors estimates for both samples were computed as described above for DEEP2, taking the volume and field geometries into account and using separate bias ($`b`$) values for All, Blue, and Red galaxies.
The resulting Schechter fits for DEEP2 are shown as the solid black lines in Figure 6. All fits use only the magnitude ranges of the data actually shown. The close match between the fitted curves and all data suggests that the Schechter formula, and in particular the assumed $`\alpha `$ values, are a good match to the luminosity function shapes over the magnitudes ranges where the data exist. The match of the Schechter form to red galaxies was explored quantitatively in Paper I and is reviewed again under errors in ยง5.
Evolutionary trends in fitted Schechter function parameters are shown in Figure 7. Besides DEEP2 and COMBO-17, this figure adds data from other recent surveys (2dF \[Norberg et al. 2002, Madgwick et al. 2002\]; SDSS \[Blanton et al. 2003, Bell et al. 2003\]; VVDS \[Ilbert et al. 2004\]; FDF \[Gabasch et al. 2004\]; DEEP1 \[Im et al. 2002\]). For reference, the parameters from these surveys are tabulated in Table 5. Since the various surveys use different values for $`\alpha `$, changing them to the same values used by DEEP2 and COMBO-17 would cause small shifts in $`M_B^{}`$ and $`\varphi ^{}`$. For example, if local All and Blues values were corrected to match DEEP2 and COMBO-17, $`M_B^{}`$ would brighten by $``$ 0.2 mag, and $`\varphi ^{}`$ would decline by $``$0.1 dex; these would act to reduce the gaps visible in Figure 7 between the local and distant values. For red galaxies, the changes are opposite: $`M_B^{}`$ would dim by $``$ 0.15 mag while $`\varphi ^{}`$ would increase by $``$0.08 dex, acting to increase the gaps. All thse corrections are small but add somewhat to the uncertainties.
Figure 7 contains the principal results of this paper. The first conclusion (from the top row) is that $`M_B^{}`$ has dimmed for all galaxies, and by roughly the same amount for All, Blue, and Red samples. COMBO-17 (red circles) agrees well with DEEP2 (black squares) in all three color bins, and VVDS and FDF agree well with them for All galaxies (the latter do not subdivide by color). The level of agreement is impressive because the samples were selected and measured in very different ways: COMBO-17 and DEEP2 are $`R`$-band selected to $`R=24`$, VVDS is $`I`$-band selected to $`I_{AB}=24`$, and FDF is $`I`$-band selected to $`I_{AB}=26.8`$. VVDS and DEEP2 use spectroscopic redshifts, COMBO-17 uses high-precision photometric redshifts based on 17 filters, and FDF uses photo-zโs derived from photometry in 9 bands including $`J`$ and $`K`$. Despite these differences, values of $`M_B^{}`$ for all four distant surveys typically agree to within $`\pm `$0.1 mag. Agreement for the two local surveys as analyzed by Bell et al. (2003, SDSS) and Norberg et al. (2002, 2dF) is also good (though Blanton et al. (2003) find SDSS $`M_B^{}`$ dimmer by 0.4 mag). In short, a consistent picture for the evolution of $`M_B^{}`$ for all galaxies since $`z=1`$ is emerging.
The dashed grey lines in the top row are an attempt to fit straight lines to $`M_B^{}`$ versus redshift using all the data. It is not clear that this is advisable since the All data in particular seem to show a leveling out in $`M_B^{}`$ at intermediate redshifts. If this is ignored, the coefficients (in Table 6) show that the total dimming in $`M_B^{}`$ to $`z=1`$ is 1.30$`\pm `$0.20 mag for the Red sample, 1.31$`\pm `$0.14 mag for the Blue sample, and 1.37$`\pm `$0.31 mag for the All sample. (The last value is not simply a weighted mean of the first two because the functions for red and blue galaxies have different shapes.) Thus, the evolution of $`M_B^{}`$ for both red and blue galaxies appears to have been very similar since $`z=1`$.
Based on DEEP2 alone, we wondered in Paper I whether $`M_B^{}`$ for red galaxies in fact evolved very much, and indeed the slope derived from DEEP2 (black squares in Figure 7) is rather shallow. However, adding the points from COMBO-17 has steepened the slope for the high-redshift data, and this is bolstered by the addition of the local values from SDSS and 2dF. We return to this topic in ยง5 when discussing uncertainties in the red fits.
The bottom row of Figure 7 shows evolution in $`\varphi ^{}`$ for the three color classes. Agreement is again very good among the data sets, but now red and blue galaxies evolve quite differently. The number density of blue galaxies remains nearly flat to $`z=1`$, whereas the number density of red galaxies appears smaller back in time. This rise, already noted in connection with Figure 6, repeats very closely the pattern found by B04, whose data showed a gradual rise in $`\varphi ^{}`$ since $`z0.8`$ by a factor of $``$2, precededed by a steeper rise before that near $`z=1`$. The new data from DEEP2, which are completely independent, also show a steep rise near $`z=1`$ followed by a shallower rise after that.
Formal values can be calculated for the decline in $`\varphi ^{}`$ at $`z=0.8`$ and at $`z=1`$ using the new DEEP2 and COMBO-17 data together with the updated values of local $`\varphi ^{}`$ shown in Figure 7. The mean value of $`\varphi ^{}`$ at $`z=1`$ is found to be 0.95$`\times 10^3\pm `$14%, where the value comes from interpolating the DEEP2 and COMBO-17 data at $`z=0.9`$ and $`z=1.1`$ in Table 4 and the error reflects an assumed uncertainty of 20% in DEEP2 and COMBO-17 separately. The local value of $`\varphi ^{}`$ is taken to be 3.44$`\times 10^3\pm `$20%, where the error is a conservative estimate for the mean of the two measured values in Figure 7. The formal value for the rise in red $`\varphi ^{}`$ from $`z=1`$ to now is therefore 3.6$`\pm `$24% (0.56$`\pm `$0.09 dex), and the rise since $`z=0.8`$ is 2.3$`\pm `$24% (0.36$`\pm `$0.09 dex). These are formal values based on the fitted values for $`\varphi ^{}`$; potential errors and uncertainties are discussed in ยง5.
The fall in the number of red galaxies back in time measured here does not agree with the earlier result from DEEP1 by Im et al. (2002) in which $`\varphi ^{}`$ for red galaxies was claimed to have held constant since $`z1`$. The two redshift bins from Im et al. are plotted as crosses in Figure 7, where they lie both low and remain constant back in time. Im et al. applied a very stringent cut to define their sample, targeting only morphologically normal, spheroid-dominated E/S0s having red colors that are consistent with passively fading stellar populations. Their numbers therefore have to be corrected upwards in any event by $`30`$% to account for non-E/S0 contamination on the distant red sequence (Bell et al. 2004a, Weiner et al. 2005). However, the real difference between Im et al. and DEEP2 is nearly a factor of two, based on counts by DEEP2 over the identical region. The reason for this bigger discrepancy has not yet been unravelled and signals that the cut used by Im et al. to define their sample was even more stringent than thought. Finally, Im et al. compared their distant values to an earlier estimate for the local number density of spheroidal galaxies that is considerably lower than the values used here. When all of these factors are combined (and coupled with a new and somewhat larger cosmic variance estimate), it is easy to see why Im et al. reached they conclusion they did. However, this case highlights the problems introduced by selecting distant red galaxy samples in different ways, to which we return later below.
Before leaving Schechter fits, we report a further test that divided the DEEP2 and COMBO-17 blue samples into two equal halves to see whether โModerately Blueโ galaxies evolve differently from โVery Blueโ galaxies. This repeats a test reported in Paper I for DEEP2 but now adds COMBO-17. The method of division used sloping lines that ran parallel to the red-sequence color-magnitude relation and bisected the blue sample in each redshift bin into equal color halves. The lines used for DEEP2 are illustrated in Figure 4 of Paper I; those for COMBO-17 were similar. Dynamically adjusted zero points for each redshift bin were used in preference to a constant color zeropoint because the latter would yield a spurious evolution if blue galaxies were reddening with time, as suggested by the motion of the median dividing line by $`0.1`$ mag for DEEP2 in Paper I. With this approach, $`M_B^{}`$ for Moderately Blue galaxies in DEEP2 was found to average 0.7 mag brighter than for Very Blue galaxies, as expected from the sloping color-magnitude relation for blue galaxies. Apart from that, the two blue sub-samples evolve similarly, with $`\varphi ^{}`$ holding constant for both halves separately and values of $`M_B^{}`$ retaining a constant offset versus redshift. COMBO-17 confirms these conclusions. This sameness of evolution is perhaps surprisingโwe might have expected Moderately Blue galaxies to evolve in a way that is intermediate between Very Blue and Red galaxies. Work in progress shows, for example, that the clustering of Moderately Blue galaxies is indeed intermediate between the outermost color classes (A. Coil et al., in prep.). To the contrary, the data suggest that blue galaxies are evolving as a bloc in the CM diagram (apart from a possible dilation or expansion in their total color range, which cannot be tested in the present data).
We end this section by comparing to other published luminosity functions divided by color classes. The discovery of color bimodality is rather recent, and the study by Giallongo et al. (2005) is one of only two that divide distant galaxies by restframe color as we do. Unfortunately, a quantitative comparison cannot be given because no results were presented for epoch $`z1`$ specifically, and our data do not go farther than that. Plots in Giallongo et al. agree with ours in showing similar dimming for both red and blue galaxies, a constant number of blue galaxies, and a rise in the number of red galaxies since $`z=1`$. Though the Giallongo et al. sample is much smaller than ours, it goes roughly two magnitudes fainter and is therefore valuable for establishing the existence of a turnover in the faint red luminosity function at $`z1`$.
A similar turnover is seen in the second study, by Cross et al. (2004), who counted red-selected galaxies and galaxies morphologically selected to be spheroids regardless of color, based on ACS images. Their counts agree well with ours despite their small sample size of 72 galaxies. The main difference with us is an even steeper turnover in faint-end slope near $`z1`$ in their red-selected sample, for which they find $`\alpha =+0.3`$; this is reduced to $`\alpha 0.5`$ when blue spheroids are included.
Several other studies have attempted to count galaxies in various ways to see whether spheroids are disappearing back in time. Reviews can be found in Schade et al. (1999) and Im et al. (2001). Motivated by predictions of semi-analytic models, Kauffmann et al. (1996) reanalyzed CFRS data and claimed a drop in spheroid density, but their conclusions were disputed (Totani & Yoshii 1998). Schade et al. (1999) and Menanteau et al. (1999) counted morphologically normal $`R^{1/4}`$-law objects in HST images out to $`z1`$ and concluded that there was indeed no drop. Sample and field sizes were small in both cases and no color cuts were applied, with the result that half or more of all distant objects were blue. Thus, there is no contradiction with the present study though the prevalence (again) of blue spheroids raises interesting questions. We return to the topic of blue spheroids in the Discussion section.
Our final reference is to CFRS, the pioneering study that first attempted to calculate the luminosity functions of distant red and blue galaxies separately (Lilly et al. 1995b). For blue galaxies, CFRS claimed a steepening in total faint-end slope back in time to $`z=1`$. As noted, we have refrained from drawing any strong conclusions about faint-end slope evolution from our data, despite the fact that DEEP2 and COMBO-17 have many more galaxies and go 1.5 magnitudes deeper than CFRS. In retrospect, the CFRS data do not look strong enough to support that claim. For red galaxies, CFRS found no evolution in either $`M_B^{}`$ or $`\varphi ^{}`$, whereas we find a dimming of $`M_B^{}`$ by $`\stackrel{>}{}`$1 mag and a rise in $`\varphi ^{}`$ by a factor of $``$4 (since $`z=1`$). Part of the difference may be that, lacking knowledge of color bimodality, CFRS used a non-evolving color cut that did not quite hit the valley at high redshift. Regardless, CFRS projected a picture in which red galaxies have been rather static since $`z=1`$, while blue galaxies have significantly changed. The picture derived in this paper is that blue galaxies are rather constant in number (though fading) over this time interval, while red galaxies are actively being generated. The general impression in CFRS of active blue galaxies versus passive red galaxies is thus essentially opposite to what we find.
### 4.3. Luminosity Density
Luminosity density provides an estimate of the total amount of light emitted by galaxies per unit volume. The luminosity density (in Johnson $`B`$ band) in this work is obtained assuming the Schechter form of the luminosity function:
$$j_B=L\varphi (L)๐L=L^{}\varphi ^{}\mathrm{\Gamma }(\alpha +2),$$
(3)
where $`j_B`$ is calculated in solar units using $`M_B_{\mathrm{}}`$=5.48 (Binney & Merrifield 1998) and $`\mathrm{\Gamma }`$ is the Gamma function. Use of this expression entails extrapolation over faint magnitudes that are not observed, the more so at high redshifts. However, fitting a given bright-end data set assuming different values of $`\alpha `$ tends to leave the product $`L^{}\varphi ^{}`$ unchanged, which means that most of the uncertainty comes from $`\mathrm{\Gamma }`$. For example, changing $`\alpha `$ from $`1.3`$ to $`1.7`$, as suggested by VVDS for their All sample at $`z=1.1`$, changes $`\mathrm{\Gamma }`$ by 230%. This case is extreme, however. Values of $`\alpha `$ for Blue and All galaxies from nearly all other studies range between $`1.0`$ and $`1.3`$, which implies a total change in $`\mathrm{\Gamma }`$ of only 30%. Plausible red $`\alpha `$โs range in value from $`0.5`$ to $`1.0`$, which changes $`j_B`$ by only 11%. We conclude that, as long as $`\alpha `$โs remain below $`1.3`$, uncertainties on $`j_B`$ are small.
The resultant $`B`$-band luminosity densities are plotted versus redshift in Figure 8. To the previously shown local points we have added a second value for SDSS measured by Blanton et al. (2003). This latter value has been multiplied by 62% and 38% to obtain the fraction of $`B`$-band light in blue and red galaxies separately, based on fractional light contributions from Hogg et al. (2002).
Local values agree remarkably well for all three color classes. Relative to them, DEEP2, COMBO-17, and FDF show at most a mild decline in $`j_B`$ for all galaxies with time. The fall in VVDS is nearly twice as large owing to their steeper faint-end slope at high redshift; however, as noted, FDF goes 10 times fainter and does not see such steepening. If constant $`\alpha `$ is adopted, as in DEEP2, COMBO-17, and FDF, the data indicate that $`j_B`$ for all galaxies has fallen by about a factor of $``$2 since $`z=1`$. The CFRS survey (Lilly et al. 1996) is also plotted in Figure 8 and shows a much steeper decline in $`j_B`$ after $`z=1`$, like VVDS. Some of this may come from their steeper faint-end slope at high redshift, but part also comes from their adopted low local value (see figure), which they took from Loveday et al. (1992). The new value measured by 2dF and SDSS is about 30% higher.
Turning to the individual color classes, we see that the luminosity density of blue galaxies (in Figure 8$`b`$) as measured by DEEP2 and COMBO-17 evolves somewhat more than total luminosity density, falling by a factor of $``$3 since $`z`$ = 1. This is consistent with the constant value of $`\varphi ^{}`$ and the change in $`M_B^{}`$ of 1.3 mag seen above for blue galaxies. Red galaxies are shown in Figure 8$`c`$, repeating a similar figure from B04. The conclusions are the sameโthe $`B`$-band luminosity density of red galaxies has remained essentially flat since $`z=0.9`$ and was possibly rising before that. The flat section is caused by the dimming of $`M_B`$ coupled with the rise in $`\varphi ^{}`$ so that $`j_B`$ remains constant. Before $`z=0.9`$, the steep decline in $`\varphi ^{}`$ wins out, and total $`j_B`$ seems to be lower.
To summarize, DEEP2 agrees with both old and new analyses of COMBO-17 in showing that the $`B`$-band luminosity density for red galaxies has remained nearly constant since $`z=0.9`$, with a possible fall beyond that. Despite the fact that only the upper part of the function is observed at $`z1`$, DEEP2 and COMBO-17 agree within 20%, and extrapolation errors must be small because the red function is known to turn over, even at high redshift (Cross et al. 2004, Giallongo et al. 2005). Barring actual loss of galaxies from the samples (see below), the constancy of $`j_B`$ for red galaxies after $`z1`$ should therefore be well established. Since stellar mass-to-light ratios are increasing with time (see below), this constancy implies that the stellar mass bound up in red galaxies has increased markedly since $`z=1`$, as stressed by B04. Given the strong implications of this result for galaxy formation, it is advisable to go back and review the errors and assumptions, which we do in the next section. Readers not interested in these details should skip directly to ยง6.
## 5. Errors, Assumptions, and Uncertainties
This section focuses on red-sequence galaxies, although many of the conclusions apply equally well to the other color classes. A major issue is whether the apparent fall in the number density of bright red galaxies back in time is due to galaxies that are being left out in different stages of the analysis. A second issue is the extent to which the conclusions are sensitive to fitting the counts with Schechter functions having constant, non-evolving $`\alpha `$. These effects and others are discussed below.
1) Completeness of the photometric catalogues: The COMBO-17 photometric catalogue has a 5-$`\sigma `$ detection limit down to $`R_{AB}26`$, nearly two magnitudes below what is needed for the luminosity function surveys ($`R_{AB}24`$). The DEEP2 catalog is shallower but also adequate (5-$`\sigma `$ limit $`R=24.5`$, Coil et al. 2004). DEEP2 makes an additional cut in surface brightness when designing the DEIMOS masks that deletes low-surface-brightness galaxies in the last half-magnitude bin (see Paper I). However, this is largely taken into account by calculating weights as a function of color as well as magnitude. Furthermore, this cut would not affect early-type galaxies, which have high surface brightness.
Errors in star-galaxy separation may result in either too few or too many galaxies, depending on the errors. Star-galaxy separation in COMBO-17 is based on 17-color photometry and is in general highly efficient; red counts near $`z=1`$ may be $``$10% too high owing to K-star interlopers (W04), but this would tend to overestimate red galaxies. Star-galaxy separation for DEEP2 was tested in Paper I using high-resolution HST images that cover part of the DEEP2 region in the Groth Survey Strip. Misclassification of red galaxies as stars amounted to $``$10%, but these were nearly cancelled by stars misclassified as galaxies, so the net effect was nil. Finally, checks of both data sets show that almost all galaxies to $`R=24`$ have adequate photometry in all bands, and the few ($``$ 1%) DEEP2 galaxies that do not have $`B`$-band photometry (โ$`B`$-dropoutsโ) are corrected for statistically in the weights (Paper I).
2) Dividing red galaxies from blue galaxies: This is done using restframe values of $`UB`$ in DEEP2 and $`UV`$ in COMBO-17. Errors in the zero points of these systems do not matter even if they vary as a function of redshift, since the dividing line is adjusted empirically to fit the color valley in each redshift bin. Division of local samples into red and blue galaxies has been done in different ways, but results are not sensitive to the method used. Madgwick et al. (2002) separated 2dF galaxies by spectral type, whereas Bell et al. (2003) separated SDSS galaxies based on concentration and optical color. Since there is very high correlation among these properties for local galaxies, it is not surprising that the results agree well, as shown in Table 5.
3) Errors in $`M_B`$: Weiner et al. (2005) checked DEEP2 restframe values of $`M_B`$ against values derived from GIM2D photometry of Groth Strip HST images by Simard et al. (1999). GIM2D fitted model bulge+disk profiles to $`V`$ and $`I`$ images to find total magnitudes, whereas the DEEP2 Hawaii CFHT photometry approximates each object by a Gaussian profile on the $`BRI`$ ground-based images. Despite these different methods plus uncertainties in HST WFPC2 photometric zero points and charge-transfer-efficiency corrections, the zero points of both $`M_B`$ systems agreed to 0.07 mag. Furthermore, any mismatch in the magnitude systems for distant and local surveys would cause only an error in the evolution of $`M_B`$, not $`\varphi ^{}`$. The COMBO-17 luminosities have never been independently checked. However, the detailed SED information allows a precise calculation of the rest-frame B-band luminosity at all $`z<1`$ without extrapolation. The main source of error is the photo-z error, which translates into a distance error. Most objects should have luminosity errors between 10% and 20%. Local survey values of $`M_B`$ are claimed to be accurate to 0.1 to 0.02 mag (for photographic 2dF and CCD SDSS magnitudes respectively).
4) $`R`$-band selection effect: The use of the $`R`$-band for selecting DEEP2 and COMBO-17 corresponds to restframe 3300 ร
at $`z=1`$, where the SEDs of red galaxies are rather dim. It might be thought that red galaxies are being โmissedโ on that account. In practice, this is completely allowed for by calculating limiting absolute magnitudes at each redshift as a function of both redshift and color using CM diagrams like those illustrated in Figure 4. The limiting $`M_B^{}`$ magnitude to which the counts are complete at each redshift and color is well understood.
5) Redshift completeness and accuracy: Redshift completeness has been simulated for COMBO-17 using Monte Carlo methods (W01, W03, W04). From these, it appears that redshifts are highly complete for red galaxies in COMBO-17 but substantially incomplete for blue galaxies in the last magnitude bin. This is consistent with the finding that nearly all failed galaxies are faint blue galaxies. Testing the completeness model independently is difficult. However, we have predicted total galaxy number counts from the best-fit luminosity functions, including extrapolations to the faint end and to somewhat higher redshifts, and find remarkable consistency of the prediction with the observed galaxy number counts in COMBO-17. Of course, the power of this test to assess the completeness of a sub-sample in any particular redshift bin is extremely limited.
Redshift incompleteness in DEEP2 was discussed in Paper I. Using the minimal versus the average model for failed redshifts typically results in no change in $`M_B`$ and a change in $`\varphi ^{}`$ of 10-20%, a small effect compared to the total measured evolution in $`\varphi ^{}`$ out to $`z=1`$. Our preferred choice for red galaxies results in higher values of $`\varphi ^{}`$, which minimizes the observed evolution. For red galaxies, Paper I also considered a third, extreme model in which all failed red galaxies were assumed to be located in whatever redshift bin was under consideration. It is possible to do this without redshifts because red-sequence galaxies near $`z=0.71.1`$ have apparent $`RI>1.25`$ and show up as a well defined ridge in the apparent CM diagram (see Figure 1 of Paper I). This test amounts to counting all possible red galaxies and dumping all of them with unknown redshifts into a single redshift bin. Even this extreme approach hardly affects results out to $`z=0.9`$ (though it does increase both counts and $`j_B`$ at $`z=1.1`$).
Errors for DEEP2 redshifts used here are negligible ($`<10^4`$ in $`z`$); catastrophic errors are at the level of 1%. The accuracy of COMBO-17 photo-zโs has been studied using simulations, yielding an estimated rms error of 0.03, which agrees with the spectroscopic cross-check in W04. The effect of such errors on the red luminosity function was simulated by B04 and shown to be small.
6) Formal Schechter fit errors and cosmic variance: The errors in $`\varphi ^{}`$ for red galaxies in Table 4 include Poisson noise and cosmic variance. In the two most distant bins, these errors are comparable and give an rms error in number density of about 20% per survey, or 14% for the two together. The error in the local zero point of $`\varphi ^{}`$ for red galaxies is estimated to be 20%, yielding a formal rms error for the total difference between near and far samples of 24%, or 0.09 dex. This is small compared to the formal rise in $`\varphi ^{}`$ since $`z=1.0`$ of 3.6, or 0.56 dex.
This completes the list of possible observational errors and selection effects. We turn now to various theoretical and model assumptions.
7) The assumption of a constant Schechter-function shape at all redshifts: In practice, this means 1) that the Schechter formula is a good match to the bright end of the luminosity function, and 2) that the shape parameter $`\alpha `$ can be assumed to hold constant with redshift. A breakdown in either one of these assumptions will produce a mismatch between the data and the model, causing errors in both $`M_B`$ and $`\varphi ^{}`$. If the shape of the real function is constant with redshift but is not well fitted by the model, the fitted parameters will drift with $`z`$ as the data are limited to progressively brighter magnitudes at higher redshift. Any real evolution in shape may cause additional errors.
Inspection of Figure 6 suggests that our Schechter model (with the adopted value of $`\alpha =0.5`$) looks like a good fit for red galaxies. Paper I tested this quantitatively by truncating DEEP2 data in nearer bins at brighter magnitudes corresponding to the cutoffs in more distant bins. For red galaxies, a drift of $`M_B^{}`$ of $`0.1`$ mag toward fainter values was seen with more truncation, whereas the measured evolution is a brightening of $`M_B^{}`$ by 1.3 back in time. The quantity $`\varphi ^{}`$ drifted upwards by $``$0.10 dex, whereas the measured evolution is a fall of 0.36 dex to $`z=0.8`$ and 0.56 dex to $`z=1.0`$ and $`z=1.0`$. The measured evolutions are therefore if anything an underestimate owing to errors in assumed Schechter function shape. High-redshift bins cannot be tested in the same way, but visual inspection indicates that the match between data and model remains good.
However, certain recent data hint that the shape of the red luminosity function may be evolving with time, which would invalidate the assumption of strictly constant $`\alpha `$. For example, de Lucia et al. (2004) see a deficit in the number of faint red galaxies in rich clusters at $`z=0.8`$ compared to Coma, and Kodama et al. (2004) detect a similar deficit of faint red galaxies in overdense field regions at $`z1`$. As noted above, Cross et al. (2004) report a turnover in the counts of distant field red galaxies at $`z1`$ that is stronger than reported for red galaxies locally. These studies at high redshift all go 1-2 mag fainter than DEEP2 and COMBO-17 and are thus better determinants of $`\alpha `$ in distant samples. Added to this is the potential flattening of faint-end slope seen by COMBO-17 in its nearest redshift bin (see Figure 6), which resembles the flattish faint-end slope seen in local SDSS data by Bell et al. (2003) (though DEEP2 and 2DF \[Madgwick et al. 2002\] disagree; see Table 5).
A flatter faint-end slope with time might indicate that smaller red-sequence galaxies formed later than larger ones. Further support for this are findings by McIntosh et al. (2005) that smaller red-sequence galaxies evolve faster in surface brightness back in time than larger ones, by van der Wel et al. (2005) that the zeropoint of the fundamental plane evolves faster back in time for small galaxies than for brighter ones, by Bundy et al. (2005) that the crossover mass between spheroids and disk galaxies was larger in the past, by Treu et al. (2005a,b) that small spheroidal galaxies arrived later on the red sequence than large ones, and by Im et al. (2001) and Cross et al. (2004) that distant blue spheroidal galaxies are significantly smaller than red ones and would preferentially populate the bottom end of the red sequence if they were fading towards it.
In short, recent data may be pointing toward a mild flattening in the faint-end slope of the luminosity function for red galaxies, with small spheroidal galaxies forming later than large ones. However, even if this is occurring, the effect on Schechter parameters is not large. Suppose for example that $`\alpha `$ is evolving from $`0.5`$ at $`z1`$ to $`1.0`$ locally, which is the largest change conceivably allowed by the data. Values for DEEP2 and COMBO-17 use $`\alpha =0.5`$, while local surveys find the average value $`\alpha =0.65`$ (Table 5). As an experiment, we have refitted local data using $`\alpha =1.0`$, <sup>1</sup><sup>1</sup>1For example, using sample data available at http://www.mpia-hd.mpg.de/homes/bell/data/glfearlycol.out, and find that $`M_B^{}`$ brightens by 0.3 mag, $`\varphi ^{}`$ declines by 0.18 dex, and $`j_B`$ declines by 2%. These changes are all small compared to the claimed evolution. We also noted above that $`\mathrm{\Gamma }`$ itself varies by only 11% over the entire range $`\alpha =0.5`$ to $`1.0`$. The conclusion is that faint-end slope is sufficiently flat for red galaxies that total luminosity density is determined quite well by data down to $`L^{}`$, as is the case at all redshifts here.
8) Using color as a surrogate for morphological type: In focusing on red galaxies, we are implicitly assuming that restframe color is a good way of finding spheroid-dominated E/S0 types at high redshifts. The method clearly works well at low redshifts, where only 15-20% of nearby red-sequence galaxies have Hubble types later than S0, being mostly edge-on and dust-reddened (Strateva et al. 2001, Weiner et al. 2005). However, contamination by non-spheroidal galaxies is larger at higher redshifts, amounting to 30% at $`z0.75`$ (Bell et al. 2004a, Weiner et al. 2005), and may increase beyond that (Cimatti et al. 2002a, 2003; Yan & Thompson 2003; Gilbank et al. 2003; Moustakas et al. 2004). Because contamination appears to be larger back in time, our measured rise in the number of red galaxies is a lower limit to the rise of morphologically normal E/S0s. Assuming that contamination on the red sequence at $`z=1`$ amounts to 30% would increase the rise in normal E/S0s by 0.06 dex (to 0.62 dex) since that time. The true correction could be larger since contamination at $`z=1`$ may be higher than at $`z=0.75`$.
9) Uncertainties in evolving stellar mass-to-light ratios: These come into play when converting luminosity density into the more fundamental quantity stellar mass, as shown below. As noted, $`j_B`$ for bright red galaxies above $`L_B^{}`$ is nearly constant out to $`z=0.9`$, and may be lower before that. Since stellar mass-to-light ratios are increasing with time, this means that the stellar mass bound up in red-sequence galaxies must also increase. But by how much? B04 investigated this question using single-burst, passively evolving models, but these are only one option. We have investigated further possibilities such as $`\tau `$ models, โfrostingโ models with a continuing low level of star formation (e.g., Gebhardt et al. 2003), and โquenchedโ models in which star formation is shut down abruptly at some epoch (J. Harker et al., in prep.). Models are set up to match the average color of red-sequence galaxies today and the relatively small amount of color evolution that is seen since $`z=1`$ ($`\mathrm{\Delta }(UB)=0.150.25`$ mag; Bell et al. 2004b, Weiner et al. 2005, Koo et al. 2005). Recipes that satisfy these constraints all yield fading in $`M_B^{}`$ between 1 and 2 mag.
Observationally measured brightenings are consistent with these model estimates. The fundamental plane zeropoint brightens by 1-2 mag (van Dokkum et al. 2000, Gebhardt et al. 2003, van Dokkum & Ellis 2003, Treu et al. 2005a,b, van der Wel et al. 2005), the magnitude-radius zeropoint brightens by 1-1.6 mag (Trujillo & Aguerri 2004; McIntosh et al. 2005), and $`M_B^{}`$ for red galaxies brightens by 1.3 mag (this paper). These observational shifts do not necessarily represent the fading of stellar populations if galaxies are merging or otherwise evolving in radius or $`\sigma `$ and changing their structure. Nevertheless, it is striking that the amount of evolution from the various scaling laws is very close to the evolution seen in $`L_B^{}`$, and this in turn is near the middle of the range predicted by the stellar population models. Combining all results together, we adopt 1.0 mag (0.4 dex) as the minimum increase in the mass-to-light ratio of a typical massive red galaxy since $`z=1`$. Since luminosity density $`j_B`$ has remained constant, this is also the minimum increase in red stellar mass over the same period.
We collect together the following potential corrections to the above estimates of $`\mathrm{\Delta }\varphi ^{}`$ for E/S0 galaxies from $`z=1`$ to now. A positive sign means that the previously estimated rise in $`\varphi ^{}`$ would be even bigger. The factors are: a possible mismatch between the adopted Schechter function shape and the actual bright end of the luminosity function, +0.10 dex; a possible nearby flattening of $`\alpha `$ from โ0.5 to โ1.0, โ0.18 dex; and contamination by distant non-E/S0s, +0.06 dex. Each effect is small, two of the three are hypothetical, and collectively they tend to cancel. For these reasons, we do not apply any corrections to the above-measured rise in $`\varphi ^{}`$, namely, 0.36 dex since $`z=0.8`$ and 0.56 dex since $`z=1`$.
We have not thus far uncovered any โsmoking gunโ as to why our counts of red galaxies in DEEP2 or COMBO-17 should be seriously in error. Nevertheless, there is a worrisome feature of the data, and that is the fact that the two surveys, DEEP2 and COMBO-17, do not show much internal evolution in $`\varphi ^{}`$ over most of their well-measured range. This point was mentioned in Paper I in connection with DEEP2, and it is visible again in Figure 7, which plots both DEEP2 and COMBO-17. In both data sets, there is a jump in $`\varphi ^{}`$ of $`0.2`$ dex from the distant surveys to the local surveys, and another jump of $`0.3`$ dex between $`z=0.9`$ and $`z=1.1`$. In between, $`\varphi ^{}`$ tends to plateau. This stagnation is illustrated another way in Figure 9, which overplots $`1/V_{max}`$ data points from the lowest bin at $`z=0.3`$ from both surveys on top of the data points for distant bins. As previously noted, the red counts at bright magnitudes tend to translate parallel to themselves, and one might even conclude that no evolution in the luminosity function, and thus in number density, has occurred. Both DEEP2 and COMBO-17 are similar in this regard. This degeneracy could be broken by having fainter data, but our two surveys do not go deep enough to permit this.
At this point, the argument involving stellar mass-to-light-ratios assumes great importance. Imagine replotting Figure 9 versus stellar mass instead of $`M_B`$. To account for the evolution in mass-to-light ratio, based on the discussion immediately above, the counts at $`z1`$ would have to be shifted over to the right by at least 1 mag, which would produce a vertical offset with respect to the low-redshift counts by about a factor of four (0.6 dex) near $`M_B=22`$, where all curves superimpose. This is nearly identical to our previously measured fall-off of 0.56 dex based on $`\varphi ^{}`$. Thus, once mass-to-light ratio evolution is allowed for, the number of massive red galaxies at fixed stellar mass is increasing about as fast as the formal fit for $`\varphi ^{}`$. This argument is very similar to the one applied by B04 to total luminosity density, but we apply it here to individual galaxy masses at the top end of the luminosity function. The distinction is a small one, but we wish to emphasize that the present version of the argument relies entirely on data that are observed.
In summary, the conclusion that the number density of red galaxies has risen significantly is supported by three pieces of evidence. One is the fitted results for $`\varphi ^{}`$, which yield formally significant values in number density that agree very well between DEEP2 and COMBO-17. However, these may be suspect because of difficulties in comparing to local surveys and because both surveys may be prone to unknown errors at the far edge of their range. The next two arguments are therefore important and rest on the high probability that stellar mass-to-light ratios of red-sequence galaxies have evolved by at least one magnitude since $`z=1`$. If so, the constancy of total luminosity density $`j_B`$ implies a rise in total red stellar mass by the same factor. Finally, the constancy of the luminosity function itself (at the bright end) implies that a comparable growth in the number of massive red galaxies (at fixed stellar mass) must also have occurred. In what follows, we adopt the formal values of $`\varphi ^{}`$ ($`\mathrm{\Delta }\varphi ^{}=0.36\pm 0.09`$ dex since $`z=0.8`$ and $`0.56\pm 0.09`$ dex since $`z=1.0`$) as our measures of the rise in the number of red-sequence galaxies.
## 6. Discussion
### 6.1. A โMixedโ Scenario for the Formation of Spheroidal Galaxies
The most interesting conclusion to emerge so far from the study of luminosity functions since $`z=1`$ is the growth in red galaxies over recent times. We have argued that this translates to a similar, or even steeper, rise in the number of morphologically normal, spheroid-dominated E/S0s. Barring mergers among a large and undiscovered populationโwhich would have to be tiny and/or highly obscured to avoid detection in our and other surveysโthis discovery means that the immediate precursors of most massive E/S0s must be visible in existing samples at $`z=1`$ and below. The implications of this were discussed by B04. We build on their arguments by adding data on the Blue luminosity function (measured here by DEEP2 and COMBO-17) and the properties of local spheroidal galaxies, which we will argue also have strong implications for formation scenarios. Our discussion focuses on typical red galaxies at high redshift since DEEP2 and COMBO-17 sample all galaxies regardless of location. The red sequence in distant clusters has also been extensively studied, but in general we do not try to fold these data into the present picture at this time.
It is well established that residence on the red sequence requires that the star formation rate be quenched, or at least strongly reduced. Stellar populations become red enough to join the red sequence just 1-2 Gyr after star formation is stopped (J. Harker et al., in prep.), but, in order for them to stay there, the star formation rate must remain low. For example, Gebhardt et al. (2003) explored a โfrosting modelโ with an early high rate of star formation, followed by a slowly decaying $`\tau `$ component. Based on colors, they found that only 7% of total stellar mass could be formed in the $`\tau `$ component; more recent limits based on O II in distant red galaxies are even lower (N. Konidaris et al., in prep.). In short, as B04 pointed out, the large build-up seen in red stellar mass after $`z=1`$ could not have arisen from star formation within red galaxies themselves. Rather, the stellar mass at the bright end of the red sequence must have migrated there via one of three processes: 1) the quenching of blue galaxies, 2) the merging of less-luminous already-quenched red galaxies, or 3) some combination of the two. In the following discussion, we focus on the bright end of the red sequence because this is where the data are complete. Galaxies may of course also be migrating to the lower end of the red sequence as well.
It is helpful to visualize this mass migration as the movement of progenitor galaxies through the color-magnitude diagram or, more fundamentally, the color-mass diagram. Sample tracks are shown in Figure 10. Two parent regions are illustrated, a narrow red locus corresponding to the red sequence, and a broader blue clump, which we will call the โblue cloud.โ The rather constant morphology of the color-magnitude diagram since $`z=1`$ suggests that these parent regions are relatively stable in size and location. In reality, they are also moving as galaxies evolve, but this will not be too important if individual galaxies move through them more rapidly. With this assumption, we show the clumps as fixed and the galaxies as moving through them with time.
Each final galaxy today is represented by its most massive progenitor at any epoch. Stellar mass is migrating toward the upper left corner, where luminous red galaxies reside. For a galaxy to get there, two things must happen: the stellar mass composing the final galaxy must be assembled via gravitational collapse, and star formation must be quenched. A key question in the formation of red-sequence galaxies is therefore when mass assembly occurred relative to star-formation quenching: that is, did quenching occur early in the process of mass build-up, midway, or late? If extremely early, the pieces that would become the final galaxy migrated to the red sequence while still small, producing a large number of small galaxies on the lower red sequence that must later merge along the sequence in a series of โdry,โ purely stellar mergers. This is Track A. If extremely late, the progenitors grew in mass hierarchically while still making stars within the blue cloud. Upon quenching, the most massive of them moved to the head of the red sequence and took up residence there without any further dry mergers whatsoever. This late-stage quenching scenario (for various masses) is shown as the tracks labeled B in Figure 10. Mixed scenarios are also possible, involving moderate mass assembly during the star-forming stage, followed by quenching and continued but limited dry merging along the red sequence. These are the tracks labeled C.<sup>2</sup><sup>2</sup>2Strictly speaking, purely stellar mergers increase the stellar mass of a galaxy but leave its color unchanged. The arrows for an instantaneous merger should therefore be horizontal in Figure 10, which could add objectionably to the scatter on the red sequence if much merging occurred (Bower et al. 1992). On the other hand, the total merging process might last some time, in which case populations would age and redden as they grow in mass, causing the track to be tilted upward, which is how we have drawn it in Figure 10. Evidently, the total scatter induced by stellar merging depends on the precise timing and amount of merging. We return to this question in ยง6.3.
The tracks in Figure 10 assume that quenching is accompanied (and perhaps triggered) by a major merger. A related but different scenario would involve the pure fading of single blue galaxies without any merging at all. We do not consider this, for two reasons given by B04. First, distant blue galaxies are disk-dominated (Bell et al. 2004b, Weiner et al. 2005), and fading alone cannot transform disks into spheroidsโchanging the structural morphology requires a merger. Second, there do not seem to be enough blue galaxies in the distant color-magnitude diagram with masses comparable to those of massive red galaxies (B04; Weiner et al. 2005; Paper I). Hence, in order to boost mass and to create spheroids, a final episode involving both quenching and merging seems to be required, at least for galaxies on the upper red sequence. On the other hand, quenching of pure disks without merging may well feed the lower red sequence. Indeed, the local red sequence contains many low-luminosity S0s (Binggeli et al. 1988) whose disks were probably quenched by ram-pressure stripping or other gas-starvation processes not involving merging.<sup>3</sup><sup>3</sup>3The red โthick-disksโ of nearby spiral galaxies may be yet other examples of quenching without merging. A more likely explanation is that they are old-disk stars that have been dynamically heated over time by encounters with molecular clouds, disk instabilities, and/or small satellite galaxies, in which case no quenching at all is needed to explain them. Because our focus here is on massive red galaxies, which lack disks, a final episode of quenching plus merging is assumed.
Yet a third process by which galaxies might migrate to the red sequence is unveiling, whereby a dusty starburst is cleansed of its interstellar medium and the underlying galaxy is revealed. Such a process might cause the galaxy to brighten as dust absorption is removed, but also to redden as the starburst ages. However, there is no need to discuss this case separately because it is already subsumed under the above cases. If the starburst is an episode in the life of a single disk galaxy (e.g., Hammer et al. 2005), then the object today is a late-type spiral and is irrelevant to the red sequence. If the starburst has been induced by a merger, then the dusty phase is a temporary stage between the original blue precursor and the final red remnant, which does not alter our fundamental model of blue galaxies merging and turning into red galaxies. The arrows in Figure 10 are meant to connect initial and final states, not represent the detailed track whereby an object moves from blue to red. The only assumption that we have made concerning that transition is that merger remnants move quickly to the red sequence without lingering very long as bright blue starbursts. This is required by the fact that few if any bright blue starbursts are visible in the CM diagram (Bell et al. 2004b; Weiner et al. 2005; Paper I). It is also supported by radiative transfer models of dust in merging galaxies, which indicate that the burst itself is heavily cloaked by dust and is optically nearly invisible (Jonsson et al. 2005). Thus, starbursting galaxies are hard to tell optically from non-starbursting galaxies, and both types populate the blue cloud, as we assume.
The above assembly processes could be represented equally well by tracks in the color-magnitude diagram as in the color-mass diagram, and the former would be closer to existing data. However, mass is the more fundamental parameter, its behavior under merging is easier to predict than light (because dust and starbursts are not a problem), and mass estimates for samples of distant galaxies are growing and soon will be a standard tool (e.g., Drory et al. 2004, Fontana et al. 2004, Drory et al. 2005, Bundy et al. 2005). With mass as the size variable, the motions of galaxies moving onto the red sequence are described by vectors moving both upward (redder) and to the left (more massive). The slopes of the vectors in Figure 10 correspond to equal-mass mergers (i.e., mass doubling); unequal mergers would have more vertical vectors.
Yet another perturbation to the model is the possibility that the most-massive progenitor might take up residence on the red sequence and then later merge with a smaller gas-rich galaxy. The resultant starburst could briefly move the remnant back to the blue cloud, followed by subsequent decay back onto the red sequence (e.g.,. Charlot & Silk 1994). However, such events (while they last) would create massive blue galaxies, which we have argued are rare. The events must therefore be short-lived and should not greatly distort our basic assumption that, once the most-massive progenitor galaxy enters the red sequence, the galaxy remains there permanently.
We make three generic points before considering Tracks A, B, and C further. First, since the number of massive spheroidal galaxies (and their associated stellar mass) has been growing over time, the makeup of the population is not stable, and mean properties such as average color, stellar age, etc., are constantly being skewed by recent arrivals (the so-called โprogenitor biasโ phenomenon of van Dokkum & Franx (2001)). The population as a whole therefore cannot be modeled using classic, single-burst, monolithic-collapse models (e.g., Eggen, Lynden-Bell, & Sandage 1962, Larson 1975), even though certain properties, such as $`L_B^{}`$ and color evolution, seem to be well fit by such modelsโthis similarity is a coincidence and these models must be abandoned. A related point pertains to what we mean by the โageโ of a galaxy. In the monolithic picture, the age of spheroidal galaxies corresponds to the epoch at which the mass collapsed and the stars were formed (both were the same). In the new picture, each spheroidal galaxy has at least three characteristic agesโthe epoch of major mass-assembly, the epoch of major star formation, and the epoch of quenchingโall of which can be different.<sup>4</sup><sup>4</sup>4Our use of the word age here is not meant to obscure the fact that mass assembly and star formation are both prolonged processes, so that any particular โageโ must be the mean of events that may have lasted billions of years.
The final point is the importance of local E/S0s in constraining formation models. Four local properties are relevant. The first is the basic fact that morphology correlates closely with average stellar ageโ disk stars are younger and blue, whereas spheroidal populations are older and red (e.g., Baade & Gaposchkin 1963). This fundamental datum motivates our basic assumption that the same process quenching star-formation also altered the morphology from disk-like to spheroidal. This process is believed to involve mergers with other similar-sized galaxies (Toomre & Toomre 1972; Toomre 1977; Mihos & Hernquist 1994, 1996; Barnes & Hernquist 1996). Many local merger remnants are known whose properties are consistent with their evolving into spheroidal galaxies once the acute merger phase is over (e.g., Schweizer 1982, Schweizer 1986, Hibbard 1995). The incidence of spheroid-dominated galaxies is also higher in groups and clusters of galaxies (e.g., Dressler 1980, Postman & Geller 1984, Hogg et al. 2003, Balogh et al. 2004, Baldry et al. 2004), where mergers were more frequent.
Second, local E/S0 galaxies populate a rather tight โfundamental planeโ linking radius, luminosity, and velocity dispersion (Faber et al. 1987, Dressler et al. 1987, Djorgovski & Davis 1987). The tightness of this plane implies that stellar mass-to-light ratio cannot scatter by more than $`\pm `$15% at any point on the plane. Two other relations, the color-magnitude relation (Faber 1973, Sandage & Visvanathan 1978, Bower, Lucey & Ellis 1992) and the Mg-$`\sigma `$ relation (e.g., Bender, Burstein & Faber 1992, Bernardi et al. 1998, Colless et al. 1999, Worthey and Collobert 2003), are also quite narrow and further link the properties of stellar populations to those of their parent galaxies.
Third, the mean light-weighted stellar-population ages of local Es scatter widely, from over 10 Gyr down to just a few Gyr (e.g., Gonzalez 1993, Trager et al. 2000a, Jรธrgensen 1999, Terlevich & Forbes 2002), with most being younger than classic single-burst models would predict (11.4 Gyr if $`z_{form}=3`$). This large number of young ages allows room for the late quenching that is required by the luminosity function data. On the other hand, there is at most a weak trend in stellar age with mass or $`\sigma `$ along the red sequence (Trager et al. 2001, Terlevich & Forbes 2002, Bernardi et al. 2005), so that stellar ages and metallicities scatter substantially at every point on all three relations. To keep the relations tight, Worthey et al. (1995) posited that an anti-correlation must exist between age and metallicity at constant mass and/or $`\sigma `$, as later verified by Jรธrgensen (1999), Trager et al. (2000b), and Bernardi et al. (2005).
The fourth and final point is that nearby Es populate a structure sequence, in which small objects rotate strongly, are flattened by rotation, and have disky isophotes and steep central surface-brightness profiles, whereas massive objects rotate weakly, are flattened by anisotropic velocity dispersions, and have boxy isophotes and core-type central profiles (e.g., Davies et al. 1983, Bender, Burstein & Faber 1992, Faber et al. 1997). At the low-mass end, these properties connect smoothly with S0s and, through them, to the remainder of the Hubble sequence (Kormendy & Bender 1996). Several authors have suggested that this structure sequence can be explained broadly by assuming that small spheroidals were produced via mergers of gas-rich, โwetโ progenitors, while more massive spheroidals were produced by progressively โdry,โ purely stellar mergers of smaller spheroidals (e.g., Bender, Burstein & Faber 1992, Kormendy & Bender 1996, Faber et al. 1997). This scenario implies that the most massive of todayโs spheroidals were created mainly by purely stellar mergers.
With these points as background, we return to the tracks in Figure 10. The early-quenching scenario (Track A) has most of its mass-assembly occurring in dry mergers along the red sequence. This can be ruled out on two grounds. First, to produce the large amount of stellar mass bound up in massive red-sequence galaxies would require a huge reservoir of small, faint galaxies on the lower red sequence. This excess is not detected at any redshiftโthe local red-sequence luminosity function is at most flat ($`\alpha 1.0`$) and if anything turns over more steeply at higher redshifts (Cross et al., 2004, Kodama et al. 2004, Giallongo et al. 2005, Figure 6 here). The required reservoir of small red galaxies therefore does not exist. Second, building up massive red galaxies from purely dry mergers along the red sequence would yield stellar populations whose metallicities are uncorrelated with stellar age and whose ages and metallicities would converge to a single value at high masses after many mergers had occurred. This fails to match the fundamental plane, color-magnitude, and Mg-$`\sigma `$ relations, which indicate that age and $`Z`$ must be anti-correlated at each location, nor does it match the strong spread in age and $`Z`$ that exists even among massive galaxies (Trager et al. 2001, Terlevich & Forbes 2002). (The upward tilt that we have placed on the arrow representing dry mergers on Track A reflects the probable increase in mean stellar age during dry merging, not an increase in mean metallicity.)
The late-quenching scenario (Track B) is extreme in the opposite sense of having no dry merging at all along the red sequence. In this picture, massive present-day spheroidals were formed via a single merger of two very massive gas-rich progenitors. The main reason for ruling out this scenario is the structure sequence among local Es, which, as noted, implies that massive Es were formed by dry, gas-poor mergers. The signatures of such mergers are distinctive because the precursors are dynamically hot, yielding fuzzy tidal tails without sharp boundaries; many examples of such dry mergers can be seen in local catalogues (e.g., Arp 1966), so it is clear that they are occurring.
The mixed scenario (Track C) involves early mass assembly and star formation, followed by quenching and further (but limited) dry merging. This scenario seems most naturally to explain the properties of local E/S0s. For example, the final mergers making small spheroidals would be mostly gas-rich, while later mergers along the red sequence would be progressively more gas-poor, as required by the structure sequence. Furthermore, the break point between boxy and disky galaxies is an upper limit on the masses of blue galaxies that have migrated onto the red sequence recently. That break point today is in the range $`M_B`$ = โ20 to โ21, where boxy and disky galaxies coexist (Faber et al. 1997, Lauer et al. in progress). With mean spheroidal $`/L_B`$ 6 (from Gebhardt et al. 2003, adjusted to the $`B`$-band and $`H_0=70`$), this translates to stellar masses in the range 1-2$`\times 10^{11}`$ M$`_{_{\mathrm{}}}`$, or blue progenitor masses of 0.5-1$`\times 10^{11}`$ M$`_{_{\mathrm{}}}`$ for equal-mass mergers. These are at the upper end of blue masses today (Bell et al. 2003), as expected.
With one more quite natural assumption, the mixed scenario might even be able to explain the narrowness of the local Fundamental Plane, color-magnitude relation, and Mg-$`\sigma `$ relations. Consider a selection of galaxies at a fixed mass today on the red sequence. In the mixed scenario, these galaxies will have arrived there via different routesโsome will have been produced by recent gas-rich mergers of two blue galaxies, while others will have quenched earlier and evolved along the red sequence via dry mergers for a longer time. In general, however, we expect there to be a broad correlation between the mass of a spheroidal galaxy today and the mass of its latest blue progenitor, massive spheroidals tending to come from more massive blue galaxies, and vice versa.
This โmemoryโ of progenitor mass then helps to shape the final metallicity of a galaxy via the mass-metallicity relation among star-forming galaxies. This relation is strong among nearby galaxies (Tremonti et al. 2004) and apparently extended well into the past to beyond $`z=1`$ (Kobulnicky et al. 2003). If the amount of dry merging on the red sequence is limited, the original mass-metallicity relation of the progenitors will survive to form the backbone of the red-sequence scaling relations seen today. On the other hand, this backbone will be blurred by the different star-formation histories of galaxiesโgalaxies that quenched early from low-mass blue progenitors and grew later via dry mergers will have rather low metallicities (reflecting their small progenitors), but their average stellar age will be high (since multiple dry mergers take time). In contrast, galaxies that quenched late and arrived on the red sequence near their present mass will have higher metallicities (reflecting more massive progenitors), but their average stellar ages will be younger because they quenched recently. Thus, the multiplicity of routes that is inherent in the mixed scenario might help to account naturally for the anti-correlation between age and metallicity that is needed to explain the scaling relations. However, the more dry merging that takes place along the red sequence, the more the underlying mass-metallicity correlation of the blue progenitors will be erased, to be replaced by age scatter. The amount of scatter in age and $`Z`$ at fixed mass is therefore an indicator of the amount of dry merging that could have occurred, which might be determined through observations and modeling.
### 6.2. Quenching and Downsizing
This section briefly discusses quenching and the related concept of โdownsizing.โ To be effective, quenching requires the removal of essentially all cold gas within galaxies, and the prevention of any more falling in. The key question is what triggers quenching: why and when does it happen? Several processes are probably involved. We adopt as a starting point the standard view that mergers of gas-rich galaxies can trigger powerful starbursts (e.g., Mihos & Hernquist 1994, 1996; Sanders & Mirabel 1996). One major source of gas removal is therefore consumption of gas in the starburst itself. The burst also generates internal stellar-driven feedback that can remove and/or heat gas, such as photoionization, O-star winds, supernovae, and radiation-driven winds operating on dust (Murray, Quataert & Thompson 2005). An additional source of gas-heating is orbital energy injected during the merger, which can drive gas out in a galactic wind (Cox et al. 2005). Gas can also be removed by external processes such as ram-pressure stripping, tidal stripping, and โharassmentโ (Moore et al. 1996), which operate more effectively in the dense environments frequented by spheroidal galaxies.
Despite this abundance of potential gas-removal mechanisms, it has been suggested that these alone may not be adequate to keep spheroidal galaxies gas-freeeโthe energy requirements seem too large, and gas in models continues to fall in, creating large numbers of massive blue galaxies that are not seen, especially at the centers of clusters (e.g., Benson et al. 2003). For example, to match the sharp turndown at high masses in the galaxy mass function, Kauffmann et al. (1999) found it necessary in semi-analytic models to truncate gas cooling arbitrarily in all halos above a circular velocity of $`V_{circ}=350`$ km s<sup>-1</sup>. Similar ad hoc recipes are being tested in smaller halos in order to see if they can produce color bimodality (e.g., R. Somerville et al., in prep.). Considerable evidence, both theoretical and empirical, suggests that feedback from AGNs might be the missing trigger for quenching (e.g., Granato et al. 2004; Dekel & Birnboim 2005; Springel, Di Matteo & Hernquist 2005), which is plausible since spheroids are precisely those galaxies that possess massive black holes (Kormendy & Richstone 1995, Magorrian et al. 1998, Tremaine et al. 2002).
Properties of local spheroids may shed further light on quenching. A major clue, as mentioned, is that local galaxies scatter significantly in age at any location on the red sequence, which suggests that galaxies of the same mass have arrived there via different routes and different quenching histories. It is also observed that red galaxies co-exist with blue galaxies over more than one order of magnitude in total stellar massโa cross-over point exists near $`3\times 10^{10}`$ M$`_{_{\mathrm{}}}`$ where the numbers of red and blue galaxies are equal (Kauffmann et al. 2003a,b), but the transition is gradual (Bell et al. 2003). Both of these facts suggest that the trigger for quenching is not simply total stellar mass or any variable closely related to it (such as luminosity or rotation speed) because any of those would produce a division in mass between spheroidals and non-spheroidals that is too abrupt. Rather, we seek one (or more) variables that are broadly related to total mass but with considerable scatter, and perhaps also having some relation to the presence and size of the central black hole, given the possible need for AGN feedback. A natural variable satisfying these requirements is stellar spheroid mass, which increases generally with total galaxy mass but scatters greatly with respect to it. Spheroid mass is also closely linked to the mass of the central black hole (Kormendy & Richstone 1995, Hรคring & Rix 2004), and thus perhaps its energy output. Finally, spheroid mass increases discontinuously during a major merger, and its rise over some threshold might be the specific trigger for quenching. This scenario would meet the requirements of Balogh et al. (2004), who concluded that the blue-to-red transition must be driven primarily by internal properties rather than by environment. But environment could play a smaller role and, through its variation from galaxy to galaxy, contribute to the scatter seen in stellar ages at each point on the scaling relations.
As a final topic, we consider the matter of โdownsizing.โ The basic concept of downsizing was introduced by Cowie et al. (1996) to explain their finding that actively star-forming galaxies at low redshift are smaller in mass than actively star-forming galaxies at high redshift, which suggested that star formation is stronger at late times in smaller galaxies than large ones. The essence of this idea was already latent in the literature. For example, it was known that early-type galaxies are more luminous and more massive than later-type galaxies (de Vaucouleurs 1977, Binggeli & Sandage 1985) and that their stellar populations are on average older (e.g., Tinsley 1968, Searle, Sargent & Bagnuolo 1973). Color and gas fraction were known to vary systematically along the Hubble sequence (de Vaucouleurs 1977; Roberts 1969), indicating progressively slower, less-efficient star formation in later Hubble types. Finally, the blue end of the Hubble sequence had been shown explicitly to be a mass sequence (van den Bergh 1976, de Vaucouleurs 1977), with low-mass Irr Is at the bottom having the smallest fraction of stars and the highest proportion of gas (Roberts 1969). All evidence together thus indicated (even then) that massive galaxies made most of their stars early, whereas small galaxies formed theirs relatively later. Recent analyses of the star-forming histories of both local galaxies from SDSS (Heavens et al. 2004) and distant galaxies from the Gemini Deep Deep survey (Juneau et al. 2005) have confirmed this basic picture.
Blumenthal et al. (1984) offered a reason for this mass-dependent sequence by suggesting that early-type galaxies arose from higher-$`\sigma `$ galaxy-sized perturbations in a cold-dark-matter universe. Such perturbations would collapse first and start making stars early. Moreover, because of the non-white nature of the CDM power spectrum, high-$`\sigma `$ perturbations of galaxy mass are embedded preferentially within larger high-$`\sigma `$ perturbations (Bardeen et al. 1986), which causes them to merge more and eventually wind up in groups and clusters. The accelerated growth of high-$`\sigma `$ perturbations was demonstrated in early hydrodynamical simulations by Cen & Ostriker (1993), which showed the first galaxies collapsing at the intersections of filaments, forming stars rapidly, and assembling later into groups and clusters. They identified these early-forming objects with E/S0s.<sup>5</sup><sup>5</sup>5It has sometimes been said in the recent literature that CDM predicts that massive galaxies form late and should therefore have younger stars, which is opposite to what the Hubble sequence actually shows. This remark demonstrates confusion between the formation of galaxies and their dark-matter halos. Massive halos indeed form late, but they are making clusters of galaxies today, not galaxies. Baryonic dissipation has reduced the collisional cross-sections of galaxies to the point that galaxies are merging much more slowly now than their parent dark-matter halos. Indeed, it is this reduced merging of galaxies at recent times that enables the overall number density of galaxies, $`\varphi ^{}`$(All), to remain constant since $`z=1`$, despite the fact that halos are continuing to merge. The larger point is that events on the scale of galaxies are becoming progressively more decoupled from events on the scale of halos, including the processes of star formation and baryonic mass assembly (see also Heavens et al. 2004). But the influence of parent halos (what is often termed โenvironmentโ) has not yet declined to zero even now. The halo occupation distribution (HOD) and related statistics (see review by Cooray & Sheth 2002) are emerging as powerful tools for unravelling the relationship between galaxies and their dark-matter halos, which is crucial to understanding galaxy formation.
The above concept of downsizing refers to the mean epoch of star formation, which was clearly earlier in massive galaxies than smaller ones. However, as noted, spheroidal galaxies have a second star-formation timescale, namely, that of quenching. These two timescales might vary differently with mass, thereby generating (in principle) two different kinds of downsizing (or even upsizing). When speaking of downsizing, it is important to clarify which timescale is meant.
The remainder of this discussion focuses on the down(or up)sizing of quenching, asking whether the typical entry mass onto the red sequence has increased or decreased with time. This amounts to asking whether it is easier (or harder) to keep galaxies of a given stellar mass free of cold gas at late epochs. Three factors suggest that keeping galaxies gas-free should get easier with time. First, gaseous infall generally declines with time in the Universe, and gas within galaxies gets converted to stars; both of these mean that there is less gas overall that needs to be removed or fended off. Second, galaxies cluster more and move more rapidly within clusters, both of which promote environmentally-driven gas-removal processes such as stripping and harassment. Third, parent dark-matter halos have higher dynamical temperatures and lower densities so that any leftover gas outside galaxies is hotter and less likely to cool.
The sum of these factors suggests that it is easier to keep galaxies of a given mass gas-free at late times. When this is coupled with the observed fact (not yet fully explained) that quenched galaxies are larger, we are led to hypothesize that the typical entry mass onto the red sequence may be decreasing with time, such that progressively smaller galaxies can find their way onto the red sequence at later epochs. We have already mentioned certain observations that point in this direction, including the deficit of small red-sequence galaxies at high redshift (Kodama et al. 2004, Cross et al., 2004, Giallongo et al. 2005) followed by possible later infill (e.g., Bell et al. 2003; COMBO-17 here), the late arrival of small spheroids on the fundamental plane (Treu et al. 2005a,b), the more rapid evolution back in time of the surface-brightness of smaller spheroids (McIntosh et al. 2005), the more rapid evolution back in time of the fundamental plane zeropoint for smaller galaxies (van der Wel et al. 2005), and the possible decrease with time in the crossover mass between spheroids and disk galaxies (Bundy et al. 2005). In aggregate, the evidence may be pointing to a downsizing of quenching, on top of the downsizing of star formation in galaxies as a whole.
### 6.3. Related Topics
The finding that the majority of spheroidal stellar mass was quenched only after $`z=1`$ amounts to a paradigm shift with wide repercussions over a range of issues in galaxy formation. This section lists some important questions that are raised by the late-quenching picture.
First, does the increase of mass on the red sequence cause a problem for the observed mass budget of blue galaxies? No, because the increase in red stellar mass is small in absolute terms. Red-sequence galaxies today make up 20% of the total number of bright galaxies, and 40% of the total stellar mass (Hogg et al. 2002). If their stellar has roughly tripled since $`z=1`$, it would have increased from 0.13 units to its present level of 0.40 units. If this increase came at the expense of blue galaxies, their stellar mass would have declined from 0.87 units to 0.60 units, a fall of only 0.16 dex that if translated to $`\varphi ^{}`$ would be barely detectable in Figure 7. Thus, even this extreme scenario in which all new stellar mass in red galaxies came entirely from the pre-existing stellar mass of blue galaxies is probably consistent with the data. But it is more likely that much of the new red stellar mass was born via continuing star-formation in blues after $`z=1`$, and also in final merger-generated starbursts. For example, it has been claimed that 50% of all star formation at $`z=1`$ is occurring in intermediate-mass LIRGS, many of which are mergers and could be the precursors of spheroidal galaxies (Hammer et al. 2005 and references therein). Total stellar mass in all types of galaxies has also probably increased since $`z=1`$, by between 1.4 and 2.0 (Fontana et al. 2004; Rudnick et al. 2003; Drory et al. 2005; but see also Bundy et al. 2005). Either of these increases would be large enough to maintain blue $`\varphi ^{}`$ approximately constant if the new stellar mass were appropriately distributed over red and blue galaxies. Thus, the observed constancy of blue number density in the face of rising numbers of red galaxies does not seem to be a problem.
A second point is how galaxy mass functions are predicted to evolve if the present luminosity function data are correct. We have already noted that the measured evolution in $`M_B^{}`$ for red galaxies is similar to the predicted change in their stellar $`/L_B`$ ratio. If these evolutions are identical, the characteristic mass $`^{}`$ for red galaxies has remained constant. The same rough equality probably also holds for blue galaxies. But, since red galaxies are more massive than blue ones, their rise in number density should cause the mass function to go up faster at higher masses, and thus the mass function should change shape. This does not appear to agree with published measurements, which show a more even rise over all masses (Drory et al. 2005), or perhaps even no rise at all (Bundy et al. 2005). On the other hand, existing mass functions are measured over small areas and are subject to large cosmic variance. Larger samples coming soon from DEEP2 (K. Bundy et al., in prep.) and VVDS may resolve this discrepancy.
A third issue is reconciling the rise in red-sequence galaxies with the rate of mergers needed to create them. Estimated merger rates for bright galaxies going back to $`z1`$ vary widely in the literature (see Lin et al. 2004 and references therein). Early merger rates from DEEP2 are based on optical pair counts, not morphologies, and are rather low: only 9% of $`L^{}`$ galaxies are estimated to have suffered a major merger since $`z=1.2`$ (Lin et al. 2004); Bundy et al. (2004) obtain similar rates based on K-band pair counts. In the present scenario, merged galaxies are assumed to migrate rapidly to the red sequence. If red galaxies have tripled in number since $`z=1`$ and make up 20% of all galaxies today (Hogg et al. 2002), then 2/3 of all red onesโand thus 13% of all galaxiesโmust have merged since $`z=1`$. This fraction 13% is not far from the DEEP2 merger fraction of 9%, and also does not allow for additional quenching that may have occured without mergers, such as in stripped S0s. In short, the required rate of conversion by mergers is not excessive and may even be consistent with the rather low DEEP2 rate. Higher merger rates could also be accommodated provided that the extra remnants wind up as the bulges of spiral galaxies, as is generally assumed in semi-analytic models (Kauffmann et al. 1996, Baugh, Cole & Frenk 1996, Somerville & Primack 2001) and suggested by Hammer et al. (2005) based on the observed frequency of LIRGs.
A further puzzle that may now need rethinking is the existence of non-solar abundance ratios in early-type galaxies, which display enhanced ratios of SNae Type II elements compared to elements generated in Type Iaโs (e.g., Worthey et al. 1992). The amount of enhancement correlates closely with velocity dispersion (Kuntschner 1998, Trager et al. 2000b). It has been customary to account for these non-solar ratios by appealing to very rapid star formation, which suppresses the iron-peak elements that are produced more slowly by Type Iaโs. Such early-burst scenarios were natural within the monolithic-collapse picture. However, if most massive ellipticals were quenched at or after $`z=1`$, they were probably making stars for at least several Gyr before that, and rapid-burst models may no longer apply. Perhaps a correlation existed among the precursors of red galaxies between $`\sigma `$ and average star-formation duration that is strong enough to explain the data. Alternatively, other factors such as galactic wind strength or IMF variations (see Trager et al. 2000b) might play a role.
Finally, further work is needed to understand the form and scatter of the fundamental plane, Mg-$`\sigma `$ relation, and red-sequence (color-magnitude) scaling relations. As noted, scatter about these relations places strong limits on the amount of stellar merging, which may prove problematic. The tilt of the red sequence is also a mystery: it could come from the underlying mass-metallicity relation of the blue progenitors or from aging during a long series of stellar mergersโboth possibilities have been mentioned here. In short, the scaling relations are sensitive diagnostics. We might have attempted simple estimates of their form and scatter here, but semi-analytic models seem a much better vehicle for such calculations, and we defer these to later papers.
Besides posing difficult questions, the late-quenching picture also opens up new opportunities. If red galaxies indeed emerged recently, it becomes feasible to study in detail why certain galaxies turn red. For example, we can look at the galaxy population just beyond $`z=1`$ and try to predict which galaxies are about to be quenched. Given the close correspondence between red galaxies and dense environments, it is natural to ask whether galaxies turn red owing to an increase in the amount of clustering around them, or whether pre-existing dense environments suddenly begin to โspawnโ red galaxies near $`z=1`$. Two studies within DEEP2 are underway to answer this question (B. Gerke et al., in prep.; M. Cooper et al., in prep.), both taking advantage of DEEP2โs high redshift accuracy, which provides the needed information on parent groups and clusters.
An important issue going well beyond the confines of this paper is the impact that the late-quenching picture will have on the relation between spheroids and black holes. The masses of present-day black holes correlate closely with the properties of their parent spheroids, whether with total luminosity (Kormendy & Richstone 1995, Hรคring & Rix 2004) or with velocity dispersion (Gebhardt et al. 2000, Ferrarese & Merritt 2000). As long as spheroids were thought to form early, it was possible to imagine that the relationship is ancient, with roots going back to $`z2`$ when spheroids and black holes were simultaneously forming (e.g., Richstone et al. 1998). However, if most spheroids emerged late, the relation could hardly have existed before $`z1`$, and the massive black holes that were growing before that time must somehow have โknownโ which spheroids they would eventually wind up in. Thus, the late-emergence of spheroids adds an important new twist to the black-hole/galaxy co-evolution story.
Although we have clearly come down in favor of the late-quenching picture, the conclusion is more indirect than we would like, resting heavily as it does on models for the evolution of stellar mass-to-light ratios. The subject would be on much firmer footing if further checks could be carried out. For example, it remains to be shown whether $`K`$-band counts and redshift distributions of near-IR-selected samples are consistent with the large drop in red galaxies at $`z1`$ claimed here. These counts should be reconsidered using fainter samples with photometric redshifts, and quenched rather than monolithic-collapse models. Furthermore, if the red counts could be extended just 1.5 magnitudes fainter, $`\varphi ^{}`$ could be measured directly at $`z=1`$. For definitive results, however, such data would have to cover $``$2 โก$`^{}`$ spaced over the sky in several statistically uncorrelated regions.
A related question is whether the emergence of spheroidal galaxies near $`z=1`$ is consistent with the properties of red objects earlier than this. Some studies have searched for red objects beyond $`z=2`$ (e.g., van Dokkum, et al. 2003; Franx et al. 2003), while others have targeted red objects near $`z=1.5`$ (e.g., Cimatti 2002a). The results are very different. Objects beyond $`z=2`$ contain old stars but are still star-forming vigorously; they are red in part because they are dusty (Fรถrster Schreiber et al. 2004; van Dokkum et al. 2004; Toft 2005). Although they may be future red-sequence galaxies, their numbers do not bear on the question of whether many galaxies quenched later near $`z=1`$.
In contrast, the number of red galaxies near $`z=1.5`$ is very relevant. Despite the fact that a large fraction of these objects are also dusty (e.g., Moustakas et al. 2004), a sizeable fraction also seem to be fully quenched (Longhetti et al. 2005, Daddi et al. 2005, Saracco et al. 2005). Number densities have been variously estimated between 10% (Daddi et al. 2005) and 100% (Saracco et al. 2005) of local massive spheroidals, leading different authors to conclude that massive Es are, or are not, fully quenched by $`z=1`$. However, such surveys are as yet small and are subject to large cosmic variance. Extending this work to larger areas (and distingushing dusty galaxies from quenched ones) would provide the sharpest test of the late-quenching model.
## 7. Summary
The evolution of $`B`$-band galaxy luminosity functions since $`z1`$ is determined using a total sample of 39,000 galaxies to $`R24`$ mag from the DEEP2 and COMBO-17 surveys. DEEP2 data come from Willmer et al. (2005, Paper I), while the COMBO-17 data come originally from Wolf et al. (2003) but have been substantially reworked using using improved photo-zโs and new color classes. Evolution is examined for blue and red samples separately by dividing galaxies using color bimodality; this is the first study aside from Willmer (2005) to compare blue and red galaxies in this way. Cosmic variance is reduced to 7-15% per redshift bin by combining the results of the surveys. DEEP2 counts agree remarkably well with COMBO-17 in all color classes at nearly all redshifts.
Luminosity functions of blue and red galaxies evolve differently with redshift; the blue counts shift to brighter magnitudes at fixed number density back in time, whereas red counts are nearly constant at fixed absolute magnitude. Both DEEP2 and COMBO-17 agree in this regard. Schechter function parameters are fit to the data assuming non-changing shape (constant $`\alpha `$), and results are compared to recent measurements from other distant surveys. Good agreement is found between DEEP2 and COMBO-17 at all redshifts, and between these and other large, recent surveys counting all galaxies. Results by color are not yet available from these other surveys.
Combining the distant Schechter parameters with local ones, we solve for the fading over time of characteristic luminosity $`M_B^{}`$ for All, Red, and Blue galaxies. All classes fade by nearly the same amount, showing fadings (per unit redshift) of 1.30$`\pm 0.20`$ mag for red galaxies, 1.31$`\pm 0.14`$ mag for blue galaxies, and 1.37$`\pm 0.31`$ mag for all galaxies. In contrast to $`M_B^{}`$, $`\varphi ^{}`$ evolves differently in different color classes: formal values for $`\varphi ^{}`$ hold steady for blue galaxies but rise for red galaxies by 0.36 dex$`\pm `$0.09 dex since $`z=0.8`$, and by 0.56$`\pm `$0.09 dex since $`z=1`$. The evolution of luminosity density, $`j_B`$, also differs with color; for blue galaxies it falls by 0.4 dex after $`z1`$, while for red galaxies it remains constant since $`z=0.9`$, possibly being smaller before that.
The simplest interpretation of these results is that the number density of blue galaxies has remained nearly constant since $`z=1`$, whereas the number density of red galaxies has increased. The latter conclusion is subjected to close scrutiny, which is warranted by the fact that most of the total red evolution (in both DEEP2 and COMBO-17) occurs between the local surveys and our data, and in the farthest bin of our dataโi.e., in the two places where the data are weakest. Although it is possible that our formal values of $`\varphi ^{}`$ may have unknown errors, we nevertheless conclude that substantial evolution in the number density of red galaxies has occurred, based on strong evidence for a rise of at least one magnitude in the mass-to-light ratios ($`/L_B`$) of red stellar populations since $`z=1`$. When this rise is taken into account, the observed near-constancy of red luminosity density translates to a rise in overall number density by at least one magnitude. A similar argument applied to the red counts at fixed absolute magnitude translates to a rise in number at fixed stellar mass that is also comparable to the formal rise from $`\varphi ^{}`$. Thus, both the new DEEP2 data and the reanalysis of COMBO-17 together strongly support the rise in red galaxies since $`z1`$ first found in COMBO-17 by Bell et al. (2004b). The rise in morphologically pure E/S0s is even larger if increasing contamination by non-E/S0s at higher redshifts is allowed for.
The implications of this rise for galaxy formation are examined. Barring the existence of a major, highly obscured and as-yet-unknown population of galaxies at low redshifts, the immediate precursors of most modern-day E/S0 galaxies must be visible in existing surveys near $`z=1`$. The lateness of the rise is inconsistent with classic, high-redshift single-burst collapse models for E/S0 formation, which predict constant numbers of spheroidal galaxies over this epoch. Instead, it appears that most present-day E/S0s arose from blue galaxies with ongoing star formation that were โquenchedโ at or after $`z1`$ and then migrated to the red sequence. The properties of nearby E/S0 galaxies support a โmixedโ scenario in which quenched galaxies enter the red sequence over a wide range of masses via โwet,โ gas-rich mergers, followed by a limited number of โdry,โ stellar mergers along the sequence. The most massive E/S0s are built up during the last stages of dry merging and are visible today as boxy, core-dominated ellipticals.
Some evidence points to a decline in the average entry-mass onto the red sequence with time, which would amount to a โdownsizingโ of quenching. This and other processes might change the shape of the red luminosity function, violating our assumption throughout of constant $`\alpha `$, but the changes are not large enough to invalidate our major conclusions. Furthermore, galaxies at a given mass on the red sequence will have arrived there via different merging and star-formation histories. Plausible differences among these histories may account for the anti-correlation that is seen between age and metallicity residuals on the red sequence and that is needed to account for the narrowness of the fundamental plane, Mg-$`\sigma `$, and color-magnitude relations of local E/S0s.
Finally, growing evidence seems to suggest that extra feedback from AGN activity might be a key ingredient in initiating quenching. Consistent with this is the fact that quenched stellar populations and massive central black holes are both uniquely associated with spheroidal galaxies. We speculate that the specific trigger for quenching occurs when the spheroid stellar mass, and perhaps also the black-hole mass, exceeds a threshold value during a merger.
CNAW thanks G. Galaz, S. Rauzy, M. A. Hendry and K. DโMellow for extensive discussions on the measurent of the luminosity function; E. Bell, J. Brinchmann, A. Gabasch, and G. Galaz for providing electronic versions of their data; S. Lilly for correspondence on the CFRS luminosity function, and G. Blumenthal, J. Cohen, and L. Cowie for useful discussions. SMF thanks R. Somerville, J. Primack, and T. Lauer for extensive discussions concerning the origin of spheroidal galaxies. The DEEP team thanks C. Steidel for sharing unpublished redshift data. The authors thank the Keck Observatory staff for their constant support during the several observing runs of DEEP1 and DEEP2; the W. M. Keck Foundation and NASA for construction of the Keck telescopes; and Bev Oke and Judy Cohen for their tireless work on LRIS that enabled the spectroscopic observations of DEEP1 galaxies. We also wish to recognize and acknowledge the highly significant cultural role and reverence that the summit of Mauna Kea has always had within the indigenous Hawaiian community; it is a privilege to be given the opportunity to conduct observations from this mountain. The DEEP1 and DEEP2 surveys were founded under the auspices of the NSF Center for Particle Astrophysics. The bulk of the work was supported by National Science Foundation grants AST 95-29098 and 00-71198 to UCSC and AST 00-71048 to UCB. Additional support came from NASA grants AR-05801.01, AR-06402.01, and AR-07532.01 from the Space Telescope Science Institute, which is operated by AURA, Inc., under NASA contract NAS 5-26555. The DEIMOS spectrograph was funded by NSF grant ARI92-14621 and by generous grants from the California Association for Research in Astronomy, and from UCO/Lick Observatory. HST imaging of the Groth Strip was planned, executed, and analyzed by Ed Groth and Jason Rhodes with support from NASA grants NAS5-1661 and NAG5-6279 from the WFPC1 IDT. SMF would like to thank the California Association for Research in Astronomy for a generous research grant and the Miller Institute at UC Berkeley for the support of a Visiting Miller Professorship. CW was supported by a PPARC fellowship. NPV acknowledges support from NASA grant GO-07883.01-96A and NSF grants NSF-0349155 from the Career Awards Program and NSF-0123690 via the ADVANCE Institutional Transformation Program at NMSU. KG was supported by Hubble Fellowship grant HF-01090.01-97A awarded by the Association of Universities for Research in Astronomy, Inc., for NASA under contract NAS5-26555. JAN acknowledges support from NASA through Hubble Fellowship grant HST-HF-01165.01-A awarded by the Space Telescope Science Institute, which is operated by the Association of Universities for Research in Astronomy, Inc., for NASA, under contract NAS 5-26555. Computer hardware gifts from Sun Microsystems and Quantum, Inc. are gratefully acknowledged. This research has made use of the NASA/IPAC Extragalactic Database (NED), which is operated by the Jet Propulsion Laboratory, California Institute of Technology, under contract with the National Aeronautics and Space Administration. Finally, we acknowledge NASAโs Astrophysics Data System Bibliographic Services.
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# Confinement-induced resonances for a two-component ultracold atom gas in arbitrary quasi-one-dimensional traps
## 1 Introduction
A strongly interacting ultracold atom gas displays interesting features of a correlated quantum many-body system when its dynamics is confined to one dimension . The presence of a transverse confining potential has been shown to induce characteristic resonances in the coupling constant of the two-particle s-wave scattering process , which have become known as confinement-induced resonances (CIR). The existence of the CIR has been revealed under the simplifying assumption of a transverse parabolic confinement potential with length scale $`a_{}`$ and for the case that the two scattering atoms belong to the same species . In this case, the center-of-mass (COM) and relative coordinates of the two particles can be separated, allowing to factorize the problem into single-particle problems. At low temperatures, only the COM ground state is occupied, the decoupled COM motion can be disregarded, and the two-body problem can be solved exactly within the pseudopotential approximation. The result is that there is exactly one bound state for any 3D scattering length $`a`$. In the limit of small binding energy, the particles are tightly bound in the lowest-energy transverse state and form a very elongated dimer. The appearance of such a bound state is purely due to the confinement, since for $`a<0`$ no dimer is formed in free space. In the opposite limit of large binding energies, the dimer becomes spherically symmetric. In this regime, the confinement is not effective, and the free-space result is recovered. Moreover, a unitary equivalence exists between the Hamiltonian and its projection onto those channels which are perpendicular to the one with lowest energy. As a consequence, to each bound state corresponds a bound state of the closed channels, which then causes the CIR . It occurs at a universal value of the ratio $`a_{}/a=๐=\zeta (1/2,1)1.4603`$, where $`\zeta `$ is the Hurvitz zeta function. The influence of the CIR has also been studied for the three-body and the four-body problem in the presence of confinement. In particular, the solution of the four-body problem completely determines the corresponding quasi-1D many-body BCS-BEC crossover phenomenon . Recently, the existence of the confinement-induced molecular bound state in a quasi-1D Fermionic <sup>40</sup>K atom gas confined in an optical trap has been reported . By using rf spectroscopy, the binding energy of the dimer has been measured as a function of the scattering length, with quantitative agreement to the results of Ref. . However, the existence of the CIR in the scattering states remains to be observed.
Although the analytical results for the parabolic confinement are instructive, realistic traps for matter waves frequently have non-linear potential forms, see for instance Ref. for a particular example of a trap on the nanometer scale. To give another example, for the problem of tunneling of a macroscopic number of ultracold atoms between two stable states of a trapping potential, the nonlinearity clearly is crucial. Hence generalization to the non-parabolic case is desirable and provided in this work. In addition, we consider traps with two different species of atoms. Note that sympathetic cooling techniques require to study this case. Different trap frequencies may arise for different atom species, e.g., because of different atom masses or different magnetic quantum numbers. Here we obtain general expressions for the bound-state energies and scattering resonances when the COM and the relative degrees of freedom do not decouple anymore. In the parabolic limit and for intraspecies scattering, we recover well-known results . For the general case, we show that more than one CIR may appear, and that it depends on the symmetry properties of the confining potential how many resonances occur. We apply our formalism to two experimentally relevant cases: (i) interspecies scattering in a two-species mixture of quantum degenerate Bose and Fermi gases in an optical trap, and (ii) a single species cloud in a magnetic trap, taking into account non-parabolic corrections due to a longitudinal magnetic field suppressing Majorana spin flips.
As we will discuss below in more detail, the CIR has a close similarity to the well-known Feshbach resonance , which arises if the Hilbert space can be divided into open and closed channels coupled together by a short-range interaction. Due to this small but finite coupling, two incoming particles initially in the open channel visit the closed channels during the scattering process. If a bound state with energy close to the continuum threshold exists, such a process is highly enhanced and a resonance results.
The paper is organized as follows: In Sec. 2, we present the general formalism. Section 3 presents the bound-state solution, while Sec. 4 contains the analysis for the scattering solutions, including the analogy to Feshbach resonances. In Sec. 5 we discuss the special case of harmonic confinement, and in Sec. 6 a particular example of a non-parabolic confinement is illustrated. Finally, we conclude in Sec. 7. Technical details have been delegated to two Appendices. We set $`\mathrm{}=1`$ throughout this paper.
## 2 The two-body problem
Let us consider the general case of two different atomic species with mass $`m_1`$ and $`m_2`$. We denote the particle coordinates by $`๐ฑ_i=(๐ฑ_{,i},z_i)`$ and their momenta by $`๐ฉ_i=(๐ฉ_{,i},p_{,i})`$. Different atoms may experience a different transversal confinement potential $`V_i(๐ฑ_{,i})`$. For ultracold atoms, only low-energy s-wave scattering is relevant, and the interaction between unlike atoms (and similarly, also the interaction between the same atoms) can be described by a Fermi-Huang pseudopotential $`V(|๐ฑ_1๐ฑ_2|)`$. Then the relevant Hamiltonian for two different atoms is given by
$$H=\frac{๐ฉ_1^2}{2m_1}+\frac{๐ฉ_2^2}{2m_2}+V_1(๐ฑ_{,1})+V_2(๐ฑ_{,2})+V(|๐ฑ_1๐ฑ_2|).$$
(1)
The pseudopotential has the standard form
$$V(๐ซ)=\frac{2\pi a}{\mu }\delta (๐ซ)\frac{}{r}r,$$
(2)
where $`\mu =m_1m_2/(m_1+m_2)`$ is the reduced mass and $`a`$ the 3D scattering length. This allows to characterize the two-body interaction by the parameter $`a`$ only. The validity of the pseudopotential approach has been verified numerically for finite-range potentials in Ref. . For further convenience, we transform to the relative/COM coordinates and momenta given by $`๐ซ=(๐ซ_{},z)`$, $`๐=(๐_{},Z)`$ and $`๐ฉ=(๐ฉ_{},p_{})`$, $`๐=(๐_{},P_{})`$, respectively. This can be done by the canonical transformation
$$\left(\begin{array}{c}๐\\ ๐ซ\\ ๐\\ ๐ฉ\end{array}\right)=\frac{1}{M}\left(\begin{array}{cccc}m_1& m_2& 0& 0\\ M& M& 0& 0\\ 0& 0& M& M\\ 0& 0& m_2& m_1\end{array}\right)\left(\begin{array}{c}๐ฑ_1\\ ๐ฑ_2\\ ๐ฉ_1\\ ๐ฉ_2\end{array}\right),$$
(3)
where $`M=m_1+m_2`$. Since the confinement is assumed to be purely transversal, the longitudinal COM coordinate $`Z`$ is free and decouples from the other degrees of freedom. Hence we eliminate it by transforming into the longitudinal COM rest frame, where the state $`|\mathrm{\Psi }`$ of the system is determined by the set of coordinates $`(๐ฑ_{,1},๐ฑ_{,2},z)`$ or, alternatively, by $`(๐_{},๐ซ)=(๐_{},๐ซ_{},z)`$. The transformed Hamiltonian takes the form
$$H=H_{}+H_{,1}+H_{,2}+V,$$
(4)
where
$$H_{}=\frac{p_{}^2}{2\mu },H_{,i}=\frac{๐ฉ_{,i}^2}{2m_i}+V_i(๐ฑ_{,i}).$$
(5)
For a more compact notation, we introduce the non-interacting Hamiltonian $`H_0=HV`$ and denote its eigenstates by
$$|k,\lambda _1,\lambda _2=e^{ikz}\psi _{\lambda _1}^{(1)}(๐ฑ_{,1})\psi _{\lambda _2}^{(2)}(๐ฑ_{,2}),$$
(6)
where $`\psi _{\lambda _i}^{(i)}`$ are single-particle eigenstates of $`H_{,i}`$ for eigenvalue $`E_{\lambda _i}^{(i)}`$. Correspondingly, the two-particle Schrรถdinger equation is given by
$$\left(H_0E\right)\mathrm{\Psi }(๐_{},๐ซ)=V(๐ซ)\mathrm{\Psi }(๐_{},๐ซ).$$
(7)
The pseudopotential (2) can be enforced by the Bethe-Peierls boundary condition
$$\mathrm{\Psi }(๐_{},๐ซ0)\frac{f(๐_{})}{4\pi r}\left(1\frac{r}{a}\right),$$
(8)
leading to the inhomogeneous Schrรถdinger equation
$$\left(H_0E\right)\mathrm{\Psi }(๐_{},๐ซ)=\frac{f(๐_{})}{2\mu }\delta (๐ซ).$$
(9)
The solution of this equation can be formally obtained in terms of a solution of the homogeneous Schrรถdinger equation, $`(H_0E)\mathrm{\Psi }_0=0`$, and the Greenโs function $`G_E=(H_0E)^1`$,
$$\mathrm{\Psi }(๐_{},๐ซ)=\mathrm{\Psi }_0(๐_{},๐ซ)+๐๐_{}^{}G_E(๐_{},๐ซ;๐_{}^{},0)\frac{f(๐_{}^{})}{2\mu }.$$
(10)
To determine $`f(๐_{})`$, we substitute Eq. (10) into Eq. (8) and find the integral equation
$$\frac{f(๐_{})}{4\pi a}=\mathrm{\Psi }_0(๐_{},0)+๐๐_{}^{}\zeta _E(๐_{},๐_{}^{})f(๐_{}^{}),$$
(11)
where we have defined the regularized integral kernel
$$\zeta _E(๐_{},๐_{}^{})=\underset{r0}{lim}\frac{1}{2\mu }\left(G_E(๐_{},๐ซ;๐_{}^{},0)\delta (๐_{}๐_{}^{})\frac{\mu }{2\pi r}\right).$$
(12)
In Eqs. (10) and (11), $`\mathrm{\Psi }_0`$ can be expressed as a superposition of single-particle eigenstates $`|k,\lambda _1,\lambda _2`$ with
$$\frac{k^2}{2\mu }+E_{\lambda _1}^{(1)}+E_{\lambda _2}^{(2)}=E.$$
(13)
We refer to the set of states with the same transverse occupation numbers $`\lambda _i`$ but arbitrary longitudinal relative momentum as a scattering channel or, simply, channel. Each channel has a minimum energy given by $`E_{\lambda _1}^{(1)}+E_{\lambda _2}^{(2)}`$. Since the interaction is short-ranged, only states fulfilling Eq. (13) appear in the asymptotic solution. For each open channel, $`E>E_{\lambda _1}^{(1)}+E_{\lambda _2}^{(2)}`$, such that there are (at least) two such states having opposite momenta. For $`E`$ just above $`E_0^{(1)}+E_0^{(2)}`$, there exists one open channel only. The corresponding solution given by Eq. (10) describes the scattering of two particles initially occupying the transverse ground-state. During the scattering process, the particles populate closed channels, but afterwards return into the single available open channel (quasi-1D picture). For $`E<E_0^{(1)}+E_0^{(2)}`$, all channels are closed and only bound-state solutions are possible. These are given by Eq. (10) with $`\mathrm{\Psi }_0(๐_{},๐ซ)=0`$. In the following, we consider both classes of solutions in more detail.
## 3 Bound-state solutions
Let us consider the situation when all channels are closed and only bound states may occur. We define the binding energy of the bound states as
$$E_B=E_0E>0,$$
(14)
where $`E_0=E_0^{(1)}+E_0^{(2)}`$ is the ground-state energy of $`H_0`$. To find bound states, we diagonalize the operator $`\zeta _E(๐_{},๐_{}^{})`$ defined in Eq. (12), where Eq. (11) yields the condition
$$\frac{f(๐_{})}{4\pi a}=๐๐_{}^{}\zeta _E(๐_{},๐_{}^{})f(๐_{}^{}).$$
(15)
For given $`a`$, bound states with binding energy $`E_B=E_0E`$ follow as solution of this eigenvalue problem. The bound-state wave function follows by inserting the corresponding eigenvector $`f(๐_{})`$ into Eq. (10) with $`\mathrm{\Psi }_0(๐_{},๐ซ)=0`$. In order to find a representation of $`\zeta _E(๐_{},๐_{}^{})`$ allowing for straightforward analytical or numerical diagonalization, we use
$$G_E(๐_{},๐ซ;๐_{}^{},0)=_0^{\mathrm{}}๐te^{Et}G_t(๐_{},๐ซ;๐_{}^{},0),$$
(16)
with the imaginary-time evolution operator
$$G_t(๐_{},๐ซ;๐_{}^{},0)=๐_{},๐ซ|\mathrm{exp}[H_0t]|๐_{}^{},0$$
(17)
for $`H_0`$. The time evolution operator $`\mathrm{exp}[H_0t]`$ can be factorized into the product $`\mathrm{exp}[H_{}t]\mathrm{exp}[H_{1,}t]\mathrm{exp}[H_{2,}t]`$. The corresponding factors in $`G_t`$ are
$$z|\mathrm{exp}[H_{}t]|z^{}=\left(\frac{\mu }{2\pi t}\right)^{1/2}e^{(zz^{})^2\mu /2t}$$
(18)
for the relative longitudinal coordinates and
$$๐ฑ_{i,}|\mathrm{exp}[H_{i,}t]|๐ฑ_{i,}^{}=\underset{\lambda }{}e^{E_\lambda ^{(i)}t}\psi _\lambda ^{(i)}(๐ฑ_{i,})\overline{\psi }_\lambda ^{(i)}(๐ฑ_{i,}^{})$$
(19)
for the transverse coordinates (the bar denotes complex conjugation). Thus $`G_t`$ can be expressed in terms of the set of coordinates $`(๐ฑ_{,1},๐ฑ_{,2},z)`$ as
$$G_t(๐_{},๐ซ;๐_{}^{},0)=\sqrt{\frac{\mu }{2\pi t}}e^{z^2\mu /2t}\underset{i=1,2}{}\underset{\lambda }{}e^{E_\lambda ^{(i)}t}\psi _\lambda ^{(i)}(๐ฑ_{,i})\overline{\psi }_\lambda ^{(i)}(๐ฑ_{,i}^{}).$$
(20)
This equation illustrates that for large imaginary times the integrand in Eq. (16) decays as $`\mathrm{exp}[E_Bt]`$. Notice that this representation is valid for $`E_B>0`$. By using
$$\frac{\mu }{2\pi r}=_0^{\mathrm{}}๐t\left(\frac{\mu }{2\pi t}\right)^{3/2}e^{r^2\mu /2t}$$
(21)
we find
$$\zeta _E(๐_{},๐_{}^{})=_0^{\mathrm{}}\frac{dt}{2\mu }\left[e^{Et}G_t(๐_{},0;๐_{}^{},0)\left(\frac{\mu }{2\pi t}\right)^{3/2}\delta (๐_{}๐_{}^{})\right].$$
(22)
To show that the integral in Eq. (22) converges also for small $`t`$, we expand $`G_t(๐_{},0;๐_{}^{},0)`$ with respect to $`t`$, see Appendix A. We find
$`\underset{t0}{lim}G_t(๐_{},0;๐_{}^{},0)`$ $`=`$ $`\left({\displaystyle \frac{\mu }{2\pi t}}\right)^{3/2}\delta (๐_{}๐_{}^{})t^{1/2}\left({\displaystyle \frac{\mu }{2\pi }}\right)^{3/2}`$ (23)
$`\times \left[{\displaystyle \frac{๐_{}^2}{2M}}+V_1(๐_{})+V_2(๐_{})\right].`$
Thus $`\zeta _E`$ can be regarded as a regular operator acting on the space $`^2`$ of square-integrable functions. We note in passing that if the two single-particle transverse Hamiltonians $`H_{,i}`$ commute with the angular momentum operators $`L_z`$, then also $`\zeta _E`$ commutes with $`L_z`$. This follows by observing that in this case we can choose for the eigenbasis $`\{\psi _\lambda ^{(i)}\}`$ a set of eigenvectors of $`L_z`$, and the product of two eigenvectors of $`L_z`$ is still an eigenvector of $`L_z`$. Hence $`\zeta _E(๐_{},๐_{}^{})`$ can be written as a sum of projectors onto states with definite angular momentum. A similar conclusion can be drawn regarding parity symmetry, when considering non-cylindrical confining potentials that obey this symmetry.
Using the transverse non-interacting ground state,
$$\psi _0(๐_{},๐ซ_{})=\psi _0^{(1)}\left(๐_{}+\frac{\mu }{m_1}๐ซ_{}\right)\psi _0^{(2)}\left(๐_{}\frac{\mu }{m_2}๐ซ_{}\right),$$
(24)
provided the overlap integral $`๐๐_{}^{}\overline{\psi }_0(๐_{}^{},0)f(๐_{}^{})0`$, the integrand in Eq. (10) decays as $`\mathrm{exp}[E_Bt]`$. This defines a spatial scale $`a_B`$ for the longitudinal size of the corresponding bound state, $`a_B=1/\sqrt{\mu E_B}.`$ For large $`E_B`$, $`a_B`$ is small and we have very tight pairs. This constitutes the dimer limit. On the other hand, for small $`E_B`$, atom pairs are very elongated. This regime is termed BCS limit. In the following, we investigate both limits in greater detail.
### 3.1 Dimer limit
For large binding energies, the atom-atom interaction dominates over the confinement. Due to the exponential factors in Eq. (22), only small imaginary times contribute significantly to the integral, and we can substitute $`G_t`$ with the short-time expansion (23) as derived in the Appendix A, yielding
$`\zeta _E(๐_{},๐_{}^{})`$ $``$ $`\left({\displaystyle \frac{\mu }{2\pi }}\right)^{3/2}{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{dt}{2\mu }}(t^{3/2}(e^{Et}1)\delta (๐_{}๐_{}^{})`$ (25)
$`e^{Et}t^{1/2}๐_{}|H_{}|๐_{}^{}).`$
Hence, the operator $`\zeta _E`$ now shares eigenfunctions with $`H_{}=๐_{}^2/2M+V_1(๐_{})+V_2(๐_{})`$. For (identical) parabolic confinement potentials, $`H_{}`$ is exactly the decoupled COM Hamiltonian. Let us denote the eigenfunctions and eigenenergies of $`H_{}`$ as $`\varphi _\lambda (๐_{})`$ and $`E_\lambda ^{(\varphi )}`$, respectively. Substituting $`\varphi _\lambda (๐_{})`$ into Eq. (15) yields after some algebra
$$\frac{1}{4\pi a}=\frac{\sqrt{2\mu |E|}}{4\pi }\left(1+\frac{E_\lambda ^{(\varphi )}}{2|E|}\right)\frac{\sqrt{2\mu E_B}}{4\pi }.$$
(26)
In the second relation, we have used Eq. (14). From this, we directly obtain the binding energy in the dimer limit $`a0^+`$ as
$$E_B\frac{1}{2\mu a^2},$$
(27)
which coincides with the result obtained in free (3D) space without confinement.
### 3.2 BCS limit
The scattering channel with lowest energy, corresponding to the transverse non-interacting ground state $`\psi _0`$, opens at the energy threshold $`E=E_0`$. For $`E_B0^+`$, as the energy approaches this threshold, the term with $`\lambda _1=\lambda _2=0`$ dominates in Eq. (20), and yields in Eq. (22) the contribution
$$\sqrt{\frac{1}{8\mu E_B}}\psi _0(๐_{},0)\overline{\psi }_0(๐_{}^{},0),$$
(28)
which diverges for $`E_B0^+`$. All other channels are still closed at $`E=E_0`$ and give finite contributions in Eq. (20). This observation suggests a useful separation of the total Hilbert space into a part $`_\mathrm{o}`$ corresponding to the open channel (or lowest-energy scattering channel) and a part $`_\mathrm{e}`$ perpendicular to that. With this separation, terms yielding a finite contribution at $`E_B0^+`$ can be summarized in the Greenโs function
$$\stackrel{~}{G}_t(๐_{},๐ซ;๐_{}^{},0)=๐_{},๐ซ|\mathrm{exp}[\stackrel{~}{H}_0t]|๐_{}^{},0,$$
(29)
where $`\stackrel{~}{H}_0`$ is the projection of $`H_0`$ onto the Hilbert subspace $`_e`$. We then define a new integral kernel,
$$\stackrel{~}{\zeta }_E(๐_{},๐_{}^{})=_0^{\mathrm{}}\frac{dt}{2\mu }\left[e^{Et}\stackrel{~}{G}_t(๐_{},0;๐_{}^{},0)\left(\frac{\mu }{2\pi t}\right)^{3/2}\delta (๐_{}๐_{}^{})\right],$$
(30)
which is also well-defined for energies above the threshold $`E=E_0`$.
For small $`E_B`$, Eq. (15) is most conveniently solved by expanding $`f(๐_{})`$ in an orthonormal basis $`|j`$ according to
$$|f=\underset{j}{}f_j|j,f_j=๐๐_{}j|๐_{}f(๐_{}),$$
(31)
where the basis state $`|\mathrm{\hspace{0.17em}0}`$ corresponds to
$$๐_{}|\mathrm{\hspace{0.17em}0}=c\psi _0(๐_{},0),$$
(32)
with normalization constant $`c`$. Although $`\psi _0(๐_{},๐ซ)`$ is a normalized element of the two-particle Hilbert space, this does not imply that $`\psi _0(๐_{},0)`$ is an element of the COM Hilbert space with norm unity. In fact, the normalization constant $`c`$ has to be computed explicitly and generally depends on the particular confinement. In this basis, Eq. (15) assumes the compact form
$$\frac{|f}{4\pi a}=\zeta _E|f=\left(\sqrt{\frac{1}{8\mu E_B}}\frac{|00|}{c^2}+\stackrel{~}{\zeta }_E\right)|f.$$
(33)
$`|0`$ is an approximate eigenstate for small $`E_B`$, since all the matrix elements are finite apart from $`0|\zeta _E|0`$ which diverges according to
$$0|\zeta _E|0\sqrt{\frac{1}{8\mu E_B}}\frac{1}{c^2}+0|\stackrel{~}{\zeta }_E|0.$$
(34)
Substituting this in Eq. (33) yields
$$\frac{1}{4\pi a}\sqrt{\frac{1}{8\mu E_B}}\frac{1}{c^2}+0|\stackrel{~}{\zeta }_E|0.$$
(35)
Neglecting the last term, the relation for the binding energy $`E_B`$ is solved in the BCS limit $`a0^{}`$,
$$E_B\frac{2a^2\pi ^2}{\mu c^4}.$$
(36)
## 4 Scattering solutions
In this section, we focus on scattering solutions at low energies $`E`$ slightly above $`E_0`$, where exactly one transverse channel is open. Then the incoming state is given by
$$\mathrm{\Psi }_0=e^{ikz}\psi _0^{(1)}(๐ฑ_{1,})\psi _0^{(2)}(๐ฑ_{2,}),$$
(37)
which describes two incoming atoms with (small) relative longitudinal momentum $`k=\sqrt{2m(EE_0)}`$ in the (transverse) single-particle ground states $`\psi _0^{(1)}`$ and $`\psi _0^{(2)}`$, respectively.
### 4.1 One-dimensional scattering length $`a_{1\mathrm{D}}`$
As done in Sec. 3, we split off the contribution from the open channel,
$$G_E(๐_{},๐ซ;๐_{}^{},0)=\psi _0(๐_{},๐ซ_{})\overline{\psi }_0(๐_{}^{},0)\frac{i\mu }{k}e^{ik|z|}+_0^{\mathrm{}}๐te^{Et}\stackrel{~}{G}_t(๐_{},๐ซ;๐_{}^{},0),$$
(38)
where $`\stackrel{~}{G}_t(๐_{},z;๐_{}^{},0)`$ is the Greenโs function restricted to $`_e`$, which is well-defined also above $`E_0`$. Inserting Eq. (38) into Eq. (10) yields for $`|z|\mathrm{}`$ the standard scattering solution,
$$\mathrm{\Psi }(๐,๐ซ)=\psi _0(๐_{},๐ซ_{})(e^{ikz}+f_e(k)e^{ik|z|}),$$
(39)
with scattering amplitude
$$f_e(k)=\frac{i}{2k}๐๐_{}^{}\overline{\psi }_0(๐_{}^{},0)f(๐_{}^{}),$$
(40)
whereas for short distances, also the term $`๐๐_{}^{}_0^{\mathrm{}}๐te^{Et}\stackrel{~}{G}_t(๐_{},๐ซ;๐_{}^{},0)f(๐_{}^{})`$ appears in the scattering solution. Since the energy is well below the continuum threshold for the closed channels, this must be regarded as a sum over localized states. Enforcing the boundary condition (8) then leads to an integral equation for $`f(๐_{})`$,
$`{\displaystyle \frac{f(๐_{})}{4\pi a}}`$ $`=`$ $`{\displaystyle ๐๐_{}^{}\stackrel{~}{\zeta }_E(๐_{},๐_{}^{})f(๐_{}^{})}`$ (41)
$`+\psi _0(๐_{},0)+{\displaystyle \frac{i\psi _0(๐_{},0)}{2k}}{\displaystyle ๐๐_{}^{}\overline{\psi }_0(๐_{}^{},0)f(๐_{}^{})}.`$
This integral equation is most conveniently solved by again expanding $`f(๐_{})`$ in the orthonormal basis $`\{|j\}`$ introduced in the previous section. Thereby, we can express Eq. (41) in compact notation,
$$\frac{|f}{4\pi a}=\frac{|0}{c}+\frac{i}{2k}\frac{|0}{c^2}0|f+\stackrel{~}{\zeta }_E|f,$$
(42)
which is formally solved by
$$|f=\frac{1/c}{1i/(ka_{1\mathrm{D}})}\left(\stackrel{~}{\zeta }_E+\frac{1}{4\pi a}\right)^1|0.$$
(43)
The parameter $`a_{1\mathrm{D}}`$ follows in the form
$$a_{1\mathrm{D}}=\frac{2c^2}{0|[\stackrel{~}{\zeta }_E+1/(4\pi a)]^1|0}.$$
(44)
From Eq. (40), $`f_e(k)=1/(1+ika_{1\mathrm{D}})`$, which allows to identify $`a_{1\mathrm{D}}`$ with the 1D scattering length. Having introduced this parameter, the 1D atom-atom interaction potential can then be written in an effective form according to
$$V_{1\mathrm{D}}(z,z^{})=g_{1\mathrm{D}}\delta (zz^{}),$$
(45)
with interaction strength $`g_{1\mathrm{D}}=1/(\mu a_{1\mathrm{D}})`$ . For very low energies, $`k0`$, we can now formally set $`E=E_0`$ in Eq. (44). For a confining trap, $`\stackrel{~}{\zeta }_{E_0}`$ is an Hermitian operator with discrete spectrum $`\{\lambda _n\}`$ and eigenvectors $`|e_n`$, which eventually have to be determined for the particular Hamiltonian. Thus we find
$$g_{1\mathrm{D}}=\frac{1}{2\mu c^2}\underset{n}{}\frac{|0|e_n|^2}{\lambda _n+1/(4\pi a)}.$$
(46)
This result has interesting consequences for the two-body interaction. The denominator can become singular for particular values of $`a`$, thereby generating a CIR. Every eigenvalue $`\lambda _n`$ corresponds to a different CIR, unless the overlap $`0|e_n`$ vanishes due to some underlying symmetry of the Hamiltonian. We anticipate that for identical parabolic confinement potentials, the decoupling of the COM motion implies that only one resonance is permitted. For confining potentials with cylindrical symmetry, there is a resonance for each eigenvector of $`\stackrel{~}{\zeta }_{E_0}`$ with zero angular momentum. For confining potentials obeying parity symmetry, the eigenstates $`|e_n`$ must be even. These two symmetries allow in principle for infinitely many resonances. In practice, however, only few of them can be resolved because the resonances become increasingly sharper when $`|0|e_n|^20`$, making them difficult to detect.
### 4.2 Interpretation of the CIR as Feshbach resonances
A very simple and illuminating analysis, similar to that for standard Feshbach resonances , is also possible for the CIR. The two-particle Schrรถdinger equation can be rewritten as an effective Schrรถdinger equation for the scattering states in the open channel, $`(EH_{\mathrm{eff}})๐ซ|\mathrm{\Psi }=0`$, with the effective Hamiltonian
$$H_{\mathrm{eff}}=H_{\mathrm{open}}+๐ซH\frac{1}{EH_{\mathrm{closed}}}H๐ซ.$$
(47)
Here, $`H_{\mathrm{open}}=๐ซH๐ซ`$ and $`H_{\mathrm{closed}}=H`$, where $`๐ซ`$ and $``$ are projectors to open and closed channels, respectively. This equation can be expressed in terms of the closed-channel eigenstates $`|\mathrm{\Phi }_n`$,
$$H_{\mathrm{eff}}=H_{\mathrm{open}}+๐ซH\underset{n}{}\frac{|\mathrm{\Phi }_n\mathrm{\Phi }_n|}{EE_n}H๐ซ,$$
(48)
with $`H_{\mathrm{closed}}|\mathrm{\Phi }_n=E_n|\mathrm{\Phi }_n`$. This implies that a Feshbach-like resonance is possible at zero momentum if two conditions are fulfilled. First, there exists a solution of $`(E_0H_{\mathrm{closed}})|\mathrm{\Phi }=0`$, i.e., $`|\mathrm{\Phi }`$ is a bound state of $`H_{\mathrm{closed}}`$ with energy $`E=E_0`$. Second, $`|\mathrm{\Phi }`$ must be coupled to the open channel, $`๐ซH|\mathrm{\Phi }0`$.
Within the pseudopotential approximation, the equation $`(E_0H_{\mathrm{closed}})|\mathrm{\Phi }=0`$ is solved in terms of the Greenโs function
$$G_{E_0}(๐_{},๐ซ;๐_{}^{},0)=_0^{\mathrm{}}๐te^{E_0t}\stackrel{~}{G}_t(๐_{},๐ซ;๐_{}^{},0)$$
(49)
by the state
$$\mathrm{\Phi }(๐_{},๐ซ)=๐๐_{}^{}G_{E_0}(๐_{},๐ซ;๐_{}^{},0)\frac{f(๐_{}^{})}{2\mu },$$
(50)
together with boundary condition
$$\mathrm{\Phi }(๐_{},๐ซ0)\frac{f(๐_{})}{4\pi r}\left(1\frac{r}{a}\right).$$
(51)
This leads to the eigenvalue equation
$$\frac{|f}{4\pi a}=\stackrel{~}{\zeta }_{E_0}|f,$$
(52)
which is solved by the eigenvectors $`|e_n`$ introduced above. This yields $`a=1/(4\pi \lambda _n)`$, implying that there is a bound state $`|\mathrm{\Phi }`$ of $`H_{\mathrm{closed}}`$ with energy equal to the energy of the incoming wave, corresponding to the resonances found in the previous subsection. The CIR is then in complete analogy to a zero-momentum Feshbach resonance. Due to the small but finite coupling to the closed channels, two incoming particles initially in the open channel visit the closed channels during the scattering process. This process is strongly intensified when a bound-state exists whose energy is close to the continuum threshold. Then, a scattering resonance results. Note that such a bound state can be occupied only virtually by two particles during the scattering process. Hence from now on we will refer to such a bound state as a virtual bound state.
It is also possible to recover the overlap condition $`0|e_n0`$ in this framework. In fact,
$$๐ซH\mathrm{\Phi }(๐_{},๐ซ)=๐ซV(๐ซ)\mathrm{\Phi }(๐_{},๐ซ)=๐ซ\frac{๐_{}|e_n}{2\mu }\delta (๐ซ)=\psi _0(๐_{},๐ซ_{})\delta (z)\frac{0|e_n}{2\mu c},$$
(53)
since $`|\mathrm{\Phi }`$ fulfills Eq. (51) with $`f(๐)=๐_{}|e_n`$. Hence, the two overlap conditions
$$๐ซH|\mathrm{\Phi }00|e_n0$$
(54)
are equivalent. When they are not fulfilled, there exists a virtual bound state with energy $`E_0`$, but it is not coupled to the incoming wave.
## 5 Special case of harmonic confinement
In the previous sections, we have formulated the theory for a general confining potential and for two different atomic species. As a simple illustration, we now consider the case of harmonic confinement, $`V_i(๐ฑ_i)=m_i\omega _i^2๐ฑ_i^2/2`$. In COM and relative coordinates,
$`V_{\mathrm{conf}}(๐_{},๐ซ_{})`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left(m_1\omega _1^2+m_2\omega _2^2\right)|๐_{}|^2+{\displaystyle \frac{1}{2}}\left({\displaystyle \frac{\mu ^2}{m_1}}\omega _1^2+{\displaystyle \frac{\mu ^2}{m_2}}\omega _2^2\right)|๐ซ_{}|^2`$ (55)
$`+\mu \left(\omega _1^2\omega _2^2\right)๐ซ_{}๐_{}.`$
In general, the COM and the relative coordinate do not decouple, and in order to find the scattering and bound-state solutions, we have to follow the procedure outlined in the previous sections. To that end, we label the single-particle transverse states by quantum numbers $`\lambda =\{m,n\}`$, where $`m`$ is the integer angular momentum and $`n`$ the integer radial quantum number. The eigenenergies and -states of the 2D harmonic oscillator
$$E_\lambda ^{(i)}=\omega _iฯต_{n,m},\psi _\lambda ^{(i)}=\frac{1}{a_i}\psi _{n,m}\left(\frac{๐ฑ_{}}{a_i}\right),$$
with the oscillator lengths $`a_i=(m_i\omega _i)^{1/2},i=1,2`$, can be expressed in terms of the quantities
$$ฯต_{n,m}=2n+|m|+1\mathrm{and}\psi _{n,m}(๐ฑ_{})=e^{im\varphi }R_{n,m}(|๐ฑ_{}|),$$
where
$$R_{n,m}(\rho )=\frac{1}{\sqrt{\pi }}\left(\frac{n!}{(n+|m|)!}\right)^{1/2}e^{\rho /2}\rho ^{|m|}L_n^{|m|}(\rho ^2),$$
with $`L_n^{|m|}(x)`$ being the standard Laguerre polynomials. A convenient choice for the orthonormal basis $`|j`$ introduced in Eq. (31) is then given by
$$๐_{}|j=๐_{}|m,n=\frac{1}{a_M}\psi _{n,m}\left(\frac{|๐_{}|}{a_M}\right),$$
(56)
with the length scale $`a_M=(m_1\omega _1+m_2\omega _2)^{1/2}`$. In particular, we find for $`|0=|0,0`$ that $`๐_{}|0`$ fulfills Eq. (32) with $`c=\sqrt{\pi }a_1a_2/a_M`$.
The single-particle imaginary-time propagator for a 2D harmonic oscillator with length scale $`a_0`$ and frequency $`\omega `$ is given by
$`{\displaystyle \underset{\lambda }{}}e^{\omega ฯต_\lambda t}{\displaystyle \frac{1}{a_0^2}}\psi _\lambda \left({\displaystyle \frac{๐ฑ_{}}{a_0}}\right)\overline{\psi }_\lambda \left({\displaystyle \frac{๐ฑ_{}^{}}{a_0}}\right)={\displaystyle \frac{1}{\pi a_0^2}}{\displaystyle \frac{e^{\omega t}}{1e^{2\omega t}}}\mathrm{exp}\left[{\displaystyle \frac{๐ฑ_{}^2+๐ฑ_{}^2}{2a_0^2}}\mathrm{coth}(\omega t)+{\displaystyle \frac{๐ฑ_{}๐ฑ_{}^{}}{a_0^2\mathrm{sinh}(\omega t)}}\right].`$
(57)
Inserting this into Eq. (20) with $`๐ฑ_{,i}=๐_{}`$, $`๐ฑ_{,i}^{}=๐_{}^{}`$ and $`z=0`$, we find
$`G_t(๐_{},0;๐_{}^{},0)`$ $`=`$ $`\sqrt{{\displaystyle \frac{\mu }{2\pi t}}}{\displaystyle \frac{\beta (1\beta )}{\pi ^2a_M^4}}{\displaystyle \frac{e^{\omega _1t}}{1e^{2\omega _1t}}}{\displaystyle \frac{e^{\omega _2t}}{1e^{2\omega _2t}}}`$ (58)
$`\times \mathrm{exp}\left[{\displaystyle \frac{๐_{}^2+๐_{}^2}{2a_M^2}}f(t)+{\displaystyle \frac{๐_{}๐_{}^{}}{a_M^2}}g(t)\right],`$
where we have introduced $`\beta =a_M^2/a_1^2`$ and
$`f(t)`$ $`=`$ $`\beta \mathrm{coth}(\omega _1t)+(1\beta )\mathrm{coth}(\omega _2t),`$
$`g(t)`$ $`=`$ $`\beta \mathrm{sinh}^1(\omega _1t)+(1\beta )\mathrm{sinh}^1(\omega _2t).`$ (59)
In order to compute explicitly the operators $`\zeta _E`$ and $`\stackrel{~}{\zeta }_E`$, we still have to project onto the discrete basis $`\{|j\}`$ and to perform the imaginary-time integral for each matrix element. In general, this cannot be achieved analytically, and one has to resort to a numerical evaluation. Only for $`\omega _1=\omega _2`$, a complete analytical solution is possible. Since the COM degrees of freedom separate, this solution is a trivial extension of Ref. . Nonetheless, along with the general analysis of the previous section, it provides a physical picture for weak interaction between the COM and the relative degrees of freedom.
### 5.1 Identical frequencies
For $`\omega _1=\omega _2=\omega `$, the COM and relative coordinates separate, $`H=H_{\mathrm{rel}}+H_{\mathrm{COM}}`$, with
$$H_{\mathrm{rel}}=\frac{๐ฉ^2}{2\mu }+\frac{1}{2}\mu \omega ^2๐ซ_{}^2+V(๐ซ),H_{\mathrm{COM}}=\frac{๐^2}{2M}+\frac{1}{2}M\omega ^2๐_{}^2.$$
(60)
In this case, we can consider the two-particle system being (asymptotically) in the ground state of the decoupled COM Hamiltonian, and just solve the relative problem . Moreover, with $`f(t)=\mathrm{coth}(\omega t)`$ and $`g(t)=\mathrm{sinh}^1(\omega t)`$, the Greenโs function (58) simplifies to
$$G_t(๐_{},0;๐_{}^{},0)=\sqrt{\frac{\mu }{2\pi t}}\frac{\beta (1\beta )}{\pi a_M^2}\frac{e^{\omega t}}{1e^{2\omega t}}\underset{n,m}{}e^{\omega ฯต_{n,m}t}\frac{1}{a_M^2}\psi _{n,m}\left(\frac{๐_{}}{a_M}\right)\overline{\psi }_{n,m}\left(\frac{๐_{}^{}}{a_M}\right).$$
(61)
In this case, $`|n,m`$ is an eigenstate of the decoupled Hamiltonian $`H_{\mathrm{COM}}`$, and describes the COM motion also for finite $`๐ซ`$. Moreover, $`a_M=(M\omega )^{1/2}`$ and $`a_\mu =a_M/(\beta (1\beta ))=(\mu \omega )^{1/2}`$ are the characteristic lengths associated with $`H_{\mathrm{COM}}`$ and $`H_{\mathrm{rel}}`$, respectively. Inserting Eq. (61) into Eq. (22) and rescaling $`t`$ by $`2\omega `$, we obtain
$$\zeta _E=\underset{n,m}{}\frac{|n,mn,m|}{4\pi a_\mu }_0^{\mathrm{}}\frac{dt}{(\pi t)^{1/2}}\left(\frac{e^{\mathrm{\Omega }_{n,m}(E)t}}{1e^t}\frac{1}{t}\right),$$
(62)
with $`\mathrm{\Omega }_{n,m}(E)=(1+ฯต_{n,m}E/\omega )/2`$. The integral on the rhs of Eq. (62) is related to the integral representation of the Hurvitz zeta function $`\zeta (1/2,\mathrm{\Omega }_{n,m})`$ .
#### 5.1.1 Bound states
The condition given in Eq. (15) for a bound state with transverse configuration $`|n,m`$ translates into
$$\zeta (\frac{1}{2},\mathrm{\Omega }_{n,m})=\frac{a_\mu }{a}.$$
(63)
The zeta function is monotonic, and has the asymptotic scaling behavior
$$\zeta \left(\frac{1}{2},\mathrm{\Omega }1\right)\mathrm{\Omega }^{1/2},\zeta \left(\frac{1}{2},\mathrm{\Omega }1\right)2\sqrt{\mathrm{\Omega }}.$$
(64)
Inverting Eq. (63), we recover the bound-state energy found in Ref. . The corresponding result is plotted in Fig. 1. As an immediate consequence of the decoupling of the COM degrees of freedom, the $`ฯต_\lambda `$-fold degenerate energies corresponding to excited transverse configurations follow from the COM transverse ground state by a shift along the ordinate in steps of $`\omega `$. This is indicated by the dotted curves in Fig. 1. Notice that for energies above $`E_0=2\omega `$, corresponding to $`E_B=2\omega \mathrm{\Omega }_{0,0}(E)<0`$, there exists an open channel, but the solutions associated with COM excited states are orthogonal to it. For this reason, the relevant condition for a bound state to exist with transverse configuration $`|n,m`$ is $`\mathrm{\Omega }_{n,m}(E)>0`$. ยฟFrom the scaling behaviors in Eq. (64), we find the limiting behaviors of the energy of the bound state at $`|a_\mu /a|1`$ as
$`E_{B,n,m}`$ $``$ $`{\displaystyle \frac{1}{2\mu a^2}}\mathrm{for}a>0,`$
$`E_{B,n,m}`$ $``$ $`{\displaystyle \frac{2a^2}{\mu a_\mu ^4}}\mathrm{for}a<0,`$ (65)
see Eqs. (27) and (36), with $`c=\sqrt{\pi }a_\mu `$ and $`E_{B,n,m}=\omega \mathrm{\Omega }_{n,m}`$. Hence, in this highly degenerate case, there is exactly one bound state for each transverse configuration and each scattering length $`a`$.
#### 5.1.2 Scattering states
In order to identify resonant bound states of the closed channel, and the corresponding zero-momentum CIR, we subtract the contribution of the lowest-energy scattering channel in Eq. (62), and obtain
$`\stackrel{~}{\zeta }_E`$ $`=`$ $`\zeta _E{\displaystyle \frac{|0,00,0|}{4\pi a_\mu }}{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{dt}{(\pi t)^{1/2}}}e^{\mathrm{\Omega }_0(E)t}`$ (66)
$`=`$ $`{\displaystyle \underset{n,m}{}}{\displaystyle \frac{|n,mn,m|}{4\pi a_\mu }}\zeta (1/2,\stackrel{~}{\mathrm{\Omega }}_{n,m}(E)),`$
with $`\stackrel{~}{\mathrm{\Omega }}_{0,0}(E)=\mathrm{\Omega }_{0,0}(E)+1`$ and $`\stackrel{~}{\mathrm{\Omega }}_{n,m}(E)=\mathrm{\Omega }_{n,m}(E)`$ for $`n+|m|>0`$. Hence the curve corresponding to the COM ground state is shifted vertically by $`2\omega `$ and coincides with the curve corresponding to the excited states $`|1,0`$, $`|0,2`$ and $`|0,2`$. Moreover, the coupling condition in Eq. (54) becomes $`0,0|n,m0`$, and is fulfilled only for $`n=m=0`$. Though there are in principle infinitely many closed-channel bound states with energy $`2\omega `$ (one for each curve), only one scattering resonance exists, since only one of them is coupled to the incoming scattering wave. Inserting Eq. (66) into Eq. (46) we recover for the 1D interaction strength $`g_{1\mathrm{D}}`$ the well known result
$$g_{1\mathrm{D}}=2\omega a_\mu \left(\frac{a_\mu }{a}๐\right)^1.$$
(67)
#### 5.1.3 Physical picture for the weakly interacting case
When $`\omega _1\omega _2`$ but $`\omega _1\omega _2`$, a weak coupling to the COM degrees of freedom is generated, with two important consequences: (i) the degeneracies of the bound-state energies are lifted, and (ii) the coupling to the other higher-lying bound states is non-zero. Since the operators $`\stackrel{~}{\zeta }_E`$ and $`\zeta _E`$ commute with the $`z`$-component $`L_z`$ of the angular momentum, the bound states are still labeled by the quantum numbers $`\{n,m\}`$. As far as the scattering solutions are concerned, the incoming wave is coupled only to states with angular momentum quantum number $`m=0`$. Since $`0,0|\stackrel{~}{\zeta }_E|0,01,0|\stackrel{~}{\zeta }_E|1,0`$, a small off-diagonal element $`0,0|\stackrel{~}{\zeta }_E|1,0`$ is sufficient to couple the bound state with $`\{n,m\}=\{1,0\}`$ to the incoming wave, yielding an additional CIR. As far as bound states are concerned, solutions with $`E>E_0`$ and $`m=0`$ leak into the open channel, and cannot be regarded as localized bound states. Hence, for $`|a_\mu /a|1`$ and $`a<0`$, there is only one bound state with zero angular momentum. In the opposite dimer limit, however, we encounter many dimer bound states.
### 5.2 The case $`\omega _1\omega _2`$: Relation to experiments
The case $`\omega _1\omega _2`$ is relevant for experiments involving two different atom species trapped in magnetic or optical traps . For instance, in optical traps the confining potential depends on the detuning $`\mathrm{\Delta }=\omega _{\mathrm{las}}hc/\lambda `$ of the laser frequency $`\omega _{\mathrm{las}}`$ from the characteristic frequency $`hc/\lambda `$ associated with the optical transition $`nsnp`$, and is therefore different for two different atom species. This conclusion also applies to magnetic traps if the atoms are confined in hyperfine states with different projection of the magnetic moment along the magnetic field. As a concrete example, let us consider a mixture of Bosonic <sup>87</sup>Rb atoms and Fermionic <sup>40</sup>K atoms. Sympathetic cooling has allowed to create an ultracold mixture of these two elements. By loading such a gas into a dipole trap and sweeping an external magnetic field, it has been possible to identify three heteronuclear Feshbach resonances and to measure the 3D interspecies scattering length $`a=14`$ nm. It seems feasible to tune the magnetic field near a Feshbach resonance and to observe the interspecies CIR. It is hence very interesting to know how many of them can be expected and to study their locations.
The confining potential for a neutral atom in a standing optical wave $`๐(๐ซ,t)=๐_0(๐ซ)\mathrm{Re}[\mathrm{exp}(i\omega _{\mathrm{las}}t)]`$ is $`V_{\mathrm{conf}}(๐ซ)=(\epsilon _0/4)\alpha ^{}๐_0^2(๐ซ)`$, where $`\alpha ^{}=e^2/(2m_e\omega _{\mathrm{las}}\epsilon _0\mathrm{\Delta })`$ is the real part of the polarizability . Let us consider a red-detuned laser field corresponding to $`\mathrm{\Delta }<0`$ and $`\alpha ^{}>0`$. In this configuration, the atoms are trapped around the maximum of the electric field. For a mixture of two species, each species experiences its own detuning $`\mathrm{\Delta }_\mathrm{K}`$ ($`\mathrm{\Delta }_{\mathrm{Rb}}`$) given by the two transition wavelengths $`\lambda _\mathrm{K}=767`$ nm and $`\lambda _{\mathrm{Rb}}=780`$ nm. Within a parabolic approximation for the potential around its minimum, the ratio $`\omega _\mathrm{K}/\omega _{\mathrm{Rb}}`$ of trap frequencies for K and Rb atoms becomes
$$\frac{\omega _\mathrm{K}}{\omega _{\mathrm{Rb}}}=\left(\frac{\mathrm{\Delta }_{\mathrm{Rb}}}{\mathrm{\Delta }_\mathrm{K}}\frac{m_{\mathrm{Rb}}}{m_\mathrm{K}}\right)^{1/2}.$$
(68)
Let us estimate this ratio for typical parameters. In order to suppress spontaneous emission, we assume an average detuning of $`\mathrm{\Delta }=(\mathrm{\Delta }_\mathrm{K}+\mathrm{\Delta }_{\mathrm{Rb}})/2=0.1\omega _{\mathrm{las}}`$, yielding $`\omega _{\mathrm{las}}=5hc(\lambda _\mathrm{K}^1+\lambda _{\mathrm{Rb}}^1)/11`$ and $`\mathrm{\Delta }_{\mathrm{Rb}}/\mathrm{\Delta }_\mathrm{k}=(5\lambda _\mathrm{K}^16\lambda _{\mathrm{Rb}}^1)/(5\lambda _{\mathrm{Rb}}^16\lambda _\mathrm{K}^1)=0.84`$. Taking also into account the mass ratio $`m_{\mathrm{Rb}}/m_\mathrm{K}=87/40`$, we have $`\omega _\mathrm{K}/\omega _{\mathrm{Rb}}=1.35`$, indicating a substantial coupling of COM and relative degrees of freedom.
Using Eq. (58), we can project the Greenโs function $`G_t(๐_{},0;๐_{}^{},0)`$ on the appropriate basis defined in Eq. (31) and then compute numerically $`\stackrel{~}{\zeta }_E`$ by performing the imaginary-time integration, see Appendix B. Then $`\stackrel{~}{\zeta }_{E_0}`$ can be diagonalized, and the effective interspecies 1D interaction constant $`g_{1\mathrm{D}}`$ follows according to Eq. (46). The results are shown in the upper viewgraph of Fig. 2 in terms of the characteristic length $`a_\mu =\sqrt{2/(\mu (\omega _\mathrm{K}+\omega _{\mathrm{Rb}}))}`$. We find two resonances, indicating that the discussion of Sec. 5.1.3 applies to this particular case. In order to illustrate the interpretation of the CIR in terms of Feshbach-type resonances with bound states of the closed channels, we also plot in the lower viewgraph of Fig. 2 the dimensionless binding energy $`\mathrm{\Omega }=2(E\omega _\mathrm{K}+\omega _{\mathrm{Rb}})/(\omega _\mathrm{K}+\omega _{\mathrm{Rb}})`$ of the corresponding bound state. As expected, the resonances occur at those values of $`a_\mu /a`$ for which the energy of the bound state of the closed channels coincides with the continuum threshold of the open channel.
## 6 Non-parabolic confining potentials
Describing the potential created by an optical or a magnetic guide as parabolic is to some extent a simplification which has to be verified. In fact, even though the lower-energy transverse states can rather well be approximated by the eigenfunctions of a 2D harmonic oscillator, in every real trap the confinement is to some degree non-parabolic. For resonant scattering, we expect to have a virtual occupation of many non-parabolic transverse states. As a consequence, the location of the CIR will be slightly moved, and new resonances could be created. This can already be seen from an analysis similar to the one in Sec. 5.1.3 for small non-parabolic corrections. In order to tackle the problem quantitatively, a full numerical treatment is required since no analytical expression for the Greenโs function is in general available, in contrast to Sec. 5.
As an example, we consider the small non-parabolicity due to the presence of a longitudinal magnetic bias field $`B_z`$ in a magnetic waveguide containing a single-species gas. This is necessary to avoid Majorana spin flips and the subsequent escape of atoms out of the trap. A magnetic trapping potential is formed according to $`V_{\mathrm{conf}}(๐ฑ)=\mu _m|๐(๐ฑ)|`$, where $`๐(๐ฑ)`$ is the applied magnetic field and $`\mu _m=m_Fg_F\mu _B`$, with $`m_F`$ being the magnetic quantum number of the atom in the hyperfine state $`|F,m_F`$, $`g_F`$ the Landรฉ factor and $`\mu _B`$ the Bohr magneton. Assuming that apart from the longitudinal bias field, the remaining magnetic fields create a parabolic and isotropic confinement in the transverse direction, the total confinement is given by
$$\chi V_{\mathrm{conf}}(๐ฑ)=\sqrt{1+2\chi (x^2+y^2)},$$
(69)
where we have scaled energy in units of the parabolic trapping frequency $`\omega `$ and length in units of $`a_\mu =(\mu \omega )^{1/2}`$. The parameter $`\chi =\omega /(\mu _mB_z)`$ is related to the Majorana spin flip rate $`\mathrm{\Gamma }_{\mathrm{loss}}`$ . The 1D effective interaction strength $`g_{1\mathrm{D}}`$ can be calculated following our general approach. We compute $`\stackrel{~}{\zeta }_{E_0}`$ numerically as outlined in Appendix B. The results are shown in Fig. 3 for $`\chi =0.067`$, which corresponds to $`\mathrm{\Gamma }_{\mathrm{loss}}=10^6\omega `$. We find two resonances reflecting the cylindrical symmetry of the potential (69) and the weakness of non-parabolic corrections. The degeneracy of the parabolic case (shown in Fig. 3) is lifted and the original CIR is split into two nearby resonances. As expected, the effect of the non-parabolic transverse states shows up only in the deep resonant region, making the parabolic solution a very good approximation away from the resonant region. In turn, this requires a good experimental resolution in order to observe the two CIR.
## 7 Conclusions
To conclude, we have presented the general solution for two-body s-wave scattering in a two-component ultracold atom gas longitudinally confined to one dimension by an arbitrary trapping potential. The underlying key property is that the center-of-mass and the relative degrees of freedom of the two-particle problem do not decouple, as it is the case for a one-component gas and a pure parabolic confinement. Thus, no reduction to an effective single-particle problem is possible and the full coupled system has to be solved. In the framework of the pseudopotential approach, we derive the energy of the bound state when all transverse channels are closed. Simple analytical results were obtained in the limiting cases of the dimer as well as the BCS limit. Moreover, scattering solutions have been obtained when just one transverse channel is open. The effective 1D interaction constant $`g_{1\mathrm{D}}`$ can be calculated after diagonalizing a reduced Greenโs function. This can be achieved analytically for the special case of parabolic confinement, where the well-known confinement-induced resonance is recovered. For a two-component gas, as well as for a non-parabolic confinement, more than one CIR occur, which reflect the symmetry properties of the confining potential. These findings were illustrated by applying our formalism to experimentally relevant questions. We are confident that once the CIR has been verified experimentally, also the effects of a non-parabolic trapping potential will be discerned.
## Acknowledgments
We thank A. Gogolin and A. Gรถrlitz for discussions. This work has been supported by the DFG-SFB/TR 12.
## Appendix A: Short-time Greenโs function
In this Appendix, we illustrate how to expand the Greenโs function
$$G_t(๐;๐^{})=๐|\mathrm{exp}[(K(๐ท)+U(๐))t]|๐^{}$$
(70)
with respect to $`t`$ yielding the expression in Eq. (23) for $`G_t(๐_{},0;๐_{}^{},0)`$. In order to simplify the notation, we have introduced the five-dimensional vectors $`๐=\{๐_{},๐ซ\}`$ and $`๐ท=\{๐_{},๐ฉ\}`$ and the functions $`K(๐ท)=๐_{}^2/2M+๐ฉ^2/2\mu `$ and $`U(๐)=V_1(๐_{}+\mu ๐ซ_{}/m_1)+V_2(๐_{}\mu ๐ซ_{}/m_2)`$ for the kinetic and the potential energy, respectively. First, we expand the Greenโs function around the free solution given by
$$๐|\mathrm{exp}[K(๐ท)t]|๐^{}=\frac{M}{2\pi t}\mathrm{exp}\left[\frac{(๐_{}๐_{}^{})^2M}{2t}\right]\left(\frac{\mu }{2\pi t}\right)^{3/2}\mathrm{exp}\left[\frac{(๐ซ๐ซ^{})^2\mu }{2t}\right].$$
(71)
In order to justify such an expansion, note that for $`t0^+`$
$$๐|K(๐ท)\mathrm{exp}[K(๐ท)t]|๐^{}=\frac{d}{dt}๐|\mathrm{exp}[K(๐ท)t]|๐^{}\delta (๐๐^{})\frac{1}{t},$$
(72)
whereas
$`๐|U(๐)\mathrm{exp}[K(๐ท)t]|๐^{}`$ $`=`$ $`U(๐)๐|\mathrm{exp}[K(๐ท)t]|๐^{}U(๐)\delta (๐๐^{}).`$ (73)
Since the kinetic energy in Eq. (72) diverges whereas the potential energy in Eq. (73) remains finite, the latter can be regarded as a small perturbation. This expansion yields
$$G_t(๐;๐^{})(1tU(๐))๐|\mathrm{exp}[K(๐ท)t]|๐^{}.$$
(74)
Let us now set $`๐_0=\{๐_{}^{},0\}`$ in Eq. (71) and expand with respect to $`t`$:
$`๐_{},0|\mathrm{exp}[K(๐ท)t]|๐_{}^{},0`$
$`=`$ $`\left({\displaystyle \frac{\mu }{2\pi t}}\right)^{3/2}{\displaystyle \frac{M}{2\pi t}}\mathrm{exp}\left[{\displaystyle \frac{(๐_{}๐_{}^{})^2M}{2t}}\right]`$
$`=`$ $`\left({\displaystyle \frac{\mu }{2\pi t}}\right)^{3/2}{\displaystyle \frac{d^2๐_{}}{(2\pi )^2}\mathrm{exp}[i๐_{}(๐_{}๐_{}^{})\frac{๐_{}^2}{2M}t]}`$
$``$ $`\left({\displaystyle \frac{\mu }{2\pi t}}\right)^{3/2}{\displaystyle \frac{d^2๐_{}}{(2\pi )^2}\left(1\frac{๐_{}^2}{2M}t\right)\mathrm{exp}[i๐_{}(๐_{}๐_{}^{})]}`$
$`=`$ $`\left({\displaystyle \frac{\mu }{2\pi t}}\right)^{3/2}\left(\delta \left(๐_{}๐_{}^{}\right)t{\displaystyle \frac{๐_{}^2}{2M}}\right).`$
In the last line, the operator $`๐_{}^2`$ stands for $`(2\pi )^2d^2๐_{}๐_{}|๐_{}๐_{}^2๐_{}|๐_{}^{}`$. Inserting $`๐_0`$ into Eq. (74), we finally obtain Eq. (23).
## Appendix B: Evaluation of the operators $`\zeta _E`$ and $`\stackrel{~}{\zeta }_E`$
In this Appendix, we outline the evaluation of the kernels $`\zeta _E`$ and $`\stackrel{~}{\zeta }_E`$ given in Eqs. (22) and (30), respectively.
### B1: Parabolic confinement, $`\omega _1\omega _2`$
First, let us consider the special case of parabolic confinement, but the two species may experience different trap frequencies. For this confinement, the Greenโs function $`G_t(๐_{},0;๐_{}^{},0)`$ is given in Eq. (58). The first step is to project this operator onto the appropriate orthonormal basis $`\{|j\}`$ defined in Eq. (31). Note that this definition allows an arbitrary choice of the basis, apart from properly fixing the vector $`|0`$. One possibility is introduced in Eq. (56). This is a natural option because it reflects the cylindrical symmetry of the problem. However, this choice would not permit further analytical progress. For this reason, we employ an alternative basis defined by
$$๐_{}|j=๐_{}|n_x,n_y=\frac{1}{a_M}\psi _{n_x}\left(\frac{x}{a_M}\right)\psi _{n_y}\left(\frac{y}{a_M}\right),$$
(75)
where $`\psi _n(x)`$ is the eigenfunction for the 1D oscillator in dimensionless units, $`\psi _n(x)=(\sqrt{\pi }2^{n_x}n!)^{1/2}\mathrm{exp}(x^2/2)H_n\left(x\right)`$, with $`H_n\left(x\right)`$ being Hermite polynomials. Note that the $`x`$ and $`y`$ directions factorize in the Greenโs function (58), allowing to perform the $`x`$ and $`y`$ integrals separately. For convenience, we introduce dimensionless coordinates $`xx/a_M`$ and find
$`[G(t)]_{๐ง,๐ฆ}`$ $`=`$ $`n_x,n_y|G_t(๐_{},0;๐_{}^{},0)|m_x,m_y`$ (76)
$`=`$ $`\sqrt{{\displaystyle \frac{\mu }{2\pi t}}}{\displaystyle \frac{\beta (1\beta )}{\pi ^2a_M^2}}{\displaystyle \frac{e^{\omega _1t}}{1e^{2\omega _1t}}}{\displaystyle \frac{e^{\omega _2t}}{1e^{2\omega _2t}}}[F(t)]_{n_x,m_x}[F(t)]_{n_y,m_y}`$
with
$$[F(t)]_{n,m}=๐x๐x^{}\overline{\psi }_n\left(x\right)\mathrm{exp}\left[\frac{x^2+x^2}{2}f(t)+xx^{}g(t)\right]\psi _m\left(x^{}\right).$$
(77)
The functions $`f(t)`$ and $`g(t)`$ are defined in Eq. (5). We perform the first integration by using the identity
$$๐ze^{(zz^{})^2}H_n(\alpha z)=\pi ^{1/2}(1\alpha ^2)^{n/2}H_n\left(\frac{\alpha z^{}}{(1\alpha ^2)^{1/2}}\right),$$
(78)
with $`\alpha =\alpha (t)=[(1+f(t))/2]^{1/2}`$, $`z=x/\alpha (t)`$ and $`z^{}=g(t)\alpha (t)x^{}/2`$, yielding
$`[F(t)]_{n,m}=(2^{n+m}m!n!)^{1/2}\alpha (t)(1\alpha (t)^2)^{n/2}`$
$`\times {\displaystyle }dx^{}\mathrm{exp}[x^2(\alpha ^2(t){\displaystyle \frac{g(t)\alpha ^2(t)}{4}})]H_n\left({\displaystyle \frac{g(t)\alpha ^2(t)}{2(1\alpha ^2(t))^{1/2}}}x^{}\right)H_m(x^{}).`$ (79)
By substituting Eq. (76) together with Eq. (B1: Parabolic confinement, $`\omega _1\omega _2`$) into $`\zeta _E`$ defined in Eq. (22), and by introducing the dimensionless time $`t^{}=\sqrt{t(\omega _1+\omega _2)}`$, we get
$`[\zeta _E]_{๐ง,๐ฆ}`$ $`=`$ $`{\displaystyle \frac{1}{4\pi a_\mu }}{\displaystyle _0^{\mathrm{}}}dt^{}\{Ah_E(t^{})\left[F\left({\displaystyle \frac{t^2}{\omega _1+\omega _2}}\right)\right]_{m_x,n_x}`$ (80)
$`\times \left[F\left({\displaystyle \frac{t^2}{\omega _1+\omega _2}}\right)\right]_{m_y,n_y}{\displaystyle \frac{2}{\pi ^{1/2}t^2}}\delta _{๐ง,๐ฆ}\}`$
with the dimensionless parameter $`A=2\pi ^{3/2}\beta (1\beta )a_\mu ^2/a_M^2`$ and
$`h_E(t^{})`$ $`=`$ $`\mathrm{exp}\left[{\displaystyle \frac{(\omega _1+\omega _2E)t^2}{\omega _1+\omega _2}}\right]\left(1\mathrm{exp}\left[{\displaystyle \frac{2\omega _1t^2}{\omega _1+\omega _2}}\right]\right)^1`$
$`\times \left(1\mathrm{exp}\left[{\displaystyle \frac{2\omega _2t^2}{\omega _1+\omega _2}}\right]\right)^1.`$
It is now possible to evaluate the matrix elements of $`[\zeta _E]_{๐ง,๐ฆ}`$ by numerically computing the double integrals in Eq. (80). Note that the integrand does not suffer from any singularity due to the rescaling of the integration variable. Moreover, the convergence of the $`x^{}`$ integral (B1: Parabolic confinement, $`\omega _1\omega _2`$) is exponentially fast. The first term in the integrand of the $`t^{}`$ integral decays exponentially at large times. Hence for large times, only the second term yields a contribution, where the integration can be performed analytically in this region. For the case of interspecies scattering of Rb and K in an optical trap, all the parameters entering in $`A,h_E(t)`$ and $`[F(t)]_{n,m}`$ can be expressed in terms of the ratios $`m_{\mathrm{Rb}}/m_\mathrm{K}`$ and $`\mathrm{\Delta }_{\mathrm{Rb}}/\mathrm{\Delta }_\mathrm{K}`$. The generalization to determine $`\stackrel{~}{\zeta }_E`$ is straightforward and not detailed further.
### B2: Non-parabolic confinement
A numerical evaluation of the operator $`\zeta _E`$ and $`\stackrel{~}{\zeta }_E`$ is less straightforward when the Greenโs function $`G_t(๐_{},0;๐_{}^{},0)`$ cannot be computed analytically. In this case, $`G_t(๐_{},0;๐_{}^{},0)`$ should be computed by numerical diagonalization of the $`H_{,i}`$ and inserting their eigenvalues and eigenfunctions into Eq. (20). For large $`t`$, this is feasible because only a small number of eigenfunctions contribute to the sum. However, for $`t0`$, the number of eigenvectors required to cancel the divergence in Eq. (22) quickly proliferates. This practical limitation can fortunately be circumvented by the following trick. Let us formally rewrite Eq. (22) as
$`\zeta _E(๐_{},๐_{}^{})`$ $`=`$ $`{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{dt}{2\mu }}\{e^{Et}[G_t(๐_{},0;๐_{}^{},0)G_t^0(๐_{},0;๐_{}^{},0)]`$ (82)
$`+e^{Et}G_t^0(๐_{},0;๐_{}^{},0)\left({\displaystyle \frac{\mu }{2\pi t}}\right)^{3/2}\delta (๐_{}๐_{}^{})\}`$
$`=`$ $`{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{dt}{2\mu }}e^{Et}[G_t(๐_{},0;๐_{}^{},0)G_t^0(๐_{},0;๐_{}^{},0)]`$
$`+\zeta _E^0(๐_{},๐_{}^{}),`$
where $`G_t^0(๐_{},0;๐_{}^{},0)`$ and $`\zeta _E^0(๐_{},๐_{}^{})`$ are the Greenโs function and the integral kernel, respectively, for an arbitrary reference confining potential $`V_0(๐ฑ_{})`$. If $`G_t^0(๐_{},0;๐_{}^{},0)`$ is known analytically, we can deal with $`\zeta _E^0(๐_{},๐_{}^{})`$ as in the previous section. For confining potentials close to the parabolic case, we choose a parabolic $`V_0(๐ฑ_{})`$.
Regarding Eq. (82), we proceed as follows. We restrict the infinite-dimensional Hilbert space to the $`๐ฉ`$ lowest eigenstates of the potential $`V_0(๐ฑ_{})`$, and diagonalize the original Hamiltonian in this $`๐ฉ`$-dimensional Hilbert space. With the eigenfunctions at hand, the Greenโs function can be computed using Eq. (20). Then, the sum in Eq. (20) is exchanged with the $`t`$-integration and the latter is performed. Next, we project the Greenโs function onto a known single-particle basis $`\{|m\}`$. To achieve numerical convergence, we increase the Hilbert space dimension $`๐ฉ`$ until the result does not change anymore. We emphasize that the overall result converges to the exact result although obviously not all the single-particle states used in computing the Greenโs function are reliable on very long distances (comparable to the numerical system size) because higher-lying energy states are increasingly inaccurate. Nevertheless, the central part (in position space) of the eigenfunctions โ which corresponds to the kinetic energy and does not feel the confinement โ is accurate enough to cancel the divergence stemming from the kinetic part. In order to compute the scattering solution, we compute $`\stackrel{~}{\zeta }_{E_0}`$ with an analogous procedure, diagonalize $`\stackrel{~}{\zeta }_{E_0}`$ numerically, and insert the result into Eq. (46). For the non-parabolic confinement in Sec. 6, a parabolic $`V_0(๐ฑ_{})`$ is appropriate. In this case, we use for $`\{|m\}`$ the orthonormal basis defined in Eq. (56). Then $`\zeta _E^0`$ is diagonal and given by Eq. (66).
## References
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# Search for First-Generation Scalar Leptoquarks in ๐โข๐ฬ collisions at โ๐ =1.96 TeV
CDF Collaboration
## Abstract
We report on a search for pair production of first-generation scalar leptoquarks ($`LQ`$) in $`p\overline{p}`$ collisions at $`\sqrt{s}`$=1.96 TeV using an integrated luminosity of 203 $`pb^1`$ collected at the Fermilab Tevatron collider by the CDF experiment. We observe no evidence for $`LQ`$ production in the topologies arising from $`LQ\overline{LQ}eqeq`$ and $`LQ\overline{LQ}eq\nu q`$, and derive 95$`\%`$ C.L. upper limits on the $`LQ`$ production cross section. The results are combined with those obtained from a separately reported CDF search in the topology arising from $`LQ\overline{LQ}\nu q\nu q`$ and 95% C.L. lower limits on the LQ mass as a function of $`\beta =BR(LQeq)`$ are derived. The limits are 236, 205 and 145 GeV/c<sup>2</sup> for $`\beta `$ = 1, $`\beta `$ = 0.5 and $`\beta `$ = 0.1, respectively.
The remarkable symmetry between quarks and leptons in the Standard Model (SM) suggests that some more fundamental theory may exist, which allows for new interactions between them. Such interactions are mediated by a new type of particle, the leptoquark (LQ)LQ and are predicted in many extensions of the SM, e.g. grand unification, technicolor, and supersymmetry with R-parity violationmodels . A LQ carries both lepton and baryon number, is a color triplet boson with spin 0 or 1, and has fractional charge. Usually it is assumed that LQs couple to fermions of the same generation to accomodate experimental constraints on flavor changing neutral currents and helicity suppressed decays.
Previous experimental limits on LQ production are summarized in dis2002 . The H1 and ZEUS experiments at the $`e^\pm p`$ collider HERA publishedd02 lower limits on the mass of a first generation LQ that depend on the unknown LQ $`lq`$ Yukawa coupling $`\lambda `$. At the LEP collider, pair production of LQs can occur in $`e^+e^{}`$ collisions via a virtual $`\gamma `$ or $`Z`$ boson in the $`s`$channel and lower limits have been presented in d03 . At the Fermilab Tevatrond04 ; d05 ; prl7200 LQ would be predominantly pair produced through $`q\overline{q}`$ annihilation and $`gg`$ fusion. Since the production is mediated via the strong interaction it is independent of $`\lambda `$, in contrast to the searches at e-p machines. The coupling strength to gluons is determined by color charges of the particles, and is model-independent in the case of scalar LQs. The production of vector LQ pairs depends on additional assumptions on LQ coupling to gluons and its cross section is typically larger than the cross section for scalar LQs production. Since the acceptance for vector and scalar LQ detection is similar, limits on the vector LQ mass will be more stringent.
In this Letter, we focus on a search for first-generation scalar LQ pairs produced in $`p\overline{p}`$ collisions at $`\sqrt{s}`$=1.96 TeV. A search for scalar LQ pairs decaying into $`\nu \nu qq`$, resulting in jets and missing transverse energy topology has been presented in prl7200 . Here we study alternative final state signatures, with LQs decaying into $`eejj`$ and a final state consisting of two electrons and two jets and LQs decaying in $`e\nu jj`$ and the final state consisting of an electron, two jets, and missing transverse energy . The results for three channels are combined and presented as a function of $`\beta `$, the LQ branching fraction into an electron and a quark.
CDF is a generalโpurpose detector built to study the physics of $`p\overline{p}`$ collisions at the Tevatron accelerator at Fermilab and it is described in detail in CDF4 . The data used in the analysis were collected during the 2002-2003 Tevatron Run II. The integrated luminosity for this data sample is 203 $`\pm `$ 12.2 pb<sup>-1</sup>. Events are selected if they pass the high $`E_T`$ electron trigger, requiring one electromagnetic trigger tower to be above threshold and a set of identification cuts on the electromagnetic cluster, track and shower profile. The efficiency of the trigger combinations used in the $`eejj`$ and $`e\nu jj`$ analyses has been measured using $`Zee`$ datazpaper ; muge and it is $`100\%`$. Electrons are reconstructed offline as calorimeter electromagnetic clusters matching a track in the central-tracking system (central electrons, $`|\eta |<1.0`$eta ) or as calorimeter electromagnetic clusters only in the forward region ($`1|\eta |3`$). Electromagnetic clusters are identified by the characteristics of their energy deposition in the calorimeter: cuts are applied on the fraction of the energy in the electromagnetic calorimeter and the isolation of the cluster. The identification efficiency for a pair of central electrons is $`92.4\%\pm 0.4`$ and for a pair of central-forward electrons is $`80\%\pm 0.4`$. The coordinate of the lepton (also assumed to be the event coordinate ) along the beamline must fall within $`60cm`$ of the center of the detector ( $`z_{vertex}`$ cut) to ensure a good energy measurement in the calorimeter. This cut has an efficiency of $`95\%\pm 0.1(stat)\pm 0.5(sys)`$, and it was determined from studies with minimum bias events. The efficiencies of the identification cuts, the trigger selection and the vertex cut, measured using $`Zee`$ data were taken into account when evaluating the signal acceptance and background estimate. Jets are reconstructed using a cone of fixed radius $`R=\sqrt{(\mathrm{\Delta }\eta ^2+\mathrm{\Delta }\varphi ^2)}=0.7`$ and required to have $`|\eta |<`$ 2.0. Jets have been calibrated as a function of $`\eta `$ and $`E_T`$ and their energy is corrected to the parton leveltopPRD . Neutrinos produce missing transverse energy, $`\text{/}E_T`$, which is measured by balancing the calorimeter energy in the transverse plane.
In the analyses we describe here, the signal selection criteria are set according to the kinematic distribution (e.g. $`E_T`$ of the electrons and $`E_T`$ of the jets) of decay products determined from Monte Carlo (MC) studies, optimized to eliminate background with a minimal loss of signal eventsfedericaThesis ; danThesis . In the dielectron and jets topology, we select events with two reconstructed isolated electrons with $`E_T>`$ 25 GeV. At least one electron is required to be central, while the other can be central or forward. Events are further selected if there are at least two jets with $`E_T>`$ 30 and 15 GeV. The dataset selected above is dominated by QCD production of $`Z`$ bosons in association with jets and $`t\overline{t}`$ production where both the $`W`$โs from top decay into an electron and neutrino. To reduce these backgrounds the following cuts are applied: i) veto of events whose reconstructed dilepton mass falls in the window $`76<m_{ee}<110`$ GeV/c<sup>2</sup> to remove the most of the $`Z`$ \+ jets contribution, ii) $`E_T(j_1)+E_T(j_2)>85`$ GeV and $`E_T(e_1)+E_T(e_2)>85`$ GeV, iii) $`\sqrt{(E_T(j_1)+E_T(j_2))^2+(E_T(e_1)+E_T(e_2))^2}>`$ 200 GeV to remove the remaining $`Z`$ \+ jets and top contributions. We studied the properties of the physics backgrounds by generating the process $`Z`$ \+ 2 jets with ALPGENalpgen \+ HERWIG HERWIG (to perform parton showering) and $`t\overline{t}`$ with PYTHIA pythia , then passing them through a complete simulation of the CDF II detector based on GEANTgeant and full event reconstruction. Other backgrounds from $`b\overline{b}`$, $`Z\tau \overline{\tau }`$, $`WW`$ are negligible due to the electron isolation and large electron and jet transverse energy requirements. To normalize the number of simulated events to data we used the theoretical cross sections for $`t\overline{t}`$ from mlm and for $`\gamma /Zee`$ \+ 2 jets from mcfm . The expected number of $`Z`$ \+ 2 jets events is $`1.9\pm 0.4`$. The expected number of $`t\overline{t}`$ events is $`0.35\pm 0.06`$ events. The background arising from multijet events where a jet is mismeasured as an electron (fake) is calculated using data, for both this analysis and the one that follows. The method used relies on the assumption that the fake electron produced by a jet will be accompanied by other particles produced by the fragmentation of the jet; thus the isolation fraction of the fake electron will generally be larger than the one corresponding to a real electron. The isolation fraction is defined here as: $`(E_T^{cone}E_T^{cluster})/E_T^{cluster}`$ where $`E_T^{cone}`$ is the sum of the electromagnetic and hadronic transverse energies measured in all towers in a radius $`R=\sqrt{(\mathrm{\Delta }\varphi ^2+\mathrm{\Delta }\eta ^2)}=0.4`$ around the electron and $`E_T^{cluster}`$ is the transverse electromagnetic energy of the electron. The phase space corresponding to the two electron isolation fractions ($`eejj`$) or to one electron isolation fraction and the $`\text{/}E_T`$ ($`e\nu jj`$) is divided in different regions. We assume that there is no correlation between the isolation of the two electrons ($`eejj`$) and the isolation of the electron and $`\text{/}E_T`$ ($`e\nu jj`$). In the region where both electrons have large isolation fraction ($`eejj`$), or where the $`\text{/}E_T`$ is small and the isolation fraction of the electron is large ($`e\nu jj`$) the $`LQ`$ contribution is expected to be negligible. We call these background-dominated regions. With these assumptions from the ratio of the number of events in the background-dominated regions we can extrapolate the contribution in the signal region. We estimate $`0_0^{+0.7}`$ fake events in the central-central category and $`4.0\pm 2.0`$ in the central-forward category. The final background estimate is $`6.2\pm 2.2`$ events. We checked the prediction of our background sources with data in a control region defined by requiring two electrons with $`E_T>`$ 25 GeV, 2 jets with $`E_T>`$ 30 GeV and $`66<m_{ee}<110`$ GeV/c<sup>2</sup>. We observe 107 events in agreement with 113 $`\pm `$ 15 predicted from SM processes.
The efficiency to detect our signal was obtained from MC simulated LQ (PYTHIA) to account for kinematical and geometrical acceptances and it is reported in Table I for various LQ mass values.
The following systematic uncertainties are considered when calculating signal acceptance and background predictions: luminosity (6%), choice of parton distribution functions (2.1%), statistical uncertainty of MC ($`<`$ 1%), jet energy scale ($`<`$ 1%), statistics of $`Ze^+e^{}`$ sample (0.8%) and $`z_{vertex}`$ cut (0.5%). After all selection cuts, 4 events are left in the $`eejj`$ channel data.
In the search in the electron and neutrino plus two jets topology, we select events with one reconstructed isolated electron with $`E_T>`$ 25 GeV. The electron is required to be central ($`|\eta |1.0`$). We veto events with a second central or forward electron to be orthogonal to the previous analysis. We then select events where there is a large missing transverse energy, $`\text{/}E_T>60`$ GeV and at least two jets with E$`{}_{T}{}^{}>`$ 30 GeV in the range $`|\eta |2`$. This time the selected dataset is dominated by QCD production of $`W`$ bosons in association with jets and top quark pairs, where either both the $`W`$โs from the top pair decay into $`l\nu `$ and one lepton is not identified, or one of the $`Ws`$ decays leptonically and the other hadronically. A small source of background is represented by $`Z`$ \+ 2 jets, where one of the electrons is not identified. To reduce these backgrounds the following cuts are applied: i)$`\mathrm{\Delta }\varphi (\text{/}E_Tjet)>10^o`$ to veto events where the transverse missing energy is mis-measured due to a jetpointing to a non instrumented region of the calorimeter , ii)$`E_T(j_1)+E_T(j_2)>80GeV`$, iii) transverse mass of electron-neutrino system, $`M_T(e\nu )>120GeV/c^2`$ to reduce the $`W`$ \+ 2 jets contribution. We studied the properties of the $`W`$ \+ jets, $`t\overline{t}`$ and $`Z`$ \+ 2 jets backgrounds using MC simulated events (ALPGEN + HERWIG and PYTHIA). The background from $`W\nu _\tau \tau `$ \+ 2 jets (ALPGEN+HERWIG) is negligible after the final mass window cut (see below), as as is the background coming from misidentified leptons and false $`\text{/}E_T`$. Our final cut consists in selecting events falling in a mass window defined around the LQ mass in the following way. We calculate the invariant mass of the electron-jet system and the transverse mass of the neutrino-jet system. Given the decay of the two LQs, there are two possible mass combinations for the electron and the neutrino with the two leading jets. We choose the combinations that minimize the difference between the electron-jet mass and the neutrino-jet transverse mass. We fit the peak of the e-jet distribution with a Gaussian, to obtain an estimate of the spread of the distribution in the signal region ($`\sigma _e`$), as well as the $`\nu jet`$ transverse mass distribution, to obtain $`\sigma _\nu `$. In the kinematic plane of $`m(ejet)`$ vs $`m_T(\nu jet)`$ we define the sides of rectangular boxes centered around various nominal LQ mass as $`3\times \sigma _{e,\nu }`$. For each LQ mass, events are accepted if they fall inside the rectangular box. The overall selection efficiency for various LQ masses is given in Table I. We checked the simulation prediction of our background sources with data in a control region defined by requiring one electron with $`E_T>`$ 25 GeV, $`\text{/}E_T>`$ 35 GeV and 2 jets with $`E_T>`$ 30 GeV. We observe 536 events in agreement with 503 $`\pm `$ 22 predicted from SM processes.
The efficiency to detect our signal was obtained from MC simulated LQ data (PYTHIA). The following systematic uncertainties are considered when calculating signal acceptance and background predictions: luminosity (6%), choice of the parton distribution functions (2.1%), statistics of MC ($`<`$ 1.0%), jet energy scale ($`<`$ 1%), electron identification (0.6%), $`z_{vertex}`$ cut (0.5%), initial and final state radiation (1.7%). The number of events in each mass region, after all selection cuts, compared with the background expectations is reported in Table II.
In the analyses described above the number of events passing the selection cuts is consistent with the expected number of background events. The conclusion of the two searches is that there is no LQ signal: hence we derive an upper limit on the LQ production cross section at 95% confidence level. We use a Bayesian approachbayes with a flat prior for the signal cross section and Gaussian priors for acceptance and background uncertainties. The cross section limits are tabulated in Table I and the mass limits are tabulated in Table III. To compare our experimental results with the theoretical expectation, we use the next-to-leading order (NLO) cross-section for scalar LQ pair production from kramer with CTEQ4 PDFcteq . The theoretical uncertainties correspond to the variations from $`M_{LQ}/2`$ to $`2M_{LQ}`$ of the renormalization scale $`\mu `$ used in the calculation. To set a limit on the LQ mass we compare our 95% CL upper experimental limit to the theoretical cross section for $`\mu =2M_{LQ}`$, which is conservative as it corresponds to the lower value of the theoretical cross section. We find lower limits on M(LQ) at 235 GeV/c<sup>2</sup> ($`\beta =1`$) and 176 GeV/c<sup>2</sup> ($`\beta =0.5`$). To obtain the best limit, we have combined the results from the two decay channels just described with the result of a search for LQ in the case where the LQ pair decays to neutrino and quark with branching ratio $`Br(LQ\nu q)=`$1.0prl7200 . The individual channel analyses are in fact optimized for fixed values of $`\beta `$ (1,0.5,0) while in the combined analysis, due to the contributions of the different decay channels, the signal acceptance can be naturally expressed as a function of $`\beta `$. As for the treatment of uncertainties, the searches in the $`eejj`$ and e$`\nu jj`$ channel use common criteria and sometimes apply the same kind of requirements so the uncertainties in the acceptances have been considered correlated. When calculating the limit combination including the $`\nu \nu jj`$ channel the uncertainties have been considered uncorrelated. For each $`\beta `$ value a 95% C.L. upper limit on the expected number of events is returned for each mass, and by comparing this to the theoretical expectation, lower limits on the LQ mass are set. The combined limit as a function of $`\beta `$ is shown in Figure 1, together with the individual channel limits. The combined mass limits are also tabulated in Table III.
In conclusion, we have performed a search for pair production of first generation scalar LQs using 203 pb<sup>-1</sup> of proton-antiproton collision data recorded by the CDF experiment during Run II of the Tevatron. The results from all the final state signatures are combined and no evidence of LQs production is observed. Assuming that a scalar LQ decays to electron and quark with variable branching ratio $`\beta `$ we exclude LQs with masses below 236 GeV/c<sup>2</sup> for $`\beta `$ = 1, 205 GeV/c<sup>2</sup> for $`\beta `$ = 0.5 and 145 GeV/c<sup>2</sup> for $`\beta `$ = 0.1.
We thank the Fermilab staff and the technical staffs of the participating institutions for their vital contributions. This work was supported by the U.S. Department of Energy and National Science Foundation; the Italian Istituto Nazionale di Fisica Nucleare; the Ministry of Education, Culture, Sports, Science and Technology of Japan; the Natural Sciences and Engineering Research Council of Canada; the National Science Council of the Republic of China; the Swiss National Science Foundation; the A.P. Sloan Foundation; the Bundesministerium fuer Bildung und Forschung, Germany; the Korean Science and Engineering Foundation and the Korean Research Foundation; the Particle Physics and Astronomy Research Council and the Royal Society, UK; the Russian Foundation for Basic Research; the Comisiรณn Interministerial de Ciencia y Tecnologรญa, Spain; and in part by the European Communityโs Human Potential Programme under contract HPRN-CT-20002, Probe for New Physics.
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# 1 Introduction
## 1 Introduction
One of the central philosophical points of modern Differential Geometry is in the idea that the existence of special metrics should imply strong geometric conditions on the manifold. An important example of special metrics are the extremal metrics of a Kรคhler manifold, introduced by Calabi \[Cal\]. By definition extremal Kรคhler metrics of a positive $`(1,1)`$-cohomology class are Kรคhler metrics with minimal $`L^2`$-norm of the scalar curvature with respect to all Kรคhler metrics in the same $`(1,1)`$-cohomology class. Calabi proved (see \[Cal\]) that extremal metrics are exactly the Kรคhler metrics with holomorphic gradient of the scalar curvature. To any positive $`(1,1)`$-cohomology class Calabi and Futaki \[Fut\] associated an invariant such that its vanishing implies that the extremal metrics of the given class are exactly the metrics with constant scalar curvature. Moreover the vanishing of the Calabi-Futaki invariant is a necessary condition for the existence of metrics with constant scalar curvature. In the case of Fano manifolds, a positive multiple of the first Chern class admits a metric with constant scalar curvature if and only if the Kรคhler metric is Einstein. The problem of existence of Kรคhler-Einstein metrics over Fano manifolds has been the subject of intense study over the last two decades. Inspired from the work of Donaldson, Mabuchi introduced the K-energy functional \[Mab\]. In a joint work with Bando, \[Ba-Ma\] they proved that the existence of a Kรคhler-Einstein metric implies a lower bound of the K-energy and the uniqueness of the Kรคhler-Einstein metric modulo the action of automorphisms of the manifold. Recently it has been shown by Chen and Tian \[Ch-Ti3\] that the existence of an extremal metric in a given class implies the lower boundedness of the K-energy and the uniqueness of the extremal metric modulo the action of automorphisms. In the special case of a multiple of the first Chern class we have a very powerfull tool which is the Kรคhler-Ricci flow. Using a similar computation as in \[Yau\] it can be proved (see \[Cao\]) that the Kรคhler-Ricci flow admit always a solution for all times. In \[Cao\] Cao proved that the solution converges to a Kรคhler-Einstein metric if the first Chern class is non positive. In this way he re-proved the celebrated Calabi-Yau theorem \[Yau\]. In the case of Fano manifolds the Kรคhler-Ricci flow does not allways converge, in fact there are Fano manifolds which do not admit Kรคhler-Einstein metrics \[Fut\], \[Tia\]. In \[Ch-Ti1\] Chen and Tian give a proof which shows that the existence of a Kรคhler-Einstein metric over a Fano manifold with nonnegative bisectional curvature implies the convergence of the Kรคhler-Ricci flow. Their hypothesis on the nonnegativity of the bisectional curvature is just a temporary technical assumption. What they really need is that the positivity of the Ricci tensor is preserved under the Kรคhler-Ricci flow. In order to prove the convergence Chen and Tian introduced in \[Ch-Ti2\] a family of energy functionals $`E_k,k=0,\mathrm{},n1`$. Until now, only $`E_0`$ and $`E_1`$ plays a key role in the convergence of the Kรคhler-Ricci flow. The functional $`E_0`$ is just the K-energy. One of the crucial points of their proof is that, over Kรคhler-Einstein manifolds admiting a metric $`\omega _0`$ such that the nonegativity of the Ricci tensor is preserved under the Kรคhler-Ricci flow with initial metric $`\omega _0`$, the energy functional $`E_1`$ is lower bounded from below (so bounded) under the flow. In fact the energy functional $`E_1`$ is always decreasing along the flows which preserve a certain uniform lower bound of the Ricci tensors of the evolved metrics. Very recently Song and Weinkove \[So-We\] have answered a question of Chen \[Che\], showing that over Kรคhler-Einstein manifolds the energy functional $`E_1`$ is bounded on the full space of potentials. We have the following general result.
###### Theorem 1
Let $`X`$ be a Fano manifold such that the K-energy functional of the canonical class $`2\pi c_1(X)`$ is bounded from below. Then the Chen-Tian energy functional $`E_1`$ of the canonical class $`2\pi c_1(X)`$ is allso bounded from below.
This is an immediate consequence of the following theorem.
###### Theorem 2
Let $`(X,\omega )`$ be a polarised Fano manifold with $`\omega 2\pi c_1(X)`$ and let $`\nu _\omega `$ and $`E_{1,\omega }`$ be respectively the the K-energy and the Chen-Tian energy functionals with reference metric $`\omega `$. Then there exist a constant $`C_\omega `$ depending only on $`\omega `$ such that for any Kรคhler-Ricci flow $`(\omega _t)_{t[0,+\mathrm{})}2\pi c_1(X)`$ we have the inequality $`E_{1,\omega }(\omega _t)2\nu _\omega (\omega _t)+C_\omega `$ for all $`t[0,+\mathrm{})`$.
We consider it necessary to write this note because our proof of theorem 2 is drastically simple. Moreover there are indications that our result will be usefull for a new existence criteria of Kรคhler-Einstein metrics.
Acknowledgments. The result presented in this note has been reported during the visit of the author in Princeton University. The author is very grateful to Professors Jean-Pierre Demailly, Joseph Kohn and Gang Tian who made possible his long visit in this institution. The author is especially grateful to Professor Gang Tian for bringing to his attention problems related with the convergence of the Kรคhler-Ricci flow. The author is allso very grateful to the referees for their nice suggestions which have contribute to improve the original version of this note.
## 2 Energy functionals
Let $`X`$ be a Fano manifold of complex dimension $`n`$ and let $`\omega 2\pi c_1(X)`$. Let $`๐ฆ_{2\pi c_1}:=\{\omega >0,\omega 2\pi c_1(X)\},`$ be the space of Kรคhler metrics in the class $`2\pi c_1(X)`$. We consider also the corresponding space of potentials $`๐ซ_\omega :=\{\phi (X,)|i\overline{}\phi >\omega \}`$ and we define $`\omega _\phi :=\omega +i\overline{}\phi `$ for every $`\phi ๐ซ_\omega `$. We will use also the notation $`_X:=(_X\omega ^n)^1_X`$ for the average operator. We remind that the scalar curvature $`\mathrm{Sc}(\omega )(X,)`$ of $`\omega `$ is defined by the formula
$$\mathrm{Sc}(\omega ):=\mathrm{Trace}__\omega (\mathrm{Ric}(\omega ))=\frac{2n\mathrm{Ric}(\omega )\omega ^{n1}}{\omega ^n},$$
where $`\mathrm{Ric}(\omega )`$ is the Ricci form. This definition coincides with the usual definition of scalar curvature of the Riemannian metric $`g=\omega (,J)`$. We define the Laplacian of a function $`f`$ by the formula
$$\mathrm{\Delta }__\omega f:=\text{Trace}__\omega (i\overline{}f)=\frac{2ni__J\overline{}__Jf\omega ^{n1}}{\omega ^n}.$$
Our Laplacian differs by a minus sign from the usual Laplace-Beltrami operator associated to the Riemannian metric $`g=\omega (,J)`$. We remind also that an energy functional $`๐ผ`$ is a continuous map $`๐ผ:๐ฆ_{2\pi c_1}\times ๐ฆ_{2\pi c_1}`$, which satisfies the conditions
$`๐ผ(\omega _1,\omega _3)`$ $`=`$ $`๐ผ(\omega _1,\omega _2)+๐ผ(\omega _2,\omega _3)`$
$`๐ผ(\omega _1,\omega _1)`$ $`=`$ $`0`$
for every $`\omega _1,\omega _2,\omega _3๐ฆ_{2\pi c_1}`$.
The generalized energy functional.
The generalized energy functional $`J_\omega :๐ซ_\omega [0,+\mathrm{})`$ is defined by the formula
$`J_\omega (\phi )`$ $`:=`$ $`{\displaystyle \underset{k=0}{\overset{n1}{}}}{\displaystyle \frac{k+1}{n+1}}{\displaystyle \underset{X}{}}i\phi \overline{}\phi \omega ^k\omega _\phi ^{nk1}`$
$`=`$ $`{\displaystyle \underset{X}{}}\phi \omega ^n{\displaystyle \frac{1}{n+1}}{\displaystyle \underset{k=0}{\overset{n}{}}}{\displaystyle \underset{X}{}}\phi \omega ^k\omega _\phi ^{nk}.`$
If $`(\phi _t)_{t(\epsilon ,\epsilon )}๐ซ_\omega `$ is a $`๐^{\mathrm{}}`$ patht then we have the important formula
$`{\displaystyle \frac{d}{dt}}{\displaystyle \underset{k=0}{\overset{n}{}}}{\displaystyle \underset{X}{}}\phi _t\omega ^k\omega _t^{nk}=(n+1){\displaystyle \underset{X}{}}\dot{\phi }_t\omega _t^n,`$ (1)
where $`\dot{\phi }_t:=\frac{}{t}\phi _t`$ and $`\omega _t:=\omega _{\phi _t}`$. In fact consider the equalities
$`{\displaystyle \frac{d}{dt}}{\displaystyle \underset{k=0}{\overset{n}{}}}{\displaystyle \underset{X}{}}\phi _t\omega ^k\omega _t^{nk}=`$
$`=`$ $`{\displaystyle \underset{k=0}{\overset{n}{}}}{\displaystyle \underset{X}{}}\left(\dot{\phi }_t\omega ^k\omega _t^{nk}+(nk)\phi _ti\overline{}\dot{\phi }_t\omega ^k\omega _t^{nk1}\right)`$
$`=`$ $`{\displaystyle \underset{k=0}{\overset{n}{}}}{\displaystyle \underset{X}{}}\left(\dot{\phi }_t\omega ^k\omega _t^{nk}+(nk)\dot{\phi }_t(\omega _t\omega )\omega ^k\omega _t^{nk1}\right)`$
$`=`$ $`{\displaystyle \underset{k=0}{\overset{n}{}}}(nk+1){\displaystyle \underset{X}{}}\dot{\phi }_t\omega ^k\omega _t^{nk}{\displaystyle \underset{k=0}{\overset{n}{}}}(nk){\displaystyle \underset{X}{}}\dot{\phi }_t\omega ^{k+1}\omega _t^{nk1}`$
$`=`$ $`(n+1){\displaystyle \underset{X}{}}\dot{\phi }_t\omega _t^n.`$
If we set $`J(\omega ,\omega _\phi ):=J_\omega (\phi )`$ for any $`\omega ๐ฆ_{2\pi c_1}`$ then the formula (1) implies that $`J`$ is an energy functional.
The K-energy functional of the canonical class $`2\pi c_1`$.
Consider the Kรคhler-Ricci flow $`(\omega _t)_t,`$
$`{\displaystyle \frac{d}{dt}}\omega _t=\omega _t\mathrm{Ric}(\omega _t)`$ (2)
with initial metric $`\omega _0๐ฆ_{2\pi c_1}`$. It was proved in \[Cao\] that the Kรคhler-Ricci flow $`(\omega _t)_t`$ exists for all $`t[0,+\mathrm{})`$ and $`(\omega _t)_t๐ฆ_{2\pi c_1}`$. This is because to solve the equation (2) is sufficient to solve the equation in terms of potentials
$`\dot{\phi }_t=\mathrm{log}{\displaystyle \frac{\omega _t^n}{\omega ^n}}+\phi _th_\omega ,`$ (3)
where $`\phi _t๐ซ_\omega ,\omega _t=\omega +i\overline{}\phi _t`$ and $`h_\omega (X,)`$ is the the real Smooth function defined by the conditions
$$\mathrm{Ric}(\omega )=\omega +i\overline{}h_\omega ,\underset{X}{}(e^{h_\omega }1)\omega ^n=0.$$
We remark that to find $`\phi ๐ซ_\omega `$ solution of Einstein the equation $`\mathrm{Ric}(\omega _\phi )=\omega _\phi ,`$ is equivalent to solve the equation
$$0=\mathrm{log}\frac{\omega _\phi ^n}{\omega ^n}+\phi h_\omega ,$$
which is also equivalent to the constant scalar curvature equation $`\mathrm{Sc}(\omega _\phi )=2n`$. This last equation is the Euler-Lagrange equation of the K-energy functional $`\nu _\omega :๐ซ_\omega `$
$$\nu _\omega (\phi ):=\underset{X}{}(\mathrm{log}\frac{\omega _\phi ^n}{\omega ^n}+\phi h_\omega )\omega _\phi ^n\frac{1}{n+1}\underset{k=0}{\overset{n}{}}\underset{X}{}\phi \omega ^k\omega _\phi ^{nk}+\underset{X}{}h_\omega \omega ^n.$$
In fact for every $`๐^{\mathrm{}}`$ path $`(\phi _t)_{t(\epsilon ,\epsilon )}๐ซ_\omega `$ we have the identity
$`{\displaystyle \frac{d}{dt}}\nu _\omega (\phi _t)={\displaystyle \frac{1}{2}}{\displaystyle \underset{X}{}}\dot{\phi }_t\left(\mathrm{Sc}(\omega _t)2n\right)\omega _t^n.`$ (4)
We prove now this identity. Set $`\mathrm{\Delta }_t:=\mathrm{\Delta }_{\omega _t}`$ and consider the derivative
$`{\displaystyle \frac{d}{dt}}{\displaystyle \underset{X}{}}\left(\mathrm{log}{\displaystyle \frac{\omega _t^n}{\omega ^n}}h_\omega \right)\omega _t^n=`$
$`=`$ $`{\displaystyle \underset{X}{}}\left({\displaystyle \frac{d}{dt}}\mathrm{log}{\displaystyle \frac{\omega _t^n}{\omega ^n}}\right)\omega _t^n+n{\displaystyle \underset{X}{}}\left(\mathrm{log}{\displaystyle \frac{\omega _t^n}{\omega ^n}}h_\omega \right)i\overline{}\dot{\phi }_t\omega _t^{n1}`$
$`=`$ $`2^1{\displaystyle \underset{X}{}}\mathrm{\Delta }_t\dot{\phi }_t\omega _t^n+n{\displaystyle \underset{X}{}}\dot{\phi }_ti\overline{}\left(\mathrm{log}{\displaystyle \frac{\omega _t^n}{\omega ^n}}h_\omega \right)\omega _t^{n1}`$
$`=`$ $`n{\displaystyle \underset{X}{}}\dot{\phi }_ti\overline{}\left(\mathrm{log}{\displaystyle \frac{\omega _t^n}{\omega ^n}}h_\omega \right)\omega _t^{n1}`$
$`=`$ $`{\displaystyle \underset{X}{}}\dot{\phi }_t\mathrm{\hspace{0.17em}2}^1\mathrm{Sc}(\omega _t)\omega _t^n+n{\displaystyle \underset{X}{}}\dot{\phi }_t\omega \omega _t^{n1}.`$
Moreover
$`{\displaystyle \frac{d}{dt}}{\displaystyle \underset{X}{}}\phi _t\omega _t^n`$ $`=`$ $`{\displaystyle \underset{X}{}}\dot{\phi }_t\omega _t^n+n{\displaystyle \underset{X}{}}\phi _ti\overline{}\dot{\phi }_t\omega _t^{n1}`$ (5)
$`=`$ $`{\displaystyle \underset{X}{}}\dot{\phi }_t\omega _t^n+{\displaystyle \underset{X}{}}\dot{\phi }_t(\omega _t\omega )\omega _t^{n1}`$
$`=`$ $`(n+1){\displaystyle \underset{X}{}}\dot{\phi }_t\omega _t^nn{\displaystyle \underset{X}{}}\dot{\phi }_t\omega \omega _t^{n1}.`$
Combining this two equalities with the identity (1) we obtain the identity (4). If we set $`\nu (\omega ,\omega _\phi ):=\nu _\omega (\phi )`$ for every $`\omega ๐ฆ_{2\pi c_1}`$ then the formula (4) implies that $`\nu `$ is an energy functional. We remark that under the Kรคhler-Ricci flow we have the identity
$$\mathrm{Sc}(\omega _t)=2n\mathrm{\Delta }_{\omega _t}\dot{\phi }_t.$$
Then using the identity (4) we deduce the inequality
$$\frac{d}{dt}\nu _\omega (\phi _t)=2^1\underset{X}{}\dot{\phi }_t\mathrm{\Delta }_{\omega _t}\dot{\phi }_t\omega _t^n=n\underset{X}{}i\dot{\phi }_t\overline{}\dot{\phi }_t\omega _t^{n1}0,$$
which shows that the K-energy decreases under the Kรคhler-Ricci flow. We remind now that the Futaki invariant $`f_{2\pi c_1}:H^0(X,T_{_{X,J}})`$ of the Kรคhler class $`2\pi c_1(X)`$ is defined by the formula
$$f_{2\pi c_1}(\xi ):=\underset{X}{}\xi .h_\omega \omega ^n$$
and the definition is independent of the choice of $`\omega ๐ฆ_{2\pi c_1}`$. We set by $`\text{Aut}__J^0(X)`$ the identity component of the group of $`J`$-holomorphic automorphisms of $`X`$. We have the obvious action $`\text{Aut}__J^0(X)\times ๐ฆ_{2\pi c_1}๐ฆ_{2\pi c_1}`$ given by the pull back. We remind the following well known fact
###### Lemma 1
Let $`(X,\omega )`$ be a compact Kรคhler manifold with $`\omega ๐ฆ_{2\pi c_1}`$.
1$`)`$ If the K-energy $`\nu _\omega `$ is bounded from below then the Futaki invariant $`f_{2\pi c_1}`$ is zero.
2$`)`$ If the Futaki invariant $`f_{2\pi c_1}`$ is zero then the K-energy $`\nu _\omega `$ is $`\text{Aut}__J^0(X)`$-invariant.
It will be usefull to write the K-energy functional under the following synthetic form
$$\nu _\omega (\phi ):=\underset{X}{}\left(\mathrm{log}\frac{\omega _\phi ^n}{\omega ^n}h_\omega \right)\omega _\phi ^n+\underset{k=0}{\overset{n}{}}\frac{a_k}{n+1}\underset{X}{}\phi \omega ^k\omega _\phi ^{nk}+C_{0,\omega },$$
where $`a_0=n,a_k=1`$ for $`k1`$ and $`C_{0,\omega }:=\underset{X}{}h_\omega \omega ^n`$.
The Chen-Tian energy functional $`E_1`$ of the canonical class $`2\pi c_1`$.
We define the energy functional $`E_{1,\omega }:๐ซ_\omega `$ of the canonical class $`2\pi c_1`$ by the formula
$`E_{1,\omega }(\phi )`$ $`:=`$ $`{\displaystyle \underset{X}{}}\left(\mathrm{log}{\displaystyle \frac{\omega _\phi ^n}{\omega ^n}}h_\omega \right)\left(\mathrm{Ric}(\omega _\phi )+\omega \right)\omega _\phi ^{n1}`$
$`+`$ $`{\displaystyle \underset{k=0}{\overset{n}{}}}{\displaystyle \frac{b_k}{n+1}}{\displaystyle \underset{X}{}}\phi \omega ^k\omega _\phi ^{nk}+C_{1,\omega },`$
where $`b_0=b_1=n1,b_k=2`$ for $`k2`$ and
$$C_{1,\omega }:=\underset{X}{}h_\omega (\mathrm{Ric}(\omega )+\omega )\omega ^{n1}.$$
This definition coincides with the Chen-Tian definition of the energy functional $`E_{1,\omega }`$ of the canonical class $`2\pi c_1`$, (see \[Ch-Ti2\]). In fact all that we need to show is that for every $`๐^{\mathrm{}}`$ path $`(\phi _t)_{t(\epsilon ,\epsilon )}๐ซ_\omega `$ we have the identity
$`{\displaystyle \frac{d}{dt}}{\displaystyle \underset{k=0}{\overset{n}{}}}{\displaystyle \frac{b_k}{n+1}}{\displaystyle \underset{X}{}}\phi _t\omega ^k\omega _t^{nk}=(n1){\displaystyle \underset{X}{}}\dot{\phi }_t(\omega _t^2\omega ^2)\omega _t^{n2}.`$ (6)
Let prove this identity. The formulas (1) and (5) implies the equality
$`{\displaystyle \frac{d}{dt}}\left[{\displaystyle \underset{X}{}}\phi _t\omega _t^n{\displaystyle \underset{k=0}{\overset{n}{}}}{\displaystyle \underset{X}{}}\phi _t\omega ^k\omega _t^{nk}\right]=n{\displaystyle \underset{X}{}}\dot{\phi }_t\omega \omega _t^{n1}.`$ (7)
Moreover we have the equality
$`{\displaystyle \frac{d}{dt}}{\displaystyle \underset{X}{}}\phi _t\omega \omega _t^{n1}`$ $`=`$ $`{\displaystyle \underset{X}{}}\dot{\phi }_t\omega \omega _t^{n1}+(n1){\displaystyle \underset{X}{}}\phi _ti\overline{}\dot{\phi }_t\omega \omega _t^{n2}`$
$`=`$ $`{\displaystyle \underset{X}{}}\dot{\phi }_t\omega \omega _t^{n1}+(n1){\displaystyle \underset{X}{}}\dot{\phi }_t(\omega _t\omega )\omega \omega _t^{n2}`$
$`=`$ $`n{\displaystyle \underset{X}{}}\dot{\phi }_t\omega \omega _t^{n1}(n1){\displaystyle \underset{X}{}}\dot{\phi }_t\omega ^2\omega _t^{n2}.`$
Combining this last equality with the equality (7), we find the identity
$$\frac{d}{dt}\underset{k=2}{\overset{n}{}}\underset{X}{}\phi _t\omega ^k\omega _t^{nk}=(n1)\underset{X}{}\dot{\phi }_t\omega ^2\omega _t^{n2}.$$
Then using the identity (1) we find that the derivative of the quantity
$`{\displaystyle \frac{n1}{n+1}}{\displaystyle \underset{k=0}{\overset{n}{}}}{\displaystyle \underset{X}{}}\phi _t\omega ^k\omega _t^{nk}{\displaystyle \underset{k=2}{\overset{n}{}}}{\displaystyle \underset{X}{}}\phi _t\omega ^k\omega _t^{nk}=`$
$`={\displaystyle \frac{n1}{n+1}}[{\displaystyle \underset{X}{}}\phi _t\omega _t^n+{\displaystyle \underset{X}{}}\phi _t\omega \omega _t^{n1}]{\displaystyle \frac{2}{n+1}}{\displaystyle \underset{k=2}{\overset{n}{}}}{\displaystyle \underset{X}{}}\phi _t\omega ^k\omega _t^{nk}=`$
$`={\displaystyle \underset{k=0}{\overset{n}{}}}{\displaystyle \frac{b_k}{n+1}}{\displaystyle \underset{X}{}}\phi _t\omega ^k\omega _t^{nk},`$
give us the required formula (6). We remind also that \[Ch-Ti2\] for every $`๐^{\mathrm{}}`$ path $`(\phi _t)_{t(\epsilon ,\epsilon )}๐ซ_\omega `$ we have the important identity
$`{\displaystyle \frac{d}{dt}}E_{1,\omega }(\phi _t)`$ $`=`$ $`{\displaystyle \underset{X}{}}\mathrm{\Delta }_t\dot{\phi }_t\mathrm{Ric}(\omega _t)\omega _t^{n1}`$ (8)
$``$ $`(n1){\displaystyle \underset{X}{}}\dot{\phi }_t(\mathrm{Ric}(\omega _t)^2\omega _t^2)\omega _t^{n2},`$
which shows in particular that $`E_1(\omega ,\omega _\phi ):=E_{1,\omega }(\phi )`$ is an energy functional.
We are now in position to prove the theorem 2.
## 3 Proof of the theorem 2
$`Proof`$. Under the Kรคhler-Ricci flow $`\dot{\phi }_t=\mathrm{log}\frac{\omega _t^n}{\omega ^n}+\phi _th_\omega ,`$ we have the following expression for the K-energy functional
$$\nu _\omega (\phi _t):=\underset{X}{}(\dot{\phi }_t\phi _t)\omega _t^n+\underset{k=0}{\overset{n}{}}\frac{a_k}{n+1}\underset{X}{}\phi _t\omega ^k\omega _t^{nk}+C_{0,\omega }.$$
The identity $`\mathrm{Ric}(\omega _t)=\omega _ti\overline{}\dot{\phi }_t=\omega +i\overline{}(\phi _t\dot{\phi }_t)`$ implies the following expressions for the energy functional $`E_1`$ under the Kรคhler-Ricci flow
$`E_{1,\omega }(\phi _t)`$ $`=`$ $`{\displaystyle \underset{X}{}}(\dot{\phi }_t\phi _t)i\overline{}(\dot{\phi }_t\phi _t)\omega _t^{n1}+2{\displaystyle \underset{X}{}}(\dot{\phi }_t\phi _t)\omega \omega _t^{n1}`$
$`+`$ $`{\displaystyle \underset{k=0}{\overset{n}{}}}{\displaystyle \frac{b_k}{n+1}}{\displaystyle \underset{X}{}}\phi _t\omega ^k\omega _t^{nk}+C_{1,\omega }`$
$`=`$ $`{\displaystyle \underset{X}{}}i(\dot{\phi }_t\phi _t)\overline{}(\dot{\phi }_t\phi _t)\omega _t^{n1}+2{\displaystyle \underset{X}{}}(\dot{\phi }_t\phi _t)\omega _t^n`$
$``$ $`2{\displaystyle \underset{X}{}}(\dot{\phi }_t\phi _t)i\overline{}\phi _t\omega _t^{n1}+{\displaystyle \underset{k=0}{\overset{n}{}}}{\displaystyle \frac{b_k}{n+1}}{\displaystyle \underset{X}{}}\phi _t\omega ^k\omega _t^{nk}+C_{1,\omega }`$
$`=`$ $`{\displaystyle \underset{X}{}}i(\dot{\phi }_t\phi _t)\overline{}(\dot{\phi }_t\phi _t)\omega _t^{n1}+2{\displaystyle \underset{X}{}}(\dot{\phi }_t\phi _t)\omega _t^n`$
$``$ $`2{\displaystyle \underset{X}{}}i\phi _t\overline{}\phi _t\omega _t^{n1}+2{\displaystyle \underset{X}{}}i\phi _t\overline{}\dot{\phi }_t\omega _t^{n1}`$
$`+{\displaystyle \underset{k=0}{\overset{n}{}}}{\displaystyle \frac{b_k}{n+1}}{\displaystyle \underset{X}{}}\phi _t\omega ^k\omega _t^{nk}+C_{1,\omega }.`$
Expanding the integral
$`{\displaystyle \underset{X}{}}i(\dot{\phi }_t\phi _t)\overline{}(\dot{\phi }_t\phi _t)\omega _t^{n1}={\displaystyle \underset{X}{}}i\phi _t\overline{}\phi _t\omega _t^{n1}`$
$`2{\displaystyle \underset{X}{}}i\phi _t\overline{}\dot{\phi }_t\omega _t^{n1}+{\displaystyle \underset{X}{}}i\dot{\phi }_t\overline{}\dot{\phi }_t\omega _t^{n1},`$
we find
$`E_{1,\omega }(\phi _t)`$ $`=`$ $`2{\displaystyle \underset{X}{}}(\dot{\phi }_t\phi _t)\omega _t^n+{\displaystyle \underset{X}{}}i\dot{\phi }_t\overline{}\dot{\phi }_t\omega _t^{n1}`$
$``$ $`{\displaystyle \underset{X}{}}i\phi _t\overline{}\phi _t\omega _t^{n1}+{\displaystyle \underset{k=0}{\overset{n}{}}}{\displaystyle \frac{b_k}{n+1}}{\displaystyle \underset{X}{}}\phi _t\omega ^k\omega _t^{nk}+C_{1,\omega }.`$
Then the trivial identity
$`{\displaystyle \frac{n1}{n+1}}[{\displaystyle \underset{X}{}}\phi _t\omega _t^n+{\displaystyle \underset{X}{}}\phi _t\omega \omega _t^{n1}]{\displaystyle \underset{X}{}}i\phi _t\overline{}\phi _t\omega _t^{n1}=`$
$`={\displaystyle \frac{1}{n+1}}\left[2n{\displaystyle \underset{X}{}}\phi _t\omega _t^n2{\displaystyle \underset{X}{}}\phi _t\omega \omega _t^{n1}\right],`$
implies the remarkable inequality
$$E_{1,\omega }(\phi _t)=2\nu _\omega (\phi _t)+\underset{X}{}i\dot{\phi }_t\overline{}\dot{\phi }_t\omega _t^{n1}+C_\omega 2\nu _\omega (\phi _t)+C_\omega ,$$
where $`C_\omega `$ is a constant depending only on $`\omega `$. $`\mathrm{}`$
Nefton Pali
E-mail address: npali@Math.Princeton.EDU
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# K3 surfaces with Picard number one and infinitely many rational points
## 1 Introduction
K3 surfaces are the two-dimensional analogues of elliptic curves in the sense that their canonical sheaf is trivial. However, as opposed to elliptic curves, little is known about the arithmetic of K3 surfaces in general. It is for instance an open question if there exists a K3 surface $`X`$ over a number field such that the set of rational points on $`X`$ is neither empty, nor dense. We will answer a longstanding question regarding the Picard group of a K3 surface. The Picard group of a K3 surface $`X`$ over a field $`k`$ is a finitely generated free abelian group, the rank of which is called the Picard number of $`X`$. The Picard number of $`\overline{X}=X\times _k\overline{k}`$, where $`\overline{k}`$ denotes an algebraic closure of $`k`$, is called the geometric Picard number of $`X`$. We will give the first known examples of explicit K3 surfaces shown to have geometric Picard number $`1`$.
Bogomolov and Tschinkel \[BT00\] showed an interesting relation between the geometric Picard number of a K3 surface $`X`$ over a number field $`K`$ and the arithmetic of $`X`$. They proved that if the geometric Picard number is at least $`2`$, then in most cases the rational points on $`X`$ are potentially dense, which means that there exists a finite field extension $`L`$ of $`K`$ such that the set $`X(L)`$ of $`L`$-rational points is Zariski dense in $`X`$, see \[BT00\]. However, it is not yet known whether there exists any K3 surface over a number field and with geometric Picard number $`1`$ on which the rational points are potentially dense. Neither do we know if there exists a K3 surface over a number field and with geometric Picard number $`1`$ on which the rational points are not potentially dense!
In December 2002, at the AIM workshop on rational and integral points on higher-dimensional varieties in Palo Alto, Swinnerton-Dyer and Poonen asked a related question. They asked whether there exists a K3 surface over a number field and with Picard number $`1`$ that contains infinitely many rational points. In this article we will show that such K3 surfaces do indeed exist. It follows from our main theorem.
###### Theorem 1.1.
In the moduli space of K3 surfaces polarized by a very ample divisor of degree $`4`$, the set of surfaces defined over $``$ with geometric Picard number $`1`$ and infinitely many rational points is Zariski dense.
Note that a polarization of a K3 surface is a choice of an ample divisor $`H`$. The degree of such a polarization is $`H^2`$. A K3 surface polarized by a very ample divisor of degree $`4`$ is a smooth quartic surface in $`^3`$. We will prove the main theorem by exhibiting an explicit family of quartic surfaces in $`_{}^3`$ with geometric Picard number $`1`$ and infinitely many rational points. Proving that these surfaces contain infinitely many rational points is the easy part. It is much harder to prove that the geometric Picard number of these surfaces equals $`1`$. It has been known since Noether that a general hypersurface in $`_{}^3`$ of degree at least $`4`$ has geometric Picard number $`1`$. A modern proof of this fact was given by Deligne, see \[DK73\], Thm. XIX.1.2. Despite this fact, it has been an old challenge, attributed to Mumford and disposed of in this article, to find even one explicit quartic surface, defined over a number field, of which the geometric Picard number equals $`1`$. Deligneโs result does not actually imply that such surfaces exist, as โgeneralโ means โup to a countable union of closed subsets of the moduli space.โ A priori, this could exclude all surfaces defined over $`\overline{}`$. Terasoma and Ellenberg have proven independently that such surfaces do exist. The following theorems state their results.
###### Theorem 1.2 ((Terasoma, 1985)).
For any positive integers $`(n;a_1,\mathrm{},a_d)`$ not equal to $`(2;3),(n;2)`$, or $`(n;2,2)`$, and with $`n`$ even, there is a smooth complete intersection $`X`$ over $``$ of dimension $`n`$ defined by equations of degrees $`a_1,\mathrm{},a_d`$ such that the middle geometric Picard number of $`X`$ is $`1`$.
###### Proof.
See \[Te85\]. โ
###### Theorem 1.3 ((Ellenberg, 2004)).
For every even integer $`d`$ there exists a number field $`K`$ and a polarized K3 surface $`X/K`$ of degree $`d`$, with geometric Picard number $`1`$.
###### Proof.
See \[El04\]. โ
The proofs of Terasoma and Ellenberg are ineffective in the sense that they do not give explicit examples. In principle it might be possible to extend their methods to test whether a given explicit K3 surface has geometric Picard number $`1`$. In practice however, it is an understatement to say that the amount of work involved is not encouraging. The explicit examples we will give to prove the main theorem also prove the case $`(n;a_1,\mathrm{},a_d)=(2;4)`$ of Theorem 1.2 and the case $`d=4`$ of Theorem 1.3.
Shioda did find explicit examples of surfaces with geometric Picard number $`1`$. In fact, he has shown that for every prime $`m5`$ the surface in $`^3`$ given by
$$w^m+xy^{m1}+yz^{m1}+zx^{m1}=0$$
has geometric Picard number $`1`$, see \[Sh81\]. However, for $`m=4`$ this equation determines a K3 surface with maximal geometric Picard number $`20`$, i.e., a singular K3 surface.
Before we prove the main theorem in Section 3, we will recall some definitions and results.
## 2 Prerequisites
A lattice is a free $``$-module $`L`$ of finite rank, endowed with a symmetric, bilinear, nondegenerate map $`\underset{ยฏ}{},\underset{ยฏ}{}:L\times L`$, called the pairing of the lattice. A sublattice of $`L`$ is a submodule $`L^{}`$ of $`L`$, such that the induced bilinear pairing on $`L^{}`$ is nondegenerate. The Gram matrix of a lattice $`L`$ with respect to a given basis $`x=(x_1,\mathrm{},x_n)`$ is $`I_x=(x_i,x_j)_{i,j}`$. The discriminant of $`L`$ is defined by $`discL=detI_x`$ for any basis $`x`$ of $`L`$. For any sublattice $`L^{}`$ of finite index in $`L`$ we have $`discL^{}=[L:L^{}]^2discL`$. The image of $`discL`$ and $`discL^{}`$ in $`^{}/_{}^{}{}_{}{}^{2}`$ is the discriminant of the inner product space $`L_{}`$, where the inner product is induced by the pairing of $`L`$.
Let $`X`$ be a smooth, projective, geometrically integral surface over a field $`k`$ and set $`\overline{X}=X\times _k\overline{k}`$, where $`\overline{k}`$ denotes an algebraic closure of $`k`$. The Picard group $`PicX`$ of $`X`$ is the group of line bundles on $`X`$ up to isomorphism, or equivalently, the group of divisor classes modulo linear equivalence. The divisor classes that become algebraically equivalent to $`0`$ over $`\overline{k}`$ (see \[Ha77\], exc. V.1.7) form a subgroup $`Pic^0X`$ of $`PicX`$. The quotient is the Nรฉron-Severi group $`NS(X)=PicX/Pic^0X`$, which is a finitely generated abelian group, see \[Ha77\], exc. V.1.7โ8, or \[Mi80\], Thm. V.3.25, for surfaces or \[Gr71\], Exp. XIII, Thm. 5.1 in general. The intersection pairing endows the group $`NS(X)/NS(X)_{tors}`$ with the structure of a lattice. Its rank is called the Picard number of $`X`$. The Picard number of $`\overline{X}`$ is called the geometric Picard number of $`X`$.
By definition a smooth, projective, geometrically integral surface $`X`$ is a K3 surface if the canonical sheaf $`\omega _X`$ on $`X`$ is trivial and $`H^1(\overline{X},๐ช_{\overline{X}})=0`$. Examples of K3 surfaces are smooth quartic surfaces in $`^3`$. The Betti numbers of a K3 surface are $`b_0=1`$, $`b_1=0`$, $`b_2=22`$, $`b_3=0`$, and $`b_4=1`$.
###### Lemma 2.1.
If $`X`$ is a K3 surface, then $`Pic^0X`$ is trivial, the Nรฉron-Severi group $`NS(X)PicX`$ is torsion free, and the intersection pairing on $`NS(X)`$ is even.
###### Proof.
See \[BPV84\], p. 21 and Prop. VIII.3.2. โ
For any scheme $`Z`$ over $`๐ฝ_q`$ with $`q=p^r`$ and $`p`$ prime and any prime $`lp`$, we define
$$H_{\text{รฉt}}^2(Z,_l)=\left(\underset{}{lim}H_{\text{รฉt}}^2(Z,/l^n)\right)__l_l,$$
see \[Ta65\], p. 94. Furthermore, for every integer $`m`$ and every vector space $`H`$ over $`_l`$ with the Galois group $`G(\overline{๐ฝ_q}/๐ฝ_q)`$ acting on it, we define the twistings of $`H`$ to be the $`G(\overline{๐ฝ_q}/๐ฝ_q)`$-spaces $`H(m)=H__lW^m`$, where
$$W=_l__l(\underset{}{lim}\mu _{l^n})$$
is the one-dimensional $`l`$-adic vector space on which $`G(\overline{๐ฝ_q}/๐ฝ_q)`$ operates according to its action on the group $`\mu _{l^n}\overline{๐ฝ_q}`$ of $`l^n`$-th roots of unity. Here we use $`W^0=_l`$ and $`W^m=Hom(W^m,_l)`$ for $`m<0`$. For a surface $`Z`$ over $`\overline{๐ฝ_q}`$ the cup-product gives $`H_{\text{รฉt}}^2(Z,_l)(m)`$ the structure of an inner product space for all integers $`m`$.
Proposition 2.2 describes the behavior of the Nรฉron-Severi group under good reduction. Its corollary will be used to show that the geometric Picard number of a certain surface is equal to $`1`$.
###### Proposition 2.2.
Let $`A`$ be a discrete valuation ring of a number field $`L`$ with residue field $`k๐ฝ_q`$. Let $`S`$ be an integral scheme with a morphism $`SSpecA`$ that is projective and smooth of relative dimension $`2`$. Assume that the surfaces $`\overline{S}=S_{\overline{L}}`$ and $`\stackrel{~}{S}=S_{\overline{k}}`$ are integral. Let $`lq`$ be a prime number. Then there are natural injective homomorphisms
$$NS(\overline{S})_lNS(\stackrel{~}{S})_lH_{\text{รฉt}}^2(\stackrel{~}{S},_l)(1)$$
(1)
of finite dimensional inner product spaces over $`_l`$. The second injection respects the Galois action of $`G(\overline{k}/k)`$.
###### Proof.
See \[VL04\], Proposition 6.2. โ
Recall that for any scheme $`Z`$ over $`๐ฝ_q`$ with $`q=p^r`$ and $`p`$ prime, the absolute Frobenius $`F_Z:ZZ`$ of $`Z`$ acts as the identity on points, and by $`ff^p`$ on the structure sheaf. Set $`\mathrm{\Phi }_Z=F_Z^r`$ and $`\overline{Z}=Z\times \overline{๐ฝ}_q`$. Let $`\mathrm{\Phi }_Z^{}`$ denote the automorphism on $`H_{\text{รฉt}}^2(\overline{Z},_l)`$ induced by $`\mathrm{\Phi }_Z\times 1`$ acting on $`Z\times \overline{๐ฝ}_q=\overline{Z}`$.
###### Corollary 2.3.
With the notation as in Proposition 2.2, the ranks of $`NS(\stackrel{~}{S})`$ and $`NS(\overline{S})`$ are bounded from above by the number of eigenvalues $`\lambda `$ of $`\mathrm{\Phi }_{S_k}^{}`$ for which $`\lambda /q`$ is a root of unity, counted with multiplicity.
###### Proof.
By Proposition 2.2 any upper bound for the rank of $`NS(\stackrel{~}{S})`$ is an upper bound for the rank of $`NS(\overline{S})`$. Let $`\sigma `$ denote the $`q`$-th power Frobenius map, i.e., the canonical topological generator of $`G(\overline{k}/k)`$. For any positive integer $`m`$, let $`\sigma ^{}`$ and $`\sigma ^{}(m)`$ denote the automorphisms induced on $`NS(\stackrel{~}{S})_l`$ and $`H_{\text{รฉt}}^2(\stackrel{~}{S},_l)(m)`$ respectively. As all divisor classes are defined over some finite extension of $`k`$, some power of Frobenius acts as the identity on $`NS(\stackrel{~}{S})`$, so all eigenvalues of $`\sigma ^{}`$ acting on $`NS(\stackrel{~}{S})`$ are roots of unity. It follows from Proposition 2.2 that the rank of $`NS(\stackrel{~}{S})`$ is bounded from above by the number of roots of $`\sigma ^{}(1)`$ that are a root of unity. As the eigenvalues of $`\sigma ^{}(0)`$ differ from those of $`\sigma ^{}(1)`$ by a factor of $`q`$, this equals the number of roots $`\lambda `$ of $`\sigma ^{}(0)`$ for which $`\lambda q`$ is a root of unity. The Corollary follows from the fact that $`\mathrm{\Phi }_{S_k}^{}`$ acts on $`H_{\text{รฉt}}^2(\overline{Z},_l)`$ as the inverse of $`\sigma ^{}(0)`$. See also \[VL04\], Corollary 6.3. โ
###### Remark 1.
Tateโs conjecture states that the upper bound mentioned is actually equal to the rank of $`NS(\stackrel{~}{S})`$, see \[Ta65\]. Tateโs conjecture has been proven for ordinary K3 surfaces over fields of characteristic $`p5`$, see \[NO85\], Thm. 0.2.
To find the characteristic polynomial of Frobenius as in Corollary 2.3, we will use the following lemma.
###### Lemma 2.4.
Let $`V`$ be a vector space of dimension $`n`$ and $`T`$ a linear operator on $`V`$. Let $`t_i`$ denote the trace of $`T^i`$. Then the characteristic polynomial of $`T`$ is equal to
$$f_T(x)=det(xIdT)=x^n+c_1x^{n1}+c_2x^{n2}+\mathrm{}+c_n,$$
with the $`c_i`$ given recursively by
$$c_1=t_1\text{and}kc_k=t_k+\underset{i=1}{\overset{k1}{}}c_it_{ki}.$$
###### Proof.
This is Newtonโs identity, see \[Bo95\], p. 5. โ
## 3 Proof of the main theorem
First we will give a family of smooth quartic surfaces in $`^3`$ with Picard number $`1`$. Let $`R=[x,y,z,w]`$ be the homogeneous coordinate ring of $`_{}^3`$. Throughout the rest of this article, for any homogeneous polynomial $`hR`$ of degree $`4`$, let $`๐_h`$ denote the scheme in $`_{}^3`$ given by
$$wf_1+2zf_2=3g_1g_2+6h,$$
(2)
with $`f_1,f_2,g_1,g_2R`$ equal to
| $`f_1=`$ | $`x^3x^2yx^2z+x^2wxy^2xyz+2xyw+xz^2+2xzw+y^3+`$ |
| --- | --- |
| | $`+y^2zy^2w+yz^2+yzwyw^2+z^2w+zw^2+2w^3,`$ |
| $`f_2=`$ | $`xy^2+xyzxz^2yz^2+z^3,`$ |
| $`g_1=`$ | $`z^2+xy+yz,`$ |
| $`g_2=`$ | $`z^2+xy.`$ |
Its base extensions to $``$ and $`\overline{}`$ are denoted $`X_h`$ and $`\overline{X}_h`$ respectively.
###### Theorem 3.1.
Let $`hR`$ be a homogeneous polynomial of degree $`4`$. Then the quartic surface $`X_h`$ is smooth over $``$ and has geometric Picard number $`1`$. The Picard group $`Pic\overline{X}_h`$ is generated by a hyperplane section.
###### Proof.
For $`p=2,3`$, let $`X_p/๐ฝ_p`$ denote the fiber of $`๐_hSpec`$ over $`p`$. As they are independent of $`h`$, one easily checks that $`X_p`$ is smooth over $`๐ฝ_p`$ for $`p=2,3`$. As the morphism $`๐_hSpec`$ is flat and projective, it follows that the generic fiber $`X_h`$ of $`๐_hSpec`$ is smooth over $``$ as well, cf. \[Ha77\], exc. III.10.2.
We will first show that $`X_2`$ and $`X_3`$ have geometric Picard number $`2`$. For $`p=2,3`$, let $`\mathrm{\Phi }_p`$ denote the absolute Frobenius of $`X_p`$. Set $`\overline{X}_p=X_p\times \overline{๐ฝ}_p`$ and let $`\mathrm{\Phi }_p^{}(i)`$ denote the automorphism on $`H_{\text{รฉt}}^i(\overline{X}_p,_l)`$ induced by $`\mathrm{\Phi }_p\times 1`$ acting on $`\overline{X}_p=X_p\times _{๐ฝ_p}\overline{๐ฝ}_p`$. Then by Corollary 2.3 the geometric Picard number of $`X_p`$ is bounded from above by the number of eigenvalues $`\lambda `$ of $`\mathrm{\Phi }_p^{}(2)`$ for which $`\lambda /p`$ is a root of unity. We will find the characteristic polynomial of $`\mathrm{\Phi }_p^{}(2)`$ from the traces of its powers. These traces we will compute with the Lefschetz formula
$$\mathrm{\#}X_p(๐ฝ_{p^n})=\underset{i=0}{\overset{4}{}}(1)^iTr(\mathrm{\Phi }_p^{}(i)^n).$$
(3)
As $`X_p`$ is a smooth hypersurface in $`^3`$ of degree $`4`$, it is a K3 surface and its Betti numbers are $`b_0=1`$, $`b_1=0`$, $`b_2=22`$, $`b_3=0`$, and $`b_4=1`$. It follows that $`Tr(\mathrm{\Phi }_p^{}(i)^n)=0`$ for $`i=1,3`$, and for $`i=0`$ and $`i=4`$ the automorphism $`\mathrm{\Phi }_p^{}(i)^n`$ has only one eigenvalue, which by the Weil conjectures equals $`1`$ and $`p^{2n}`$ respectively. From the Lefschetz formula (3) we conclude $`Tr(\mathrm{\Phi }_p^{}(2)^n)=\mathrm{\#}X_p(๐ฝ_{p^n})p^{2n}1`$. After counting points on $`X_p`$ over $`๐ฝ_{p^n}`$ for $`n=1,\mathrm{},11`$, this allows us to compute the traces of the first $`11`$ powers of $`\mathrm{\Phi }_p^{}(2)`$. With Lemma 2.4 we can then compute the first coefficients of the characteristic polynomial $`f_p`$ of $`\mathrm{\Phi }_p^{}(2)`$, which has degree $`b_2=22`$. Writing $`f_p=x^{22}+c_1x^{21}+\mathrm{}+c_{22}`$ we find the following table.
$$\begin{array}{cccccccccccc}p& c_1& c_2& c_3& c_4& c_5& c_6& c_7& c_8& c_9& c_{10}& c_{11}\\ & & & & & & & & & & & \\ 2& 3& 2& 12& 0& 32& 64& 128& 128& 256& 0& 2048\\ & & & & & & & & & & & \\ 3& 5& 6& 72& 27& 891& 0& 9477& 4374& 78732& 19683& 708588\end{array}$$
The Weil conjectures give a functional equation $`p^{22}f_p(x)=\pm x^{22}f_p(p^2/x)`$. As in our case (both for $`p=2`$ and $`p=3`$) the middle coefficient $`c_{11}`$ of $`f_p`$ is nonzero, the sign of the functional equation is positive. This functional equation allows us to compute the remaining coefficients of $`f_p`$.
If $`\lambda `$ is a root of $`f_p`$ then $`\lambda /p`$ is a root of $`\stackrel{~}{f}_p(x)=p^{22}f_p(px)`$. Hence, the number of roots of $`\stackrel{~}{f}_p(x)`$ that are also a root of unity gives an upper bound for the geometric Picard number of $`X_p`$. After factorization into irreducible factors, we find
| $`\stackrel{~}{f}_2=`$ | $`\frac{1}{2}(x1)^2\left(2x^{20}+x^{19}x^{18}+x^{16}+x^{14}+x^{11}+2x^{10}+x^9+x^6+x^4x^2+x+2\right)`$ |
| --- | --- |
| $`\stackrel{~}{f}_3=`$ | $`\frac{1}{3}(x1)^2(3x^{20}+x^{19}3x^{18}+x^{17}+6x^{16}6x^{14}+x^{13}+6x^{12}x^{11}+`$ |
| | $`7x^{10}x^9+6x^8+x^76x^6+6x^4+x^33x^2+x+3)`$ |
Neither for $`p=2`$ nor for $`p=3`$ the roots of the irreducible factor of $`\stackrel{~}{f}_p`$ of degree $`20`$ are integral. Therefore these roots are not roots of unity and we conclude that $`\stackrel{~}{f}_p`$ has only two roots that are roots of unity, counted with multiplicities. By Corollary 2.3 this implies that the geometric Picard number of $`X_p`$ is at most $`2`$.
Note that besides the hyperplane section $`H`$, the surface $`X_2`$ also contains the conic $`C`$ given by $`w=g_2=z^2+xy=0`$. We have $`H^2=degX_2=4`$ and $`HC=degC=2`$. As the genus $`g(C)`$ of $`C`$ equals $`0`$ and the canonical divisor $`K`$ on $`X_2`$ is trivial, the adjunction formula $`2g(C)2=C(C+K)`$ yields $`C^2=2`$. Thus $`H`$ and $`C`$ generate a sublattice of $`NS(\overline{X}_2)`$ with Gram matrix
$$\left(\begin{array}{cc}4& 2\\ 2& 2\end{array}\right).$$
We conclude that the inner product space $`NS(\overline{X}_2)_{}`$ has rank $`2`$ and discriminant $`12^{}/_{}^{}{}_{}{}^{2}`$. Similarly, $`X_3`$ contains the line $`L`$ given by $`w=z=0`$, also with genus $`0`$ and thus $`L^2=2`$. The hyperplane section on $`X_3`$ and $`L`$ generate a sublattice of $`NS(\overline{X}_3)`$ of rank $`2`$ with Gram matrix
$$\left(\begin{array}{cc}4& 1\\ 1& 2\end{array}\right).$$
We conclude that the inner product space $`NS(\overline{X}_3)_{}`$ also has rank $`2`$, and discriminant $`9^{}/_{}^{}{}_{}{}^{2}`$.
Let $`\rho `$ denote the geometric Picard number $`\rho =rkNS(\overline{X}_h)`$. It follows from Proposition 2.2 that there is an injection $`NS(\overline{X}_h)_{}NS(\overline{X}_p)_{}`$ of inner product spaces for $`p=2,3`$. Hence we get $`\rho 2`$. If equality held, then both these injections would be isomorphisms and $`NS(\overline{X}_2)_{}`$ and $`NS(\overline{X}_3)_{}`$ would be isomorphic as inner product spaces. This is not the case because they have different discriminants. We conclude $`\rho 1`$. As a hyperplane section $`H`$ on $`X_h`$ has self intersection $`H^2=40`$, we find $`\rho =1`$. Since $`NS(\overline{X}_h)`$ is a $`1`$-dimensional even lattice (see Lemma 2.1), the discriminant of $`NS(\overline{X}_h)`$ is even. The sublattice of finite index in $`NS(\overline{X}_h)`$ generated by $`H`$ gives
$$4=discH=[NS(\overline{X}_h):H]^2discNS(\overline{X}_h).$$
Together with $`discNS(\overline{X}_h)`$ being even this implies $`[NS(\overline{X}_h):H]=1`$, so $`H`$ generates $`NS(\overline{X}_h)`$, which is isomorphic to $`Pic\overline{X}_h`$ by Lemma 2.1. โ
###### Remark 2.
Corollary 2.3 was pointed out to the author by Jasper Scholten and people have used it before to bound the geometric Picard number of a surface. However, since all nonreal roots of the characteristic polynomial of Frobenius come in conjugate pairs, the upper bound has the same parity as the second Betti number of the surface. For K3 surfaces this means that the upper bound is even (and therefore at least $`2`$). The strategy of the proof of Theorem 3.1 allows us to sharpen such an upper bound. If the reductions modulo two different primes give the same upper bound $`r`$, but the corresponding Nรฉron-Severi groups have discriminants that do not differ by a square factor, then in fact $`r1`$ is an upper bound.
Kloosterman has used our method to construct an elliptic K3 surface with Mordell-Weil rank $`15`$ over $`\overline{}`$, see \[Kl05\]. In the proof of Theorem 3.1 we were able to compute the discriminant up to squares of the Nรฉron-Severi lattice of $`\overline{X}_p`$ because we knew a priori a sublattice of finite index. Kloosterman realized that it is not always necessary to know such a sublattice. For an elliptic surface $`Y`$ over $`\overline{๐ฝ_p}`$, the image in $`^{}/_{}^{}{}_{}{}^{2}`$ of the discriminant of the Nรฉron-Severi lattice can also be deduced from the Artin-Tate conjecture, which has been proved for ordinary K3 surfaces in characteristic $`p5`$, see \[NO85\], Thm. 0.2, and \[Mi75\], Thm. 6.1. It allows one to compute the ratio $`discNS(Y)\mathrm{\#}Br(Y)/(NS(Y)_{tors}^2)`$ from the characteristic polynomial of Frobenius acting on $`H_{\text{รฉt}}^2(Y,_l)`$. For an elliptic surface the Brauer group has square order, so this ratio determines the same element in $`^{}/_{}^{}{}_{}{}^{2}`$ as $`discNS(Y)`$.
###### Remark 3.
In the proof we counted points over $`๐ฝ_{p^n}`$ for $`p=2,3`$ and $`n=1,\mathrm{},11`$ in order to find the traces of powers of Frobenius up to the $`11`$-th power. We could have got away with less counting. In both cases $`p=2`$ and $`p=3`$ we already know a $`2`$-dimensional subspace $`W`$ of $`NS(\overline{X}_p)__lH_{\text{รฉt}}^2(\overline{X}_p,_l)(1)`$, generated by the hyperplane section $`H`$ and another divisor class. Therefore it suffices to find out the characteristic polynomial of Frobenius acting on the quotient $`V=H_{\text{รฉt}}^2(\overline{X}_p,_l)(1)/W`$. This implies it suffices to know the traces of powers of Frobenius acting on $`V`$ up to the $`10`$-th power.
An extra trick was used for $`p=3`$. The family of planes through the line $`L`$ given by $`w=z=0`$ cuts out a fibration of curves of genus $`1`$. We can give all nonsingular fibers the structure of an elliptic curve by quickly looking for a point on it. There are efficient algorithms available in for instance Magma to count the number of points on these elliptic curves.
Using these few speed-ups we let a computer run to compute the characteristic polynomial of several random surfaces given by an equation of the form $`wf_1=zf_2`$ over $`๐ฝ_3`$ or $`wf_1=g_1g_2`$ over $`๐ฝ_2`$, as in (2). If the middle coefficient of the characteristic polynomial was zero, no more effort was spent on trying to find the sign of the functional equation (see proof of Theorem 3.1) and the surface was discarded. After one night two examples over $`๐ฝ_3`$ were found with geometric Picard number $`2`$ and one example over $`๐ฝ_2`$. With the Chinese Remainder Theorem this allows us to construct two families of surfaces with geometric Picard number $`1`$. One of these families consists of the surfaces $`X_h`$. A program written in Magma that checks the characteristic polynomial of Frobenius on $`X_2`$ and $`X_3`$ is electronically available from the author upon request.
###### Remark 4.
For $`p=2,3`$, let $`A_pNS(\overline{X}_p)`$ denote the lattice as described in the proof of Theorem 3.1, i.e., $`A_2`$ is generated by a hyperplane section and a conic, and $`A_3`$ is generated by a hyperplane section and a line. Then in fact $`A_p`$ equals $`NS(\overline{X}_p)`$ for $`p=2,3`$. Indeed, we have $`discA_p=[NS(\overline{X}_p):A_p]^2discNS(\overline{X}_p)`$. For $`p=2`$ this implies $`discNS(\overline{X}_2)=12`$ or $`discNS(\overline{X}_2)=3`$. The latter is impossible because modulo $`4`$ the discriminant of an even lattice of rank $`2`$ is congruent to $`0`$ or $`1`$. We conclude $`discNS(\overline{X}_2)=12`$, and therefore $`[NS(\overline{X}_2):A_2]=1`$, so $`A_2=NS(\overline{X}_2)`$.
For $`p=3`$ we find $`discNS(\overline{X}_3)=9`$ or $`discNS(\overline{X}_3)=1`$. Suppose the latter equation held. By the classification of even unimodular lattices we find that $`NS(\overline{X}_3)`$ is isomorphic to the lattice with Gram matrix
$$\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right).$$
By a theorem of Van Geemen this is impossible, see \[VG04\], 5.4. From this contradiction we conclude $`discNS(\overline{X}_3)=9`$ and thus $`[NS(\overline{X}_3):A_3]=1`$, so $`A_3=NS(\overline{X}_3)`$.
Since there are $`\left(\genfrac{}{}{0.0pt}{}{4+3}{3}\right)=35`$ monomials of degree $`4`$ in $`[x,y,z,w]`$, the quartic surfaces in $`_{}^3`$ are parametrized by the space $`_{}^{34}`$, which we will denote by $`M`$. Let $`M^{}^{27}M`$ denote the subvariety of those surfaces $`X`$ for which the coefficients of the monomials $`x^4`$, $`x^3y`$, $`x^3z`$, $`y^4`$, $`y^3x`$, $`y^3z`$, and $`x^2z^2`$ in the defining polynomial of $`X`$ are all zero. Note that the vanishing of the coefficients of the first six of these monomials is equivalent to the tangency of the plane $`H_w`$ given by $`w=0`$ to the surface $`X`$ at the points $`P=[1:0:0:0]`$ and $`Q=[0:1:0:0]`$. Thus, the vanishing of these coefficients yields a singularity at $`P`$ and $`Q`$ in the plane curve $`C_X=H_wX`$. If the singularity at $`P`$ in $`C_X`$ is not worse than a double point, then the vanishing of the coefficient of $`x^2z^2`$ is equivalent to the fact that the line given by $`y=w=0`$ is one of the limit-tangent lines to $`C_X`$ at $`P`$.
###### Proposition 3.2.
There is a nonempty Zariski open subset $`UM^{}`$ such that every surface $`XU`$ defined over $``$ is smooth and has infinitely many rational points.
###### Proof.
The singular surfaces in $`M^{}`$ form a closed subset of $`M^{}`$. So do the surfaces $`X`$ for which the intersection $`H_wX`$ has worse singularities than just two double points at $`P`$ and $`Q`$. Leaving out these closed subsets we obtain an open subset $`V`$ of $`M^{}`$. Let $`XV`$ be given. The plane quartic curve $`C_X=XH_w`$ has two double points, so the geometric genus $`g`$ of the normalization $`\stackrel{~}{C}_X`$ of $`C_X`$ equals $`p_a2`$, where $`p_a`$ is the arithmetic genus of $`C_X`$, see \[Ha77\], exc. IV.1.8. As we have $`p_a=\frac{1}{2}(41)(42)=3`$, we get $`g=1`$. Now assume $`X`$ is defined over $``$. One of the limit-tangents to $`C_X`$ at $`P`$ is given by $`w=y=0`$. Its slope, being rational, corresponds to a rational point $`P^{}`$ on $`\stackrel{~}{C}_X`$ above $`P`$. Fixing this point as the unit element $`๐ช=P^{}`$, the curve $`\stackrel{~}{C}_X`$ obtains the structure of an elliptic curve. Let $`DPic^0(\stackrel{~}{C}_X)`$ be the pull back under normalization of the divisor $`PQPic^0(C_X)`$. By the theory of elliptic curves there is a unique point $`T`$ on $`\stackrel{~}{C}_X`$ such that $`D`$ is linearly equivalent to $`T๐ช`$, see \[Si86\], Prop. III.3.4. As $`D`$ is defined over $``$, so is $`T`$. By Mazurโs theorem (see \[Si86\], Thm. III.7.5 for statement, \[Ma77\], Thm. 8 for a proof), the point $`T`$ has finite order if and only if $`mT=๐ช`$ for some $`m\{1,2,\mathrm{},10,12\}`$. Note that we have $`lcm(1,2,\mathrm{},10,12)=2520`$. Take for $`U`$ the complement in $`V`$ of the closed subset of those $`X`$ for which we have $`2520T=๐ช`$ for the corresponding point $`T`$ on $`\stackrel{~}{C}_X`$. Then each $`XU`$ contains an elliptic curve with infinitely many rational points. By choosing a Weierstrass equation, one verifies easily that if we take $`X_0=X_h`$ with $`h=0`$, then the corresponding point $`T`$ on $`\stackrel{~}{C}_{X_0}`$ satisfies $`mT๐ช`$ for $`m\{1,2,\mathrm{},10,12\}`$. Therefore, we find $`X_0U`$, so $`U`$ is nonempty. โ
###### Remark 5.
If $`\stackrel{~}{C}_X`$ is the normalization of $`C_X`$ as in the proof of Proposition 3.2, then generically there is another rational point $`P^{\prime \prime }`$ on $`\stackrel{~}{C}_X`$ above $`P`$, besides $`P^{}`$. Generically this point also has infinite order and the Mordell-Weil rank of $`\stackrel{~}{C}_X`$ is at least $`2`$ with independent points $`P^{\prime \prime }`$ and $`T`$ as in the proof of Proposition 3.2. For $`X=X_h`$ with $`h=0`$ however, the curve $`\stackrel{~}{C}_X`$ is given by
$$3x^2y^2+xy^2z+4xyz^2+2xz^3+5yz^3+z^4=0.$$
As the point $`P=[1:0:0]`$ is a cusp, there is only one point above $`P`$ on $`\stackrel{~}{C}_X`$ here. The conductor of this elliptic curve equals $`686004`$. Both points on $`\stackrel{~}{C}_X`$ above $`Q=[0:1:0]`$ are rational and we have an extra rational point $`[1:1:1]`$. These generate the full Mordell-Weil group of rank $`3`$.
###### Remark 6.
By requiring other coefficients to vanish than is required for $`M^{}`$, we can find quartic surfaces $`Y`$ for which the plane $`H_w`$ given by $`w=0`$ is tangent at $`[1:0:0:0]`$, $`[0:1:0:0]`$, and $`[0:0:1:0]`$. Then the intersection $`H_wY`$ has geometric genus $`0`$ and if its normalization has a point defined over $``$, then this intersection is birational to $`^1`$. The quartic surface $`Z`$ given by
$$w(x^3+y^3+z^3+x^2z+xw^2)=3x^2y^24x^2yz+x^2z^2+xy^2z+xyz^2y^2z^2$$
(4)
is an example of such a surface. As in the proof of Theorem 3.1, modulo $`3`$ the surface $`Z`$ contains the line $`z=w=0`$. Also, the reduction of $`Z`$ at $`p=2`$ contains a conic again, as the right-hand side of (4) factors over $`๐ฝ_4`$ as $`(xy+xz+\zeta yz)(xy+xz+\zeta ^2yz)`$, with $`\zeta ^2+\zeta +1=0`$. An argument very similar to the one in the proof of Theorem 3.1 then shows that $`Z`$ also has geometric Picard number $`1`$ with the Picard group generated by a hyperplane section. The only difference is that Frobenius does not act trivially on the conic $`w=xy+xz+\zeta yz=0`$. The hyperplane section $`H_wZ`$ is a curve of geometric genus $`0`$, parametrized by
$$[x:y:z:w]=[(t^2+t1)(t^2t3):2(t+2)(t^2+t1):2(t+2)(t^2t3):0].$$
The Cremona transformation $`[x:y:z:w][yz:xz:xy]`$ gives a birational map from this curve to a nonsingular plane curve of degree $`2`$. It turns out that the curve on $`Z`$ given by $`x=0`$ has a triple point at $`[0:0:0:1]`$, so it is birational to $`^1`$ as well. It can be parametrized by
$$[x:y:z:w]=[0:1+t^3:t(1+t^3):t^2].$$
From the local and global Torelli theorem for K3 surfaces, see \[PS71\], one can find a very precise description of the moduli space of polarized K3 surfaces in general, see \[Be85\]. A polarization of a K3 surface $`Z`$ by a very ample divisor of degree $`4`$ gives an embedding of $`Z`$ as a smooth quartic surface in $`^3`$ with the very ample divisor corresponding to a hyperplane section. An isomorphism between two smooth quartic surfaces in $`^3`$ that sends one hyperplane section to an other hyperplane section comes from an automorphism of $`^3`$. As any two hyperplane sections are linearly equivalent, we conclude that the moduli space of K3 surfaces polarized by a very ample divisor of degree $`4`$ is isomorphic to the open subset in $`M=^{34}`$ of smooth quartic surfaces modulo the action of $`PGL(4)`$ by linear transformations of $`^3`$. We are now ready to prove the main theorem of this article.
###### Proof Theorem 1.1.
By the description of the moduli space of K3 surfaces polarized by a very ample divisor of degree $`4`$ given above, it suffices to prove that the set $`SM()`$ of smooth surfaces with geometric Picard number $`1`$ and infinitely many rational points is Zariski dense in $`M`$. We will first show that $`SM^{}`$ is dense in $`M^{}`$. Note that the coefficients of the monomials $`x^4`$, $`x^3y`$, $`x^3z`$, $`y^4`$, $`y^3x`$, $`y^3z`$, and $`x^2z^2`$ in $`wf_1+2zf_23g_1g_2`$ in (2) are zero, so if the coefficients of these monomials in a homogeneous polynomial $`hR`$ of degree $`4`$ are all zero, then $`X_h`$ is contained in $`M^{}`$. It follows that the set
$$T=M^{}\{X_h:hR,h\text{ homogeneous of degree }4\}$$
is dense in $`M^{}`$. Let $`U`$ be as in Proposition 3.2. Then $`U`$ is a dense open subset of $`M^{}`$, so $`TU`$ is also dense in $`M^{}`$. By Theorem 3.1 and Proposition 3.2 every surface in $`TU`$ has geometric Picard number $`1`$ and infinitely many rational points. Thus we have an inclusion $`TUSM^{}`$, so $`SM^{}`$ is dense in $`M^{}`$ as well.
Let $`W`$ denote the vector space of $`4\times 4`$โmatrices over $``$ and let $`T`$ denote the dense open subset of $`(W)`$ corresponding to elements of $`PGL(4)`$. Let $`\phi :T\times M^{}M`$ be given by sending $`(A,X)`$ to $`A(X)`$. Note that $`T()\times (SM^{})`$ is dense in $`T\times M^{}`$ and $`\phi `$ sends $`T()\times (SM^{})`$ to $`S`$. Hence, in order to prove that $`S`$ is dense in $`M`$, it suffices to show that $`\phi `$ is dominant, which can be checked after extending to the algebraic closure. A general quartic surface in $`^3`$ has a one-dimensional family of bitangent planes, i.e., planes that are tangent at two different points. This is closely related to the theorem of Bogomolov and Mumford, see the appendix to \[MM83\]. In fact, for a general quartic surface $`Y^3`$, there is such a bitangent plane $`H`$, for which the two tangent points are ordinary double points in the intersection $`HY`$. Let $`Y`$ be such a quartic surface and $`H`$ such a plane, say tangent at $`P`$ and $`Q`$. Then there is a linear transformation that sends $`H`$, $`P`$, and $`Q`$ to the plane given by $`w=0`$, and the points $`[1:0:0:0]`$ and $`[0:1:0:0]`$. Also, one of the limit-tangent lines to the curve $`YH`$ at the singular point $`P`$ can be sent to the line given by $`y=w=0`$. This means that there is a linear transformation $`B`$ that sends $`Y`$ to an element $`X`$ in $`M^{}`$. Then $`\phi (B^1,X)=Y`$, so $`\phi `$ is indeed dominant. โ
###### Remark 7.
The explicit polynomials $`f_1,f_2,g_1,g_2`$ for $`X_h`$ in (2) were found by letting a computer pick random polynomials modulo $`p=2`$ and $`p=3`$ such that the surface $`X_h`$ with $`h=0`$ is contained in $`M^{}`$ as in Proposition 3.2. The computer then computed the characteristic polynomial of Frobenius and tested if there were only $`2`$ eigenvalues that were roots of unity, see Remark 3.
###### Remark 8.
In finding the explicit surfaces $`X_h`$ not much computing power was needed, as we constructed the surface to have good reduction at small primes $`p`$ so that counting points over $`๐ฝ_{p^n}`$ was relatively easy. Based on ideas of for instance Alan Lauder, Daqing Wan, Kiran Kedlaya, and Bas Edixhoven, it should be possible to develop more efficient algorithms for finding characteristic polynomials of (K3) surfaces. Together with these algorithms, the method used in the proof of Theorem 3.1 becomes a strong tool in finding Picard numbers of K3 surfaces over number fields.
## 4 Open problems
We end with the remark that still very little is known about the arithmetic of K3 surfaces, especially those with geometric Picard number $`1`$. We reiterate three questions that remain unsolved.
###### Question 1.
Does there exist a K3 surface over a number field such that the set of rational points is neither empty nor dense?
###### Question 2.
Does there exist a K3 surface over a number field with geometric Picard number $`1`$, such that the set of rational points is potentially dense?
###### Question 3.
Does there exist a K3 surface over a number field with geometric Picard number $`1`$, such that the set of rational points is not potentially dense?
###### Acknowledgements.
The author thanks the American Institute of Mathematics (Palo Alto) and the Institut Henri Poincarรฉ (Paris) for inspiring working conditions. The author also thanks Bjorn Poonen, Arthur Ogus, Jasper Scholten, Bert van Geemen, and Hendrik Lenstra for very useful discussions, and Brendan Hassett for pointing out a mistake in the first version of this article.
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# Runaway electromagnetic cascade in shear flows and high energy radiation of astrophysical jets
## 1 Introduction
Traditionally, the dissipation of a relativistic bulk motion into radiation is assumed to be associated with the shock acceleration of charged particles. The latter should certainly play some role, however, at least in a relativistic case their role is limited by a number of factors$`^{\mathrm{?},\mathrm{?}}`$. On the other hand, in some cases (blazars or gamma-ray bursts) the energy release in gamma-rays can be comparable to the estimated total kinetic energy of the bulk motion in the source and one has to search for a more efficient mechanism.
Derishev et al.$`^\mathrm{?}`$ and Stern$`^\mathrm{?}`$ independently suggested that interacting neutral particles can convert bulk kinetic energy into radiation more efficiently than this can be done by charge particles. Indeed, they easily cross the shock front or the boundary of the shear layer and can be converted into charge particles (e.g. via $`e^+e^{}`$ pair productions by two photons) inside the relativistic fluid. Then, secondary charged particles gain a factor $`\mathrm{\Gamma }^2`$ (where $`\mathrm{\Gamma }`$ is the Lorentz factor of the fluid) in energy due to gyration in magnetic field associated with the fluid and can, in turn, emit new high energy photons which can leave the fluid and interact in the external environment, producing new pairs which โreflectโ a fraction of energy towards the fluid.
In this work we consider this scenario for a shear flow taking place in astrophysical jets. We study the mechanism numerically using nonlinear Large Particle Monte-Carlo Code developed by Stern et al.$`^\mathrm{?}`$. The complete solution of the problem should account for the feedback of particles on the fluid dynamics and requires detailed numerical treatment of the hydrodynamical part of the problem. This objective is beyond the scope of this work and we restrict this first study to a simple model and try to answer question whether the supercritical runaway regime exists at reasonable conditions.
## 2 Qualitative consideration
Very schematically, the process can be split into five steps (see Fig.1).
Step 1. A high energy external photon (which origin is not important) enters the jet and interacts with a soft photon producing $`e^+e^{}`$ pair.
Step 2. The pair (originally been produced in upstream direction in the fluid frame) turns around due to gyration in the magnetic field of the jet, gaining the rest frame energy by factor $`\mathrm{\Gamma }^2`$, where $`\mathrm{\Gamma }`$ is the Lorentz factor of the jet.
Step 3. The pair Comptonizes soft photons up to high energies.
Step 4. Some of these photons leave the jet and produce pairs in the external environment.
Step 5. Pairs in the external environment Comptonizes soft photons more or less isotropically and some of Comptonized photons enter the get again. This step completes the cycle.
Each step can be characterized by its โenergy transmission coefficientโ $`c_i`$ defined as average ratio of total energy of particles with energy above pair production threshold before and after the step. $`c_2`$ is large ($`\mathrm{\Gamma }^2`$), others are smaller than 1. If $`c_1\times c_2\times c_3\times c_4\times c_5>1`$ then the regime is supercritical: each cycle produces more particles than the previous one and their number grows exponentially. In this case we deal with particle breeding rather with particle acceleration: the spectrum of particles changes slowly (and can evolve to a softer state), but the number and the total energy of particles rapidly grows.
Steps 1, 3, 4 and 5 require a field of soft photons to provide the conversion of high energy photons into pairs and production of new high energy photons through inverse Compton scattering. There are many possible sources of soft photons which were already considered in the literature: blackbody and X-ray radiation of the accretion disk; disk radiation scattered or reprocessed in the broad line region $`^{\mathrm{?},\mathrm{?},\mathrm{?}}`$; external IR radiation of dust; direct and scattered synchrotron radiation of high energy electrons in the jet $`^\mathrm{?}`$.
Next condition is the presence of a transversal or chaotic magnetic field, both in the jet (to provide step 2) and external environment (to provide isotropization of Comptonized photons at step 5).
## 3 Monte-Carlo simulations and their results
The jet was represented as a cylinder of radius $`R_j`$ and length $`20R_j`$ at distance of $`20R_j`$ from the central black hole. Actually the jet should be a cone rather than a cylinder, we neglect the jet divergence for simplicity and adopt constant physical conditions along the jet. In the course of the simulation the jet undergoes the differential deceleration: we split the jet into 500 cylindric shells, calculate the momentum transferred to each shell and decelerate each shell independently from others. This is a very rough simplification: the deceleration should depend on $`z`$ and this will lead to a complicated situation including formation of internal shocks.
The trajectories of electrons and positrons in the magnetic field were simulated directly assuming transversal geometry of the field $`H_j`$ in the jet and $`H_e`$ in the external matter. There is a primary constant soft photon field filling the whole space.
At the start of simulation the space at $`0.9<r<1`$ and $`20<z<30`$ (where $`r`$ and $`z`$ are the distances from the jet axis and from the black hole in units of $`R_j`$) is filled by seed isotropic high energy photons whose energy density is several orders of magnitude smaller than the energy density of the jet. After that there is no injection of external photons and all particles participating in the further simulation are descenders of these seed photons.
### 3.1 Example 1. Weak magnetic field and a โminimalโ seed radiation
In this example we try the simplest variant of the external photon field: two-component emission of accretion disk. First component is a blackbody spectrum integrated over disk radius with the maximum temperature 5 eV and luminosity $`L_d0.510^{45}`$erg s<sup>-1</sup>. The second component is the power-law with photon index $`\alpha =1.7`$ and cutoff at 50 keV with luminosity $`L_x0.510^{44}`$ erg s<sup>-1</sup>. The jet radius was taken $`R_j=10^{16}`$ cm, the distance from the source $`R=210^{17}`$ cm. The fraction of scattered disk radiation was taken $`\xi =0.05`$. Magnetic field in the jet was $`H_j=0.6`$ G, which corresponds to the Poynting flux $`10^{43}`$ erg s<sup>-1</sup>. The jet energy flux can be much higher, therefore such magnetic field implies matter dominated jet. The low magnetic field in this run is not our arbitrary choice since a stronger field reduces the criticality of the system. At $`H_j=1`$G the system is still supercritical, but the exponential cascade breeding is too slow. We adopt the total power $`L=10^{45}`$ erg s<sup>-1</sup>.
The total energy release as a function of time is shown in Fig. 2. We observe a reasonably fast breeding with e-folding time $`\tau 1.37`$ (see Fig. 2). The active layer is rather thin: a half of the energy release is concentrated within $`\delta r0.02`$ from the jet boundary. At $`t12`$ the regime changes: the external shell decelerates and the active layer gets wider ($`\delta r=0.05`$ at $`t=15`$ and $`\delta r=0.07`$ at $`t=20`$). The cascade breeding slows down as the photon path length through the cycle (see ยง2) increases. The total energy release into photons reached 20% of the total jet energy at the end of simulation at $`t=20`$. Our simplified model with a limited number of large particles can not follow up the evolution of the system for a longer time correctly.
While the jet decelerates in our model, the external environment is fixed at rest. Actually it undergoes a radiative acceleration. Our assumption is reasonable if the density of external medium is $`n>10^5`$ cm<sup>-3</sup>: then we actually can neglect its acceleration.
The evolution of the photon spectrum is shown in Fig. 3a. Early spectrum demonstrate two distinct components: TeV Comptonization peak (mainly Comptonized emission of the disk UV photons scattered in the broad-line region) and a synchrotron maximum. Qualitatively this spectrum resembles some BL Lac spectra, but in our simulation the synchrotron maximum is harder than usually observed.
After $`t10`$ the spectrum changes: the main peak moves to lower energies and the synchrotron peak declines. The reason for such evolution is evident: the system enters the nonlinear stage because soft synchrotron radiation of the cascade exceeds the initial soft photon field. Comptonization losses increase while synchrotron losses do not change.
### 3.2 Example 2. Strong magnetic field.
In this example we take $`R=210^{16}`$ cm, $`R_j=10^{15}`$ cm, $`H=100`$ G, $`\mathrm{\Gamma }=10`$. The Poynting flux at such parameters is $`S10^{45}`$erg s<sup>-1</sup>, which is sufficient for a moderately powerful magnetically dominated jet. The disk luminosity is $`L_d=10^{45}`$erg s<sup>-1</sup>, $`L_x=10^{44}`$erg s<sup>-1</sup>. The ratio of Compton to synchrotron losses at such parameters is $`U_{\mathrm{rad}}/U_\mathrm{B}\times \mathrm{\Gamma }^21`$. However, a particle with comoving Lorentz factor $`\gamma >10^5`$ interacts with scattered disk radiation in the Klein-Nishina regime and synchrotron losses dominate. With this reason the system remains subcritical and our trials have shown no runaway regime at such conditions.
The situation changes if we add a preexisting synchrotron radiation of the jet similar to that proposed by Ghisellini & Madau $`^\mathrm{?}`$. Then we have more soft (UV and optical) photons for Comptonization and more X-rays for conversion of lower energy gamma-rays into pairs. The seed synchrotron component was taken as a power law spectrum with photon index $`\alpha =1.7`$ with an exponential cutoff at 50 keV and the total power $`0.005`$ of the jet power.
At the start, the electromagnetic cascade breeds very rapidly (see Fig.2) with $`\tau 0.17`$. The time constant is so small due to a short free path of high energy photons moving in transversal direction as the soft photon density is much higher than in the previous example. The active layer is very thin, $`\delta r210^3`$, and the breeding cycle is short. Evidently, such regime can not last for a long time and at $`t2`$ the active layer decelerates, the photon spectrum gets softer (see Fig. 3b) and the breeding slows down (Fig. 2). At the end of the simulation run at $`t=21`$ the energy release reached 12% of the total jet energy and is probably underestimated (see Discussion).
The hard to soft evolution of high energy peak of the photon spectrum shown in Fig. 3 spans almost all range of peak energies observed in blazars. The latest spectrum peaks in MeV range as in MeV blazars, where the observed spectra have a much sharper maximum. One also can see a hint on IRโradio synchrotron component observed in blazars. This component here is less prominent than in blazars, or, in other terms, intermediate 10 eV - 100 keV luminosity in our simulations has a higher level than the observed one.
## 4 Discussion
We have demonstrated that the supercritical runaway cascade does develop at reasonable conditions and can convert into radiation at least $``$ 20% of the jet kinetic energy. This is certainly not an ultimate value: with our simplified model we are able to reproduce only the initial stage of the evolution. Actually we can expect a wealth of interesting nonlinear phenomena with jet deceleration, formation of internal shocks, non stationary behaviour producing flares and moving bright features. A more realistic model should include a detailed treatment of fluid hydrodynamics coupled with large particle electromagnetic cascade.
The model can reproduce at least the high energy component of blazar radiation. On the other hand, examples presented in this work do not reproduce the low energy synchrotron components as prominent as observed in blazars. The general impression is that our simulated spectra are qualitatively similar, but flatter and smoother than observed.
The reason why our simulated spectra are flatter (in $`E^2N_E`$ scale) and more featureless than the observed spectra is probably a too high maximal Lorentz factor of pairs produced in the jet, $`\gamma _{\mathrm{max}}`$. At early stages, $`\gamma _{\mathrm{max}}10^8`$ in Example 1 and $`\gamma _{\mathrm{max}}=10^6`$ in Example 2 (comoving values). The synchrotron radiation energy is 1000 MeV and 15 MeV (rest frame), respectively, while the observed synchrotron peak energy in blazars varies in the range $`10^60.1`$ MeV. A high $`\gamma _{\mathrm{max}}`$ leads to a many generation electromagnetic cascade which does produce a flat spectrum$`^\mathrm{?}`$.
The reason why $`\gamma _{\mathrm{max}}`$ is so high in our simulations is a short scale of the breeding cycle at the early stage: a high energy photon despite a large opacity can cross the jet boundary, if the latter is sharp. At later stages, the transition becomes smooth, the opacity increases due to generation on new soft photons and highest energy photons cannot cross a region with large $`\mathrm{\Gamma }`$ increment. Therefore $`\gamma _{\mathrm{max}}`$ drops and we can see this in Fig. 3. Moreover, in both cases we can see the formation of a poorly developed synchrotron component in a right place in late spectra. Probably, if we were able to follow up the system evolution for a longer time and with a more realistic treatment, we would obtain spectra in a better agreement with observations.
## ACKNOWLEDGMENTS
The work is supported by the RFBR grant 04-02-16987, Academy of Finland, Jenny and Antti Wihuri Foundation, Vilho, Yrjรถ and Kalle Vรคisรคlรค Foundation, and the NORDITA Nordic project in High Energy Astrophysics.
## References
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# Deformations of Fuchsian Systems of Linear Differential Equations and the Schlesinger System.
## NOTATION
* $``$ stands for the complex plane.
* $`^1`$ stands for the extended complex plane ($`=`$ the Riemann sphere):
$$^1=\mathrm{}.$$
* $`^n`$ stands for the $`n`$-dimensional complex space.
* In the coordinate notation, a point $`๐^n`$ will be written as $`๐=(t_1,\mathrm{},t_n).`$
* $`_{}^n`$ is the set of points $`๐^n,`$ whose coordinates $`t_1,\mathrm{},t_n`$ are pairwise different:
$$_{}^n=^n\underset{\begin{array}{c}1i,jn\\ ij\end{array}}{}\{๐:t_i=t_j\}.$$
* $`๐_k`$ stands for the set of all $`k\times k`$ matrices with complex entries.
* $`[,]`$ denotes the commutator: for $`A,B๐_k,[A,B]=ABBA`$.
* $`I`$ stands for the identity matrix of an appropriate dimension.
## 0. Introduction
The systematic study of linear differential equations in the complex plane with coefficients dependent on parameters has been started by Lazarus Fuchs in the late eighties of the nineteenth century<sup>1</sup><sup>1</sup>1L.Fuchs died in 1902.\[FuL1\], \[FuL2\], \[FuL3\], \[FuL4\]. In particular, L.Fuchs investigated the equations whose monodromy does not depend on such parameters. These investigations were continued in the beginning of the twentieth century by L.Schlesinger<sup>2</sup><sup>2</sup>2A student of L.Fuchs., \[Sch2\], \[Sch3\], \[Sch4\], R.Fuchs<sup>3</sup><sup>3</sup>3The son of L.Fuchs., \[FuR1\], \[FuR2\], and R. Garnier, \[Gar\].
L.Schlesingerโs research was closely related to the Hilbert 21st problem (a.k.a. the Riemann - Hilbert monodromy problem), which requires to construct a Fuchsian system with prescribed monodromy (for the explanation of terminology see Section 1 of the present article). In the paper \[Sch2\], which appeared exactly one hundred years ago โ in 1905, L.Schlesinger proposed the idea that it would be very fruitful to study the deformations of Fuchsian systems
(0.1)
$$\frac{dY}{dx}=\underset{1jn}{}\frac{Q_j(๐)}{xt_j}Y,$$
where the residues $`Q_j`$ depend holomorphically on the pole loci $`๐=(t_1,\mathrm{},t_n)`$, and investigate the dependence of the solution $`Y`$ on $`๐`$, as well as on $`x`$.
Emphasizing this idea, L.Schlesinger explained that he was guided by the analogy with the theory of algebraic functions, where he had studied algebraic functions as functions of both the โmain variableโ and the loci of ramification points considered as parameters (see \[Sch1, pp. 287 - 288\]).
Also in the paper \[Sch2\], the system of PDEs
(0.2)
$$\{\begin{array}{ccc}\hfill \frac{Q_j}{t_k}& =\frac{[Q_j,Q_k]}{t_jt_k},\hfill & \hfill 1j,kn,kj,\\ \hfill \frac{Q_j}{t_j}& =\underset{\begin{array}{c}1kn\\ kj\end{array}}{}\frac{[Q_j,Q_k]}{t_jt_k},\hfill & \hfill 1jn,\end{array}$$
which is now known as the Schlesinger system, was introduced and the statement that the holomorphic deformation (0.1) is isomonodromic if and only if its coefficients $`Q_j(๐)`$ satisfy the system (0.2) was formulated. (See page 294 of \[Sch2\], four bottom lines of this page.) This formulation was repeated in the book \[Sch3\], pp. 328-329, and later in the paper \[Sch4\], p. 106. (References to the earlier paper \[Sch2\] are relatively rare. Usually, one refers to the more recent paper \[Sch4\].)
Over the years the Schlesinger system and the isomonodromic deformations of Fuchsian systems were extensively studied; we would like to mention in particular the papers of T. Miwa \[Miwa\] and of B. Malgrange \[Ma1\], where it was proved that the Schlesinger system enjoys the Painlevรฉ property (its solutions are meromorphic functions in the universal covering space over $`_{}^n`$), and the book by A.R. Its and V.Yu. Novokshenov \[ItNo\], where the connections between the isomonodromic deformations and the transcendents of Painlevรฉ were revealed.
However, in the 1990s the famous negative solution to the Riemann - Hilbert monodromy problem due to A.A. Bolibrukh (see \[Bol1\], \[Bol2\]) gave a strong motivation for the revision of the classical results, concerning the isomonodromic deformations of the Fuchsian systems.
For example, it should be noted that in the original works \[Sch2\]\[Sch4\] of L. Schlesinger no assumptions concerning the non-resonance of the matrices $`Q_j`$ are made. In such generality the above-cited statement of L. Schlesinger fails. If a holomorphic deformation (0.1) of Fuchsian systems is such that the residues $`Q_j`$ satisfy the Schlesinger system (0.2), then this deformation is isomonodromic, but the converse statement is not true, in general.
It was also A.A. Bolibrukh who constructed the first explicit example of the non-Schlesinger isomonodromic deformation. In this example the monodromy is non-trivial and the residues $`Q_j(๐)`$ are rational functions of $`๐`$ (see \[Bol3\] and \[Bol4\] โ in both papers the same example appears as Example 2; see also <sup>4</sup><sup>4</sup>4Unfortunately, to our best knowledge this review has not yet been translated into English. Section 3 of the review paper \[Bol5\], where this example appears as Example 3).
At the same time it was shown independently in \[Kats1\] that almost every isomonodormic deformation of Fuchsian systems with generic rational solutions<sup>5</sup><sup>5</sup>5In particular, with trivial monodromy. is non-Schlesinger (for more details see Remark 5.1 in Section 5 of this article).
Thus the isomonodromic property of the deformation (0.1) implies the Schlesinger system for the residues $`Q_j(๐ญ)`$ under the non-resonance condition, but not in general.
Unfortunately, careless treatment of the non-resonance condition is very common in the history of problems related to the monodromy of Fuchsian systems. It can also be found in some works of V. Volterra and of G.D. Birkhoff (see \[Gant, Chapter XV, ยง9\] for details). This tradition continues to certain extent in the above-mentioned paper \[Miwa\] of T. Miwa on the Painlevรฉ property of isomonodromic deformations: in this paper the non-resonance condition appears as the equation (2.22), but is omitted both in the introduction and in the formulation of the main result.
Without the assumption of non-resonance the main result of \[Miwa\] does not hold: there exist isomonodromic deformations of the form (0.1), where the residues $`Q_j(๐)`$ are not meromorphic in the universal covering space over $`_{}^n`$. (The appropriate example is presented in Section 5 of this article. We note that this phenomenon does not occur in the above-mentioned example of the non-Schlesinger isomonodromic deformation due to A.A. Bolibrukh: in that example the residues $`Q_j(๐)`$ are rational functions of $`๐`$.)
The main goal of the present work is to answer the following question: how to describe the class of holomorphic deformations (0.1) with the property that the residues $`Q_j(๐ญ)`$ satisfy the Schlesinger system (0.2), when one omits the non-resonance assumption?
The presentation of our results is organized as follows.
In the first section after this introduction we recall the basic notions concerning the Fuchsian system and introduce a certain canonical multiplicative decomposition of the fundamental solution in a neighborhood of its singular point $`t_j`$. This is the so-called regular-principal factorization: the fundamental solution is represented as the product of a regular factor, holomorphic and invertible at the point $`t_j,`$ and a principal factor, holomorphic (multi-valued) and invertible everywhere except at $`t_j.`$ This principal factor is the multiplicative analogue of the principal part in the Laurent decomposition: it contains the information about the nature of the singularity.
In Section 2 we introduce the main notion of the present article (Definition 2.9): the so-called isoprincipal<sup>6</sup><sup>6</sup>6Iso- (from $`\stackrel{..}{\iota }`$$`\sigma o\varsigma `$ \- equal - in Old Greek) is a combining form. families of Fuchsian systems. These are the holomorphic families (0.1) with the property that all the principal factors of a suitably normalized fundamental solution $`Y(x,๐)`$ are, in a certain sense, preserved. We show that every isoprincipal family is also isomonodromic and that the converse is true, when the non-resonance condition is in force.
In Section 3 we formulate and prove our main result (Theorem 3.1) that the family (0.1) is isoprincipal if and only if the residues $`Q_j(๐ญ)`$ satisfy the Schlesinger system (0.2). This result holds in the general case, without the assumption of non-resonance.
In the next section we discuss the isoprincipal deformation of a given Fuchsian system. Using our Theorem 3.1, we also outline how to establish the Painlevรฉ property of the Schlesinger system and indicate possible generalizations.
Finally, in Section 5 we illustrate the general theory with explicit examples of the isoprincipal and isomonodromic deformations. In particular, we give an example of the isomonodromic deformation which is not related to the Schlesinger system and does not possess the Panlevรฉ property. This example is based on the theory of the isoprincipal families of Fuchsian systems with generic rational solutions, developed in \[Kats1\], \[Kats2\] and \[KaVo2\].
## 1. Fuchsian differential systems
### 1.1. Fuchsian differential systems
A Fuchsian differential system is a linear system of ordinary differential equations of the form
(1.1)
$$\frac{dY}{dx}=\left(\underset{1jn}{}\frac{Q_j}{xt_j}\right)Y,$$
where $`Q_j`$, $`1jn,`$ are square matrices of the same dimension, say $`Q_j๐_k`$, and $`t_1,\mathrm{},t_n`$ are pairwise distinct points of the complex plane $``$. The variable $`x`$ โlivesโ in the punctured Riemann sphere $`^1\{t_1,\mathrm{},t_n\}`$, the โunknownโ $`Y`$ is an $`๐_k`$-valued matrix function of $`x`$. Under the condition
(1.2)
$$\underset{1jn}{}Q_j=0$$
the point $`x_0=\mathrm{}`$ is a regular point for the system (1.1). If this condition is satisfied (which we always assume in the sequel), then in a neighborhood of the point $`x_0=\mathrm{}`$ there exists a fundamental solution $`Y=Y(x)`$ of (1.1) satisfying the initial condition
(1.3)
$$Y(x)|_{x=\mathrm{}}=I.$$
This solution $`Y`$ can be analytically continued into the multi-connected domain $`^1\{t_1,\mathrm{},t_n\}`$. However, for $`x^1\{t_1,\mathrm{},t_n\}`$ the value of $`Y`$ at the point $`x`$ depends, in general, on the path $`\alpha `$ from $`x_0=\mathrm{}`$ to $`x`$, along which the analytic continuation is performed:
$$Y=Y(x,\alpha ).$$
More precisely, $`Y`$ depends not on the path $`\alpha `$ itself, but on its homotopy class in $`^1\{t_1,\mathrm{},t_n\}`$. Thus $`Y`$ is a multi-valued holomorphic function in the punctured Riemann sphere $`^1\{t_1,\mathrm{},t_n\}`$ or, better to say, $`Y`$ is a singled-valued holomorphic function on the universal covering surface of the punctured Riemann sphere $`^1\{t_1,\mathrm{},t_n\}`$ with the distinguished point $`\mathrm{}`$.
### 1.2. Universal covering spaces
Recall (see \[Fo, Chapter 1, Sections 3 - 5\] if need be) that the universal covering space $`\mathrm{cov}(๐ณ;x_0)`$ of an arcwise connected topological space $`๐ณ`$ with the distinguished point $`x_0๐ณ`$ is the set of pairs $`(x,\alpha ),`$ where $`x`$ is a point in $`๐ณ`$ and $`\alpha `$ is a homotopy class of continuous mappings
$$\alpha :\{s:\mathrm{\hspace{0.17em}0}s1\}๐ณ,\alpha (0)=x_0,\alpha (1)=x.$$
Such a mapping $`\alpha `$ is called a path in $`๐ณ`$ from $`x_0`$ to $`x.`$ A path in $`๐ณ`$ from $`x_0`$ to $`x_0`$ is called a loop with the distinguished point $`x_0`$.
The product $`\beta \alpha `$ of two paths $`\alpha ,\beta `$ in $`๐ณ`$, where $`\alpha `$ is a path from $`a`$ to $`b`$ and $`\beta `$ is a path from $`b`$ to $`c`$, is defined as the path from $`a`$ to $`c`$, obtained by going first along $`\alpha `$ from $`a`$ to $`b`$ and then along $`\beta `$ from $`b`$ to $`c`$:
$$(\beta \alpha )(s)=\{\begin{array}{cc}\hfill \alpha (2s),& 0s\frac{1}{2},\hfill \\ \hfill \beta (2s1),& \frac{1}{2}s1.\hfill \end{array}$$
With respect to this product, the homotopy classes of loops in $`๐ณ`$ with the distinguished point $`x_0`$ form the so-called fundamental group $`\pi (๐ณ;x_0)`$ of the space $`๐ณ`$ with the distinguished point $`x_0`$. The fundamental group $`\pi (๐ณ;x_0)`$ acts on the universal covering space $`\mathrm{cov}(๐ณ;x_0)`$ (on the right) as the group of deck transformations
(1.4)
$$(x,\alpha )(x,\alpha \gamma ),(x,\alpha )\mathrm{cov}(๐ณ;x_0),\gamma \pi (๐ณ;x_0).$$
### 1.3. Monodromy
Let $`Y=Y(x,\alpha )`$ be the solution of (1.1) โ (1.3) defined on the universal covering surface $`\mathrm{cov}(^1\{t_1,\mathrm{},t_n\};\mathrm{}).`$ For each loop $`\gamma \pi (^1\{t_1,\mathrm{},t_n\};\mathrm{})`$ let us consider the function $`Y_\gamma `$ defined on $`\mathrm{cov}(^1\{t_1,\mathrm{},t_n\};\mathrm{})`$ by
(1.5)
$$Y_\gamma (x,\alpha )\stackrel{\text{def}}{=}Y(x,\alpha \gamma ).$$
The expression (1.5) means that the value of $`Y_\gamma `$ at the point $`(x,\alpha )`$ of the universal covering surface $`\mathrm{cov}(^1\{t_1,\mathrm{},t_n\};\mathrm{})`$ is obtained by the analytic continuation of the solution $`Y`$ of (1.1) โ (1.3): first along the loop $`\gamma `$ from the distinguished point $`x_0=\mathrm{}`$ to itself, then along the path $`\alpha `$ from $`x_0`$ to $`x`$.
Thus $`Y_\gamma `$ is also a fundamental solution of the linear system (1.1) and, therefore, there exists a unique invertible constant matrix $`M_\gamma ๐_k`$ such that
(1.6)
$$Y_\gamma (x,\alpha )Y(x,\alpha )M_\gamma ,(x,\alpha )\mathrm{cov}(^1\{t_1,\mathrm{},t_n\};\mathrm{}).$$
###### Definition 1.1.
Let $`Y=Y(x,\alpha )`$ be the solution of (1.1) โ (1.3) on the universal covering surface $`\mathrm{cov}(^1\{t_1,\mathrm{},t_n\};\mathrm{})`$ and let $`\gamma \pi (^1\{t_1,\mathrm{},t_n\};\mathrm{})`$.
The constant (with respect to $`x`$) matrix $`M_\gamma ๐_k`$, which appears in the identity (1.6), is said to be the monodromy matrix of the solution $`Y`$, corresponding to the loop $`\gamma `$.
Note that for a pair of loops $`\gamma _1,\gamma _2\pi (^1\{t_1,\mathrm{},t_n\})`$ and the cooresponding monodromy matrices $`M_{\gamma _1},M_{\gamma _2}`$ of the solution $`Y`$ it holds that
$$Y(x,\alpha \gamma _1\gamma _2)=\left(YM_{\gamma _2}\right)(x,\alpha \gamma _1)=Y(x,\alpha \gamma _1)M_{\gamma _2}=Y(x,\alpha )M_{\gamma _1}M_{\gamma _2}.$$
Therefore, the monodromy matrices of $`Y`$ satisfy the following multiplicative identity:
(1.7)
$$M_{\gamma _1\gamma _2}=M_{\gamma _1}M_{\gamma _2}\gamma _1,\gamma _2\pi (^1\{t_1,\mathrm{},t_n\}).$$
This means that the mapping $`\gamma M_\gamma `$ is a linear representation of the fundamental group $`\pi (^1\{t_1,\mathrm{},t_n\};\mathrm{}).`$
###### Definition 1.2.
Let $`Y`$ be the solution of (1.1) โ (1.3) on the universal covering surface $`\mathrm{cov}(^1\{t_1,\mathrm{},t_n\};\mathrm{}).`$ The linear representation of the fundamental group $`\pi (^1\{t_1,\mathrm{},t_n\};\mathrm{})`$
(1.8)
$$\gamma M_\gamma ,\gamma \pi (^1\{t_1,\mathrm{},t_n\};\mathrm{}),$$
where $`M_\gamma `$ denotes the monodromy matrix of the solution $`Y`$, corresponding to the loop $`\gamma `$, is called the monodromy representation of the solution $`Y`$.
### 1.4. The regular-principal factorization for a fundamental solution of a Fuchsian system: single-valued case
Each of the points $`t_j,`$ $`1jn,`$ is a singularity of the solution $`Y`$. This means that at least one of the two functions $`Y`$ and $`Y^1`$ is not holomorphic at $`t_j`$. More information about the nature of the singularity at $`t_j`$ can be obtained from a certain multiplicative decomposition of the solution $`Y`$ near the point $`t_j`$, which is called the regular-principal factorization.
In order to explain the idea of the regular-principal factorization, let us assume for the moment that the solution $`Y=Y(x)`$ is single-valued in the domain $`^1\{t_1,\mathrm{},t_n\}`$ โ that is, the monodromy representation of $`Y`$ is trivial:
$$M_\gamma =I\gamma \pi (^1\{t_1,\mathrm{},t_n\};\mathrm{}).$$
For $`1jn`$ let $`๐ฑ_j`$ be an open simply connected neighborhood of $`t_j`$ in $``$, such that $`t_k๐ฑ_j`$ for $`kj`$. Then it follows, for instance, from G.D. Birkhoffโs results on factorization of matrix functions holomorphic in the annulus (see \[Birk1, ยง7\]) that in the punctured neighborhood $`๐ฑ_j\{t_j\}`$ the solution $`Y(x)`$ admits a factorization of the form
$$Y(x)=H_j(x)P_j(x),x๐ฑ_j\{t_j\},$$
where the function $`H_j(x)`$ is holomorphic and invertible in the entire (non-punctured) neighborhood $`๐ฑ_j`$ and the function $`P_j(x)`$ is holomorphic, single-valued and invertible in the punctured plane $`\{t_j\}`$. The factors $`H_j(x)`$ and $`P_j(x)`$ are said to be, respectively, the regular factor and the principal factor of the solution $`Y`$ at its singular point $`t_j`$.
In the general case, when the monodromy representation of $`Y`$ may be non-trivial, the regular-principal factorization of $`Y`$ is more involved.
Indeed, on the one hand the solution $`Y`$ is normalized at the distinguished point $`x_0=\mathrm{}`$. In order to consider $`Y`$ in a neighborhood of $`t_j`$, we have to choose a homotopy class of paths, connecting the distinguished point $`x_0=\mathrm{}`$ with this neighborhood of $`t_j`$, and such a choice is not unique.
On the other hand, even in the single-valued case $`x_0=\mathrm{}`$ is, in general, a singular point of the principal factor $`P_j`$. In the general case $`P_j`$ will have to be considered as a function on a universal covering surface of the punctured plane $`\{t_j\}`$. Therefore, $`P_j`$ needs to be normalized at some distinguished point in $`\{t_j\}`$ (obviously different from $`x_0=\mathrm{}`$) and continued analytically from there.
Thus, before we can present the regular-principal factorization of $`Y`$ in the general case, we need some preparation.
### 1.5. Branches of the solution of a Fuchsian system in a neighborhood of the singular point
We propose the following terminology:
###### Definition 1.3.
Let $`Y`$ be the solution of (1.1) โ (1.3) on the universal covering surface $`\mathrm{cov}(^1\{t_1,\mathrm{},t_n\};\mathrm{}).`$
For $`1jn`$ assume that:
1. $`๐ฑ`$ is a domain in $`^1\{t_1,\mathrm{},t_n\}`$;
2. $`p`$ is a point in the domain $`๐ฑ`$;
3. $`\alpha `$ is a path in $`^1\{t_1,\mathrm{},t_n\}`$ from the distinguished point $`x_0=\mathrm{}`$ to the point $`p`$.
Define the function $`Y_\alpha `$ on the universal covering surface $`\mathrm{cov}(๐ฑ;p)`$ by the analytic continuation of the solution $`Y`$ first along the path $`\alpha `$, then inside the domain $`๐ฑ`$:
(1.9)
$$Y_\alpha (x,\beta )\stackrel{\text{def}}{=}Y(x,\beta \alpha ),$$
where $`x๐ฑ`$ and $`\beta `$ is a path in $`๐ฑ`$ from $`p`$ to $`x`$.
Then the function $`Y_\alpha `$, holomorphic in $`\mathrm{cov}(๐ฑ;p)`$, is said to be the branch of the solution$`Y`$ in the domain $`๐ฑ`$, corresponding to the path $`\alpha `$.
In what follows we shall be mostly dealing with the branches of the solution $`Y`$ in domains of the form $`๐ฑ_j\{t_j\}`$, where $`๐ฑ_j`$ is a simply connected neighborhood of the singular point $`t_j`$, such that $`t_k๐ฑ_j`$ for $`kj`$. In this case the universal covering surface $`\mathrm{cov}(๐ฑ_j\{t_j\};p_j)`$ has a simple structure: the fundamental group $`\mathrm{cov}(๐ฑ_j\{t_j\};p_j)`$ is cyclic, generated by the loop with the distinguished point $`p_j`$ which makes one positive circuit of $`t_j`$ in $`๐ฑ_j\{t_j\}`$.
###### Definition 1.4.
For $`1jn`$ assume that:
1. $`๐ฑ_j`$ is an open simply connected neighborhood of the singular point $`t_j`$, such that $`t_k๐ฑ_j`$ for $`kj`$;
2. $`p_j`$ is a point in the punctured neighborhood $`๐ฑ_j\{t_j\}`$;
3. $`\alpha _j`$ is a path in $`^1\{t_1,\mathrm{},t_n\}`$ from the distinguished point $`x_0=\mathrm{}`$ to the point $`p_j`$.
Let $`\beta _j`$ be the loop in the punctured neighborhood $`๐ฑ_j\{t_j\}`$ with the distinguished point $`p_j`$ which makes one positive circuit of $`t_j`$, and let $`\gamma _j`$ be the loop in the punctured sphere $`^1\{t_1,\mathrm{},t_n\}`$ with the distinguished point $`x_0=\mathrm{}`$, defined by
(1.10)
$$\gamma _j\stackrel{\text{def}}{=}\alpha _j^1\beta _j\alpha _j$$
(that is, the loop $`\gamma _j`$ goes from $`x_0=\mathrm{}`$ to $`p_j`$ along $`\alpha _j`$, then makes one positive circuit of $`t_j`$ along the small loop $`\beta _j`$, then goes again along $`\alpha _j`$, but in the opposite direction: from $`p_j`$ to $`x_0=\mathrm{}`$).
Then:
* the loop $`\beta _j`$ is said to be the small loop around $`t_j`$ in the punctured neighborhood $`๐ฑ_j\{t_j\}`$;
* the loop $`\gamma _j`$ is said to be the big loop around $`t_j`$, corresponding to the path $`\alpha _j`$.
###### Remark 1.5.
Note that for a suitable choice of the paths $`\alpha _1,\mathrm{},\alpha _n`$ the corresponding big loops $`\gamma _1,\mathrm{},\gamma _n`$ generate the fundamental group $`\pi (^1\{t_1,\mathrm{},t_n\};\mathrm{})`$. These generators are not free: choosing $`\alpha _1,\mathrm{},\alpha _n`$ carefully we can ensure, for example, that $`\gamma _1\mathrm{}\gamma _n=1.`$
We observe that the surface $`\mathrm{cov}(๐ฑ_j\{t_j\};p_j)`$ in Definition 1.3 is naturally embedded into the universal covering surface $`\mathrm{cov}(\{t_j\};p_j)`$, which is isomorphic to the Riemann surface of the logarithm $`\mathrm{ln}\zeta `$.
Although the basic properties of the function $`\mathrm{ln}\zeta `$ are very well-known, we shall discuss them in some detail, because they are important for our future considerations.
### 1.6. The Riemann surface of $`\mathrm{ln}\zeta `$
For each fixed $`\zeta \{0\}`$ the equation
$$e^\lambda =\zeta $$
has a countable set of solutions $`\lambda =\lambda (\zeta )\stackrel{\text{def}}{=}\mathrm{ln}\zeta `$. These solutions can be parameterized as
(1.11)
$$\mathrm{ln}\zeta =\mathrm{ln}|\zeta |+i\mathrm{arg}\zeta ,\text{ where }\mathrm{ln}|\zeta |$$
and $`\mathrm{arg}\zeta `$ is an equivalence class of real numbers (the values of $`\mathrm{arg}\zeta `$) modulo addition by $`2\pi `$.
Now we explain how the function $`\mathrm{ln}\zeta `$ can be defined as a single-valued holomorphic function on the universal covering surface $`\mathrm{cov}(\{0\};1)`$ of the punctured plane $`\{0\}`$ with the distinguished point $`\zeta _0=1`$.
Let us choose a point $`\zeta \{0\}`$ and a path $`\vartheta `$ in $`\{0\}`$ from $`\zeta _0=1`$ to $`\zeta `$. Then there exists a unique $`\theta `$ such that, up to homotopy in $`\{0\}`$, the path $`\vartheta `$ can be parameterized as follows:
(1.12)
$$\vartheta (s)=e^{(\mathrm{ln}|\zeta |+i\theta )s},0s1$$
(here $`\mathrm{ln}|\zeta |`$ is the real-valued logarithm). The real number $`\theta `$ is said to be the value of $`\mathrm{arg}\zeta `$ corresponding to the path $`\vartheta `$. In this manner we establish a $`1`$-to-$`1`$ correspondence between the values of $`\mathrm{arg}\zeta `$ and the homotopy classes of paths in $`\{0\}`$ from $`\zeta _0=1`$ to $`\zeta `$.
Thus the function $`\mathrm{arg}\zeta `$ is defined as a single-valued continuous function on $`\mathrm{cov}(\{0\};1)`$. Accordingly, the function $`\mathrm{ln}\zeta `$ is defined by (1.11) as a single-valued holomorphic function on the universal covering surface $`\mathrm{cov}(\{0\};1)`$. In the sequel we shall often refer to the surface $`\mathrm{cov}(\{0\};1)`$ as the Riemann surface of $`\mathrm{ln}\zeta `$.
The fundamental group $`\pi (\{0\};1)`$ is cyclic, generated by the loop with the distinguished point $`\zeta _0=1`$ which makes one turn counterclockwise around the origin. The corresponding deck transformation of $`\mathrm{cov}(\{0\};1)`$ is denoted by
(1.13)
$$\zeta \zeta e^{2\pi i},$$
so that the following monodromy relations hold:
(1.14)
$$\mathrm{arg}(\zeta e^{2\pi i})=\mathrm{arg}\zeta +2\pi ,\mathrm{ln}(\zeta e^{2\pi i})=\mathrm{ln}\zeta +2\pi i.$$
### 1.7. Transplants of functions defined on the Riemann surface of $`\mathrm{ln}\zeta `$
Let us consider the universal covering surface $`\mathrm{cov}(\{t_j\};p_j)`$, where $`t_j`$ is some point in the complex plane $``$ and $`p_j`$ is a distinguished point in the punctured plane $`\{t_j\}`$. Let us choose some value $`\theta _j`$ of $`\mathrm{arg}(p_jt_j)`$ and let $`\vartheta _j`$ denote the corresponding path (see (1.12)) in $`\{0\}`$ from $`\zeta _0=1`$ to $`\zeta =p_jt_j`$. Let us define a mapping from the universal covering surface $`\mathrm{cov}(\{t_j\};p_j)`$ into the Riemann surface of $`\mathrm{ln}\zeta `$ as follows.
To each point $`(x,\alpha )\mathrm{cov}(\{t_j\};p_j),`$ where $`\alpha `$ is a path in $`\{t_j\}`$ from $`p_j`$ to $`x`$, we associate the point $`(xt_j,\alpha _{t_j}\vartheta _j)\mathrm{cov}(\{0\};1)`$, where the path $`\alpha _{t_j}`$, which leads in $`\{0\}`$ from $`p_jt_j`$ to $`xt_j`$, is obtained by the parallel translation of the path $`\alpha `$:
$$\alpha _{t_j}(s)=\alpha (s)t_j,0s1.$$
This mapping is an isomorphism between $`\mathrm{cov}(\{t_j\};p_j)`$ and the Riemann surface of $`\mathrm{ln}\zeta `$. It will be denoted by
(1.15)
$$x\underset{\mathrm{arg}\left(p_jt_j\right)=\theta _j}{\overset{}{}}xt_j,x\mathrm{cov}(\{t_j\};p_j),xt_j\mathrm{cov}(\{0\};1),$$
and the inverse mapping will be denoted by
(1.16)
$$x\underset{\mathrm{arg}\left(p_jt_j\right)=\theta _j}{\overset{}{}}x+t_j,x\mathrm{cov}(\{0\};1),x+t_j\mathrm{cov}(\{t_j\};p_j).$$
Accordingly, we shall denote the deck transformation of $`\mathrm{cov}(\{t_j\};p_j)`$, corresponding to the loop with the distinguished point $`p_j`$ which makes one positive circuit around $`t_j`$ in $`\{t_j\}`$, by
(1.17)
$$xt_j+(xt_j)e^{2\pi i},x\mathrm{cov}(\{t_j\};p_j).$$
With this notation we can consider the function $`\mathrm{ln}(xt_j)`$ as a function of $`x`$, holomorphic on the universal covering surface $`\mathrm{cov}(\{t_j\};p_j)`$ and such that (see (1.14))
(1.18)
$$\mathrm{ln}\left((xt_j)e^{2\pi i}\right)=\mathrm{ln}(xt_j)+2\pi i.$$
More generally, we propose the following
###### Definition 1.6.
Let a function $`E(\zeta )`$ be defined on the Riemann surface of $`\mathrm{ln}\zeta `$. Let $`t_j`$, let $`p_j\{t_j\}`$ and let us choose a value $`\theta _j`$ of $`\mathrm{arg}(p_jt_j)`$. For each $`x\mathrm{cov}(\{t_j\};p_j)`$ let $`xt_j`$ denote the image of $`x`$ in the Riemann surface of $`\mathrm{ln}\zeta `$ under the isomorphism (1.15).
The function $`E(xt_j)`$, defined as a function of $`x`$ on the universal covering surface $`\mathrm{cov}(\{t_j\};p_j)`$, is said to be the transplant of the function $`E(\zeta )`$ into $`\mathrm{cov}(\{t_j\};p_j)`$, corresponding to the value $`\theta _j`$ of $`\mathrm{arg}(p_jt_j)`$.
### 1.8. The regular-principal factorization for a fundamental solution of a Fuchsian system: general case
###### Theorem 1.7.
Let $`Y`$ be the solution of the Fuchsian system (1.1) โ (1.2), satisfying the initial condition (1.3). For $`j=1,\mathrm{},n`$ assume that:
1. $`๐ฑ_j`$ is an open simply connected neighborhood of the singular point $`t_j`$, such that $`t_k๐ฑ_j`$ for $`kj`$;
2. $`p_j`$ is a point in the punctured neighborhood $`๐ฑ_j\{t_j\}`$;
3. $`\alpha _j`$ is a path in $`^1\{t_1,\mathrm{},t_n\}`$ from the distinguished point $`x_0=\mathrm{}`$ to the point $`p_j`$.
Then for each $`j`$, $`1jn,`$ the branch $`Y_{\alpha _j}`$ of the solution $`Y`$ in the punctured neighborhood $`๐ฑ_j\{t_j\}`$, corresponding to the path $`\alpha _j`$, admits the factorization
(1.19)
$$Y_{\alpha _j}(x)=H_{\alpha _j}(x)P_{\alpha _j}(x),x\mathrm{cov}(๐ฑ_j\{t_j\};p_j),$$
where the factors possess the following properties:
* the matrix function $`H_{\alpha _j}(x)`$ is holomorphic and invertible in the entire <sup>7</sup><sup>7</sup>7In particular, the function $`H_{\alpha _j}(x)`$ is single-valued in $`๐ฑ_j`$. (non-punctured) neighborhood $`๐ฑ_j`$;
* the matrix function $`P_{\alpha _j}(x)`$ is holomorphic and invertible on the universal covering surface $`\mathrm{cov}(\{t_j\};p_j)`$.
###### Definition 1.8.
The factorization (1.19), where the factors $`H_{\alpha _j}`$ and $`P_{\alpha _j}`$ possess the properties (R) and (P), is said to be the regular-principal factorization of the branch $`Y_{\alpha _j}`$ of the solution $`Y`$ in a punctured neighborhood of the singular point $`t_j.`$ The factors $`H_{\alpha _j}`$ and $`P_{\alpha _j}`$ are said to be, respectively, the regular factor and the principal factor of the branch $`Y_{\alpha _j}`$.
###### Proof of Theorem 1.7.
For $`j=1,\mathrm{},n`$ let $`\gamma _j`$ be the big loop around $`t_j`$, corresponding to the path $`\alpha _j`$ (see Definition 1.4), and let $`M_{\gamma _j}`$ be the corresponding monodromy matrix of $`Y`$.
Then, according to Definitions 1.1 and 1.3, the monodromy matrix $`M_{\gamma _j}`$ is given by
(1.20)
$$M_{\gamma _j}=Y_{\alpha _j}^1(x)Y_{\alpha _j}(t_j+(xt_j)e^{2\pi i}),x\mathrm{cov}(๐ฑ_j\{t_j\};p_j),$$
where $`Y_{\alpha _j}`$ is the branch of $`Y`$ in $`๐ฑ_j\{t_j\}`$, corresponding to the path $`\alpha _j`$.
Since the matrix $`M_{\gamma _j}`$ is invertible, there exists a matrix denoted by $`\mathrm{ln}M_{\gamma _j}`$, such that <sup>8</sup><sup>8</sup>8Here we refer to \[Gant, Chapter VIII, Section 8\]. Such a matrix $`\mathrm{ln}M_{\gamma _j}`$ is not unique, of course, but for our purposes any choice of $`\mathrm{ln}M_{\gamma _j}`$ will do.
$$e^{\mathrm{ln}M_{\gamma _j}}=M_{\gamma _j}.$$
Let us choose a transplant $`\mathrm{ln}(xt_j)`$ of the function $`\mathrm{ln}\zeta `$ into $`\mathrm{cov}(\{t_j\};p_j)`$. Then, according to (1.18), the matrix function
$$(xt_j)^{{\displaystyle \frac{1}{2\pi i}}\mathrm{ln}M_{\gamma _j}}\stackrel{\text{def}}{=}e^{\mathrm{ln}(xt_j){\displaystyle \frac{1}{2\pi i}}\mathrm{ln}M_{\gamma _j}},$$
which is holomorphic and invertible on $`\mathrm{cov}(\{t_j\};p_j)`$, satisfies the relation
$$((xt_j)e^{2\pi i})^{{\displaystyle \frac{1}{2\pi i}}\mathrm{ln}M_{\gamma _j}}=(xt_j)^{{\displaystyle \frac{1}{2\pi i}}\mathrm{ln}M_{\gamma _j}}M_{\gamma _j},x\mathrm{cov}(\{t_j\};p_j).$$
Hence, in view of (1.20), the branch $`Y_{\alpha _j}`$ of the solution $`Y`$ in $`๐ฑ_j\{t_j\}`$ has the form
$$Y_{\alpha _j}(x)=\mathrm{\Phi }(x)(xt_j)^{{\displaystyle \frac{1}{2\pi i}}\mathrm{ln}M_{\gamma _j}},x\mathrm{cov}(\{t_j\};p_j),$$
where $`\mathrm{\Phi }(x)`$ is a matrix function, holomorphic, invertible and single-valued in the punctured neighborhood $`๐ฑ_j\{t_j\}`$.
Now, according to \[Birk1, ยง7\], we can factorize the function $`\mathrm{\Phi }(x)`$ as
$$\mathrm{\Phi }(x)=\mathrm{\Phi }_+(x)\mathrm{\Phi }_{}(x),x๐ฑ_j\{t_j\}$$
where $`\mathrm{\Phi }_+(x)`$ and $`\mathrm{\Phi }_{}(x)`$ are matrix functions, single-valued, holomorphic and invertible in, respectively, the entire (non-punctured) neighborhood $`๐ฑ_j`$ and the punctured plane $`\{t_j\}`$. We set
$`H_{\alpha _j}(x)`$ $`\stackrel{\text{def}}{=}\mathrm{\Phi }_+(x),`$
$`P_{\alpha _j}(x)`$ $`\stackrel{\text{def}}{=}\mathrm{\Phi }_{}(x)(xt_j)^{{\displaystyle \frac{1}{2\pi i}}\mathrm{ln}M_{\gamma _j}}`$
and obtain the desired factorization (1.19) with the properties (R) and (P). This completes the proof. โ
###### Remark 1.9.
Of course, the principal and regular factors of the branch $`Y_{\alpha _j}`$ of the solution $`Y`$ in a punctured neighborhood of the singularity $`t_j`$ are determined only up to the transformation
(1.21)
$$P_{\alpha _j}(x)T(x)P_{\alpha _j}(x),H_{\alpha _j}(x)T(x)^1H_{\alpha _j}(x),$$
where $`T(x)`$ is an invertible entire matrix function. However, once the choice of, say, the regular factor $`H_{\alpha _j}`$ is fixed, the principal factor $`P_{\alpha _j}`$ is uniquely determined.
Moreover, if we choose a different path from $`x_0`$ to $`p_j`$, say $`\alpha _j^{}`$, then, in view of Definitions 1.1 and 1.3, the branches $`Y_{\alpha _j}`$ and $`Y_{\alpha _j^{}}`$ of the solution $`Y`$ are related by
$$Y_{\alpha _j^{}}(x)=Y_{\alpha _j}(x)M_{\alpha _j^1\alpha _j^{}},x\mathrm{cov}(\{t_j\};p_j),$$
where $`M_{\alpha _j^1\alpha _j^{}}`$ is the monodromy matrix of $`Y,`$ corresponding to the loop $`\alpha _j^1\alpha _j^{}\pi (^1\{t_1,\mathrm{},t_n\};\mathrm{})`$. Hence the branch $`Y_{\alpha _j^{}}`$ admits the regular-principal factorization
$$Y_{\alpha _j^{}}(x)=H_{\alpha _j^{}}(x)P_{\alpha _j^{}}(x),x\mathrm{cov}(\{t_j\};p_j),$$
where
$`H_{\alpha _j}(x)`$ $`=H_{\alpha _j^{}}(x),`$ $`x`$ $`๐ฑ_j,`$
$`P_{\alpha _j^{}}(x)`$ $`=P_{\alpha _j}(x)M_{\alpha _j^1\alpha _j^{}},`$ $`x`$ $`\mathrm{cov}(\{t_j\};p_j).`$
Thus the regular factor at the singular point $`t_j`$ can be chosen independently of the choice of the path $`\alpha _j`$ and will be denoted simply by $`H_j(x)`$.
###### Remark 1.10.
Up to now, we have made no use of the fact that the system (1.1) is Fuchsian (that is, each singularity is a simple pole for the coefficients of the system). In particular, the regular - principal factorization (1.19) of the fundamental solution of a linear differential system also takes place in a neighborhood of the isolated singularity, where the coefficients of the system have a higher order pole or even an essential singular point.
However, in the special case when the system is Fuchsian, the general form of the fundamental solution in a neighborhood of its singular point is quite well known (see, for instance, \[Gant, Chapter XV, ยง10\]) and thus much more precise statements concerning the principal factors of the solution of the system (1.1) can be made.
If the matrix $`Q_j`$ is non-resonant <sup>9</sup><sup>9</sup>9 A square matrix $`Q`$ is said to be non-resonant if distinct eigenvalues of $`Q`$ do not differ by integers or, in other words, if the spectra of the matrices $`Q+nI`$ and $`Q`$ are disjoint for every $`n0`$., then the principal factor $`P_{\alpha _j}`$ can be chosen in the form
(1.22)
$$P_{\alpha _j}(x)=(xt_j)^{A_{\alpha _j}},$$
where $`A_{\alpha _j}`$ is a matrix, similar to the matrix $`Q_j`$:
(1.23)
$$A_{\alpha _j}=C_{\alpha _j}^1Q_jC_{\alpha _j},$$
where $`C_{\alpha _j}`$ is an invertible matrix. In the general case (without the assumption that the matrix $`Q_j`$ is non-resonant) the principal factor $`P_{\alpha _j}`$ can be chosen in the form
(1.24)
$$P_{\alpha _j}(x)=(xt_j)^{Z_{\alpha _j}}(xt_j)^{A_{\alpha _j}},$$
where $`Z_{\alpha _j}`$ is a diagonalizable matrix with integer eigenvalues $`l_1,\mathrm{},l_k`$ and $`A_{\alpha _j}`$ is a non-resonant matrix, whose eigenvalues $`\widehat{\lambda }_p`$ are related to the eigenvalues $`\lambda _p`$ of the matrix $`Q_j`$ by the equations
(1.25)
$$\widehat{\lambda }_p=\lambda _pl_p,1pk.$$
The matrices $`Z_{\alpha _j},A_{\alpha _j}`$ also possess certain additional properties, but we shall not go into further details, because in the sequel our considerations will be mostly based on the existence of the regular - principal factorization (1.19) described in Theorem 1.7 rather than on the specific form of the factors.
## 2. Holomorphic Families of Fuchsian Differential Systems
### 2.1. Families of Fuchsian systems, parameterized by the pole loci
In the present paper we consider a family of linear differential systems of the form
(2.1)
$$\frac{dY}{dx}=\left(\underset{1jn}{}\frac{Q_j(๐)}{xt_j}\right)Y.$$
The variable $`x`$ โlivesโ in the punctured Riemann sphere $`^1\{t_1,\mathrm{},t_n\}`$, where $`t_1,\mathrm{},t_n`$ are pairwise distinct points of the complex plane $``$. However, now $`t_1,\mathrm{},t_n`$ are not fixed but serve as the parameters of the family. The string $`๐=(t_1,\mathrm{},t_n)`$ is considered as a point of $`_{}^n`$ and the residue matrices $`Q_j`$ are assumed to depend on the the parameter $`๐`$. The โunknownโ $`Y`$ is a square matrix function depending both on the โmain variableโ $`x`$ and on the parameter $`๐`$.
###### Definition 2.1.
Assume that the matrix functions $`Q_j(๐),`$ $`1jn,`$ are defined and holomorphic for $`๐๐`$, where $`๐`$ is a domain in $`_{}^n.`$ Then the family (2.1) is said to be a holomorphic family of Fuchsian systems, parameterized by the pole loci.
In what follows we assume that the condition
(2.2)
$$\underset{1jn}{}Q_j(๐)0,๐๐,$$
holds. Thus for every fixed $`๐๐`$ the point $`x=\mathrm{}`$ is a regular point for the system (2.1), considered as a differential system with respect to $`x`$.
For each fixed $`๐๐`$, the fundamental solution $`Y=Y(x,๐)`$ of the differential system (2.1) with the initial condition<sup>10</sup><sup>10</sup>10In view of (2.2), for each fixed $`๐๐`$ the point $`x=\mathrm{}`$ is regular for the system (2.1) and hence the initial condition (2.3) can be posed.
(2.3)
$$Y(x,๐)|_{x=\mathrm{}}=I$$
is defined and holomorphic as a function of $`x`$ on the universal covering surface $`\mathrm{cov}(^1\{t_1,\mathrm{},t_n\};\mathrm{})`$.
In the present section our goal is to compare the properties of $`Y(x,๐)`$, such as the monodromy representation or the principal factors, for different $`๐`$. More precisely, we would like to understand what does it mean that โthe monodromy representation or the principal factors of $`Y(x,๐)`$ are the same for different $`๐`$โ? In the previous section these notions were defined in terms of homotopy classes of paths on the punctured Riemann sphere $`^1\{t_1,\mathrm{},t_n\}`$. However, for different $`๐=(t_1,\mathrm{},t_n)`$ the domains $`^1\{t_1,\mathrm{},t_n\}`$ are different. Thus one ought to explain how to consider โthe same paths for different $`๐`$โ. This can be done because we can confine ourselves to local considerations.
### 2.2. Cylindrical neighborhoods of points in $`^n`$
###### Definition 2.2.
Let $`๐ฒ_j,`$ $`1jn`$, be subsets of the complex plane $``$.
The Cartesian product
$$๐ฆ=๐ฒ_1\times \mathrm{}\times ๐ฒ_n^n$$
is said to be the cylindrical set with the bases $`๐ฒ_j,`$ $`1jn`$.
###### Definition 2.3.
Let $`๐=(t_1,\mathrm{},t_n)`$ be a point in $`^n`$. For $`j=1,\mathrm{},n`$ let $`๐ฒ_j`$ be an open neighborhood of $`t_j`$ in $``$, such that:
1. the set $`๐ฒ_j`$ is simply connected;
2. the set $`\overline{๐ฒ_j}`$ is connected.
Denote by $`๐ฆ`$ the cylindrical set with the bases $`๐ฒ_j`$, $`1jn`$.
The set $`๐ฆ`$ is said to be an open cylindrical neighborhood of the point $`๐`$ in $`^n`$.
###### Definition 2.4.
Let subsets $`๐ข,๐`$ of $`^n`$ be such that:
1. the set $`๐`$ is open;
2. the closure $`\overline{๐ข}`$ is compact;
3. $`\overline{๐ข}๐`$.
Then we say that the set $`๐ข`$ is compactly included in $`๐`$ and denote this relation by
$$๐ข๐.$$
### 2.3. Isomonodromic families of Fuchsian systems
Let $`๐^0`$ be a point in the domain<sup>11</sup><sup>11</sup>11 Recall that $`๐`$ is the domain where the residue matrices $`Q_j(๐)`$ from (2.1) are defined and holomorphic. Since $`๐_{}^n`$, the coordinates $`t_1^0,t_2^0,\mathrm{},t_n^0`$ of every point $`๐^\mathrm{๐}๐`$ are pairwise distinct. $`๐`$ and let $`๐ฆ`$ be a cylindrical neighborhood of $`๐^0`$, such that $`๐ฆ๐`$. Then
(2.4)
$$\overline{๐ฒ_p}\overline{๐ฒ_q}=\mathrm{}\mathrm{.\; 1}p,qn,pq,$$
hence<sup>12</sup><sup>12</sup>12Here we refer to a relatively delicate result from general topology: if $`K_1,K_2`$ are two disjoint compact subsets of $`^m`$ and each of two sets $`^mK_1`$, $`^mK_2`$ is connected, then the set $`^m\{K_1K_2\}`$ is connected, as well. (See for example \[HW\], Corollary of Theorem VI.10.) Actually we do not need this result in full generality. For our goal it is enough to consider only the case of $`m=2`$ and some very special compact sets $`K^2`$, such as finite unions of disks, etc. In this particular case, the above stated result is elementary. the set $`_k\overline{๐ฒ}_k`$ is connected.
For a fixed $`๐๐ฆ`$ each homotopy class of loops in the punctured sphere $`^1\{t_1,\mathrm{},t_n\}`$ with the distinguished point $`x_0=\mathrm{}`$ has a representative which is a loop in the perforated sphere $`^1_k\overline{๐ฒ}_k`$.
On the other hand, if $`\gamma `$ is a loop in the perforated sphere $`^1_k\overline{๐ฒ}_k`$ with the distinguished point $`x_0=\mathrm{}`$, then for each fixed $`๐๐ฆ`$ the loop $`\gamma `$ serves as a path of the analytic continuation with respect to $`x`$ for the solution $`Y(x,๐)`$ of (2.1) โ (2.3) and one can consider the corresponding monodromy matrix <sup>13</sup><sup>13</sup>13See Definition 1.1. $`M_\gamma (๐).`$
Although the loop $`\gamma `$ does not depend on $`๐`$, the corresponding monodromy matrix $`M_\gamma (๐)`$ does, in general. We distinguish the following special case:
###### Definition 2.5.
Let (2.1) be a holomorphic family of Fuchsian systems, where the residue matrices $`Q_j(๐)`$ are holomorphic in a domain $`๐_{}^n`$ and satisfy (2.2). For each $`๐๐`$, let $`Y(x,๐)`$ be the solution of (2.1) โ (2.3).
The family (2.1) is said to be an isomonodromic family of Fuchsian systems with the distinguished point $`x_0=\mathrm{}`$ if for every $`๐^0๐`$ there exists a cylindrical open neighborhood $`๐ฆ๐`$ of $`๐^0`$, such that the following holds.
* For every loop $`\gamma `$ in the perforated sphere $`^1_k\overline{๐ฒ_k}`$ with the distinguished point $`x_0=\mathrm{}`$ and every pair of points $`๐^{\mathbf{}},๐^{\mathbf{\prime \prime }}๐ฆ`$ the monodromy matrices $`M_\gamma (๐^{\mathbf{}}),M_\gamma (๐^{\mathbf{\prime \prime }})`$ of the solutions $`Y(x,๐^{\mathbf{}})`$, $`Y(x,๐^{\mathbf{\prime \prime }})`$, which correspond to this loop $`\gamma `$, are equal:
(2.5)
$$M_\gamma (๐^{\mathbf{}})=M_\gamma (๐^{\mathbf{\prime \prime }})$$
###### Remark 2.6.
Note that if the family (2.1) is isomondromic with the distinguished point $`x_0=\mathrm{}`$ and $`๐^0๐`$, then the monodromy matrices of $`Y`$ are constant with respect to $`๐`$ in every cylindrical open neighborhood $`๐ฆ`$ of $`๐^0`$, such that $`๐ฆ๐`$.
### 2.4. Isoprincipal families: the informal definition
Now we introduce the notion of the isoprincipal family of Fuchsian systems, which is the central notion in the present article.
Again, we assume that (2.1) is a holomorphic family of Fuchsian systems, where the residue matrices $`Q_j(๐)`$ are holomorphic in a domain $`๐_{}^n`$ and satisfy (2.2), and consider the solution $`Y(x,๐)`$ of (2.1) โ (2.3).
According to Proposition 1.7, for each fixed $`๐๐`$ a branch $`Y_{\alpha _j}`$ of the solution $`Y`$ in a neighborhood of $`t_j`$ admits the regular-principal factorization, but now both the regular factor $`H_j`$ and the principal factor $`P_{\alpha _j}`$ may depend on $`๐`$:
(2.6)
$$Y_{\alpha _j}(x,๐)=H_j(x,๐)P_{\alpha _j}(x,๐).$$
For example, if the principal factor is chosen in the form (1.24), then $`Z_{\alpha _j}=Z_{\alpha _j}(๐),`$ $`A_{\alpha _j}=A_{\alpha _j}(๐)`$ and
(2.7)
$$P_{\alpha _j}(x,๐)=(xt_j)^{Z_{\alpha _j}(๐)}(xt_j)^{A_{\alpha _j}(๐)}.$$
Roughly speaking, the family (2.1) is isoprincipal if for every $`j`$, $`1jn`$, the matrices $`Z_{\alpha _j}`$ and $`A_{\alpha _j}`$ in (2.7) do not depend on $`๐`$: $`Z_{\alpha _j}(๐)Z_{\alpha _j},A_{\alpha _j}(๐)A_{\alpha _j}`$, and
(2.8)
$$P_{\alpha _j}(x,๐)=(xt_j)^{Z_{\alpha _j}}(xt_j)^{A_{\alpha _j}}.$$
If the matrices $`Z_{\alpha _j}`$ and $`A_{\alpha _j}`$ in (2.7) do not depend on $`๐`$, then the principal factor $`P_{\alpha _j}(x,๐)`$ of the form (2.8) possesses the following property: it depends only on the difference $`xt_j`$.
For our goals, the specific form (2.8) of the principal factors is of no importance. We just need each principal factor $`P_{\alpha _j}(x,๐)`$ to depend only on the difference $`xt_j.`$
This means that there exist functions $`E_{\alpha _j},`$ such that
(2.9)
$$P_{\alpha _j}(x,๐)=E_{\alpha _j}(xt_j),$$
or, in the language of differential equations,
(2.10a) $`{\displaystyle \frac{P_{\alpha _j}}{t_{\mathrm{}}}}`$ $`=0,\mathrm{}j,`$
(2.10b) $`{\displaystyle \frac{P_{\alpha _j}}{x}}`$ $`={\displaystyle \frac{P_{\alpha _j}}{t_j}}={\displaystyle \frac{dE_{\alpha _j}}{d\zeta }}|_{\zeta =xt_j}.`$
The formal definition of the isoprincipal family of Fuchsian systems (see Definition 2.9 below) is more involved, since for each $`๐`$ the branch $`Y_{\alpha _j}(x,๐)`$ and the principal factor $`P_{\alpha _j}(x,๐)`$ depend on the choice of a path $`\alpha _j`$ in the punctured sphere $`^1\{t_1,\mathrm{},t_n\}`$, which connects the distinguished point $`x_0=\mathrm{}`$ with a neighborhood of $`x=t_j`$.
Moreover, for each $`๐`$ the principal factor $`P_{\alpha _j}(x,๐)`$ should be a transplant <sup>14</sup><sup>14</sup>14See Definition 1.6. of a function $`E_{\alpha _j}(\zeta )`$, holomorphic on the Riemann surface of $`\mathrm{ln}\zeta `$, into a universal covering surface over $`\{t_j\}`$ and these transplants should be defined coherently with respect to $`๐`$.
### 2.5. Isoprincipal families: the formal definition
###### Definition 2.7.
Let a function $`E(\zeta )`$ be defined on the Riemann surface of $`\mathrm{ln}\zeta `$. Let $`๐ฒ_j`$ be a simply connected domain, compactly included in $``$. Let $`p_j\overline{๐ฒ_j}`$ and let us choose a branch of $`\mathrm{arg}(p_jt_j)`$, continuous with respect to $`t_j`$ in $`๐ฒ_j`$ (since $`๐ฒ_j`$ is simply connected and $`๐ฒ_j\{p_j\}`$, such a choice can be made). For every $`t_j๐ฒ_j`$ let us specify the value of $`\mathrm{arg}(p_jt_j)`$ in this manner and consider the corresponding transplant $`E(xt_j)`$ of the function $`E(\zeta )`$ into $`\mathrm{cov}(\{t_j\};p_j).`$
Then the family of transplants $`\{E(xt_j)\}_{t_j๐ฒ_j}`$ is said to be coherent with respect to $`t_j`$ in $`๐ฒ_j`$.
###### Definition 2.8.
Let $`๐^0`$ be a point in a domain $`๐_{}^n`$ and let $`๐ฆ,๐ฅ`$ be a pair of open cylindrical neighborhoods of the point $`๐^0`$, such that
$$๐ฅ_{}^n\text{ and }๐ฆ\left\{๐ฅ๐\right\}.$$
Then the pair $`๐ฆ,๐ฅ`$ is said to be a nested pair of open cylindrical neighborhoods of the point $`๐^0`$ in $`๐.`$
###### Definition 2.9 (Formal definition of the isoprincipal family).
Let (2.1) be a holomorphic family of Fuchsian systems, where the residue matrices $`Q_j(๐)`$ are holomorphic in a domain $`๐_{}^n`$ and satisfy (2.2). For $`๐๐`$ let $`Y(x,๐)`$ be the solution of (2.1) โ (2.3).
The family (2.1) is said to be an isoprincipal family of Fuchsian systems with the distinguished point $`x_0=\mathrm{}`$ if for every $`๐^0๐`$ there exists a nested pair of open cylindrical neighborhoods of $`๐^0`$:
$$๐ฅ=๐ฑ_1\times \mathrm{}\times ๐ฑ_n,๐ฆ=๐ฒ_1\times \mathrm{}\times ๐ฒ_n,๐ฆ๐ฅ,$$
such that the following holds.
* For every path $`\alpha _j`$ in the perforated sphere $`^1_k\overline{๐ฒ_k}`$ from the distinguished point $`x_0=\mathrm{}`$ to a point $`p_j๐ฑ_j\overline{๐ฒ_j}`$, $`1jn`$, there exists a coherent family of transplants $`\{E_{\alpha _j}(xt_j)\}_{t_j๐ฒ_j}`$ of a function $`E_{\alpha _j}(\zeta )`$, holomorphic and invertible on the Riemann surface of $`\mathrm{ln}\zeta `$, such that for each $`๐๐ฆ`$ the branch $`Y_{\alpha _j}(x,๐)`$ of the solution $`Y(x,๐)`$ in the punctured domain $`๐ฑ_j\{t_j\}`$ admits the representation
(2.11)
$$Y_{\alpha _j}(x,๐)=H_j(x,๐)E_{\alpha _j}(xt_j),x\mathrm{cov}(๐ฑ_j\{t_j\};p_j),$$
where $`H_j(x,๐)`$ is a function, holomorphic (with respect to $`x`$) and invertible in the entire domain $`๐ฑ_j.`$
###### Remark 2.10.
In view of Definition 1.8, Definition (2.9) means that the family (2.1) is isoprincipal with the distinguished point $`x_0=\mathrm{}`$ if every branch of the solution $`Y(x,๐)`$ in a neighborhood of each singular point $`x=t_j`$ admits the regular-principal factorization (2.11), where the principal factor is the appropriately shifted copy of a function $`E_{\alpha _j}(\zeta )`$, which is holomorphic and invertible on the Riemann surface of $`\mathrm{ln}\zeta `$ and does not depend on $`๐`$. The meaning of the words โappropriately shifted copyโ is made precise in Definition 2.7: this is what we call a coherent family of transplants of $`E_{\alpha _j}(\zeta )`$.
Thus Definition 2.9 is a formal interpretation of the informal definition in Section 2.4.
We would also like to note that it suffices to consider only the branches corresponding to a certain choice of the paths $`\alpha _1,\mathrm{},\alpha _n`$ โ the same choice as the one mentioned in Remark 2.13 below.
### 2.6. Every isoprincipal family is an isomonodromic one.
###### Theorem 2.11.
Let (2.1) be a holomorphic family of Fuchsian systems, where the residue matrices $`Q_j(๐ญ)`$ are holomorphic in a domain $`๐_{}^n`$ and satisfy (2.2).
Assume that the family (2.1) is isoprincipal with the distinguished point $`x_0=\mathrm{}`$.
Then this family is isomonodromic with the distinguished point $`x_0=\mathrm{}`$.
Before we turn to the proof of Theorem 2.11, let us introduce the following โ$`๐`$-dependentโ counterpart of Definition 1.4:
###### Definition 2.12.
Let $`๐^0`$ be a point in the domain $`๐`$ and let $`๐ฆ๐ฅ`$ be a nested pair of open cylindrical neighborhoods of $`๐^0`$ in $`๐`$. For $`1jn`$ let $`\alpha _j`$ be a path in the perforated sphere $`^1_k\overline{๐ฒ_k}`$ from the distinguished point $`x_0=\mathrm{}`$ to a point $`p_j๐ฑ_j\overline{๐ฒ_j}.`$
Furthermore, assume that $`\beta _j`$ is the loop in the annulus $`๐ฑ_j\overline{๐ฒ_j}`$ with the distinguished point $`p_j`$ which makes one positive circuit of the set $`๐ฒ_j`$, and let $`\gamma _j`$ be the loop in the perforated sphere $`^1_k\overline{๐ฒ_k}`$ with the distinguished point $`x_0=\mathrm{}`$, defined by
(2.12)
$$\gamma _j\stackrel{\text{def}}{=}\alpha _j^1\beta _j\alpha _j.$$
Then:
* the loop $`\beta _j`$ is said to be the small loop around the set $`๐ฒ_j`$ in the annulus $`๐ฑ_j\overline{๐ฒ_j}`$;
* the loop $`\gamma _j`$ is said to be the big loop around the set $`๐ฒ_j`$, corresponding to the path $`\alpha _j`$.
###### Remark 2.13.
Similarly to the case of a fixed $`๐`$ (see Remark 1.5), for a suitable choice of the paths $`\alpha _1,\mathrm{},\alpha _n`$ the corresponding big loops $`\gamma _1,\mathrm{},\gamma _n`$ generate the fundamental group $`\pi (^1_k\overline{๐ฒ_k};\mathrm{})`$.
###### Proof of Theorem 2.11.
Let $`๐^0`$ be a point in $`๐`$ and let $`๐ฆ๐ฅ`$ be a nested pair of open cylindrical neighborhoods of $`๐^0`$ as in Definition 2.9.
In view of Remark 2.13, it suffices to prove that if $`\alpha _j`$ is a path in $`^1_k\overline{๐ฒ_k}`$ from the distinguished point $`x_0=\mathrm{}`$ to a point $`p_j๐ฑ_j\overline{๐ฒ_j}`$ and $`\gamma _j`$ is the corresponding big loop around $`๐ฒ_j,`$ then the monodromy matrix $`M_{\gamma _j}(๐)`$ of $`Y(x,๐)`$ does not depend on $`๐`$:
(2.13)
$$M_{\gamma _j}(๐)=\text{const},๐๐ฆ.$$
In view of (1.20), for each fixed $`๐๐ฆ`$ the monodromy matrix $`M_{\gamma _j}(๐)`$ is given by
$$M_{\gamma _j}(๐)=Y_{\alpha _j}^1(x,๐)Y_{\alpha _j}(t_j+(xt_j)e^{2\pi i},๐),x\mathrm{cov}(๐ฑ_j\{t_j\};p_j),$$
where $`Y_{\alpha _j}(x,๐)`$ is the branch of $`Y(x,๐)`$ in $`๐ฑ_j\{t_j\}`$, corresponding to the path $`\alpha _j`$.
Substituting the expression (2.6) for $`Y_{\alpha _j}`$ into the above identity and taking into account that the factor $`H_j(x,๐)`$ is a single-valued function of $`x`$, we obtain
$$\begin{array}{c}M_{\gamma _j}(๐)=E_{\alpha _j}^1(xt_j)E_{\alpha _j}((xt_j)e^{2\pi i})\hfill \\ \hfill =\left(E_{\alpha _j}^1(\zeta )E_{\alpha _j}(\zeta e^{2\pi i})\right)|_{\zeta =xt_j},x\mathrm{cov}(๐ฑ_j\{t_j\};p_j).\end{array}$$
Thus the function $`E_{\alpha _j}^1(\zeta )E_{\alpha _j}(\zeta e^{2\pi i})`$, holomorphic on the Riemann surface of $`\mathrm{ln}\zeta `$, is constant with respect to $`\zeta `$ on a certain non-empty open subset of this surface. Therefore, this function is identically constant on the Riemann surface of $`\mathrm{ln}\zeta `$ and we write
(2.14)
$$M_{\gamma _j}(๐)=E_{\alpha _j}^1(\zeta )E_{\alpha _j}(\zeta e^{2\pi i})๐๐ฆ,\zeta \mathrm{cov}(\{0\};1).$$
But the right-hand side of the last identity does not depend on $`๐`$, hence we obtain (2.13). โ
### 2.7. Every non-resonant isomonodromic family is an isoprincipal one.
The converse to Theorem 2.11 is only conditionally true: it holds under the assumption that all the matrices $`Q_j`$ are non-resonant (see footnote $`^\text{9}`$). In general, however, an isomonodromic family can be non-isoprincipal. The appropriate counterexample will be presented in Section 5 of this paper.
###### Lemma 2.14.
Let (2.1) be a holomorphic family of Fuchsian systems, where the residue matrices $`Q_j(๐ญ)`$ are holomorphic in a domain $`๐_{}^n`$ and satisfy (2.2). Assume that this family is isomonodromic with the distinguished point $`x_0=\mathrm{}`$.
Then the family (2.1) is isospectral in the following sense: for every pair of points $`๐ญ^{\mathbf{}},๐ญ^{\mathbf{\prime \prime }}๐`$ and each $`j`$, $`1jn`$, the spectra $`\mathrm{spec}Q_j(๐ญ^{\mathbf{}})`$ and $`\mathrm{spec}Q_j(๐ญ^{\mathbf{\prime \prime }})`$ are equal <sup>15</sup><sup>15</sup>15As usual, the spectra are considered โwith multliplicitiesโ.:
(2.15)
$$\mathrm{spec}Q_j(๐^{\mathbf{}})=\mathrm{spec}Q_j(๐^{\mathbf{\prime \prime }}),๐^{\mathbf{}},๐^{\mathbf{\prime \prime }}๐,1jn.$$
###### Theorem 2.15.
Let (2.1) be a holomorphic family of Fuchsian systems, where the residue matrices $`Q_j(๐ญ)`$ are holomorphic in a domain $`๐_{}^n`$ and satisfy (2.2).
Assume that this family satisfies the following conditions:
1. It is isomonodromic with the distinguished point $`x_0=\mathrm{}`$;
2. At least for one point $`๐๐`$, each of the matrices $`Q_j(๐),j=1,\mathrm{},n,`$ is non-resonant.
Then the family (2.1) is isoprincipal with the distinguished point $`x_0=\mathrm{}`$.
###### Remark 2.16.
It should be noted that, unlike the rest of our considerations in the present article, the proofs of Lemma 2.14 and Theorem 2.15 utilize the explicit form of the principal factors mentioned in Remark 1.10.
###### Proof of Lemma 2.14.
Let $`๐^0`$ be a point in $`๐`$ and let $`๐ฆ๐ฅ`$ be a nested pair of open cylindrical neighborhoods of $`๐^0`$ in $`๐`$ as in Definition 2.5. For $`j=1,\mathrm{},n`$ let us choose a path $`\alpha _j`$ in the perforated sphere $`^1_k\overline{๐ฒ_k}`$ from the distinguished point $`x_0=\mathrm{}`$ to a point $`p_j๐ฑ_j\overline{๐ฒ_j}`$. As usual, for each $`๐๐ฆ`$ we denote by $`Y_{\alpha _j}(x,๐)`$ the branch of the solution $`Y(x,๐)`$ of (2.1) โ (2.3) in the punctured domain $`๐ฑ_j\{t_j\}`$, corresponding to this path $`\alpha _j`$.
Then, in view of Remark 1.10 (see (1.24), (1.25)), the branch $`Y_{\alpha _j}(x,๐)`$ admits the regular-principal factorization
$$Y_{\alpha _j}(x,๐)=H_j(x,๐)(xt_j)^{Z_{\alpha _j}(๐)}(xt_j)^{A_{\alpha _j}(๐)},๐๐ฆ,x\mathrm{cov}(๐ฑ_j\{t_j\};p_j),$$
where
$$\mathrm{spec}\left(e^{2\pi iA_{\alpha _j}(๐)}\right)=\mathrm{spec}\left(e^{2\pi iQ_j(๐)}\right).$$
Let $`\gamma _j`$ be the big loop around $`๐ฒ_j`$, corresponding to the path $`\alpha _j`$. Then the monodromy matrix $`M_{\gamma _j}(๐)`$ of $`Y(x,๐)`$, corresponding to the loop $`\gamma _j`$ is given by
(2.16)
$$M_{\gamma _j}(๐)=Y_{\alpha _j}^1(x,๐)Y_{\alpha _j}(t_j+(xt_j)e^{2\pi i},๐)=e^{2\pi iA_{\alpha _j}(๐)}.$$
Therefore,
$$\mathrm{spec}(M_{\gamma _j}(๐))=\mathrm{spec}\left(e^{2\pi iQ_j(๐)}\right),๐๐ฆ.$$
Since the family (2.1) is isomonodromic, we have
$$M_{\gamma _j}(๐)=M_{\gamma _j}(๐^\mathrm{๐}),๐๐ฆ,$$
hence
$$\mathrm{spec}\left(e^{2\pi iQ_j(๐)}\right)=\mathrm{spec}\left(e^{2\pi iQ_j(๐^\mathrm{๐})}\right),๐๐ฆ.$$
This means that the spectra $`\mathrm{spec}Q_j(๐)`$ and $`\mathrm{spec}Q_j(๐^\mathrm{๐})`$ coincide modulo integers. But the function $`Q_j(๐)`$ is continuous with respect to $`๐`$ in $`๐ฆ,`$ hence
$$\mathrm{spec}Q_j(๐)=\mathrm{spec}Q_j(๐^\mathrm{๐}),๐๐ฆ.$$
Since the above identity holds for all $`๐`$ in a neighborhood of every point $`๐^\mathrm{๐}๐`$, we obtain the identity (2.15). โ
###### Proof of Theorem 2.15.
Let $`๐^0`$ be a point in $`๐`$ and let $`๐ฆ๐ฅ`$ be a nested pair of open cylindrical neighborhoods of $`๐^0`$ in $`๐`$ as in Definition 2.5.
As in the proof of Lemma 2.14, for $`j=1,\mathrm{},n`$ let us choose a path $`\alpha _j`$ in the perforated sphere $`^1_k\overline{๐ฒ_k}`$ from the distinguished point $`x_0=\mathrm{}`$ to a point $`p_j๐ฑ_j\overline{๐ฒ_j}`$ and consider for each fixed $`๐๐ฆ`$ the branch $`Y_{\alpha _j}(x,๐)`$ of the solution $`Y(x,๐)`$ of (2.1) โ (2.3) in the punctured domain $`๐ฑ_j\{t_j\}`$, corresponding to this path $`\alpha _j`$.
Since by Lemma 2.14 the family (2.1) is isospectral, the matrix $`Q_j(๐)`$ is non-resonant for every $`๐๐`$. Hence, in view of Remark 1.10 (see the expressions (1.22), (1.23)), the branch $`Y_{\alpha _j}`$ admits the regular - principal factorization
$$Y(x,๐)=H_j(x,๐)(xt_j)^{A_{\alpha _j}(๐)},x\mathrm{cov}(๐ฑ_j\{t_j\};p_j),$$
where the matrix $`A_{\alpha _j}(๐)`$ is similar to the matrix $`Q_j(๐)`$. Therefore,
(2.17)
$$\mathrm{spec}(A_{\alpha _j}(๐))=\mathrm{spec}(A_{\alpha _j}(๐^\mathrm{๐})),๐๐ฆ.$$
In view of Definition 2.9, it remains to prove that the matrix $`A_{\alpha _j}(๐)`$ does not actually depend on $`๐`$.
Let $`\gamma _j`$ be the big loop around $`๐ฒ_j`$, corresponding to the path $`\alpha _j`$. Then the monodromy matrix $`M_{\gamma _j}(๐)`$ of $`Y(x,๐)`$, corresponding to the loop $`\gamma _j`$, is given by (see (2.16))
$$M_{\gamma _j}(๐)=e^{2\pi iA_{\alpha _j}(๐)},๐๐ฆ.$$
Since the family (2.1) is isomonodromic, we have
$$M_{\gamma _j}(๐)=M_{\gamma _j}(๐^\mathrm{๐}),๐๐ฆ,$$
hence
$$e^{2\pi iA_{\alpha _j}(๐)}=e^{2\pi iA_{\alpha _j}(๐^\mathrm{๐})},๐๐ฆ.$$
But, in view of (2.17), the last identity implies
$$A_{\alpha _j}(๐)=A_{\alpha _j}(๐^\mathrm{๐})=\text{const},๐๐ฆ.$$
## 3. Isoprincipal Families of Fuchsian Systems and the Schlesinger System
### 3.1. The Schlesinger system
A natural question arises: how to express the property of a family of Fuchsian systems
(3.1)
$$\frac{dY}{dx}=\left(\underset{1jn}{}\frac{Q_j(๐)}{xt_j}\right)Y,$$
to be isoprincipal in terms of the residues matrix functions $`Q_j(๐)`$? Here, as always, we assume that the residue matrices $`Q_j(๐)`$ are holomorphic in a domain $`๐_{}^n`$ and satisfy
(3.2)
$$\underset{1jn}{}Q_j(๐)0,๐๐.$$
It turns out that an answer to this question is given by the so-called Schlesinger system of PDEs:
(3.3)
$$\{\begin{array}{ccc}\hfill \frac{Q_i}{t_j}& =\frac{[Q_i,Q_j]}{t_it_j},\hfill & \hfill 1i,jn,ij,\\ \hfill \frac{Q_i}{t_i}& =\underset{\begin{array}{c}1jn\\ ji\end{array}}{}\frac{[Q_i,Q_j]}{t_it_j},\hfill & \hfill 1in.\end{array}$$
The following theorem is the main result of the present article.
###### Theorem 3.1 (The main result).
Let (3.1) be a holomorphic family of Fuchsian systems, where the residue matrices $`Q_j(๐ญ)`$ are holomorphic in a domain $`๐_{}^n`$ and satisfy (3.2).
Then the family (3.1) is isoprincipal with the distinguished point $`x_0=\mathrm{}`$ if and only if the residue matrices $`Q_j(๐ญ)`$ satisfy with respect to $`๐ญ`$ the Schlesinger system (3.3) in the domain $`๐`$.
###### Remark 3.2.
Note that (3.3) implies
$$\frac{}{t_j}\underset{1in}{}Q_i=0,1jn,$$
that is, $`_iQ_i(๐)`$ is a first integral of the Schlesinger system.
In particular, if functions $`Q_j(๐)`$, $`1jn,`$ satisfy the Schlesinger system (3.3) in the domain $`๐`$ and at some point $`๐^\mathrm{๐}๐`$ it holds that
$$\underset{1jn}{}Q_j(๐^\mathrm{๐})=0,$$
then these functions $`Q_j(๐)`$ satisfy the relation (3.2).
###### Remark 3.3.
We would like to stress that in Theorem 3.1 no assumptions on the spectra of the residue matrices $`Q_j(๐)`$ are made. Thus Theorem 3.1 for the isoprincipal families of Fuchsian systems can be viewed as an amended version of L. Schlesingerโs statement, concerning the isomonodromic deformations (see \[Sch2\] and the introduction of the present article).
In the case of the isomonodromic families of Fuchsian systems Theorem 3.1 implies the following:
1. If the residue matrices $`Q_j(๐ญ)`$ satisfy the Schlesinger system (3.3), then the family (3.1) is isoprincipal and hence by Theorem 2.11 also isomonodromic.
2. If the family (3.1) is isomonodromic and, in addition, all the residue matrices $`Q_j(๐ญ)`$ are non-resonant (at least at some point), then by Theorem 2.15 the family (3.1) is isoprincipal and hence the residue matrices $`Q_j(๐ญ)`$ satisfy the Schlesinger system (3.3).
We remark that in the statement (ii) the assumption of non-resonance for the residues $`Q_j(๐)`$ cannot be omitted: in Section 5 we shall present an example of the isomonodromic family (3.1), where the residue matrices $`Q_j(๐)`$ are resonant and do not satisfy the Schlesinger system (3.3) (thus contradicting the statement of L. Schlesinger).
Nevertheless, our proof of the โonly ifโ part of the Theorem 3.1 largely follows the original proof of L. Schlesinger for the isomonodromic case (see also \[IKSY, Section 3.5\], where the modern adaptation of Schlesingerโs proof is presented). In particular, the overdetermined linear system (3.6), which appears in Proposition 3.6 below and is crucial in the derivation of the Schlesinger system, can be found in \[Sch2, Section II\].
The proof of Theorem 3.1 will be split into parts and presented as a series of propositions, culminating with Propositions 3.15 and 3.16.
### 3.2. The auxiliary system related to the isomonodromic family of Fuchsian systems
In order to prove Theorem 3.1, we have to study the partial derivatives of the solution $`Y(x,๐)`$, satisfying the initial condition
(3.4)
$$Y(x,๐)|_{x=\mathrm{}}=I,$$
with respect to the parameters $`t_1,\mathrm{},t_n`$.
First of all, let us choose and fix a point $`๐^\mathrm{๐}๐`$ and let $`๐ฆ๐`$ be a cylindrical open neighborhood of $`๐^\mathrm{๐}`$. Since the coefficients of the system (3.1) and the initial condition (3.4) depend holomorphically on $`๐ญ`$, the solution $`Y(x,๐ญ)`$ is holomorphic jointly in $`x`$ and $`๐ญ`$ in the Cartesian product $`\mathrm{cov}(^1_k\overline{๐ฒ_k};\mathrm{})\times ๐ฆ.`$
In particular, the partial derivatives $`{\displaystyle \frac{Y}{t_j}}(x,๐^\mathrm{๐})`$ are defined and holomorphic with respect to $`x`$ in $`\mathrm{cov}(^1_k\overline{๐ฒ_k};\mathrm{}).`$ Since these definitions of $`{\displaystyle \frac{Y}{t_j}}(x,๐^\mathrm{๐})`$ agree for various choices of $`๐ฆ`$ as long as $`๐ฆ`$ is sufficiently small, we conclude that for each fixed $`๐ญ^\mathrm{๐}๐`$ the partial derivatives $`{\displaystyle \frac{Y}{t_j}}(x,๐ญ^\mathrm{๐})`$ are defined and holomorphic as functions of $`x`$ on the the same surface as the function $`Y(x,๐ญ^\mathrm{๐})`$ itself โ the universal covering surface $`\mathrm{cov}(^1\{t_1^0\mathrm{},t_n^0\};\mathrm{})`$.
It turns out that in terms of these partial derivatives of $`Y`$ the property of the family (3.1) to be isomonodromic can be expressed as follows:
###### Proposition 3.4.
Let (3.1) be a holomorphic family of Fuchsian systems, where the residue matrices $`Q_j(๐ญ)`$ are holomorphic in a domain $`๐_{}^n`$ and satisfy (3.2).
Then the family (3.1) is isomonodromic with the distinguished point $`x_0=\mathrm{}`$ if and only if the solution $`Y(x,๐ญ)`$ of (3.1), (3.4) satisfies a linear system of the form
(3.5)
$$\{\begin{array}{cc}\hfill \frac{Y}{x}& =\underset{1jn}{}\frac{Q_j(๐)}{xt_j}Y,\hfill \\ \hfill \frac{Y}{t_j}& =T_j(x,t)Y,1jn,\hfill \end{array}$$
where for each $`๐ญ๐`$ the functions $`T_j(x,t)`$, $`1jn,`$ are single-valued holomorphic with respect to $`x`$ in $`^1\{t_1,\mathrm{},t_n\}`$.
###### Definition 3.5.
Let the family (3.1) of Fuchsian systems be isomonodromic with the distinguished point $`x_0=\mathrm{}`$.
The system (3.5) with the single-valued coefficients $`T_j(x,t)`$, $`1jn,`$ which appears in Theorem 3.4, is said to be the auxiliary linear system related to the isomonodromic family (3.1) of Fuchsian systems.
###### Proof of Proposition 3.4.
The first equation of the system (3.5) is just the Fuchsian system (3.1) itself, hence we only need to prove that for each fixed $`๐๐`$ the logarithmic derivatives
$$T_j(x,๐)\stackrel{\text{def}}{=}\frac{Y}{t_j}(x,๐)Y^1(x,๐),$$
which a priori are defined as holomorphic functions of $`x`$ on the universal covering surface $`\mathrm{cov}(^1\{t_1,\mathrm{},t_n\};\mathrm{}),`$ are single-valued in the punctured sphere $`^1\{t_1,\mathrm{},t_n\}.`$
Let us choose a point $`๐^0๐`$ and a cylindrical open neighborhood $`๐ฆ๐`$ of $`๐^0`$. Let $`\gamma \pi (^1_k\overline{๐ฒ_k};\mathrm{})`$ and let us denote by
$$xx\gamma ,x\mathrm{cov}(^1\underset{k}{}\overline{๐ฒ_k};\mathrm{})$$
the deck transformation of the universal covering surface $`\mathrm{cov}(^1_k\overline{๐ฒ_k};\mathrm{})`$, corresponding <sup>16</sup><sup>16</sup>16See (1.4). to this loop $`\gamma `$.
According to Definition 1.1, the monodromy matrix $`M_\gamma (๐)`$ of the solution $`Y(x,๐)`$, which corresponds to the loop $`\gamma `$, is given by
$$M_\gamma (๐)=Y^1(x,๐)Y(x\gamma ,๐),x\mathrm{cov}(^1\underset{1kn}{}\overline{๐ฒ_k};\mathrm{}),๐๐ฆ.$$
The monodromy matrix $`M_\gamma (๐)`$ does not depend on $`x`$ and hence is a holomorphic single-valued <sup>17</sup><sup>17</sup>17Recall (see Definition 2.2) that all the bases $`๐ฒ_k`$ of the cylindrical neighborhood $`๐ฆ`$ are simply connected. function of $`๐`$ in $`๐ฆ`$. Differentiating the equality
$$Y(x\gamma ,๐)=Y(x,๐)M_\gamma (๐)$$
with respect to $`t_j`$, we obtain
$$\frac{Y}{t_j}(x\gamma ,๐)=\frac{Y}{t_j}(x,๐)M_\gamma (๐)+Y(x,๐)\frac{M_\gamma }{t_j}(๐).$$
Therefore, the logarithmic derivative
$$T_j(x,๐)=\frac{Y}{t_j}(x,๐)Y^1(x,๐)$$
satisfies the monodromy relation
$$T_j(x\gamma ,๐)=T_j(x,๐)+Y(x,๐)\frac{M_\gamma }{t_j}(๐)Y^1(x\gamma ,๐).$$
The last equality implies that the following two statements are equivalent <sup>18</sup><sup>18</sup>18This equivalence is stronger than the statement of Proposition 3.6 in the sense that it holds for each individual loop $`\gamma `$ and each individual index $`j`$.:
1. The monodromy matrix $`M_\gamma (๐)`$ does not depend on $`t_j`$:
$$\frac{M_\gamma }{t_j}(๐)0,๐๐ฆ.$$
2. It holds that
$$T_j(x\gamma ,๐)T_j(x,๐),x\mathrm{cov}(^1\underset{1kn}{}\overline{๐ฒ_k};\mathrm{}),๐๐ฆ.$$
However, the statement 2) holds for every $`\gamma \pi (^1_k\overline{๐ฒ_k};\mathrm{})`$ if and only if for each $`๐๐ฆ`$ the function $`T_j(x,๐)`$ is a single-valued function of $`x`$ in $`^1\{t_1,\mathrm{},t_n\}`$. In view of Definition 2.5, this completes the proof. โ
### 3.3. The auxiliary system related to the isoprincipal family of Fuchsian systems
Now we turn to the case, when the family (3.1) is not only isomonodromic but, moreover, isoprincipal. We claim that in this special case the auxiliary linear system (3.5) can be written explicitly in terms of the residues $`Q_j(๐)`$.
###### Proposition 3.6.
Let (3.1) be a holomorphic family of Fuchsian systems, where the residue matrices $`Q_j(๐ญ)`$ are holomorphic in a domain $`๐_{}^n`$ and satisfy (3.2). Assume that the family (3.1) is isoprincipal with the distinguished point $`x_0=\mathrm{}`$.
Then the solution $`Y(x,๐ญ)`$ of (3.1), (3.4) satisfies the following auxiliary system:
(3.6)
$$\{\begin{array}{cc}\hfill \frac{Y}{x}& =\underset{1jn}{}\frac{Q_j(๐)}{xt_j}Y,\hfill \\ \hfill \frac{Y}{t_j}& =\frac{Q_j(๐)}{xt_j}Y,1jn.\hfill \end{array}$$
In the proof of Proposition 3.6 we shall use the following
###### Lemma 3.7.
Let $`๐ฐ,๐ฑ`$ be simply-connected domains in the complex plane $``$, such that:
1. $`๐ฐ๐ฑ`$;
2. the set $`๐ฑ\overline{๐ฐ}`$ is connected.
Let $`๐ฆ`$ be a domain in $`^n`$ and let $`H(x,๐ญ)`$ be a function of $`x๐ฑ`$ and $`๐ญ๐ฆ`$, possessing the following properties:
1. the function $`H(x,๐)`$ is holomorphic (jointly in $`x`$ and $`๐`$) in $`\{๐ฑ\overline{๐ฐ}\}\times ๐ฆ`$;
2. for each fixed $`๐๐ฆ`$ the function $`H(x,๐)`$ is holomorphic with respect to $`x`$ in the entire domain $`๐ฑ`$.
Then the function $`H(x,๐ญ)`$ is holomorphic (jointly in $`x`$ and $`๐ญ`$) in the domain $`๐ฑ\times ๐ฆ`$.
###### Remark 3.8.
Lemma 3.7 is a special case of the well-known Hartogs lemma. However, in this simple case the conclusion follows immediately from the Cauchy integral formula.
Indeed, let $`\beta `$ be a smooth loop in the annulus $`๐ฑ\overline{๐ฒ}`$ which makes one positive circuit of the set $`๐ฒ`$ and let $`\mathrm{\Delta }๐ฆ`$ be a polydisk:
$$\mathrm{\Delta }=\mathrm{\Delta }_1\times \mathrm{}\mathrm{\Delta }_n.$$
Then for each $`๐\mathrm{\Delta }`$ and $`x\overline{๐ฒ}`$ it holds that
$$\begin{array}{c}H_j(x,๐)=\frac{1}{2\pi i}\underset{\beta _j}{}\frac{H_j(\zeta ,๐)}{\zeta x}๐\zeta \hfill \\ \hfill =\frac{1}{(2\pi i)^{n+1}}\underset{\beta _j}{}\left(\underset{\mathrm{\Delta }_n}{}\mathrm{}\underset{\mathrm{\Delta }_1}{}\frac{H_j(\zeta ,\tau _1,\mathrm{}\tau _n)}{(\tau _1t_1)\mathrm{}(\tau _nt_n)}๐\tau _1\mathrm{}๐\tau _n\right)\frac{d\zeta }{\zeta x},\end{array}$$
where $`\mathrm{\Delta }_k`$ denotes the boundary of the disk $`\mathrm{\Delta }_k`$. Since the contours of integration lie in the domain $`\{๐ฑ\overline{๐ฒ}\}\times ๐ฆ`$, where $`H_j(x,๐)`$ is jointly holomorphic in $`x`$ and $`๐`$ (in particular, continuous), the integral represents a function, jointly holomorphic in $`x`$ and $`๐`$ for $`x`$ in a neighborhood of $`\overline{๐ฒ_j}`$ and $`๐\mathrm{\Delta }`$.
###### Proof of Proposition 3.6.
Let $`๐^\mathrm{๐}๐`$ and let $`๐ฆ๐ฅ`$ be a nested pair of cylindrical open neighborhoods of $`๐^\mathrm{๐}`$.
For $`\mathrm{}=1,\mathrm{},n`$ let us consider the logarithmic derivative
(3.7)
$$T_{\mathrm{}}(x,๐)=\frac{Y}{t_{\mathrm{}}}(x,๐)Y^1(x,๐).$$
Since by Theorem 2.11 the isoprincipal family (3.1) is also isomonodromic, Proposition 3.4 implies that for each fixed $`๐๐ฆ`$ the function $`T_{\mathrm{}}(x,๐)`$ is single-valued holomorphic with respect to $`x`$ in the punctured sphere $`^1\{t_1,\mathrm{},t_n\}`$. We have to prove that
(3.8)
$$T_{\mathrm{}}(x,๐)=\frac{Q_{\mathrm{}}(๐)}{xt_{\mathrm{}}}.$$
To this end let us introduce the auxiliary function
(3.9)
$$F_{\mathrm{}}(x,๐)\stackrel{\text{def}}{=}T_{\mathrm{}}(x,๐)+\frac{Q_{\mathrm{}}(๐)}{xt_{\mathrm{}}},$$
which is single-valued holomorphic with respect to $`x`$ in the punctured sphere $`^1\{t_1,\mathrm{},t_n\}.`$ We are going to show that for $`j=1,\mathrm{},n`$ the following statement holds:
* The point $`x=t_j`$ is a removable singularity of the function $`F_{\mathrm{}}(x,๐)`$.
Then, according to Liouville theorem, the function $`F_{\mathrm{}}(x,๐)`$ is constant with respect to $`x`$. Since, as follows from (3.4),(3.9)
(3.10)
$$T_{\mathrm{}}(x,๐)|_{x=\mathrm{}}=F_{\mathrm{}}(x,๐)|_{x=\mathrm{}}=0,$$
we can then conclude that
$$F_{\mathrm{}}(x,๐)0,$$
which leads to the desired result (3.8).
Now we turn to the proof of the statement $`()`$. To begin with, let us choose and fix a path $`\alpha _j`$ in the perforated sphere $`^1_k\overline{๐ฒ_k}`$ from the distinguished point $`x_0=\mathrm{}`$ to a point $`p_j๐ฑ_j\overline{๐ฒ_j}.`$ Then, according to Definition 2.9, for each $`๐๐ฆ`$ the branch $`Y_{\alpha _j}(x,๐)`$ of the solution $`Y(x,๐)`$ in $`๐ฑ_j\{t_j\}`$ admits the regular-principal factorization
(3.11)
$$Y_{\alpha _j}(x,๐)=H_j(x,๐)E_{\alpha _j}(xt_j),x\mathrm{cov}(๐ฑ_j\{t_j\};p_j),$$
where the family $`\{E_{\alpha _j}(xt_j)\}_{t_j๐ฒ_j}`$ is a coherent family of transplants of a function $`E_{\alpha _j}(\zeta )`$, holomorphic and invertible on the Riemann surface of $`\mathrm{ln}\zeta `$, and the function $`H_j(x,๐)`$ is holomorphic with respect to $`x`$ and invertible in the entire (non-punctured) domain $`๐ฑ_j`$.
Then the principal factor $`E_{\alpha _j}(xt_j)`$ is jointly holomorphic in $`x`$ and $`t_j`$ in $`\mathrm{cov}(\overline{๐ฒ_j};p_j)\}\times ๐ฒ_j,`$ and hence the regular factor
$$H_j(x,๐)=Y_{\alpha _j}^1(x,๐)E_{\alpha _j}(xt_j)$$
is jointly holomorphic (single-valued) in $`x`$ and $`๐`$ in $`\{๐ฑ_j\overline{๐ฒ_j}\}\times ๐ฆ`$. Since the function $`H_j(x,๐)`$ is also holomorphic with respect to $`x`$ in the entire domain $`๐ฑ_j`$, Lemma 3.7 implies that it is jointly holomorphic in $`x`$ and $`๐`$ in $`๐ฑ_j\times ๐ฆ`$.
Thus we can differentiate the equality 3.11 with respect to $`t_{\mathrm{}}`$, $`1\mathrm{}n`$. First we consider the case $`\mathrm{}j`$. Then, since
(3.12)
$$\frac{E_{\alpha _j}(xt_j)}{t_{\mathrm{}}}=0,\mathrm{}j,$$
we obtain for $`x๐ฑ_j\{t_j\}`$
$$T_{\mathrm{}}(x,๐)=\frac{Y_{\alpha _j}}{t_{\mathrm{}}}(x,๐)=\frac{H_j}{t_{\mathrm{}}}(x,๐)E_{\alpha _j}(xt_j)=\frac{H_j}{t_{\mathrm{}}}(x,๐)H_j^1(x,๐)Y_{\alpha _j}(x,๐).$$
Hence
$$F_{\mathrm{}}(x,๐)=\frac{H_j}{t_{\mathrm{}}}(x,๐)H_j^1(x,๐)+\frac{Q_{\mathrm{}}(๐)}{xt_{\mathrm{}}},x๐ฑ_j\{t_j\},\mathrm{}j,$$
which proves the statement $`()`$ in the case $`\mathrm{}j`$.
Next we differentiate the equality (3.11) with respect to $`t_j`$. Since
(3.13)
$$\frac{E_{\alpha _j}(xt_j)}{x}=\frac{E_{\alpha _j}(xt_j)}{t_j}=\frac{dE_{\alpha _j}(\zeta )}{d\zeta }|_{\zeta =xt_j},$$
we obtain
$$\begin{array}{c}T_j(x,๐)=\frac{H_j}{t_j}(x,๐)H_j^1(x,๐)\hfill \\ \hfill H_j(x,๐)\frac{E_{\alpha _j}(xt_j)}{x}E_{\alpha _j}^1(xt_j)H_j^1(x,๐),x\mathrm{cov}(๐ฑ_j\{t_j\};p_j).\end{array}$$
On the other hand, differentiating the equality 3.11 with respect to $`x`$, we get
$$\begin{array}{c}\frac{Y_{\alpha _j}}{x}(x,๐)Y_{\alpha _j}^1(x,๐)=\frac{H_j}{x}(x,๐)H_j^1(x,๐)+\hfill \\ \hfill +H_j(x,๐)\frac{E_{\alpha _j}(xt_j)}{x}E_{\alpha _j}^1(xt_j)H_j^1(x,๐),x\mathrm{cov}(๐ฑ_j\{t_j\};p_j),\end{array}$$
hence
$$\begin{array}{c}F_j(x,๐)=\left(\frac{H_j}{x}(x,๐)+\frac{H_j}{t_j}(x,๐)\right)H_j^1(x,๐)+\hfill \\ \hfill +\frac{Q_j(๐)}{xt_j}\frac{Y_{\alpha _j}}{x}(x,๐)Y_{\alpha _j}^1(x,๐),x\mathrm{cov}(๐ฑ_j\{t_j\};p_j).\end{array}$$
Taking into account that $`Y_{\alpha _j}(x,๐)`$ is a branch of the solution $`Y(x,๐)`$ of the Fuchsian system (3.1), we obtain
$$F_j(x,๐)=\left(\frac{H_j}{x}(x,๐)+\frac{H_j}{t_j}(x,๐)\right)H_j^1(x,๐)\underset{\begin{array}{c}1kn\\ kj\end{array}}{}\frac{Q_k(๐)}{xt_k},x๐ฑ_j\{t_j\}.$$
Thus the function $`F_j(x,๐)`$ has at $`x=t_j`$ a removable singularity, which proves the statement $`()`$ in the case $`j=\mathrm{}`$. โ
###### Remark 3.9.
Note that the relations (3.12), (3.13), which are instrumental in the proof of Proposition 3.6, are precisely the relations (2.10) in the informal definition of the isoprincipal family of Fuchsian systems in Section 2.4.
###### Remark 3.10.
We observe that the auxiliary linear system (3.6) leads to a linear system for the regular factor $`H_j(x,๐)`$ of the regular-principal factorization (3.11).
Indeed, in view of (3.12) and (3.13), we can differentiate the equality (3.11) to obtain
$`{\displaystyle \frac{H_j}{t_{\mathrm{}}}}(x,๐)H_j^1(x,๐)`$ $`={\displaystyle \frac{Y_{\alpha _j}}{t_{\mathrm{}}}}(x,๐)Y_{\alpha _j}^1(x,๐),j\mathrm{},`$
$`\left({\displaystyle \frac{H_j}{t_j}}(x,๐)+{\displaystyle \frac{H_j}{x}}(x,๐)\right)H_j^1(x,๐)`$ $`=\left({\displaystyle \frac{Y_{\alpha _j}}{t_j}}(x,๐)+{\displaystyle \frac{Y_{\alpha _j}}{x}}(x,๐)\right)Y_{\alpha _j}^1(x,๐),`$
and therefore
(3.14)
$$\{\begin{array}{cc}\hfill \frac{H_j}{t_{\mathrm{}}}& =\frac{Q_{\mathrm{}}}{xt_{\mathrm{}}}H_j,\mathrm{}j,\hfill \\ \hfill \frac{H_j}{t_j}+\frac{H_j}{x}& =\underset{\begin{array}{c}1\mathrm{}n\\ \mathrm{}j\end{array}}{}\frac{Q_{\mathrm{}}}{xt_{\mathrm{}}}H_j.\hfill \end{array}$$
Using the change of variables
(3.15) $`\zeta `$ $`=xt_j,`$
(3.16) $`L_j(\zeta ,๐)`$ $`\stackrel{\text{def}}{=}H_j(\zeta +t_j,๐),`$
one can rewrite the system (3.14) in the following form:
(3.17)
$$\{\begin{array}{cc}\hfill \frac{L_j}{t_{\mathrm{}}}& =\frac{Q_{\mathrm{}}}{\zeta +t_jt_{\mathrm{}}}L_j,\mathrm{}j,\hfill \\ \hfill \frac{L_j}{t_j}& =\underset{\begin{array}{c}1\mathrm{}n\\ \mathrm{}j\end{array}}{}\frac{Q_{\mathrm{}}}{\zeta +t_jt_{\mathrm{}}}L_j.\hfill \end{array}$$
The system (3.17) is nothing more than the system (3.6) with the constraint
$$xt_j=\zeta =\text{const}.$$
Note that, although $`x=t_j`$ is a singularity of the Fuchsian system (3.1) and the auxiliary system (3.6), the right-hand side of the system (3.17) is holomorphic with respect to $`\zeta `$ at $`\zeta =0`$ (compare with Remark 1.2 in \[Ma1\]). This occurs because the function $`L_j(\zeta ,๐)`$, defined in (3.16), is holomorphic with respect to $`\zeta `$ and invertible at $`\zeta =0.`$
### 3.4. The Frobenius theorem
The auxiliary system (3.6), related to the isoprincipal family (3.1) is an overdetermined system of PDEs. The criterion for the existence of solution of such a system is known (see \[Na, Section 2.11\]) as the Frobenius theorem:
###### Theorem 3.11 (The Frobenius theorem for Pfaffian <sup>19</sup><sup>19</sup>19A Pfaffian system is a system of the form (3.18). systems).
Let $`๐_๐ฉ`$ and $`๐_๐ช`$ be domains in, respectively, $`^p`$ and $`^q`$. Consider the following system of PDEs:
(3.18)
$$\frac{๐}{\mu _j}=\mathit{\varphi }_๐(๐,๐),1jq,$$
where $`๐(๐)`$ is an unknown $`^p`$-valued function of the variable $`๐=(\mu _1,\mathrm{},\mu _q)^q`$ and $`\mathbf{\varphi }_๐ฃ(๐,๐)=(\varphi _{1,j}(๐,๐),\mathrm{},\varphi _{p,j}(๐,๐))`$, $`1jq,`$ are given $`^p`$-valued functions, holomorphic with respect to $`๐,๐`$ in the domain $`๐_๐ฉ\times ๐_๐ช.`$
Then the following statements (i) and (ii) are equivalent:
1. For every pair of points $`๐^\mathrm{๐}๐_๐,`$ $`๐^\mathrm{๐}๐_๐`$ there exists a solution $`๐(๐)`$ of the system (3.18), holomorphic in a neighborhood of $`๐^\mathrm{๐}`$ and satisfying the initial condition <sup>20</sup><sup>20</sup>20Note that, according to the uniqueness theorem for ordinary differential equations, such a solution $`๐(๐)`$ is necessarily unique.
$$๐(๐^\mathrm{๐})=๐^\mathrm{๐}.$$
2. The $``$-valued functions $`\varphi _{i,j}(๐,๐)`$ satisfy the equations
(3.19)
$$\begin{array}{c}\frac{\varphi _{i,j}(๐,๐)}{\mu _k}+\underset{1\mathrm{}p}{}\frac{\varphi _{i,j}(๐,๐)}{\lambda _{\mathrm{}}}\varphi _{\mathrm{},k}(๐,๐)=\hfill \\ \hfill =\frac{\varphi _{i,k}(๐,๐)}{\mu _j}+\underset{1\mathrm{}p}{}\frac{\varphi _{i,k}(๐,๐)}{\lambda _{\mathrm{}}}\varphi _{\mathrm{},j}(๐,๐),\\ \hfill 1ip,1j,kq,\end{array}$$
in the domain $`๐_๐\times ๐_๐`$.
###### Definition 3.12.
The condititon (3.19), formulated in the Frobenius theorem 3.11, is said to be the compatibility condition for the overdetermined system of PDEs (3.18). The overdetermined system of PDEs (3.18) which satisfies the compatibility condition (3.19) is said to be compatible.
###### Remark 3.13.
Note that the compatibility condition (3.19) can be obtained by substituting the equations (3.18) into the identity
$$\frac{^2\lambda _i}{\mu _j\mu _k}=\frac{^2\lambda _i}{\mu _k\mu _j},1ip,1j,kq.$$
Thus the statement (ii) of the Frobenius theorem (3.11) follows from the statement (i) immediately.
In what follows, we shall often deal with linear overdetermined systems of PDEs, depending on a parameter. Because of the global existence theorem for such systems, the following stronger version of the Frobenius theorem 3.11 holds in this case:
###### Theorem 3.14 (The Frobenius theorem for linear systems with a parameter).
Let $`๐_๐ช`$ and $`๐_๐ซ`$ be domains in, respectively, $`^q`$ and $`^r`$. Consider the following linear systems of PDEs:
(3.20)
$$\frac{๐}{\mu _j}=\mathrm{\Phi }_j(๐,๐)๐,1jq,$$
where $`๐๐_๐ช`$ is the โmainโ variable, $`๐๐_๐ซ`$ is a parameter and $`\mathrm{\Phi }_j(๐,๐)`$, $`1jq,`$ are $`๐_p`$-valued functions, holomorphic with respect to $`๐`$ in the domain $`๐_๐ช`$ for each fixed $`๐๐_๐ซ`$.
Then:
1. For each fixed value of the parameter $`๐`$, say $`๐=๐^\mathrm{๐}๐_๐`$, the linear system
(3.21)
$$\frac{๐}{\mu _j}=\mathrm{\Phi }_j(๐,๐^\mathrm{๐})๐,1jq,$$
has a fundamental solution <sup>21</sup><sup>21</sup>21In other words, an $`๐_p`$-valued solution $`๐(๐,๐^\mathrm{๐})`$, such that $`det๐(๐,๐^\mathrm{๐})0`$ for $`๐๐_๐`$. $`๐(๐,๐^\mathrm{๐})`$ if and only if the functions $`\mathrm{\Phi }_j(๐,๐^\mathrm{๐})`$ satisfy the equations
(3.22)
$$\begin{array}{c}\frac{\mathrm{\Phi }_j(๐,๐^\mathrm{๐})}{\mu _k}\frac{\mathrm{\Phi }_k(๐,๐^\mathrm{๐})}{\mu _j}=[\mathrm{\Phi }_k(๐,๐^\mathrm{๐}),\mathrm{\Phi }_j(๐,๐^\mathrm{๐})],\hfill \\ \hfill ๐๐_๐,1j,kq.\end{array}$$
2. If the functions $`\mathrm{\Phi }_j(๐,๐)`$ are jointly holomorphic with respect to $`๐,๐`$ in the domain $`๐_๐\times ๐_๐`$ and satisfy the equations
(3.23)
$$\begin{array}{c}\frac{\mathrm{\Phi }_j(๐,๐)}{\mu _k}\frac{\mathrm{\Phi }_k(๐,๐)}{\mu _j}=[\mathrm{\Phi }_k(๐,๐),\mathrm{\Phi }_j(๐,๐)],\hfill \\ \hfill ๐๐_๐,๐๐_๐,1j,kq,\end{array}$$
then for every point $`๐^\mathrm{๐}๐_๐`$ and every $`๐_p`$-valued function $`๐^\mathrm{๐}(๐),`$ holomorphic and invertible in $`๐_๐,`$ there exists a unique fundamental solution $`๐(๐,๐)`$ of the system (3.20), jointly holomorphic with respect to $`๐,๐`$ in the domain $`๐_๐\times ๐_๐`$ and satisfying the initial condition
(3.24)
$$๐(๐,๐)|_{๐=๐^\mathrm{๐}}=๐^\mathrm{๐}(๐).$$
### 3.5. Proof of Theorem 3.1
###### Proposition 3.15.
Let (3.1) be a holomorphic family of Fuchsian systems, where the residue matrices $`Q_j(๐ญ)`$ are holomorphic in a domain $`๐_{}^n`$ and satisfy (3.2). Assume that the family (3.1) is isoprincipal with the distinguished point $`x_0=\mathrm{}`$. Then the residues $`Q_1(๐ญ),\mathrm{},Q_n(๐ญ)`$ satisfy the Schlesinger system (3.3) in the domain $`๐`$.
###### Proof.
According to Proposition 3.6, the solution $`Y(x,๐)`$ of (3.1), (3.4) satisfies the auxiliary linear system (3.6). Therefore, $`Y(x,๐)`$ is the fundamental solution of the overdetermined linear system (3.6) with the initial condition
$$Y(\mathrm{},๐^\mathrm{๐})=I,$$
where $`๐^\mathrm{๐}`$ is an arbitrary fixed point in the domain $`๐`$.
Hence, in view of Theorem 3.14, the linear system (3.6) (which is a special case of the system (3.20) with $`q=n+1`$, $`r=0`$, $`๐(๐)=Y(x,๐)`$) is compatible. The compatibility condition (3.23) takes in this case the form of the following two equations:
(3.25) $`{\displaystyle \frac{}{x}}\left({\displaystyle \frac{Q_j}{xt_j}}\right)+{\displaystyle \underset{1in}{}}{\displaystyle \frac{}{t_j}}\left({\displaystyle \frac{Q_i}{xt_i}}\right)`$ $`={\displaystyle \underset{1in}{}}{\displaystyle \frac{[Q_i,Q_j]}{(xt_i)(xt_j)}},1jn,`$
(3.26) $`{\displaystyle \frac{}{t_i}}\left({\displaystyle \frac{Q_j}{xt_j}}\right){\displaystyle \frac{}{t_j}}\left({\displaystyle \frac{Q_i}{xt_i}}\right)`$ $`={\displaystyle \frac{[Q_j,Q_i]}{(xt_i)(xt_j)}},1i,jn.`$
Since the residues $`Q_j`$ do not depend on $`x`$, from the equation (3.25) we obtain
$$\underset{1in}{}\frac{Q_i}{t_j}\frac{1}{xt_i}=\underset{1in}{}\frac{[Q_i,Q_j]}{(xt_i)(xt_j)},1jn.$$
Both sides of the last equality are rational functions of $`x`$, which are holomorphic in the punctured sphere $`^1\{t_1,\mathrm{},t_n\}`$ and have simple poles at the points $`x=t_i`$, $`1in`$. Equating the residues at each pole $`x=t_i`$, $`1in`$, we obtain the equations (3.3).
The equation (3.26) leads in a similar way to the first of the equations (3.3) and, therefore, provides no additional information. โ
###### Proposition 3.16.
Let (3.1) be a holomorphic family of Fuchsian systems, where the residue matrices $`Q_j(๐ญ)`$ are holomorphic in a domain $`๐_{}^n`$ and satisfy (3.2). Assume that the residue matrices $`Q_j(๐ญ)`$ satisfy the Schlesinger system (3.3).
Then the family (3.1) is isoprincipal with the distinguished point $`x_0=\mathrm{}`$.
###### Proof.
Let us choose $`๐^\mathrm{๐}๐`$ and a nested pair of open cylindrical neighborhoods $`๐ฆ๐ฅ`$ of $`๐^\mathrm{๐}`$ in $`๐.`$ Let $`Y(x,๐)`$ be the solution of (3.1), (3.4) and for $`1jn`$ let $`\alpha _j`$ be a path from $`x_0=\mathrm{}`$ to a point $`p_j๐ฑ_j\overline{๐ฒ_j}`$.
Then, according to Theorem 1.7 the branch $`Y_{\alpha _j}(x,๐^\mathrm{๐})`$ of the solution $`Y(x,๐^\mathrm{๐})`$ in $`๐ฑ_j\{t^0\}`$, corresponding to the path $`\alpha _j`$, admits the regular-principal factorization
(3.27)
$$Y_{\alpha _j}(x,๐^\mathrm{๐})=H_j(x,๐^\mathrm{๐})P_{\alpha _j}(x,๐^\mathrm{๐}),x\mathrm{cov}(๐ฑ_j\{๐^\mathrm{๐}\};p_j),$$
where the factors $`H_j(x,๐^\mathrm{๐})`$ and $`P_{\alpha _j}(x,๐^\mathrm{๐})`$ are holomorphic with respect to $`x`$ and invertible in, respectively, $`๐ฑ_j`$ and $`\mathrm{cov}(\{t_j^0\};p_j).`$
ยฟFrom here we proceed in four steps:
#### Step 1
Firstly, we construct a function $`E_{\alpha _j}(\zeta )`$, holomorphic and invertible on the Riemann surace of $`\mathrm{ln}\zeta `$ and such that the principal factor $`P_{\alpha _j}(x,๐^\mathrm{๐})`$ is a transplant of $`E_{\alpha _j}(\zeta )`$ into $`\mathrm{cov}(\{t_j^0\};p_j)`$:
$$P_{\alpha _j}(x,๐^\mathrm{๐})=E_{\alpha _j}(xt_j^0).$$
To this end we choose some value $`\theta _j^0`$ of $`\mathrm{arg}(p_jt_j^0)`$ and consider the corresponding isomorphism from the Riemann surace of $`\mathrm{ln}\zeta `$ onto $`\mathrm{cov}(\{t_j^0\};p_j)`$ (see (1.16)):
(3.28)
$$\zeta \underset{\mathrm{arg}\left(p_jt_j^0\right)=\theta _j^0}{\overset{}{}}\zeta +t_j^0,\zeta \mathrm{cov}(\{0\};1),\zeta +t_j^0\mathrm{cov}(\{t_j^0\};p_j).$$
We set
(3.29)
$$E_{\alpha _j}(\zeta )\stackrel{\text{def}}{=}P_{\alpha _j}(\zeta +t_j^0,๐^\mathrm{๐}),$$
where $`\zeta +t_j^0`$ denotes the image in $`\mathrm{cov}(\{t_j^0\};p_j)`$ of the point $`\zeta \mathrm{cov}(\{0\};1)`$ under the isomorphism (3.28).
#### Step 2
Secondly, we construct the regular factor $`H_j(x,๐)`$ for the solution $`Y(x,๐)`$ of (3.1), (3.4). In view of Remark 3.10, we consider the overdetermined linear system
(3.30)
$$\{\begin{array}{ccc}\hfill \frac{L_j(\zeta ,๐)}{t_i}& =\frac{Q_i(๐)}{\zeta +t_jt_i}L_j(\zeta ,๐),\hfill & \hfill 1in,ij,\\ \hfill \frac{L_j(\zeta ,๐)}{t_j}& =\underset{\begin{array}{c}1in\\ ij\end{array}}{}\frac{Q_i(๐)}{\zeta +t_jt_i}L_j(\zeta ,๐),\hfill \end{array}$$
with the initial condition
(3.31)
$$L_j(\zeta ,๐)|_{๐=๐^\mathrm{๐}}=H_j(\zeta +t_j^0,๐^\mathrm{๐}),$$
where $`H_j(x,๐^\mathrm{๐})`$ is the regular factor in the factorization (3.27). Here $`\zeta `$ is a small parameter:
(3.32)
$$\begin{array}{c}|\zeta |<ฯต,\text{ where }ฯต>0\text{ is such that }\hfill \\ \hfill x_1๐ฒ_j\{x:|xx_1|<ฯต\}๐ฑ_j,1jn.\end{array}$$
We claim that the overdetermined linear system (3.30), depending on the parameter $`\zeta `$, is compatible (see Calculation 1 below).
Therefore, the system (3.30) has a solution $`L_j(\zeta ,๐)`$, satisfying the initial condition (3.31). This solution $`L_j(\zeta ,๐)`$ is jointly holomorphic in $`\zeta ,๐`$ for $`|\zeta |<ฯต,๐๐ฆ`$ and invertible.
We define the function $`H_j(x,๐)`$ by
(3.33)
$$H_j(x,๐)\stackrel{\text{def}}{=}L_j(xt_j,๐),๐๐ฆ,x๐ฑ_j$$
#### Step 3
Thirdly, we consider the product
(3.34)
$$Z_{\alpha _j}(\zeta ,๐)\stackrel{\text{def}}{=}L_j(\zeta ,๐)E_{\alpha _j}(\zeta ).$$
Note that, in view of (3.31), for $`๐=๐^\mathrm{๐}`$ we have
$$Z_{\alpha _j}(\zeta ,๐^\mathrm{๐})=Y_{\alpha _j}(\zeta +t_j^0,๐^\mathrm{๐}),$$
hence the function $`Z_{\alpha _j}(\zeta ,๐^\mathrm{๐})`$ satisfies the linear system
(3.35)
$$\frac{dZ_{\alpha _j}(\zeta ,๐^\mathrm{๐})}{d\zeta }=\underset{1in}{}\frac{Q_i(๐^\mathrm{๐})}{\zeta +t_j^0t_i^0}Z_{\alpha _j}(\zeta ,๐^\mathrm{๐}).$$
Also, as follows from (3.30), we have
(3.36) $`{\displaystyle \frac{Z_{\alpha _j}(\zeta ,๐)}{t_i}}`$ $`={\displaystyle \frac{Q_i(๐)}{\zeta +t_jt_i}}Z_{\alpha _j}(\zeta ,๐),1in,ij,`$
(3.37) $`{\displaystyle \frac{Z_{\alpha _j}(\zeta ,๐)}{t_j}}`$ $`={\displaystyle \underset{\begin{array}{c}1in\\ ij\end{array}}{}}{\displaystyle \frac{Q_i(๐)}{\zeta +t_jt_i}}Z_{\alpha _j}(\zeta ,๐).`$
Furthermore, the function $`Z_{\alpha _j}(\zeta ,๐)`$ can be shown (see Calculation 2 below) to satisfy with respect to $`\zeta `$ the linear system
(3.38)
$$\frac{Z_{\alpha _j}(\zeta ,๐)}{\zeta }=\underset{1in}{}\frac{Q_i(๐)}{\zeta +t_jt_i}Z_{\alpha _j}(\zeta ,๐).$$
We note that for each fixed $`๐๐ฆ`$ the linear differential system (3.38) with respect to $`\zeta `$ has no singularities in the punctured domain
$$๐ฑ_{t_j}\stackrel{\text{def}}{=}\{\zeta :\zeta +t_j๐ฑ_j\}\{0\},$$
hence its fundamental solution $`Z_{\alpha _j}(\zeta ,๐)`$ is holomorphic with respect to $`\zeta `$ on a universal covering surface over the domain $`๐ฑ_{t_j}`$. Therefore, the function
$$L_j(\zeta ,๐)=Z_{\alpha _j}(\zeta ,๐)E_{\alpha _j}^1(\zeta )$$
is holomorphic with respect to $`\zeta `$ and invertible on this universal covering surface. On the other hand, the function $`L_j(\zeta ,๐)`$ is also single-valued holomorphic with respect to $`\zeta `$ and invertible in the open disk $`\{\zeta :|\zeta |<ฯต\}`$ (see (3.32)). Hence the function $`L_j(\zeta ,๐)`$ is single-valued holomorphic with respect to $`\zeta `$ and invertible in the non-punctured domain $`๐ฑ_{t_j}\{0\}.`$
It follows that the function $`H_j(x,๐)`$, defined in (3.33), is holomorphic (single-valued) with respect to $`x`$ and invertible in the entire domain $`๐ฑ_j`$.
#### Step 4
Finally, we consider the coherent family of transplants $`\{E_{\alpha _j}(xt_j)\}_{t_j๐ฒ_j}`$ (see Definition 2.7), corresponding to the unique branch of $`\mathrm{arg}(p_jt_j)`$ which is continuous with respect to $`t_j`$ in $`๐ฒ_j`$ and takes the value $`\theta ^0`$, chosen at Step 1, at the point $`t_j=t_j^0`$.
For each $`๐๐ฆ`$ we define the function $`Y_{\alpha _j}(x,๐)`$ by
(3.39)
$$Y_{\alpha _j}(x,๐)\stackrel{\text{def}}{=}H_j(x,๐)E_{\alpha _j}(xt_j)=Z_{\alpha _j}(xt_j,๐),x\mathrm{cov}(๐ฑ_j\{t_j\};p_j).$$
Note that in view of (3.29), (3.31) this definition agrees for $`๐=๐^\mathrm{๐}`$ with (3.27) and the notation $`Y_{\alpha _j}(x,๐^\mathrm{๐})`$ for the branch of the solution $`Y(x,๐^\mathrm{๐})`$ in $`๐ฑ_j\{t_j^0\}`$, corresponding to the path $`\alpha _j`$. Now we show that for every $`๐ญ๐ฆ`$ the function $`Y_{\alpha _j}(x,๐ญ)`$, defined in (3.39), is the branch of the solution $`Y(x,๐ญ)`$ of (3.1), (3.4) in the punctured domain $`๐ฑ_j\{t_j\}`$, corresponding to the path $`\alpha _j`$.
First of all, in view of (3.38), the function $`Y_{\alpha _j}(x,๐)`$ satisfies with respect to $`x`$ the system (3.1):
$$\frac{Y_{\alpha _j}}{x}(x,๐)=\underset{1in}{}\frac{Q_i(๐)}{xt_i}Y_{\alpha _j}(x,๐).$$
Hence for each $`๐๐ฆ`$ the function $`Y_{\alpha _j}(x,๐)`$ can be analytically continued along the path $`\alpha _j`$ in the opposite direction: from $`p_j`$ to $`\mathrm{}`$. The value of this continuation at $`x=\mathrm{}`$ will be denoted by $`\widehat{Y}(\mathrm{},๐);`$ in particular it holds that
(3.40)
$$\widehat{Y}(\mathrm{},๐^\mathrm{๐})=Y(\mathrm{},๐^\mathrm{๐})=I.$$
Furthermore, the value $`\widehat{Y}(\mathrm{},๐)`$ can be considered as the initial value at the distinguished point $`x_0=\mathrm{}`$ for a fundamental solution $`\widehat{Y}(x,๐)`$ of the Fuchsian system (3.1), defined on the universal covering surface $`\mathrm{cov}(^1\{t_1,\mathrm{},t_n\};\mathrm{})`$. The function $`Y_{\alpha _j}(x,๐)`$ is the branch of this solution $`\widehat{Y}(x,๐)`$ in the punctured domain $`๐ฑ_j\{t_j\}`$, corresponding to the path $`\alpha _j`$.
Now we note that, in view of (3.36) โ (3.37),
$$\frac{\widehat{Y}}{t_i}(x,๐)=\frac{Q_i(๐)}{xt_i}\widehat{Y}(x,๐),1in;$$
in particular,
(3.41)
$$\frac{\widehat{Y}}{t_i}(x,๐)|_{x=\mathrm{}}=0,1in.$$
Combining (3.41) with (3.40), we observe that the solution $`\widehat{Y}(x,๐)`$ satisfies the initial condition
$$\widehat{Y}(x,๐)|_{x=\mathrm{}}=I,๐๐ฆ,$$
and by the uniqueness theorem for linear differential systems coincides with the fundamental solution $`Y(x,๐)`$ of (3.1), (3.4).
Thus the function $`Y_{\alpha _j}(x,๐)`$, defined in (3.39), is the branch of the solution $`Y(x,๐)`$ of (3.1), (3.4) in the punctured domain $`๐ฑ_j\{t_j\}`$, corresponding to the path $`\alpha _j`$. The equality (3.39) itself can now be considered as the regular-principal factorization of the branch $`Y_{\alpha _j}(x,๐)`$. In view of Definition 2.9, we conclude that the family (3.1) is isoprincipal with the distinguished point $`x_0=\mathrm{}`$.
In order to complete the proof, it remains to present the calculations, omitted in the above reasonings.
#### Calculation 1
We show that the overdetermined linear system (3.30), depending on the parameter $`\zeta `$, is compatible.
In this case the compatibility condition (3.23) of Theorem (3.14) is represented by the following two equalities:
(3.42a)
$$\begin{array}{c}\frac{}{t_j}\left(\frac{Q_i}{\zeta +t_jt_i}\right)+\underset{\begin{array}{c}1kn\\ kj\end{array}}{}\frac{}{t_i}\left(\frac{Q_k}{\zeta +t_jt_k}\right)=\hfill \\ \hfill =\underset{\begin{array}{c}1kn\\ kj\end{array}}{}\frac{[Q_k,Q_i]}{(\zeta +t_jt_k)(\zeta +t_jt_i)},1in,ij,\end{array}$$
and
(3.42b)
$$\begin{array}{c}\frac{}{t_k}\left(\frac{Q_i}{\zeta +t_jt_i}\right)\frac{}{t_i}\left(\frac{Q_k}{\zeta +t_jt_k}\right)=\hfill \\ \hfill =\frac{[Q_k,Q_i]}{(\zeta +t_jt_k)(\zeta +t_jt_i)},1i,kn,i,kj.\end{array}$$
The equality (3.42a) can be simplified as follows:
$$\begin{array}{c}\frac{Q_i}{t_j}\frac{1}{\zeta +t_jt_i}+\underset{\begin{array}{c}1kn\\ kj\end{array}}{}\frac{Q_k}{t_i}\frac{1}{\zeta +t_jt_k}=\hfill \\ \hfill =\underset{\begin{array}{c}1kn\\ ki,j\end{array}}{}\frac{[Q_k,Q_i]}{t_kt_i}\left(\frac{1}{\zeta +t_jt_k}\frac{1}{\zeta +t_jt_i}\right).\end{array}$$
In view of (3.3), the right-hand side of the last equality can be rewritten as
$$\begin{array}{c}\underset{\begin{array}{c}1kn\\ ki,j\end{array}}{}\frac{Q_k}{t_i}\left(\frac{1}{\zeta +t_jt_k}\frac{1}{\zeta +t_jt_i}\right)=\hfill \\ \hfill =\frac{Q_j}{t_i}\frac{1}{\zeta +t_jt_i}+\underset{\begin{array}{c}1kn\\ kj\end{array}}{}\frac{Q_k}{t_i}\frac{1}{\zeta +t_jt_k},\end{array}$$
where we have used the fact that, in view of Remark 3.2,
$$\underset{1kn}{}\frac{Q_k}{t_i}=0,1in.$$
Since, as follows from (3.3),
$$\frac{Q_i}{t_j}=\frac{Q_j}{t_i},1i,jn,ij,$$
we conclude that the equality (3.42a), indeed, holds. The equality (3.42b) can be verified analogously.
#### Calculation 2
We show that the function $`Z_{\alpha _j}(\zeta ,๐)`$, defined in (3.34), satisfies with respect to $`\zeta `$ the equation (3.38).
This can be done as follows. We consider the auxiliary function
(3.43)
$$X_{\alpha _j}(\zeta ,๐)\stackrel{\text{def}}{=}\frac{Z_{\alpha _j}}{\zeta }(\zeta ,๐)\underset{1in}{}\frac{Q_i(๐)}{\zeta +t_jt_i}Z_{\alpha _j}(\zeta ,๐).$$
It will be shown below that $`X_{\alpha _j}(\zeta ,๐)`$ satisfies the linear system
(3.44)
$$\{\begin{array}{ccc}\hfill \frac{X_{\alpha _j}}{t_i}& =\frac{Q_i}{\zeta +t_jt_i}X_{\alpha _j},\hfill & \hfill 1in,ij,\\ \hfill \frac{X_{\alpha _j}}{t_j}& =\underset{\begin{array}{c}1in\\ ij\end{array}}{}\frac{Q_i}{\zeta +t_jt_i}X_{\alpha _j}.\hfill \end{array}$$
(Note that this system is the same as the system (3.30) for the function $`L_j`$.)
Since, in view of (3.35), the solution $`X_{\alpha _j}(\zeta ,๐)`$ of the linear system (3.44) satisfies the initial condition
$$X_{\alpha _j}(\zeta ,๐)|_{๐=๐^\mathrm{๐}}=0,$$
the uniqueness theorem for linear differential systems implies
$$X(\zeta ,๐)0.$$
Therefore, the function $`Z_{\alpha _j}(\zeta ,๐)`$ satisfies the equation (3.38).
Now let us prove that $`X_{\alpha _j}(\zeta ,๐)`$ satisfies, indeed, the linear system (3.44). In view of (3.37), we have
$$\begin{array}{c}\frac{X_{\alpha _j}}{t_j}=\underset{\begin{array}{c}1in\\ ij\end{array}}{}Q_i\frac{}{\zeta }\left(\frac{Z_{\alpha _j}}{\zeta +t_jt_i}\right)\underset{1in}{}\frac{}{t_j}\left(\frac{Q_i}{\zeta +t_jt_i}\right)Z_{\alpha _j}\hfill \\ \hfill \underset{\begin{array}{c}1i,kn\\ kj\end{array}}{}\frac{Q_iQ_k}{(\zeta +t_jt_i)(\zeta +t_jt_k)}Z_{\alpha _j}\end{array}$$
$$\begin{array}{c}=\underset{\begin{array}{c}1in\\ ij\end{array}}{}\frac{Q_i}{\zeta +t_jt_i}\frac{Z_{\alpha _j}}{\zeta }\underset{1in}{}\frac{Q_i}{t_j}\frac{1}{\zeta +t_jt_i}Z_{\alpha _j}\hfill \\ \hfill \underset{\begin{array}{c}1i,kn\\ ij\end{array}}{}\frac{Q_kQ_i}{(\zeta +t_jt_i)(\zeta +t_jt_k)}Z_{\alpha _j}.\end{array}$$
Now we substitute (3.3) to obtain
$$\begin{array}{c}\frac{X_{\alpha _j}}{t_j}=\underset{\begin{array}{c}1in\\ ij\end{array}}{}\frac{Q_i}{\zeta +t_jt_i}\frac{Z_{\alpha _j}}{\zeta }\underset{\begin{array}{c}1in\\ ij\end{array}}{}\frac{Q_iQ_j}{\zeta (\zeta +t_jt_i)}Z_{\alpha _j}\hfill \\ \hfill \underset{\begin{array}{c}1i,kn\\ i,kj\end{array}}{}\frac{Q_iQ_k}{(\zeta +t_jt_i)(\zeta +t_jt_k)}Z_{\alpha _j}\\ \hfill =\left(\underset{\begin{array}{c}1in\\ ij\end{array}}{}\frac{Q_i}{\zeta +t_jt_i}\right)\left(\frac{Z_{\alpha _j}}{\zeta }\underset{1kn}{}\frac{Q_k}{\zeta +t_jt_k}Z_{\alpha _j}\right)\\ \hfill =\underset{\begin{array}{c}1in\\ ij\end{array}}{}\frac{Q_i}{\zeta +t_jt_i}X_{\alpha _j}.\end{array}$$
The first equation of the system (3.44) is obtained analogously. โ
## 4. Isomondromic and isoprincipal deformations of Fuchsian systems
###### Definition 4.1.
Let a Fuchsian system
(4.1)
$$\frac{dY}{dx}=\left(\underset{1jn}{}\frac{Q_j^0}{xt_j^0}\right)Y,$$
where $`๐^\mathrm{๐}=(t_1^0,\mathrm{},t_n^0)_{}^n`$, $`Q_1^0,\mathrm{},Q_n^0๐_k`$ and
(4.2)
$$\underset{1jn}{}Q_j^0=0,$$
be given.
Let a holomorphic family of Fuchsian systems
(4.3)
$$\frac{dY}{dx}=\left(\underset{1jn}{}\frac{Q_j(๐)}{xt_j}\right)Y,$$
where the residue matrices $`Q_j(๐)`$ are holomorphic in a neighborhood $`๐_{}^n`$ of $`๐^\mathrm{๐}`$, be such that:
(4.4) $`{\displaystyle \underset{1jn}{}}Q_j(๐)`$ $`0,`$ $`๐`$ $`๐,`$
(4.5) $`Q_j(๐^\mathrm{๐})`$ $`=Q_j^0,`$ $`1j`$ $`n.`$
If the holomorphic family of Fuchsian systems (4.3) is isoprincipal (respectively, isomonodromic) with the distinguished point $`x_0=\mathrm{}`$, then it is said to be an isoprincipal (respectively, isomonodromic) deformation with the distinguished point $`x_0=\mathrm{}`$ of the Fuchsian system (4.1).
###### Remark 4.2.
Note that, according to Theorem 2.11, an isoprincipal deformation (4.3) of a given Fuchsian system (4.1) is also an isomonodromic one. As Theorem 2.15 implies, the converse is true under the condition that all the matrices $`Q_1^0,\mathrm{},Q_n^0`$ are non-resonant.
Now we are going to address the question of the existence of an isoprincipal deformation of a given Fuchsian system. It follows from Theorem 3.1 and Remark 3.2 that the holomorphic family (4.3) is an isoprincipal deformation with the distinguished point $`x_0=\mathrm{}`$ of the Fuchsian system (4.1) if and only if the residues $`Q_j(๐)`$ satisfy the Schlesinger system
(4.6)
$$\{\begin{array}{ccc}\hfill \frac{Q_i}{t_j}& =\frac{[Q_i,Q_j]}{t_it_j},\hfill & \hfill 1i,jn,ij,\\ \hfill \frac{Q_i}{t_i}& =\underset{\begin{array}{c}1jn\\ ji\end{array}}{}\frac{[Q_i,Q_j]}{t_it_j},\hfill & \hfill 1in.\end{array}$$
Thus the question is, whether the Cauchy problem for the Schlesinger system with the initial condition (4.5) is solvable. An answer to this question follows from the Frobenius theorem (Theorem 3.11):
###### Proposition 4.3.
Let $`๐ญ^\mathrm{๐}=(t_1^0,\mathrm{},t_n^0)_{}^n`$, $`Q_1^0,\mathrm{},Q_n^0๐_k`$ be given.
Then there exist a neighborhood $`๐_{}^n`$ of $`๐ญ^\mathrm{๐}`$ and unique matrix functions $`Q_1(๐ญ),\mathrm{},Q_n(๐ญ)`$, holomorphic in $`๐`$, which satisfy the Schlesinger system (4.6) and the initial condition (4.5).
###### Proof.
According to Theorem 3.11 and Remark 3.13, we have to to verify the identity
$$\frac{^2Q_i}{t_jt_k}=\frac{^2Q_i}{t_kt_j},1i,j,kn,$$
substituting the expressions (4.6) for the partial derivatives of $`Q_i.`$
In the case $`i=kj`$ we have
$$\begin{array}{c}\frac{^2Q_i}{t_it_j}=\frac{}{t_j}\left(\underset{\begin{array}{c}1kn\\ ki\end{array}}{}\frac{[Q_j,Q_i]}{t_it_j}\right)\hfill \\ \hfill =\frac{[Q_i,[Q_i,Q_j]][Q_i,Q_j]}{(t_it_j)^2}+\underset{\begin{array}{c}1kn\\ ki\end{array}}{}\frac{[Q_i,[Q_j,Q_k]]+[Q_k,[Q_i,Q_j]]}{(t_it_j)(t_it_k)}\\ \hfill =\frac{[Q_i,[Q_i,Q_j]][Q_i,Q_j]}{(t_it_j)^2}\underset{\begin{array}{c}1kn\\ ki\end{array}}{}\frac{[Q_j,[Q_k,Q_i]]}{(t_it_j)(t_it_k)}\\ \hfill =\frac{}{t_i}\left(\frac{[Q_i,Q_j]}{t_it_j}\right)=\frac{^2Q_i}{t_jt_i},\end{array}$$
where we have used the Jacobi identity
$$[A,[B,C]]+[B,[C,A]]+[C,[A,B]]=0,A,B,C๐_k.$$
In the case $`kij`$ the computations are analogous and will be omitted. โ
###### Remark 4.4.
Proposition 4.3 was originally proved by L. Schlesinger in \[Sch2\]. Here we would also like to mention the paper \[Boa\] by P. Boalch, where some general considerations concerning the compatibility of systems from a class, containing the Schlesinger system, are presented.
As an immediate consequence of Theorem 3.1 and Proposition 4.3, we obtain
###### Theorem 4.5.
Let a Fuchsian system (4.1) โ (4.2) be given.
Then there exists a neighborhood $`๐_{}^n`$ of $`๐ญ^\mathrm{๐}`$ and a unique isoprincipal deformation (4.3) of the Fuchsian system (4.1) with the distinguished point $`x_0=\mathrm{}`$, where the residue matrices $`Q_j(๐ญ)`$ are holomorphic in the neighborhood $`๐`$.
###### Remark 4.6.
Theorem 4.5 can also be proved in another way.
Let $`๐ฆ๐ฅ`$ be a nested pair of open cylindrical neighborhoods of $`๐^\mathrm{๐}`$ in $`_{}^n`$ and for $`1jn`$ let $`\alpha _j`$ be a path from $`x_0=\mathrm{}`$ to a point $`p_j๐ฑ_j\overline{๐ฒ_j}`$. Let $`Y(x,๐^\mathrm{๐})`$ be the solution of (4.1) with the initial condition
$$Y(x,๐^\mathrm{๐})|_{x=\mathrm{}}=I$$
and let
$$Y_{\alpha _j}(x,๐^\mathrm{๐})=H_j(x,๐^\mathrm{๐})E_{\alpha _j}(xt_j^0),x\mathrm{cov}(๐ฑ_j\{t_j^0\};p_j)$$
be the regular-principal factorization of the appropriate branch $`Y_{\alpha _j}(x,๐^\mathrm{๐})`$ in the punctured domain $`๐ฑ_j\{t_j^0\}`$ (here the function $`E_{\alpha _j}(\zeta )`$, holomorphic and invertible in the Riemann surface of $`\mathrm{ln}\zeta `$, is defined in the same way as in the proof of Proposition 3.16 โ see (3.29)).
First we assume that the family (4.3), where the residue matrices $`Q_j(๐)`$ are holomorphic in $`๐ฆ`$, is an isoprincipal deformation of the Fuchsian system (4.1) with the distinguished point $`x_0=\mathrm{}`$ and consider the solution $`Y(x,๐)`$ of (4.3) with the initial condition
(4.7)
$$Y(x,๐)|_{x=\mathrm{}}=I.$$
Since $`Y(x,๐)`$ is jointly holomorphic in $`x`$ and $`๐`$ and invertible in the domain $`\mathrm{cov}(^1_k\overline{๐ฒ_k};\mathrm{})\times ๐ฆ`$, so is the โratioโ
$$R(x,๐)\stackrel{\text{def}}{=}Y(x,๐)Y^1(x,๐^\mathrm{๐}).$$
Moreover, taking into account the initial condition (4.7) and the fact that the family (4.3) is by Theorem 2.11 isomonodromic, we reach the following conclusion:
* The function $`R(x,๐)`$ is holomorphic (single-valued) with respect to $`x,๐`$ and invertible in the domain $`\left\{^1_k\overline{๐ฒ_k}\right\}\times ๐ฆ`$ and it holds that
$$R(x,๐)|_{x=\mathrm{}}I,๐๐ฆ.$$
Similarly, considering the coherent family of transplants $`\{E_{\alpha _j}(xt_j)\}_{t_j๐ฒ_j},`$ we observe that the โratioโ
$$F_j(x,๐)\stackrel{\text{def}}{=}E_{\alpha _j}(xt_j)E_{\alpha _j}^1(xt_j^0),x\mathrm{cov}(\overline{๐ฒ_j};p_j),๐๐ฆ$$
is holomorphic (single-valued) with respect to $`x,๐`$ and invertible in $`\left\{\overline{๐ฒ_j}\right\}\times ๐ฆ.`$
Since the family (4.3) is assumed to be isoprincipal, there exist functions $`H_j(x,๐)`$ (the regular factors of $`Y(x,๐)`$), such that:
* Each function $`H_j(x,๐)`$, $`1jn`$, is holomorphic (single-valued) with respect to $`x,๐`$ and invertible in $`๐ฑ_j\times ๐ฆ;`$
* $$F_j(x,๐)H_j^1(x,๐^\mathrm{๐})=H_j^1(x,๐)R(x,๐),๐๐ฆ,x๐ฑ_j\overline{๐ฒ_j},1jn.$$
Note that the function $`F_j(x,๐)H_j^1(x,๐^\mathrm{๐})`$ on the left-hand side of the equality (Pb) is holomorphic (single-valued) with respect to $`x,๐`$ and invertible in $`\left\{\overline{๐ฒ_j}\right\}\times ๐ฆ`$. It is determined entirely in terms of the initial data $`๐^\mathrm{๐},Q_1^0,\mathrm{},Q_n^0`$. The equality (Pb) itself can be viewed as a factorization problem for this function, where one looks for the factors $`H_j(x,๐)`$, $`R(x,๐)`$, possessing the analyticity properties (H) and (R). If such factors can be found, then, reversing the reasonings above, one arrives at the isoprincipal deformation of the Fuchsian system (4.1).
The uniqueness of the solution, possessing the properties (H) and (R), for the factorization problem (Pb) follows immediately from the Liouville theorem. The existence of this solution for $`๐`$ in a sufficiently small neighborhood $`๐ฆ`$ of $`๐^\mathrm{๐}`$ can be established by elementary means, since for $`๐=๐^\mathrm{๐}`$ the solution exists (it is trivial: $`R=I`$).
However, analyzing the factorization problem (Pb) more carefully and systematically (see, for instance, \[BIK, ยง5\]) and taking into account Theorem 3.1, one can reach the stronger conclusion that the functions $`Q_j(๐),`$ which satisfy the Schlesinger system (4.6) with the initial condition (4.5), are meromorphic in the universal covering space $`\mathrm{cov}(_{}^n;,๐^\mathrm{๐}).`$ This result was obtained by B. Malgrange in \[Ma1\] and (in the non-resonant case) by T. Miwa in \[Miwa\]. Our proof, which involves the isoprincipal deformations and is outlined above, will be presented in more detail elsewhere.
###### Remark 4.7.
As was already mentioned (see Remarks 1.10 and 2.16), most of the considerations of the present article need not be restricted to the case of linear differential systems with only Fuchsian singularities.
For instance, one can consider linear systems of the form
$$\frac{dY}{dx}=\left(\underset{j=1}{\overset{n}{}}\underset{k=0}{\overset{p_j}{}}\frac{Q_{j,k}}{(xt_j)^{k+1}}\right)Y,$$
where
$$\underset{j=1}{\overset{n}{}}Q_{j,0}=0.$$
In this case the regular-principal factorization of the solution and the notion of the isoprincipal family
$$\frac{dY}{dx}=\left(\underset{j=1}{\overset{n}{}}\underset{k=0}{\overset{p_j}{}}\frac{Q_{j,k}(๐)}{(xt_j)^{k+1}}\right)Y$$
can be introduced as in Definitions 1.8, 2.9.
Furthermore, the auxiliary linear system related to the isoprincipal family takes the form (compare with (3.5))
$$\{\begin{array}{cc}\hfill \frac{Y}{x}& =\left(\underset{j=1}{\overset{n}{}}\underset{k=0}{\overset{p_j}{}}\frac{Q_{j,k}(๐)}{(xt_j)^{k+1}}\right)Y,\hfill \\ \hfill \frac{Y}{t_j}& =\left(\underset{k=0}{\overset{p_j}{}}\frac{Q_{j,k}(๐)}{(xt_j)^{k+1}}\right)Y,1jn.\hfill \end{array}$$
The compatibility condition for this overdetermined system is given by the system
$$\{\begin{array}{cc}\hfill \frac{Q_{i,k}}{t_j}& =\underset{\begin{array}{c}0\mathrm{}p_ik\\ 0qp_j\end{array}}{}(1)^l\left(\genfrac{}{}{0pt}{}{l+q}{l}\right)\frac{[Q_{i,k+l},Q_{j,q}]}{(t_it_j)^{q+l+1}},0kp_i,1i,jn,ij,\hfill \\ \hfill \frac{Q_{i,k}}{t_i}& =\underset{\begin{array}{c}1jn\\ ji\end{array}}{}\underset{\begin{array}{c}0\mathrm{}p_ik\\ 0qp_j\end{array}}{}(1)^l\left(\genfrac{}{}{0pt}{}{l+q}{l}\right)\frac{[Q_{i,k+l},Q_{j,q}]}{(t_it_j)^{q+l+1}},0kp_i,1in,\hfill \end{array}$$
which itself is compatible and contains the Schlesinger system 4.6 as a special case (corresponding to $`p_1=\mathrm{}=p_n=0`$).
## 5. An example
In order to illustrate the distinction between the isoprincipal and the isomonodromic deformations of Fuchsian systems in the case when the non-resonance condition is violated, we present the following explicit example.
Let us consider the linear differential system
(5.1)
$$\frac{dY}{dx}=\left(\begin{array}{cc}\frac{1}{x(x1)}& 0\\ 0& \frac{1}{(x2)(x3)}\end{array}\right)Y.$$
Note that the system (5.1) is of the form
$$\frac{dY}{dx}=\left(\frac{Q_0^0}{x}+\frac{Q_1^0}{x1}+\frac{Q_2^0}{x2}+\frac{Q_3^0}{x3}\right)Y,$$
where
$$Q_0^0=Q_1^0=\left(\begin{array}{cc}1& 0\\ 0& 0\end{array}\right),Q_2^0=Q_3^0=\left(\begin{array}{cc}0& 0\\ 0& 1\end{array}\right),$$
and hence it is Fuchsian and resonant. Moreover, $`x=\mathrm{}`$ is a regular point of the system (5.1), since
$$Q_0^0+Q_1^0+Q_2^0+Q_3^0=0.$$
The solution $`Y(x)`$ satisfying the initial condition $`Y_|\mathrm{}=I`$ is a rational matrix functions (in particular, single-valued); it has the following form:
$$Y(x)=\left(\begin{array}{cc}\frac{x}{x1}& 0\\ 0& \frac{x2}{x3}\end{array}\right)\stackrel{\text{def}}{=}Y^0(x).$$
The principal factors of $`Y^0`$ can be chosen <sup>22</sup><sup>22</sup>22See Remark 1.9. in the form
$`P_0(x)`$ $`=\left(\begin{array}{cc}x& 0\\ 0& 1\end{array}\right),`$ $`P_1(x)`$ $`=\left(\begin{array}{cc}{\displaystyle \frac{1}{x1}}& 0\\ 0& 1\end{array}\right),`$
$`P_2(x)`$ $`=\left(\begin{array}{cc}1& 0\\ 0& x2\end{array}\right),`$ $`P_3(x)`$ $`=\left(\begin{array}{cc}1& 0\\ 0& {\displaystyle \frac{1}{x3}}\end{array}\right).`$
Now we are going to construct the isoprincipal deformation of the system (5.1). For simplicity, we assume that the singular points $`x=1,2,3`$ are fixed while the position of one singularity (x=0) varies: $`x=t`$. Then we need to determine a holomorphic family of matrix functions $`Y(x,t)`$ such that
(5.2)
$$Y(x,t)|_{x=\mathrm{}}=I,Y(x,t)|_{t=0}=Y^0(x),$$
$`P_0(xt)`$ is the principal factor of $`Y(x,t)`$ at $`x=t`$ and for $`k=1,2,3`$ $`P_k(x)`$ is the principal factor of $`Y(x,t)`$ at $`x=k`$. Such a family is unique and consists of rational matrix functions of the form:
$$Y(x,t)=\left(\begin{array}{cc}\frac{xt}{x1}& 0\\ 0& \frac{x2}{x3}\end{array}\right).$$
The function $`Y(x,t)`$ satisfies with respect to $`x`$ the Fuchsian system
(5.3)
$$\begin{array}{c}\frac{dY}{dx}=\left(\begin{array}{cc}\frac{t1}{(xt)(x1)}& 0\\ 0& \frac{1}{(x2)(x3)}\end{array}\right)Y\hfill \\ \hfill =\left(\frac{Q_0^0}{xt}+\frac{Q_1^0}{x1}+\frac{Q_2^0}{x2}+\frac{Q_3^0}{x3}\right)Y,\end{array}$$
and the constant functions
$$Q_j(t)Q_j^0,j=0,1,2,3,$$
satisfy, of course, the Schlesinger system
$$\{\begin{array}{cc}\hfill \frac{dQ_j}{dt}& =\frac{[Q_0,Q_j]}{tj},j=1,2,3,\hfill \\ \hfill \frac{dQ_0}{dt}& =\underset{j=1}{\overset{3}{}}\frac{[Q_j,Q_0]}{tj},\hfill \end{array}$$
with the initial condition
$$Q_j(0)=Q_j^0,j=0,1,2,3.$$
However, the deformation (5.3) is not the only possible isomonodromic deformation of the system (5.1). Indeed, let us consider a family of rational functions
$$Y(x,t)=\left(\begin{array}{cc}\frac{xt}{x1}& \frac{2t(xt)h(t)}{(x1)(x3)(t3)}\\ 0& \frac{x2}{x3}\end{array}\right),$$
where $`h(t)`$ is a function, holomorphic in a neighborhood of $`t=0`$ (outside this neighborhood $`h(t)`$ may have arbitrary singularities). Such a family also satisfies the normalizing conditions (5.2); moreover, we have
$`Y(x,t)^1`$ $`=\left(\begin{array}{cc}{\displaystyle \frac{x1}{xt}}& {\displaystyle \frac{2th(t)}{(x2)(t3)}}\\ 0& {\displaystyle \frac{x3}{x2}}\end{array}\right),`$
$`{\displaystyle \frac{Y(x,t)}{x}}Y(x,t)^1`$ $`=\left(\begin{array}{cc}{\displaystyle \frac{t1}{(xt)(x1)}}& {\displaystyle \frac{2t(xt)h(t)}{(x1)(x2)(x3)(t3)}}\\ 0& {\displaystyle \frac{1}{(x2)(x3)}}\end{array}\right).`$
Thus we obtain the following deformation of the system (5.1):
(5.4)
$$\begin{array}{c}\frac{dY}{dx}=\left(\begin{array}{cc}\frac{t1}{(xt)(x1)}& \frac{2t(xt)h(t)}{(x1)(x2)(x3)(t3)}\\ 0& \frac{1}{(x2)(x3)}\end{array}\right)Y\hfill \\ \hfill =\left(\frac{Q_0(t)}{xt}+\frac{Q_1(t)}{x1}+\frac{Q_2(t)}{x2}+\frac{Q_3(t)}{x3}\right)Y,\end{array}$$
where
$$Q_0(t)Q_0^0,Q_1(t)=\left(\begin{array}{cc}1& \frac{t(t1)h(t)}{t3}\\ 0& 0\end{array}\right)$$
$$Q_2(t)=\left(\begin{array}{cc}0& \frac{2t(t2)h(t)}{t3}\\ 0& 1\end{array}\right),Q_3(t)=\left(\begin{array}{cc}0& th(t)\\ 0& 1\end{array}\right).$$
Since the monodromy of $`Y(x,t)`$ for every $`t`$ is trivial, the deformation (5.4) is isomonodromic, but if $`h(t)0`$, then it is not isoprincipal and the coefficients $`Q_j(t)`$ do not satisfy the Schlesinger system. Moreover, since we only require $`h(t)`$ to be holomorphic in a neighborhood of $`t=0`$, the behavior of the functions $`Q_j(t)`$ outside this neighborhood may be arbitrary. In particular, these functions $`Q_j(t)`$ need not be meromorphic with respect to $`t`$ in $`\{1,2,3\}`$.
###### Remark 5.1.
The example presented above is based on the theory of holomorphic families (4.3) of Fuchsian systems, whose solutions $`Y(x,๐)`$ are generic rational matrix functions of $`x`$ for every fixed $`๐`$. This theory was developed in \[Kats1\], \[Kats2\] and \[KaVo2\] (see also the electornic version of the latter work \[KaVo2-e\].) For such Fuchsian systems the number of poles $`n`$ is even, so we write $`2n`$ instead of $`n`$. All such Fuchsian systems with this property can be parameterized as follows. If $`k`$ is the dimension of the square residue matrices $`Q_j`$, then to each pole $`t_j`$ a $`k1`$-vector is related as a โfreeโ parameter. Therefore, to each $`2n`$-tuple $`๐=(t_1,\mathrm{},t_2n)`$ the total of $`(k1)\times 2n`$ complex parameters is related. To every choice of these parameters (satisfying a certain non-degeneracy condition) corresponds a different system of the form (4.3), whose solution $`Y(x,๐)`$ is a generic rational matrix function. Considering $`๐`$ as variable, we assign $`(k1)\times 2n`$ complex parameters to each $`๐_{}^{2n}`$.
Thus the families (4.3), possessing the property that $`Y(x,๐)`$ is a generic rational matrix function of $`x`$ for each fixed $`๐`$, can be parameterized by $`(k1)\times 2n`$ complex valued functions of $`๐`$, and these functional parameters are free. To obtain holomorphic families, we have to require that these complex valued functions are holomorphic. Of course, the monodromy of any rational matrix function of $`x`$ is trivial. Hence we can parameterize the class of the isomonodromic deformations (4.3) of Fuchsian systems with generic rational solutions by $`(k1)\times 2n`$ functional parameters, and these parameters can be arbitrary (up to the mentioned non-degeneracy condition) holomorphic functions of $`๐`$.
It turns out that the deformation corresponding to a given choice of the functional parameters is isoprincipal if and only if all these parameters are constant functions. In particular, the class of the isomonodromic deformations (considered without the non-resonance condition) is much richer than its subclass of the isoprincipal deformations.
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# Cosmological expansion and the uniqueness of gravitational action
## I Introduction
One of the great surprises in modern cosmology has been the accelerating expansion of the universe. The first sign of this unexpected behaviour came from distant supernova observations (SNIa) in 1998, when it was discovered that far away supernovae are dimmer than expected snia . The SNIa observations have since improved, notably by a number very distant supernova observations by the Hubble Space Telescope, and the main conclusion still stands: a critical density matter dominated universe is not consistent with the data. Such a universe, also known as the Einstein-deSitter (EdS) model, was also disfavoured by independent observations based on measuring the small fluctuations in the cosmic microwave background (CMB) cmb . Combining these two observations along with a third independent observation based on the clustering of galaxies on very large scales lss , the present day cosmological model emerges. The cosmological concordance model is a geometrically flat $`\mathrm{\Lambda }`$CDM model, which is a critical density universe dominated by cold dark matter and dark energy in the form of a cosmological constant. Furthermore, cross-correlating the observations of galaxy clusters with CMB measurements, one has seen evidence of the decay of gravitational potentials on very large scales isw , as predicted in the $`\mathrm{\Lambda }`$CDM model. Although independent measurements are statistically not very significant, combining measurements leads to a picture that is very well consistent with the concordance $`\mathrm{\Lambda }`$CDM model iswcomb .
A natural alternative to adding a mysterious new form of energy to our physical picture of the universe, is to consider modifications of general relativity. This is also motivated by the fact that we only have precision observations of gravity from sub-millimeter scales to solar system scales, which is very far from the present Hubble radius that is the scale relevant to the question of dark energy.
A particular class of models that has drawn a significant amount attention recently are the so-called $`f(R)`$ gravity models (see e.g. turner ; turner2 ; allemandi ; meng ; nojiri3 and references therein). They form a class of higher derivative gravity theories that include higher order curvature invariants in the gravitational action. Such an extension of the Einstein-Hilbert (EH) Lagrangian can be viewed natural as there is no a priori reason why the gravitational action should be linear in the Ricci scalar $`R`$. Furthermore, higher order terms can naturally appear in low energy effective Lagrangians of quantum gravity and string theory.
After $`f(R)`$ theories were proposed as an alternative solution to the dark energy problem cappo1 -nojiri2 , it was quickly realized that two major obstacles exist. Firstly, some of these theories exhibit instabilities that will render space-time unstable dolgov . A second serious but more practical hindrance has been the difficulty in actually integrating the equations of motion numerically and comparing to observational data.
Another twist in the tale is the question of frames and conformal transformations. While the Einstein and Jordan frames are mathematically equivalent on the classical level, the physical equivalence of the two frames at the perturbation and quantum level has been discussed in reference to $`f(R)`$ theories. Possible restrictions arise due to non-standard gravitational effects constrained by fifth force experiments (see confprobs and references therein).
Recently, an inverse approach to the problem of integrating the complicated equations of motion arising from a general $`f(R)`$ action has been proposed in cappo3 . Instead of specifying a particular form of $`f(R)`$ one considers the inverse problem of reconstructing $`f(R)`$ from the expansion history of the universe.
In this letter we are also concerned with the inverse problem of reconstructing $`f(R)`$, given an expansion history. We show that for any barotropic equation state, the functional form of $`f(R)`$ and hence the gravitational action, is not uniquely determined by the expansion history of the universe. Instead, for a given fluid one can always construct a class of gravitational actions that will have the same cosmological evolution as the EH action. As an example, we demonstrate this explicitly for cosmologically relevant solutions.
## II Basic formalism of $`f(R)`$ gravity
The action for $`f(R)`$ gravity is (with $`8\pi G=1`$) (see e.g. cappo3 )
$$S=d^4x\sqrt{g}\left(f(R)+_m\right),$$
(1)
form which it follows in the standard metric approach (as opposed to the so-called Palatini approach)
$$G_{\mu \nu }=R_{\mu \nu }\frac{1}{2}Rg_{\mu \nu }=T_{\mu \nu }^c+T_{\mu \nu }^m,$$
(2)
where
$`T_{\mu \nu }^c`$ $`=`$ $`{\displaystyle \frac{1}{f^{}(R)}}\{{\displaystyle \frac{1}{2}}g_{\mu \nu }(f(R)Rf^{}(R))+`$ (3)
$`+`$ $`f^{}(R)^{;\mu \nu }(g_{\alpha \mu }g_{\beta \nu }g_{\mu \nu }g_{\mu \nu })\}).`$
The standard minimally coupled stress-energy tensor, $`\stackrel{~}{T}_{\mu \nu }^m`$, from the matter Lagrangian, $`_m`$ in the action, Eq. (1), is related to $`T_{\mu \nu }^m`$ by $`T_{\mu \nu }^m=\stackrel{~}{T}_{\mu \nu }^m/f^{}(R)`$. The equations of motion arising from the action in a Friedmann-Robertson-Walker universe are
$$H^2+\frac{k}{a^2}=\frac{1}{3}\left(\rho _c+\frac{\rho _m}{f^{}(R)}\right)$$
(4)
and
$$2\dot{H}+3H^2+\frac{k}{a^2}=\left(p_c+\frac{p_m}{f^{}(R)}\right),$$
(5)
where the energy density and pressure of the curvature fluid are
$`\rho _c`$ $`=`$ $`{\displaystyle \frac{1}{f^{}(R)}}\left\{{\displaystyle \frac{1}{2}}\left(f(R)Rf^{}(R)\right)3H\dot{R}f^{\prime \prime }(R)\right\}`$
$`p_c`$ $`=`$ $`{\displaystyle \frac{1}{f^{}(R)}}\{\dot{R}^2f^{\prime \prime \prime }(R)+2H\dot{R}f^{\prime \prime }(R)+\ddot{R}f^{\prime \prime }(R)`$ (6)
$`{\displaystyle \frac{1}{2}}(f(R)Rf^{}(R))\}.`$
In addition to these, we also have a constraint equation for the curvature scalar
$$f^{\prime \prime }(R)\left\{R+6\left(\dot{H}+2H^2+\frac{k}{a^2}\right)\right\}=0$$
(7)
and we can write the continuity equation of the total fluid from Eqs. (4) and (5) as $`\dot{\rho }+3H(\rho +p)=0,`$ where $`\rho =\rho _c+\rho _m/f^{}(R)`$ and $`p=p_c+p_m/f^{}(R)`$.
From the Einsteinโs equations we see that due to the presence of the $`\ddot{R}`$ term in $`p_c`$, we have a fourth order differential equation system for the scale factor of $`a(t)`$. Given a gravitational theory, or given the form of $`f(R)`$, one could hope to solve the arising differential equation. However, even for simple choices of $`f(R)`$, such as $`f(R)=R+\mu ^4/R`$, the resulting differential equation (see eg. turner ) is complicated and one must resort to numerical methods. The numerical solutions are also problematic since one then needs information about the higher derivatives of the scale factor.
## III Uniqueness of the gravitational action
It is clear that for $`f(R)=R`$, Equations (4) and (5) reduce to the standard Friedmann equations of Einstein gravity, since then $`\rho _c=p_c=0`$. However, this does not guarantee that $`f(R)=R`$ is the only choice that reduces to the standard Friedmann equations but in general standard equations are reached if conditions
$`\rho _c+{\displaystyle \frac{\rho _m}{f^{}(R)}}`$ $`=`$ $`\rho _m,`$ (8)
$`p_c+{\displaystyle \frac{p_m}{f^{}(R)}}`$ $`=`$ $`p_m`$ (9)
hold. Assuming that the continuity equation holds independently for the ordinary matter fluid, $`\dot{\rho }_m+3H(\rho _m+p_m)=0`$, these two conditions are in fact equivalent.
The fact that one can recover the standard Friedmann equations is a consequence of higher derivatives in the Einsteinโs equations. The general Friedman equations are higher order differential equations than the standard equations. Thus the standard solutions can also be solutions of the general equations without violating the uniqueness of the solutions of differential equations. One can step even a little further and show the standard Friedmann equations supplemented with Eq. (8) are enough that both general Eqs. (4) and (5) are fulfilled.
When the ordinary Einsteinโs equations hold, the curvature scalar can be written as $`R=\rho _m+3p_m`$, so that for any barotropic equation of state (expect for the radiation dominated case), $`p_m=p_m(\rho _m)`$, the curvature scalar is a function of the matter density, $`R=R(\rho _m)`$. One can then write, at least formally, $`\rho _m=\rho _m(R)`$, e.g. for cold dark matter $`\rho _m=R`$. Therefore, the quantity $`H\dot{R}`$ appearing in the expression for $`\rho _c`$ can be written as $`H\dot{R}=(\rho _m3ka^2)(\rho _m+p_m)R^{}(\rho _m)`$. In a flat universe, the curvature fluid is hence a function of $`R`$ only, $`\rho _c=\rho _c(R)`$ and therefore the condition, Eq. (8), gives a second order homogeneous differential equation for $`g(R)R+f(R)`$:
$$3\frac{\rho _m}{\rho _m^{}}\left(\rho _m+p_m\right)g^{\prime \prime }(R)\left(\frac{1}{2}R+\rho _m\right)g^{}(R)+\frac{1}{2}g(R)=0,$$
(10)
where $`\rho _m`$ and $`p_m`$ are given as functions of $`R`$ only. Hence for any barotropic fluid, there exists a general $`f(R)`$ with two integration constants (expect if $`\rho _m+p_m=0`$ which leads to a first order differential equation) that will lead to the exactly same background evolution as in the standard EH gravity. In other words, for any barotropic equation of state $`f(R)`$ gravity constructed in this way, will always have among its solutions the same solution as the solution arising from EH gravity.
In a non-flat universe, the conclusions are unchanged but the technical details are somewhat more complicated since now one cannot write $`H\dot{R}`$ as a function of $`R`$ only due to the presence of the $`k`$-term. Instead, one needs to express everything in terms of the scale factor $`a`$, resulting in a differential equation for $`F(a)=f(R(a))`$ to be solved. After inverting $`R=R(a)`$ one obtains an expression for $`f(R)`$.
## IV Constructing equivalent gravitational actions
To illustrate the general conclusions of the previous section, we construct explicit gravitational actions that have identical background expansion as the solution arising from the EH action among their solutions.
Consider first a fluid with a constant equation of state, $`p_m=w\rho _m`$, where $`w`$ is a arbitrary constant, excluding cases $`w=1,\mathrm{\hspace{0.17em}1}/3`$. Now $`\rho _m=\rho _{m,0}a^{3(1+w)}`$, and as explained above, expressing all in terms of $`R`$, writing $`g(R)=G(u)`$, where $`u=R/(\rho _{m,0}(3w1))`$, we obtain the general equation for $`G`$:
$`0`$ $`=`$ $`6{\displaystyle \frac{w+1}{3w1}}\left(1\stackrel{~}{k}u^{\frac{3w+1}{3(w+1)}}\right)u^2G^{\prime \prime }(u)`$ (11)
$`{\displaystyle \frac{3w+1}{3w1}}uG^{}(u)+G(u),`$
where $`\stackrel{~}{k}=k/\rho _{m,0}`$. In the case of a flat universe, $`k=0`$, Eq. (11) can be solved:
$$g(R)=c_+(R)^{\alpha _+(w)}+c_{}(R)^{\alpha _{}(w)},$$
(12)
where $`\alpha _\pm (w)=(9w+7\pm \sqrt{9w^2+78w+73})/(12(w+1))`$ and $`c_\pm `$ are constants (note the $`R`$ is negative with our conventions). In contrast to the $`k0`$ case, the spatially flat solution does not depend on the (present day) density $`\rho _{m,0}`$. In the special case of a flat matter dominated universe, $`w=0`$, $`k=0`$, the solution translates to
$`f(R)=R`$ $`+`$ $`c_+(R)^{(7+\sqrt{73})/12}`$ (13)
$`+`$ $`c_{}(R)^{(7\sqrt{73})/12}.`$
Note how the other solution actually grows with $`|R|`$ i.e. its effect is larger at early times.
The radiation dominated case $`w=1/3`$ is not included to the solution given above because now $`R0`$ and the scale parameter can not be expressed in terms of $`R`$. However, this means that for radiation domination $`\dot{R}0`$ and therefore from Eqs (8) and (II) any function $`g(R)`$ with $`g(0)=g^{}(0)=0`$ will give the same radiation dominated expansion. The other excluded value $`w=1`$, corresponding to the deSitter model, gives $`R=4\rho _\mathrm{\Lambda }`$. Hence $`\dot{R}=0`$ and Eq. (10) simplifies to $`Rg^{}(R)2g(R)=0`$ giving $`g(R)=cR^2`$. This solution also coincides with the limiting value of $`\alpha _+`$ -solution: $`\alpha _+(w1)2`$.
Another interesting example is the flat $`\mathrm{\Lambda }`$CDM model
$$a(t)=\left(\frac{\mathrm{\Omega }_m}{\mathrm{\Omega }_\mathrm{\Lambda }}\right)^{1/3}\mathrm{sinh}^{2/3}(\frac{3}{2}\sqrt{\mathrm{\Omega }_\mathrm{\Lambda }}H_0t),$$
(14)
in which case the differential equation for $`g(R)`$ reads as
$$(R+3\beta )(R+4\beta )g^{\prime \prime }(R)(\frac{1}{6}R+\beta )g^{}(R)\frac{1}{6}g(R)=0,$$
(15)
where we have defined $`\beta 3\mathrm{\Omega }_\mathrm{\Lambda }H_0^2`$. This equation has a solution in terms of the hypergeometric function $`F(\alpha ,\beta ;\gamma ;z)`$ and the solution in terms of $`f(R)`$ is
$`f(R)`$ $`=`$ $`R+c_+z^{\alpha _+}F(\alpha _+,{\displaystyle \frac{3}{2}}\alpha _+;\alpha _{}\alpha _++1;{\displaystyle \frac{1}{z}})`$ (16)
$`+`$ $`c_{}z^\alpha _{}F(\alpha _{},{\displaystyle \frac{3}{2}}\alpha _{};\alpha _+\alpha _{}+1;{\displaystyle \frac{1}{z}})`$
where $`z3R/\beta `$ and $`\alpha _\pm \alpha _\pm (w=0)`$. It is easy to see that in the limit $`R\mathrm{}`$ (or $`z+\mathrm{}`$) we recover the matter dominated case (13) as expected.
## V Stability
An important question in considering $`f(R)`$ models as realistic theories of gravity is the stability of the ground state dolgov . The stability criterion can be formulated as a condition on the sign of the potential given by nojiri
$`U(R)`$ $`=`$ $`{\displaystyle \frac{1}{3}}R{\displaystyle \frac{f^{(1)}(R)f^{(3)}(R)R}{3f^{(2)}(R)^2}}{\displaystyle \frac{f^{(1)}(R)R}{3f^{(2)}(R)}},`$ (17)
$``$ $`{\displaystyle \frac{f^{(1)}(R)}{3f^{(2)}(R)}}+{\displaystyle \frac{2f(R)f^{(3)}(R)}{3f^{(2)}(R)^2}}`$
where we have assumed that the universe is homogeneous (note that we have different metric conventions compared to nojiri ). If $`U(R)>0`$ for classical, unperturbed solution, then the linear perturbations $`\delta R`$ are oscillatory without exponentially decreasing or decaying modes. In terms of function $`g(R)`$ the stability condition obtained from Eq. (17) reads as (for $`g^{\prime \prime }0`$)
$$\left(2gRg^{}\right)g^{(3)}\left(1+g^{}Rg^{\prime \prime }\right)g^{\prime \prime }>0.$$
(18)
Equation (18) shows that for the radiation dominated case gravitational stability is achieved simply if $`g^{\prime \prime }(0)<0`$ and for the deSitter case (where $`g=cR^2`$) if $`c<0`$. In order to have an understanding whether stable choices of $`f(R)`$ exist in the more general case where $`p=w\rho `$, we consider the two solutions $`f_\pm (R)=R+c_\pm (R)^{\alpha _\pm }`$ separately and leave more general considerations for future work. We find that both solutions are stable in two separate regions,
$`c_\pm `$ $`<`$ $`0,`$
$`c_\pm `$ $`>`$ $`{\displaystyle \frac{6(1+w)(R)^{1\alpha _\pm (w)}}{17+15w\sqrt{9w^2+78w+73}}}`$ (19)
and hence one can always construct stable extensions of the EH action for such a fluid.
## VI Conclusions and discussion
In the present letter we have considered the uniqueness of the gravitational action among $`f(R)`$ gravity models. We have showed that for any barotropic fluid, the evolution of the scale factor of the universe does not uniquely determine the form of $`f(R)`$, and hence nor the gravitational action. Instead, one can always construct a $`f(R)`$ theory that will have the same background evolution among its solutions. In other words, cosmology covered by standard Friedmann equations has always a generalized counterpart with the very same classical evolution independently of the actual equation of state of the ordinary matter.
As an example, we have constructed explicit forms of $`f(R)`$ in a number cosmologically interesting cases. By considering small fluctuations around the ground state, we have shown that such modified theories of gravity are also stable.
In this letter we have not directly addressed the dark energy problem. However, our results are likely to be relevant to such considerations. Assuming that we have a matter dominated universe, we now know the form of $`f(R)`$ we need in order to reproduce the matter dominated regime exactly. This will guide us in considering choices of $`f(R)`$ that will lead to late time acceleration while preserving standard matter dominated expansion at early times.
These considerations are also closely related to the models obtained by conformal transformations. Although the gravitational actions constructed in this article may be rather unintuitive, their counterparts in scalar-tensor -gravity may seem natural. Moreover from the analysis presented here we know that some highly non trivial, probably non-minimally coupled scalar-tensor theories reproduce exactly the standard evolution of the background metric.
Since the background expansion alone cannot distinguish between different choices of $`f(R)`$, one must study perturbations in such cosmologies. This will ultimately guide us in choosing the correct gravitational action which is possibly not the simple Einstein-Hilbert action.
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# Exact Quantization of Einstein-Rosen Waves Coupled to Massless Scalar Matter
(May 6, 2005)
## Abstract
We show in this letter that gravity coupled to a massless scalar field with full cylindrical symmetry can be exactly quantized by an extension of the techniques used in the quantization of Einstein-Rosen waves. This system provides a useful testbed to discuss a number of issues in quantum general relativity such as the emergence of the classical metric, microcausality, and large quantum gravity effects. It may also provide an appropriate framework to study gravitational critical phenomena from a quantum point of view, issues related to black hole evaporation, and the consistent definition of test fields and particles in quantum gravity.
Symmetry reductions of general relativity have been used as model systems to extract information about quantum gravity. They usually allow the discussion of specific problems without the difficulties present in the full theory. Some of the most popular choices in this regard (Bianchi models) have only a finite number of degrees of freedom and, hence, are not suitable to address some of the more nagging questions posed by the study of quantum gravity (diffeomorphism invariance or issues related to the presence of an infinite number of local degrees of freedom such as perturbative non-renormalizability). Fortunately there are other symmetry reductions that retain these features while still being exactly solvable both at the classical and quantum levels. Chief among them are the Einstein-Rosen waves Einstein and Rosen (1937); Kuchar (1971), obtained by requiring that space-time metrics have two commuting, spacelike, and hypersurface orthogonal Killing vector fields (one translational and the other rotational). This model has been extensively studied in the past Ashtekar and Pierri (1996); Ashtekar and Varadarajan (1994) and some intriguing results have been derived, in particular the appearance of unexpected large quantum gravity effects Ashtekar (1996) and a detailed picture of the emergence of the causal structure of space-time in the classical limit Barbero G. et al. (2003, 2003). An improvement that would increase the usefulness of this system as a toy model for quantum gravity would be the coupling of matter. The availability of a solvable model with matter would open up a host of interesting possibilities deserving a careful investigation. We show in this letter that such a model exists and discuss how it can be exactly quantized. Specifically we will consider here the quantization of Einstein-Rosen waves coupled to a cylindrically symmetric massless scalar field.
To our knowledge this system was first discussed from a classical point of view by Chandrasekhar Chandrasekhar (1986) who showed that a full solution to the Einstein field equations can be found in this case. The specific form of this solution suggests that a Hamiltonian treatment of the system would lead to a description very similar to the one found by Ashtekar, Pierri, and Varadarajan Ashtekar and Pierri (1996); Ashtekar and Varadarajan (1994) for pure gravity in the asymptotically flat case. This is a strong indication that the model is amenable to quantization by a suitable extension of known techniques. The main point of this letter is to show that this is indeed the case.
The possible applications of such a model are manifold and can be classified in several different categories. First of all there is the issue of extracting information about geometry in quantized gravity. We do not expect to find here the kind of precise geometric information offered by loop quantum gravity in the form of geometric observables such as areas or volumes. Nevertheless our approach gives some indications about the validity of a metric description in the realm of quantum gravity. This has already been considered for pure gravity by studying expectation values of metric components. Here we propose to follow a different philosophy; instead of obtaining some approximate semiclassical metric by taking expectation values of a metric operator in some quantum state we can use the scalar field and its particle-like excitations to explore space-time geometry operationally, much in the same way as one uses the geodesics followed by test particles to understand the geometrical features of a given space-time metric. The availability of an external probe allows us to discuss the microcausality of the model both from the perspective of the gravitational and scalar degrees of freedom. The agreement between both points of view, that we discuss later, is a clear indication of the usefulness of the present approach in the discussion of quantum gravitational effects and supports the results already obtained for the purely gravitational case.
A second set of questions that can possibly be addressed within this framework in the quantum regime are related to critical phenomena in gravitational collapse and problems in black hole physics<sup>1</sup><sup>1</sup>1This would require dropping the radial asymptotic flatness condition. By doing this we can have the self-similar solutions needed to discuss critical collapse Wang (1986) or escape the conclusions of Berger. et al. (2003) about the absence of compact trapped surfaces in the asymptotically flat case.. These issues have been recently considered by Wang in his study of critical collapse of a cylindrically symmetric scalar field in four dimensions Wang (1986). This system displays some rich and non-trivial behavior, in particular, the possibility of forming solutions with future, spacelike singularities by the collapse of massless scalar matter and the appearance of a critical metric separating solutions with different singular behavior. Notice that having an exact solution, backreaction effects are automatically taken into account without any approximation. Finally we want to point out other possible uses of this model such as the discussion of the validity of the usual perturbative schemes in quantum gravity, the development of new ones, the discussion of issues in QFT in curved spacetimes, and the application to other useful symmetry reductions of these type (Gowdy models) that are similar to Einstein-Rosen waves and, hence, can also be solved after coupling massless scalar fields.
Our starting point is the four dimensional action for a massless scalar $`\mathrm{\Phi }_s`$ coupled to gravity with cylindrical symmetry
$`{}_{}{}^{4}S={\displaystyle \frac{1}{16\pi G_N}}{\displaystyle _{\times I}}d^4x\sqrt{|^4g|}\left[R{\displaystyle \frac{1}{2}}{}_{}{}^{4}g_{}^{ab}_a\mathrm{\Phi }_s_b\mathrm{\Phi }_s\right]`$
$`+{\displaystyle \frac{1}{8\pi G_N}}{\displaystyle _{(\times I)}}d^3x(\sqrt{|^3h|}K\sqrt{|^3h^0|}K^0).`$
Here we have included the surface terms necessary to have a well-defined variational principle, $`I[z_1,z_2]`$ is a closed interval in the direction of the translational Killing vector $`_z`$, and $`K`$, $`K^0`$ are the extrinsic curvatures of the boundary defined by the dynamical metric $`{}_{}{}^{4}g_{ab}^{}`$ and a fiducial metric $`{}_{}{}^{4}g_{ab}^{0}`$ that we choose as Minkowski in the following (we denote the induced metrics on the boundary as $`{}_{}{}^{3}h_{ab}^{}`$ and $`{}_{}{}^{3}h_{ab}^{0}`$). The Geroch formalism (and a subsequent conformal transformation) allows us to reduce the previous action to the following three dimensional one by taking advantage of the translational symmetry
$`{}_{}{}^{3}S`$ $`=`$ $`{\displaystyle \frac{1}{16\pi G_3}}{\displaystyle _{}}d^3x\sqrt{|g|}[^3R{\displaystyle \frac{1}{2}}g^{ab}_a\varphi _g_b\varphi _g`$
$`{\displaystyle \frac{1}{2}}g^{ab}_a\varphi _s_b\varphi _s]+{\displaystyle \frac{1}{8\pi G_3}}{\displaystyle }_{}d^2x[\sqrt{|h|}K\sqrt{|h^0|}K^0].`$
Here $`g_{ab}`$ is a 3-dimensional metric and $`{}_{}{}^{3}R`$ the corresponding scalar curvature, $`\varphi _g`$ is the scalar field that encodes the local gravitational degrees of freedom of the model Ashtekar and Pierri (1996), $`\varphi _s`$ is the massless matter scalar field in three dimensions, and $`G_3`$ is the gravitational constant per unit length along the symmetry axis (with dimensions of inverse energy; in the following we choose units such that $`\mathrm{}=c=8G_3=1`$). The integration is extended to a 3-manifold $``$ with boundary $``$ with the appropriate topology. At this point it could be argued that the inclusion of the massless scalar is a rather trivial addition to the system because it just plays the role of an extra field of the same type of the gravitational scalar already present in the 2+1 dimensional description of Einstein-Rosen waves. However this is the most important and unexpected feature of (Exact Quantization of Einstein-Rosen Waves Coupled to Massless Scalar Matter) because both the gravitational and matter degrees of freedom are described by the same type of term in the three-dimensional Lagrangian in spite of their very different meaning in the original action, the Geroch reduction, and the conformal transformation used to arrive at (Exact Quantization of Einstein-Rosen Waves Coupled to Massless Scalar Matter). It is also striking that they couple only through the metric and not directly (there are no cross terms).
As we are interested in the quantization of the system it is necessary to obtain the Hamiltonian corresponding to (Exact Quantization of Einstein-Rosen Waves Coupled to Massless Scalar Matter). Although the final answer turns out to be quite simple it is not completely obvious, and it has some surprising features, so we provide some details on its derivation. To this end we choose a foliation of $``$ with timelike unit normal $`n^a`$, a radial unit vector $`\widehat{r}^a`$, and denote as $`\sigma ^a`$ the azimutal, hypersurface orthogonal, Killing field (notice that this is not a unit vector). We further introduce two additional vector fields $`t^a`$ and $`r^a`$ defined as $`t^a=Nn^a+N^r\widehat{r}^a`$ and $`r^a=e^{\gamma /2}\widehat{r}^a`$, where $`N`$ is the lapse function, $`N^r`$ the radial shift, and $`\gamma `$ is an additional field. It is possible to find conditions that ensure that $`t^a`$, $`r^a`$, and $`\sigma ^a`$ are coordinate vectors. If we define $`_\sigma \sigma ^a_a`$, $`_rr^a_a`$, and $`_tt^a_a`$ these are the following:
$`_\sigma N=_\sigma N^r=_\sigma \gamma =0;[\sigma ,\widehat{r}]^a=[\sigma ,n]^a=0;`$
$`n^a_rN+\widehat{r}^a(_rN^r_te^{\gamma /2})+Ne^{\gamma /2}[\widehat{r},n]^a=0.`$
Writing the metric as $`g_{ab}=n_an_b+\widehat{r}_a\widehat{r}_b+\frac{1}{R^2}\sigma _a\sigma _b`$ (with $`R^2g_{ab}\sigma ^a\sigma ^b`$) we obtain the line element in these coordinates
$`ds^2=(N^{r2}N^2)dt^2+2e^{\gamma /2}N^rdtdr+e^\gamma dr^2+R^2d\sigma ^2.`$
The action can be rewritten now as
$`{}_{}{}^{3}S`$ $`=`$ $`{\displaystyle _{t_1}^{t_2}}dt{\displaystyle _0^{\stackrel{~}{r}}}dr\{Ne^{\gamma /2}(\gamma ^{}R^{}2R^{\prime \prime })`$
$`{\displaystyle \frac{1}{N}}(e^{\gamma /2}\dot{\gamma }2N^r)(\dot{R}e^{\gamma /2}N^rR^{})+`$
$`+{\displaystyle \frac{R}{2N}}\left[e^{\gamma /2}\dot{\varphi }_g^22N^r\dot{\varphi }_g\varphi _g^{}+e^{\gamma /2}(N^{r2}N^2)\varphi _g^2\right]`$
$`+{\displaystyle \frac{R}{2N}}[e^{\gamma /2}\dot{\varphi }_s^22N^r\dot{\varphi }_s\varphi _s^{}+e^{\gamma /2}(N^{r2}N^2)\varphi _s^2]\}`$
$`+2{\displaystyle _{t_1}^{t_2}}๐t(Ne^{\gamma /2}R^{}1),`$
where we have denoted $`_t`$ with a dot, $`_r`$ with a prime. The Hamiltonian when we take $`\stackrel{~}{r}\mathrm{}`$ is
$`H`$ $`=`$ $`{\displaystyle _0^{\mathrm{}}}dr\{N^re^{\gamma /2}[p_RR^{}p_\gamma ^{}+p_\gamma \gamma ^{}+\varphi _g^{}p_g+\varphi _s^{}p_s]`$
$`+Ne^{\gamma /2}[2R^{\prime \prime }\gamma ^{}R^{}p_Rp_\gamma +{\displaystyle \frac{1}{2R}}p_g^2`$
$`+{\displaystyle \frac{R}{2}}\varphi _g^2+{\displaystyle \frac{1}{2R}}p_s^2+{\displaystyle \frac{R}{2}}\varphi _s^2]\}+2(1e^{\gamma _{\mathrm{}}/2}),`$
where $`p_R`$, $`p_\gamma `$, $`p_g`$, and $`p_s`$ are the momenta canonically conjugate to $`R`$, $`\gamma `$, $`\varphi _g`$, and $`\varphi _s`$ respectively, $`\gamma _{\mathrm{}}lim_r\mathrm{}\gamma (r)`$, and the fall-off of the fields, that ensures asymptotic flatness in 2+1 dimensions and implies $`N1`$ and $`R^{}1`$, is the one used in Ashtekar and Pierri (1996). All fields are chosen to be regular in the axis. From the previous expression the Hamiltonian of the system and the constraints can be immediately read. To proceed ahead we fix the gauge with the conditions $`R(r)=r`$ and $`p_\gamma (r)=0`$ (the same as in the absence of matter). It is straightforward to show that they are admissible. After fixing the gauge and solving the constraints we get
$`\gamma (R)={\displaystyle \frac{1}{2}}{\displaystyle _0^R}๐rr\left[\varphi _g^2+{\displaystyle \frac{p_g^2}{r^2}}+\varphi _s^2+{\displaystyle \frac{p_s^2}{r^2}}\right],`$
the three dimensional line element can be written as
$$ds^2=e^\gamma [e^\gamma _{\mathrm{}}dt^2+dR^2]+R^2d\sigma ^2$$
(2)
and the reduced Hamiltonian is
$$H=2(1e^{\gamma _{\mathrm{}}/2}).$$
This is a function of the sum of the Hamiltonians for two massless cylindrically symmetric fields evolving in a fictitious Minkowskian background. For every solution to the field equations $`\gamma _{\mathrm{}}`$ is a constant of motion. Taking advantage of this we can introduce an auxiliary, solution-dependent, time variable as in Ashtekar and Pierri (1996) defined according to $`T=e^{\gamma _{\mathrm{}}/2}t`$, that allows us to simplify the form of the field equations to get
$`_T^2\varphi _g\varphi _g^{\prime \prime }{\displaystyle \frac{1}{R}}\varphi _g^{}=0,_T^2\varphi _s\varphi _s^{\prime \prime }{\displaystyle \frac{1}{R}}\varphi _s^{}=0.`$ (3)
Equations (3) describe two massless, cylindrically symmetric scalar fields in 2+1 dimensions. Classically this is a time redefinition that amounts to a change of the coordinate $`t`$; once we pick a certain solution to (3) we can choose to write (2) either in terms of $`t`$ or $`T`$. Quantum mechanically the situation is more complicated because the evolution of wave packets generically involves the superposition of Hilbert space vectors with energy dependent phases so a change in the functional form of the energy completely changes the evolution of the states. It is very important to notice that the form of the Hamiltonian means that the model is not free. The two fields that appear are coupled in a non trivial way.
In order to quantize the system we define field and momenta operators $`\widehat{\varphi }_{g,s}(R)`$, $`\widehat{p}_{g,s}(R)`$ satisfying the commutation relations $`[\widehat{\varphi }_{g,s}(R),\widehat{p}_{g,s}(R^{})]=i\delta (R,R^{})`$ and introduce creation and annihilation operators as usual according to
$`\varphi _{g,s}(R)={\displaystyle \frac{1}{\sqrt{2}}}{\displaystyle _0^{\mathrm{}}}๐kJ_0(Rk)[a_{g,s}(k)+a_{g,s}^{}(k)],`$
$`p_{g,s}(R)={\displaystyle \frac{iR}{\sqrt{2}}}{\displaystyle _0^{\mathrm{}}}๐kkJ_0(Rk)[a_{g,s}^{}(k)a_{g,s}(k)],`$
with non-zero commutators given by
$`[a_g(k),a_g^{}(q)]=\delta (k,q),[a_s(k),a_s^{}(q)]=\delta (k,q).`$
These operators are defined in a Hilbert space built as a tensor product of two Fock spaces $`_g`$ and $`_s`$, $`=_g_s`$ with a vacuum state $`|\mathrm{\Omega }=|0^g|0^s`$ defined in terms of the vacua annihilated by $`a_{g,s}(k)`$. States with a fixed number of quanta of โgravitationalโ or โscalarโ type are obtained by repeated action of the corresponding creation operators $`|k_{g,s}A_{g,s}^{}(k)|\mathrm{\Omega }`$, where we have written $`A_g^{}(k)a_g^{}(k)๐_s`$, $`A_s^{}(k)๐_ga_s^{}(k)`$.
The quantum Hamiltonian in $``$ is
$`\widehat{H}=2\left(1\mathrm{exp}[{\displaystyle \frac{1}{2}}{\displaystyle _0^{\mathrm{}}}๐kk[A_g^{}(k)A_g(k)+A_s^{}(k)A_s(k)]]\right).`$
We have normal ordered the exponent to remove the zero point energy of the vacuum. This Hamiltonian is a non-linear and bounded function of the sum of the Hamiltonians for two massless, cylindrically symmetric scalar fields in 2+1 dimensions, $`H_0^g`$ and $`H_0^s`$. It is an observable of the system (the energy) and the generator of time evolution in the time variable $`t`$ (from $`t_1`$ to $`t_2`$)
$$U(t_2t_1)=\mathrm{exp}\left[2i(t_2t_1)\left(1e^{\frac{1}{2}[H_0^g+H_0^s]}\right)\right].$$
(4)
It is important to realize at this point that this is the physical evolution. The free Hamiltonians $`H_0^g`$ and $`H_0^s`$ are indeed observables but are not directly related to the time evolution of the system. It is necessary to take this fact into account in the search for semiclassical states because coherent states should display a classical behavior under the evolution given by (4) rather than under the one that would be defined by the โfreeโ Hamiltonian $`H_0^g+H_0^s`$. The expectation value of the field and momenta operators should evolve in $`t`$ according to the classical field equations in terms of $`t`$; notice that they are not (3).
The unitary evolution given by $`U(t)`$ defines the $`S`$-matrix of the system (the $`S`$-matrix in QFT is basically the evolution operator in the limit $`t\mathrm{}`$). Its matrix elements in n-particle states can be computed in a straightforward way because they are eigenstates of $`H_0^g`$ and $`H_0^s`$. The only non-zero matrix elements on states with a definite number of both quanta (i.e. gravitational and matter) are those connecting state vectors with the same number of particles of each type; hence there is no conversion of quanta of one type into the other. It is necessary at this point to stress that the split of the Hilbert space as a tensor product of two Hilbert spaces should not immediately lead us to interpret one of them as โgravitationalโ and the other as โmatterโ; in fact the classical metric depends on both the gravitational scalar and the matter scalar and semiclassical approximations of it obtained by computing expectation values of a metric operator would also depend on both the โgravitationalโ and the โmatterโ part of the state. In this sense a vector such as $`|0^g|\mathrm{\Phi }^s`$ should not be interpreted as an approximate description of a pure matter state $`|\mathrm{\Psi }^s`$ on some quantum approximation of the Minkowski metric. In fact the quantum state that most closely resembles the Minkowski metric is the vacuum $`|\mathrm{\Omega }`$. By the way, this is the only coherent state of the system that we know under the evolution given by (4).
The fact that the $`S`$-matrix on n-particle states can be found in a straightforward way on the generalized orthonormal basis $`|p^g|q^s`$ and its diagonal character does not mean that interesting quantum information on the system cannot be obtained. Quite on the contrary some significant features of quantum spacetime geometry can be seen. We will concentrate here on the discussion of microcausality. Though we have looked at this problem somewhere else Barbero G. et al. (2003, 2003) the inclusion of scalar matter in the model allows us to consider the issue both from the point of view of the metric scalar and the matter scalar. As in previous work we concentrate on the vacuum expectation value of the commutator of the fields evolved with (4) in the Heisenberg picture
$`\mathrm{\Omega }|[\widehat{\varphi }_g(R^{},t^{}),\widehat{\varphi }_g(R,t)]|\mathrm{\Omega }=\mathrm{\Omega }|[\widehat{\varphi }_s(R^{},t^{}),\widehat{\varphi }_s(R,t)]|\mathrm{\Omega }`$
$`={\displaystyle \frac{i}{2}}{\displaystyle _0^{\mathrm{}}}๐kJ_0(R^{}k)J_0(Rk)\mathrm{sin}[(t^{}t)E(k)],`$ (5)
with $`E(k)=2(1e^{k/2})`$. As can be seen the commutator on the vacuum for both types of fields is exactly the same. The microcausality of the model as described by the gravitational and matter scalars coincide. As we showed in Barbero G. et al. (2003), we can see from (5) the emergence of the sharp and well defined cylindrical light cone structure corresponding to the quantization of a cylindrical massless scalar field in a 2+1-dimensional Minkowskian background. This happens in the limit when spatial distances and time intervals are large in comparison with the length scale $`\mathrm{}G_3`$. Finally it is interesting to point out that even though the diagonal matrix elements of the commutator of $`\widehat{\varphi }_g(R_2,t_2)`$ and $`\widehat{\varphi }_s(R_1,t_1)`$ are zero, the non-diagonal ones are generically different from zero for $`t_2t_1`$. This is a clear indication that both fields interact in a non-trivial way.
The availability of a matter field allows us to explore how the quantization of gravity reflects on the geometric properties of spacetime by using its particle-like excitations as quantum test particles. Though the details of this will appear elsewhere we want to explain here how it can be done. A possible approach to the problem is to interpret $`\mathrm{\Omega }|\widehat{\varphi }_s(R_2,t_2)\widehat{\varphi }_s(R_1,t_1)|\mathrm{\Omega }`$ as the probability amplitude for a particle created at the radial distance $`R_1`$ in the instant $`t_1`$ to be found at $`(R_2,t_2)`$. As in ordinary QFT in Minkowski spacetime this interpretation is only approximate (i.e. valid only above a certain distance scale) because the $`R`$-dependent states $`\widehat{\varphi }_s(R)|\mathrm{\Omega }`$ do not constitute an orthonormal basis. One can, however, introduce the analogous of the Newton-Wigner localized states for this system, define one-particle wave functions depending on the radial coordinate, and study their time evolution under the dynamics given by (4). If we choose appropriately peaked states it should be possible to study to what extent the evolution of wave packets follows the null geodesics of some cylindrical spacetime metric and the length scales in which a metric gives an accurate description of spacetime geometry.
Finally we want to point out that once the massless case is understood it could be possible to use it as a guide to find a consistent way to introduce other types of test fields (such as massive scalars or electromagnetic fields) that further improve our ability to explore quantum geometry and quantum gravity. This will be the focus of our attention in the near future.
###### Acknowledgements.
We want to thank G. Mena and J. M. Martรญn Garcรญa for drawing our attention to references Chandrasekhar (1986); Wang (1986) and also them, A. Ashtekar, L. Garay, and M. Varadarajan for discussions. IG is supported by a Spanish Ministry of Science and Education FPU fellowship. This work is also supported by the Spanish MEC under the research project BFM2002-04031-C02-02.
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# Natural Nuclear Reactor Oklo and Variation of Fundamental Constants: Computation of Neutronics of Fresh Core
## I Introduction
The discovery of the natural nuclear reactor in Gabon (West Africa) was possibly one of the most momentous events in reactor physics since in 1942 Enrico Fermi with his team achieved an artificial self-sustained fission chain reaction. Soon after the discovery of the ancient natural reactor Oklo in Gabon (West Africa) 1 ; 2 ; 3 , one the authors of the present paper (Yu.P.) and his postgraduate student A.I. Shlyakhter realized that the โOklo phenomenonโ could be used to find the most precise limits on possible changes of fundamental constants. At that time they considered probabilistic predictions of unknown absorption cross sections based on static nuclear properties. Near the neutron binding energy ($`B_n=68`$ MeV) the resonances form a fence with a mean separation of tens of electron volts. The magnitude of the cross section depends on the proximity of the energy $`E_{th}`$ of the thermal neutron to the nearest resonance. If the energy $`E_{th}=25`$ meV falls directly on a resonance, then the cross section increases to as much $`10^510^6`$ b 4a ; 4b ; 5 . If the cross sections have changed with time, then the entire fence of resonances as a whole has shifted by a small amount $`\mathrm{\Delta }E_r`$. This shift can be established most accurately from the change of the cross section of strong absorbers (for instance of $`_{\mathrm{\hspace{0.25em}\hspace{0.25em}62}}^{149}\text{Sm}`$). The estimate of the shift $`\mathrm{\Delta }E_r`$ on the basis of experimental data for the Oklo reactor ($`\mathrm{\Delta }E_r310^{17}`$ eV/year) 6a ; 6b ; 6c ; 7 allowed us to get the most accurate estimate of a possible limit of the rate of change of fundamental constants. This estimate has remained the most accurate one for 20 years. In 1996 a paper by Damour and Dyson was published in which the authors checked and confirmed the results of Ref.6a ; 6b ; 6c . Damour and Dyson were the first to calculate the dependence of the capture cross section on the temperature $`T_C`$ of the core: $`\widehat{\sigma }_{r,Sm}(T_C)`$ 8a ; 8b . In 2000 Fujii et al. published a paper in which the authors significantly reduced the experimental error of the cross section $`\widehat{\sigma }_{r,Sm}(T_C)`$ 9 . In both papers the authors averaged the samarium cross section with a Maxwell velocity spectrum over a wide interval of the core temperature $`T_C`$.
After the publication of Refs. 8a ; 8b , one of the authors of the present paper (Yu. P.) realized that the limit on the change of the cross section can be significantly improved at least in two directions:
1. Instead of a Maxwell distribution, the samarium cross section should be averaged with the spectrum of the Oklo reactor which contains the tail of Fermi spectrum of slowing down epithermal neutrons.
2. The range of admissible core temperatures $`T_C`$ can be significantly reduced, assuming that $`T_C`$ is the equilibrium temperature at which the effective multiplication factor $`K_{eff}`$ of the reactor is equal to one. On account of the negative power coefficient (void + temperature) such a state of the reactor will be maintained for a long time until the burn-up results in a reduction of the reactivity excess and hence of $`T_C`$. Since $`_{\mathrm{\hspace{0.25em}\hspace{0.25em}62}}^{149}\text{Sm}`$ burns up about 100 times faster than $`_{\mathrm{\hspace{0.25em}\hspace{0.25em}92}}^{235}\text{U}`$, the core will contain only that amount of samarium that was generated immediately before the reactor shut down. Therefore one needs to know the reactivity excess and $`T_C`$ at the end of the cycle.
To solve this problem one must use modern neutron-physical and thermo-hydrodynamical methods of reactor calculations. We have built up a complete computer model of the Oklo reactor core RZ2 and established its material composition. We have chosen three variants of its initial composition in order to estimate its effect on the spread of results. To increase the reliability of the results we have used modern versions of two Monte Carlo codes. One of them, which has been developed at the Kurchatov Institute, is the licensed Russian code MCU-REA with the library DLC/MCUDAT-2.2 of nuclear data 10 ; the other one is the well known international code MCNP4C with library ENDF/B-VI 11 . Both codes give similar results. We have calculated the multiplication factors, reactivity and neutron flux for the fresh cores, and the void and temperature effects 12 . As expected, the reactor spectrum differs strongly from a Maxwell distribution (see below). The cross section $`\widehat{\sigma }_{r,Sm}(T_C,\mathrm{\Delta }E_r)`$, averaged with this distribution, is significantly different from the cross section averaged with a Maxwell distribution 13 . We use our result for the averaged cross section to estimate the position of resonances at the time of Oklo reactor activity. This allows us to obtain the most accurate limits on the change of the fine structure constant in the past.
The paper is organized as follows. In section I we describe briefly the history of the discovery of the natural Oklo reactor and itemize the main parameters of its cores. We consider mainly the core RZ2 . We describe in detail the neutronics of this core calculated by modern Monte Carlo codes. However simple semianalytical considerations are also useful to clarify the picture. We consider the power effect which is a sum of the temperature and void effects. At the end of the section we discuss the computational difficulties in the calculations of the unusually large core RZ2 and demonstrate that Monte Carlo methods are, in general, inadequate for the calculations of core burn-up.
The main result of Section I is the neutron spectrum in the fresh core. In Section II we apply this spectrum to obtain the averaged cross section of $`_{\mathrm{\hspace{0.25em}\hspace{0.25em}62}}^{149}\text{Sm}`$ in the past. We begin this Section with an explanation of the way of obtaining precise limits on the variation of fundamental constants using the available Oklo reactor data. We describe different approaches to the problem and relate the variation of the constants to the change in the averaged cross sections for thermal neutrons. Using our value for the cross section of $`_{\mathrm{\hspace{0.25em}\hspace{0.25em}62}}^{149}\text{Sm}`$ we obtain limits on the variation of the fine structure constant which is the best available at the moment. At the end of this Section we compare our result with the results obtained in other papers and discuss possible reasons of differences.
## II Neutronics of the fresh core
### II.1 History of the discovery and parameters of the Oklo reactor
#### II.1.1 History of the discovery of the natural reactor
The first physicist to say in May or June of 1941 that a nuclear chain reaction could have been more easily realized a billion years ago was Yakov Borisovich Zeldovich 14 . At that time he was considering the possibility of getting a fission chain reaction in a homogeneous mixture of natural uranium with ordinary water. His calculations (with Yu.B. Khariton) showed that this could be achieved with an approximately two-fold enrichment of natural uranium 15 ; 16 . A billion years ago the relative concentration of the light uranium isotope was significantly higher, and a chain reaction was possible in a mixture of natural uranium and water. โYakov Borisovich said nothing about the possibility of a natural reactor, but his thoughts directly lead us to the natural reactor discovered in Gabon in 1972โ reminisced I.I. Gurevich 14 . Later, in 1957, G. Whetherill and M. Inghram arrived at the same conclusion 17a ; 17b . Going from the present concentration of uranium in pitchblende, they concluded that about two billion years ago, when the proportion of $`_{\mathrm{\hspace{0.25em}\hspace{0.25em}92}}^{235}\text{U}`$ exceeded 3%, conditions could be close to critical. Three years later, P. Kuroda 18a ; 18b showed that, if in the distant past there was water present in such deposits, then the neutron multiplication factor ($`K_{\mathrm{}}`$) for an infinite medium could exceed unity and a spontaneous chain reaction could arise. But before 1972 no trace of a natural reactor has been found. On the 7th of June 1972, during a routine mass-spectroscopic analysis in the French Pierrelatte factory that produced enriched fuel, H. Bouzigues 1 ; 3 noticed that the initial uranium hexafluoride contains $`\zeta _5=0.717\%`$ of $`_{\mathrm{\hspace{0.25em}\hspace{0.25em}92}}^{235}\text{U}`$ atoms instead of the $`0.720`$%, which is the usual concentration in terrestrial rock, meteorites and lunar samples. The French Atomic Energy Authority (CEA) began an investigation into this anomaly. The phenomenon was named โOklo phenomenonโ. The results of this research were published in the proceedings of two IAEA symposia 2 ; 19 . The simplest hypothesis of a contamination of the uranium by depleted tails of the separation process was checked and shown to be wrong. Over a large number of steps of the production process, the anomaly was traced to the Munana factory near Franceville (Gabon) where the ore was enriched. The ore with a mean uranium concentration of ($`0.40.5`$)% got delivered there from the Oklo deposit. The isotope analysis of the uranium-rich samples showed a significant depletion of the $`_{\mathrm{\hspace{0.25em}\hspace{0.25em}92}}^{235}\text{U}`$ isotope and also a departure from the natural distribution of those rare earth isotopes, which are known as fission products 1 ; 3 ; 20a ; 20b . This served as a proof of the existence in the distant past of a spontaneous chain reaction. It had taken less than three months to produce this proof. A retrospective analysis of documents and samples of the Munane enrichment factory showed that in 1970-72 ore was delivered for processing that contained at times up to $`20`$% of uranium depleted to $`0.64`$% of isotope $`_{\mathrm{\hspace{0.25em}\hspace{0.25em}92}}^{235}\text{U}`$ 21 . Considering that the ore was mixed during mining, the uranium concentration could be even higher in some samples, and the depletion even stronger. Altogether more than 700 tons of depleted uranium has been mined that had taken part in the chain reaction. The deficit of $`_{\mathrm{\hspace{0.25em}\hspace{0.25em}92}}^{235}\text{U}`$ (that had not been noticed at first) was about 200 kg. By agreement with the Government of Gabon, the uranium ore production company of Franceville (COMUF) agreed to halt mining in the region of the natural reactor. A Franco-Gabon group headed by R. Naudet began a systematic study of the Oklo phenomenon. Numerous samples, obtained by boring, were sent for analysis to various laboratories around the world. They allowed a reconstruction of the functioning of the reactor in the Precambrian epoch.
#### II.1.2 Geological history of the Oklo deposit
As was shown by the $`U/Pb`$ analysis, the Oklo deposit with a uranium concentration of about $`0.5`$% in the sediment layer was formed about $`210^9`$ years ago 22 ; 23 ; 24 . During this epoch an important biological process was taking place: the transition from prokaryotes, i.e. cells without nucleus, to more complex unicellular forms containing a nucleus - eucaryotes. The eucaryotes began to absorb carbon oxide and hence saturate the atmosphere with oxygen. Under the influence of oxygen, the uranium oxides began transforming into forms containing more oxygen, which are soluble in water. Rains have washed them into an ancient river, forming in its mouth a sandstone sediment, rich in uranium, of 4 to 10 meters thickness and a width of 600 to 900 meters 24 . The heavier uranium particles settled more quickly to the ground of the nearly stagnant water of the river delta. As a result the sandstone layer got enriched with uranium up to 0.5% (as in an enrichment factory). After its formation, the uranium-rich layer, that was resting on a basalt bed, was covered by sediments and sank to a depth of 4 kilometers. The pressure on this layer was 100 MPa25 . Under this pressure the layer got fractured and ground water entered the clefts. Under the action of the filtered water that was subjected to a high pressure, and as a result of not completely understood processes, lenses formed with a very high uranium concentration (up to $`2060`$% in the ore) with a width of 10 to 20 meters and of the order of 1 meter thickness 26 . The chain reaction took place in these lenses. After the end of the chain reaction the deposit was raised to the surface by complicated tectonic processes and became accessible for mining. Within tens of meters six centers of reactions were found immediately, and altogether the remains of 17 cores were found 27 .
The age $`T_0`$ of the reactor was determined from the total number of $`_{\mathrm{\hspace{0.25em}\hspace{0.25em}92}}^{235}\text{U}`$ nuclei burnt up in the past, $`N_{5b}(d)`$ , and the number of nuclei existing today, $`N_5(T_0)`$(here $`N_5`$ is the density of $`_{\mathrm{\hspace{0.25em}\hspace{0.25em}92}}^{235}\text{U}`$ and $`d`$ is the duration of the chain reaction). For such a way of determining $`T_0`$ it is necessary to know the number of $`_{\mathrm{\hspace{0.25em}\hspace{0.25em}94}}^{239}\mathrm{Pu}`$ nuclei formed as a result of neutron capture by $`_{\mathrm{\hspace{0.25em}\hspace{0.25em}92}}^{238}U`$ and decayed to $`_{\mathrm{\hspace{0.25em}\hspace{0.25em}92}}^{235}\text{U}`$, and the fluence $`\mathrm{\Psi }=\mathrm{\Phi }d`$ ($`\mathrm{\Phi }`$ being the neutron flux). Another independent method consists of the determination of the amount of lead formed as a result of the decay of $`_{\mathrm{\hspace{0.25em}\hspace{0.25em}92}}^{235}\text{U}`$, assuming that it did not occur in such a quantity in the initial deposit 22 . Both methods yield $`T_0=1.81(5)10^9`$ years 7 ; 28 . Below we assume in our calculations the value of $`T_0=1.810^9`$ years.
The duration of the work of the reactor can be established from the amount of $`_{\mathrm{\hspace{0.25em}\hspace{0.25em}94}}^{239}\mathrm{Pu}`$ formed. One can separate the decayed $`_{\mathrm{\hspace{0.25em}\hspace{0.25em}94}}^{239}\mathrm{Pu}`$ from the decayed $`_{\mathrm{\hspace{0.25em}\hspace{0.25em}92}}^{235}\text{U}`$ using the different relative yields of Nd isotopes:
$`\delta _{Nd}^9`$ $`=^{150}\text{Nd}/(^{143}\text{Nd}+^{144}\text{Nd})=0.1175`$ for $`_{\mathrm{\hspace{0.25em}\hspace{0.25em}94}}^{239}\mathrm{Pu}`$ and $`\delta _{Nd}^9=0.0566`$ for $`_{\mathrm{\hspace{0.25em}\hspace{0.25em}92}}^{235}\text{U}`$ 29 . However this comparison is masked by the fission of $`_{\mathrm{\hspace{0.25em}\hspace{0.25em}92}}^{238}\text{U}`$ by fast neutrons: $`\delta _{Nd}^8=0.1336`$. Taking account of this contribution one arrives at an estimate of $`d0.6`$ million years 30 . This was the value we adopted in our calculations.
The total energy yield of the reactor has been estimated to be $`1.510^4`$ MWa 31 . Such a fission energy is obtained by two blocks of the Leningrad atomic power station with a hundred percent load in 2.3 years. Assuming a mean duration of $`d=610^5`$ a for the work of the reactor one gets a mean power output of only $`P_P=25`$ kW.
### II.2 Composition and size of the Oklo RZ2 reactor
The cores of the Oklo reactor have been numbered. The most complete data are available for core RZ2 . This core of the Oklo reactor is of the shape of an irregular rectangular plate that lies on a basalt bed at an angle of 45<sup>o</sup>. The thickness of the plate is $`H=1`$ m, its width is $`b=1112`$ m, and its length is $`l=1920`$ m (see Fig.8a in Refs.31 and 22 ). Thus the volume of the RZ2 core is about 240 m<sup>3</sup>. Since in the case of large longitudinal and transverse sizes the shape of the reactor is not essential, we have assumed as a reactor model a flat cylinder of height $`H=1`$ m and radius $`R`$ which is determined by the core burn-up. The energy yield is $`P_Pd=1.510^4`$ MWa $`5.4810^6`$ MWd. At a consumption of $`_{\mathrm{\hspace{0.25em}\hspace{0.25em}92}}^{235}\text{U}`$ of $`g=1.3`$ g/MWd32 , the total amount of burnt up fissile matter is
$$\mathrm{\Delta }M_b=gP_Pd=7.12\mathrm{tons}.$$
(1)
Taking into account that half of the burnt up $`_{\mathrm{\hspace{0.25em}\hspace{0.25em}92}}^{235}\text{U}`$ isotope is replenished from the decay of the produced $`_{\mathrm{\hspace{0.25em}\hspace{0.25em}94}}^{239}\mathrm{Pu}`$, we find the original mass of the burnt up $`_{\mathrm{\hspace{0.25em}\hspace{0.25em}92}}^{235}\text{U}`$:
$$\mathrm{\Delta }M_5=4.75\mathrm{tons}.$$
(2)
In case of a uniform burn-up, the average density of the burnt up $`_{\mathrm{\hspace{0.25em}\hspace{0.25em}92}}^{235}\text{U}`$ is
$$\mathrm{\Delta }\gamma _5(d)=\frac{\mathrm{\Delta }M_5}{\pi R^2H}=\frac{1.51}{R^2}\mathrm{g}/\mathrm{cm}^3,$$
(3)
where $`R`$ is given in meters. The relative average initial burn-up is
$$\overline{y}_5(d)=\frac{\mathrm{\Delta }\gamma _5(d)}{\gamma _{5,i}(0)}=\frac{1.51\mathrm{g}/\mathrm{cm}^3}{\gamma _{5,i}(0)R_i^2}\mathrm{and}R_i=\left[\frac{1.51\mathrm{g}/\mathrm{cm}^3}{\gamma _{5,i}(0)\overline{y}_5(d)}\right]^{1/2}.$$
(4)
Processing the data of Table 2 from Ref.31 gives a value of $`\overline{y}_5(T_0)50\%`$ for the present-day average over the core. In the past it was 1.355 times smaller (see below) on account of the higher concentration $`\gamma _{5,i}(0)`$ of uranium, i.e. $`\overline{y}_5(0)=36.9\%`$. Thus the radius is given by the following formula:
$$R_i=\left[\frac{4.09\mathrm{g}/\mathrm{cm}^3}{\gamma _{5,i}(0)}\right]^{1/2}\mathrm{m}$$
(5)
The approximate composition of the rock in core RZ2 is shown in Table 1 of Ref.7 . On the basis of these data one can calculate the elemental composition of the ore by weight (see the penultimate column of Table1). For comparison we show in the last column the composition by weight from the book of Yu. A. Shukolyukov (Table 2.1 of Ref. 33 ), which was based on early data of R. Naudet 34 . These values coincide within 10$`\%`$. Since all those elements are relatively weakly absorbing, such differences practically do not play any role.
It is much more important to know the amount of uranium and of water at the beginning of the work of the reactor. The connection between the uranium content in the core (in $`\%`$) and the density of the dehydrated core has been measured experimentally in Ref.35 (Fig.1). The content of uranium in the core varies greatly between different samples. To determine the influence of the uranium content on the reactor parameters we have chosen three initial values for the density of uranium in the dehydrated ore: $`Y_{U,i}(T_0)=`$35, 45 and 55%, taken to be constant over the reactor. The value of the ore density $`\gamma _i(T_0)`$ that corresponds to $`Y_{U,i}(T_0)`$ is shown in the second row of Table 2 (see Fig.1). In the fifth row of Table 2 we show the density of the empty rock. The density of water in the reactor is $`0.30.5`$ g/cm<sup>3</sup> 27 . This water consists of bound (crystalline) and unbound water which evaporates after 100<sup>o</sup>C. In our reactor model we assumed a total density of water of $`\gamma _{H_2O}=0.355`$ g/cm<sup>3</sup>, of which 0.155 g/cm<sup>3</sup> was taken for the density of unbound water. Assuming a porosity of about 6% one can take for the density of dry ore with water the same value as for the dehydrate ore.
The density of $`_{\mathrm{\hspace{0.25em}\hspace{0.25em}92}}^{238}`$U of the fresh core RZ2 at the epoch of the formation of the reactor was
$$\gamma _8(0)=\gamma _8(T_0)(1\zeta _5)\mathrm{exp}(+T_0/\tau _8),$$
(6)
where the lifetime of $`_{\mathrm{\hspace{0.25em}\hspace{0.25em}92}}^{238}`$U is $`\tau _8=6.4510^9`$ year. The value of $`\gamma _U^8(0)`$ increases on account of the decay of uranium into lead. The density of $`_{\mathrm{\hspace{0.25em}\hspace{0.25em}92}}^{235}\text{U}`$ ($`\gamma _{U,5}(0)`$, g/cm<sup>3</sup>) in the fresh core ($`\tau _5=1.01510^9`$ year) is
$$\gamma _{U,5}(0)=\gamma _{U,5}(T_0)\zeta _5\mathrm{exp}(+T_0/\tau _5).$$
(7)
The values of $`\gamma _{5,i}(0)`$ and $`\gamma _{8,i}(0)`$ are shown in Table 2. Also in the table are the calculated values of the densities of uranium $`\gamma _{U,i}(0)`$ and of the empty rock (without Pb), and the new fraction of uranium in the dry ore at the beginning of the cycle of the Oklo reactor. From eq.(6) and eq.(7) we get for the ratio $`\gamma _{U,i}(0)/\gamma _{U,i}(T_0)`$
$$\gamma _{U,i}(0)/\gamma _{U,i}(T_0)=(1\zeta _5)\mathrm{exp}(T_0/\tau _8)+\zeta _5\mathrm{exp}(T_0/\tau _5)=1.355,$$
(8)
independent of $`\gamma _{U,i}(T_0)`$. The initial concentration $`N_{k,i}(0)`$ of nuclei, which is needed for the calculations, was calculated from the formula
$$N_{k,i}(0)=\gamma _{k,i}(0)N_A/A_k,$$
(9)
where $`N_A=6.02210^{23}`$ mol<sup>-1</sup> is the Avogadro number and $`A_k`$ is the atomic weight. For $`_{\mathrm{\hspace{0.25em}\hspace{0.25em}92}}^{238}`$U and $`_{\mathrm{\hspace{0.25em}\hspace{0.25em}92}}^{235}\text{U}`$ we have used the data of Table 2; for the atoms of the rock we used the percentages by weight from Table 1. The oxygen content of water was added to the oxygen of the core. The composition of the fresh core RZ2 that we used in calculations with different initial content of uranium is shown in Table 3. Although the accuracy of the densities of some elements in this table is only a few percent, the values of $`N_{k,i}(0)`$ are given with four decimal places for reproducibility of results.
It follows from Table 2 that the enrichment of isotope $`_{\mathrm{\hspace{0.25em}\hspace{0.25em}92}}^{235}\text{U}`$ ($`\zeta _5(0)=\gamma _5(0)/\gamma _U(0)`$) was $`\zeta _5(0)=3.1\%`$ 1.8 billion years ago. Uranium of such enrichment is used in Russian VVER reactors of atomic power stations. Since the ratio of nuclei U/H is about equal and the sizes of both reactors are comparable, one can immediately and without any calculation say that a chain reaction was possible in Oklo 7 .
### II.3 Calculation of the fresh core
#### II.3.1 Semianalytical calculation of core RZ2
Consider first the bare reactor without reflector. Since the reactor is large compared with the neutron migration length $`M`$, one can apply the single-group diffusion theory Wein ; 32 . For a stationary neutron flux $`\mathrm{\Phi }(\stackrel{}{r})`$ the following equation holds:
$$\left[^2+\frac{1}{M^2}\right]\mathrm{\Phi }(\stackrel{}{r})=\frac{K_{\mathrm{}}}{K_{\mathrm{eff}}M^2}\mathrm{\Phi }(\stackrel{}{r}),\mathrm{\Phi }\left(\pm \frac{H}{2}\right)=\mathrm{\Phi }(R)=0.$$
(10)
The solution $`\mathrm{\Phi }(\stackrel{}{r})`$ that satisfies this equation with boundary conditions eq.(10) is
$$\mathrm{\Phi }(\stackrel{}{r})=\mathrm{\Phi }_0\mathrm{cos}\left(\frac{\pi }{H}x\right)J_0\left(\frac{2,405r}{R}\right),$$
(11)
where $`J_0(Br)`$ is the zeroth Bessel function, and the effective multiplication factor is
$$K_{\mathrm{eff}}=\frac{K_{\mathrm{}}}{1+M^2B^2};B^2=B_H^2+B_R^2;B_H=\frac{\pi }{H};B_R=\frac{2.405}{R}.$$
(12)
In Table 4 we show the two constants, $`K_{\mathrm{},i}`$ and $`M_i^2`$, calculated with codes MCNP4C and MCU-REA for three different cores 35a ; 35b . These constants are needed to calculate $`K_{\mathrm{eff},i}`$ by formula (12). The values of $`K_{\mathrm{},i}`$ calculated for one and the same composition differ by a few tenth of a per cent; the values of $`M_i^2=K_{\mathrm{},i}\tau _i+L_i^2`$ differ by a few per cent. In row 9 of Table 4 we show the values of $`K_{\mathrm{eff},i}^{(1)}`$ calculated with the approximate formula (12). They are smaller than the direct calculations using Monte Carlo code (row 1). The difference in reactivity amounts to $`\mathrm{\Delta }\rho _1=(0.20.3)\%`$. The diffusion length in the fresh core is $`L=1.62.1`$ cm, and the total migration length is $`M=67`$ cm. These lengths get less with increasing uranium concentration.
The mean neutron flux, averaged over the reactor, is
$$\overline{\mathrm{\Phi }}=\frac{1}{V}_V\mathrm{\Phi }(\stackrel{}{r})๐\stackrel{}{r}=\mathrm{\Phi }_0\frac{4}{\pi }\frac{J_1(2.405)}{2.405}$$
(13)
($`J_1(2.405)=0.51905`$). From formula (13) we get the following formula of the volume nonuniformity coefficient $`K_V`$, independent of $`R`$ and $`H`$:
$$K_V=\mathrm{\Phi }_0/\overline{\mathrm{\Phi }}=\frac{\pi }{2}\frac{2.405}{2J_1(2.405)}=3.638.$$
(14)
This formula is useful to check the accuracy of calculation of the spatial distribution $`\mathrm{\Phi }(\stackrel{}{r})`$. The absolute value of the mean neutron flux for $`P_P=2.510^2`$ MW is equal to
$$\overline{\mathrm{\Phi }}=\phi \frac{v_f}{E_f}P_P=1.8810^{15}\mathrm{n}/\mathrm{s}\phi ,$$
(15)
where $`\phi `$ is the neutron flux per cm<sup>2</sup> and one fast fission neutron which is calculated with the Monte Carlo code; $`\nu _f/E_f=7.510^{16}`$ n/MW$``$s is the number of fast neutrons per second and a power of 1 MW ($`E_f`$ is the fission energy; $`\nu _f`$ is the number of fast neutrons per fission). For thermal neutrons formula (15) holds with $`\phi _{\mathrm{th}}`$. The mean thermal neutron flux with energies $`E_n<0.625`$ eV is very small in the case of $`Y_{U,2}(0)=49.4\%`$ it is $`\overline{\mathrm{\Phi }}_{\mathrm{th}}=0.6310^8`$ n/cm$`{}_{}{}^{2}`$s. The thermal flux in the center of the core is $`\mathrm{\Phi }_{\mathrm{th}}^{\mathrm{max}}(0)=2.0010^8`$ n/cm<sup>2</sup>s. The total mean flux, integrated over all energies, is equal to $`\overline{\mathrm{\Phi }}=3.910^8`$ n/cm<sup>2</sup>s. These results were found using code MCU-REA. The results of calculations using other Monte Carlo codes are similar (see Table 5). The low neutron flux determines the specifics of the function of the reactor.
As a reflector one can assume the same core but without uranium. The analytical calculations for the reactor with reflector are more cumbersome. Therefore we used numerical methods for these calculations. The results are shown in Table 6. Both Monte Carlo programs give values of the reactivity reserve for variant 1 of the core which coincide within the statistical accuracy. The difference of reactivity is 0.26% for core RZ2 of the bare reactor and 0.46% for core 3. For the reactor with reflector the difference is smaller: 0.20% and 0.39%, respectively. Since the migration length is small, a reflector of thickness $`\mathrm{\Delta }=0.5`$ m is practically infinite: the results for a reflector of thickness $`\mathrm{\Delta }=0.5`$ m coincide with those for $`\mathrm{\Delta }=1`$ m. Compared with the bare reactor, the reflector makes a contribution of $`\delta \rho (\mathrm{\Delta })=0.8รท0.9\%`$. This contribution drops with increasing uranium content in the core. The cold reactor with a fresh core is strongly overcritical, since temperature and void effects have not yet been taken into account, also the initial strong absorbers which afterwards burn up rapidly. In Table 7 we show the neutron capture in the infinite fresh core per fast fission neutron. Capture by $`_{\mathrm{\hspace{0.25em}\hspace{0.25em}92}}^{235}\text{U}`$ amounts to 55.7% and by $`_{\mathrm{\hspace{0.25em}\hspace{0.25em}92}}^{238}`$U to 33.8%. These are followed by hydrogen (3.9%), iron (3.8%), silicon (0.8%) etc. Code MCNP4C gives similar values.
#### II.3.2 Power effect
Reactor Oklo is controlled by the core temperature $`T_C`$ 36 . During heating the water was driven out of the core until the multiplication factor was equal to one. At first the large overcriticality was compensated by the power effect, which is the sum of the temperature and void effects. In Table 8 and Fig.2 we show the dependence of the water density on the temperature for several pressures. At a pressure of 100 MPa in the Oklo reactor and $`T_C=700`$K, the density of water is 65% of its value for $`T_C=300`$K and normal pressure. In this case the difference between crystalline and free water disappears apparently. The power effects is shown in Fig.3 35a ; 35b . Near $`T_C=700`$K all the $`K_{\mathrm{eff}i}`$ become equal to one. Therefore we can assume $`T_C700`$K as the most likely temperature of the fresh active core (we neglect the small difference between the temperatures of fuel and water). For the variant of the composition of core RZ2 , the power effect is $`\mathrm{\Delta }\rho _P=11.6\%`$. The void effect at 700 K accounts for 73% of this value, and the temperature effect for 27%. These results were obtained with code MCNP4C. Code MCU-REA gives similar values. In Table 4 we show the numerical values of $`K_{\mathrm{eff}}(T_C)`$, $`K_{\mathrm{}}(T_C)`$ and $`M^2(T_C)`$, calculated with code MCU-REA for the bare reactor with core composition $`i=2`$. With increased temperature $`K_{\mathrm{}}(T_C)`$ drops and $`M^2`$ and the leakage increase (Fig.4).
To determine the spread of results depending on the uncertainty of the initial composition of the core the calculations were carried out over a wide range of the content of uranium ($`Y_{U,i}(0)=39.459.6\%`$ by weight) and of water ($`\omega _{H_2O}^0=0.3550.455`$) in the ore. For reliability the calculations for the bare reactor were carried out with two codes: MCU-REA and MCNP4C. The results were additionally controlled by the single-group formula (12) with the parameters shown in Fig.4. The calculations of the $`T_C`$ dependence of $`K_{\mathrm{eff}}`$ for variants 2 and 3 are similar (Fig.3). For variant 1 with lower uranium content ($`Y_{U1}(0)=38.4\%`$ by weight) the curve $`K_{\mathrm{eff}}(T_C)`$ lies visibly lower. For a water content of $`\omega _{H_2O}^0=0.455`$ the curve for variant 2 lies significantly higher. As a result the core temperature at which the reactor became critical was $`T_C=725`$K with a spread of $`\pm 55`$K. Taking account of the fuel burn-up and of slogging of the reactor a loss of reactivity takes place. This leads to a drop of $`T_C`$. The calculations of burn-up are continued at present.
The reactor could have worked also in a pulsating mode: when the temperature exceeded 710 K, then the unbound water was boiled away and the reactor stopped on account of the void effect. Then the water returned and it started to work again 27 ; 36 ; 37a ; 37b . However a detailed analysis of the pulsating mode of operation of the reactor is outside the scope of the present paper.
#### II.3.3 Computational problems in calculations of large reactors
When making Monte Carlo calculations of the neutron flux in large reactors, one encounters certain difficulties. In such reactors many generations are produced before a neutron that was created in the center of the reactor reaches its boundary. This time depends on the relation between the migration length and the size of the reactor 38 . At $`T_C=300`$K these values are for core $`i=2`$ equal to $`M=7`$ cm and $`R=8.1`$ m. 230 generations are needed before a centrally produced neutron reaches the boundary, detects the boundary condition and returns to the center. In order to reproduce the spatial distribution of the neutron flux with sufficient accuracy one must calculate tens of such journeys. Analyzing the solution of the time dependent diffusion equation one finds that over 6000 cycles are needed to get the fundamental harmonic with an accuracy of a few per cent. Experience with such calculations shows that one needs $`(510)10^3`$ histories per cycle in order to keep an acceptable statistical accuracy. Thus we needed $`(46)10^7`$ neutron trajectories for our calculations. For several hundred calculations we have explored an order of 10<sup>10</sup> trajectories taking up several months of continuous work of a modern PC cluster. In spite of such a large volume of calculations we could not find the volume nonuniformity coefficient $`K_V`$ of the neutron flux with good accuracy. To do these calculations we had to divide the core into tens of volume elements which led to a reduction of the statistical accuracy in each of them. As a result the value of $`K_V`$ in formula (14) was reproducible with an accuracy not better than 10%. This is obviously insufficient to carry out the calculation of the burn-up that depends on the magnitude of the absolute flux in different parts of the core. One must admit that the Monte Carlo method is not suitable for the calculation of large reactors and one must resort to different approaches.
The reactor neutron spectrum below 0.625 eV is needed in order to average the cross sections of strong absorbers (e.g. $`_{\mathrm{\hspace{0.25em}\hspace{0.25em}62}}^{149}\text{Sm}`$). The spectrum for three compositions of the fresh core without reflector, calculated with code MCNP4C for $`T_C=300`$K, is shown in Fig.5. One can see small peaks which correspond to excitations of rotational and vibrational levels of H<sub>2</sub>O. For comparison we also show in Fig. 5 the Maxwell neutron spectrum that was used by all previous authors to average the $`_{\mathrm{\hspace{0.25em}\hspace{0.25em}62}}^{149}\text{Sm}`$ cross section 6a ; 6b ; 6c ; 7 ; 8a ; 8b ; 9 . The spectra are significantly different. The Maxwell spectrum has a much higher peak but is exponentially small above 0.3 eV where the reactor spectrum is a Fermi distribution. In our calculations we have used the Nelkin model of water which automatically takes account of the chemical bond of hydrogen nuclei. Calculations at other values of $`T_C`$ yield similar results.
## III Variation of fundamental constants
The Oklo reactor is an instrument that is sensitive to the neutron cross sections in the distant past. Comparing them with current values one can estimate how constant they, and hence also the fundamental constants, are in time 6a ; 6b ; 6c ; 7 .
### III.1 Early approaches
In 1935 E. Miln posed the question: how do we know that the fundamental constants are actually constant in time 39 . He thought that the answer could be found only by experiment. A little later D. Dirac proposed that originally all constants were of one order of magnitude but that the gravitational constant dropped at a rate of $`\dot{G}/Gt_0^1`$ during the lifetime $`t_0`$ of the universe 40a ; 40b . In 1967 G. Gamov suggested that, on the contrary, the electromagnetic constant is increasing: $`\dot{\alpha }/\alpha t_0`$ 41 . Both hypothesis were wrong since they contradicted geological and paleobotanical data from the early history of the Earth. Without entering into a detailed discussion of these and many other later publications on this subject one must admit that there is a problem of the experimental limit on the rate of change of the fundamental constants (see the early review by F. Dyson 42 ).
The authors of Refs. 6a ; 6b ; 6c ; 7 noticed that the sensitivity to variations of the nuclear potential increases by several orders of magnitude if one considers neutron capture. Owing to the sharp resonances of the absorption cross section the nucleus is a finely tuned neutron receiver. A resonance shifts on the energy scale with changes of the nuclear potential similarly as the frequency of an ordinary radio receiver shifts when the parameters of the resonance circuit are changed (Fig.6) 43 . Qualitatively one can understand the absence of a significant shift of the near-threshold resonances on the grounds that all strong absorbers are highly burnt up in the Oklo reactor and weak absorbers are burnt up weakly (Fig.7) 44 ; 7 . Holes in the distributions are seen for strong absorbers: $`_{\mathrm{\hspace{0.25em}\hspace{0.25em}62}}^{149}`$Sm, $`_{\mathrm{\hspace{0.25em}\hspace{0.25em}63}}^{151}`$Eu, $`_{\mathrm{\hspace{0.25em}\hspace{0.25em}64}}^{155}`$Gd, $`_{\mathrm{\hspace{0.25em}\hspace{0.25em}64}}^{157}`$Gd. The depth of burn-up, calculated using the present absorption values, are in satisfactory agreement with experiment, particularly if one remembers that the neutron spectrum over which one must average the cross section is not known very well. Thus in the 1.8 billion years since the work of the Oklo reactor, the resonances (or, in other words, the levels of the compound nuclei) have shifted by less than $`\mathrm{\Delta }E_r\mathrm{\Gamma }\gamma /2`$ ($`\mathrm{\Gamma }_\gamma =0.1`$ eV). Therefore the average rate of the shift did not exceed $`310^{11}`$ eV/year. This value is at least three orders of magnitude less than the experimental limit on the rate of change of the transition energy in the decay of $`^{187}`$Re 42 .
At present there are no theoretical calculations giving a reliable connection between the positions of all resonances with parameters of the nuclear potential. But already the preliminary qualitative estimates allow one to reduce the limits on the rates of change of the coupling constants of the strong and electromagnetic interactions $`\dot{\overline{\alpha }}/\alpha `$ and $`\dot{\overline{\delta \alpha }}/\alpha `$. We confirm the absence of a power or logarithmic dependence on the lifetime of the universe. It is desirable to have a more detailed calculation of the influence of variations of the fundamental constants on the parameters of the neutron resonances.
### III.2 Basic formulae
#### III.2.1 Averaging the Breit-Wigner formula
When a slow neutron is captured by a nucleus of isotope $`_{\mathrm{\hspace{0.25em}\hspace{0.25em}62}}^{149}\text{Sm}`$, then a nuclear reaction takes place with formation of an excited intermediate compound nucleus and subsequent emission of $`m`$ $`\gamma `$โquanta:
$$n+_{\mathrm{\hspace{0.25em}\hspace{0.25em}62}}^{149}\text{Sm}_{\mathrm{\hspace{0.25em}\hspace{0.25em}62}}^{150}\mathrm{Sm}^{}_{\mathrm{\hspace{0.25em}\hspace{0.25em}62}}^{150}\mathrm{Sm}+m\gamma .$$
(16)
Near a strong $`S`$ resonance one can neglect the effect of the other resonances and describe the cross section with the Breit-Wigner formula
$$\sigma _{\gamma ,Sm}(E_C)=g_0\frac{\pi \text{}\text{}\text{h}^2}{2m_nE_C}\frac{\mathrm{\Gamma }_n(E_C)\mathrm{\Gamma }_\gamma }{(E_CE_r)^2+\mathrm{\Gamma }_{tot}^2/4},$$
(17)
where $`g_0=(2J+1)/(2S+1)(2I+1)`$ is the statistical factor, $`S=1/2`$ is the electron spin, $`I`$ is the nuclear spin and $`J`$ is the spin of the compound nucleus. The full width is $`\mathrm{\Gamma }_{tot}=\mathrm{\Gamma }_n(E)+\mathrm{\Gamma }_\gamma `$, where $`\mathrm{\Gamma }_n(E)`$ and $`\mathrm{\Gamma }_\gamma `$ are the neutron and $`\gamma `$ width, respectively. The neutron width is given by 45
$$\mathrm{\Gamma }_n(E_C)=\mathrm{\Gamma }_n^0\sqrt{\frac{E_C}{E_0}};E_0=1\mathrm{eV}.$$
(18)
The parameters of the lowest resonances of a number of absorbers is given in Table 9.
In formula (17) the neutron energy is given in the c.m. frame: $`E_c=\frac{1}{2}m_n\stackrel{}{V}_L\stackrel{}{V}_k^2`$. It depends on the velocities of the nucleus $`\stackrel{}{V}_k`$ and the neutron $`\stackrel{}{V}_L`$ in the lab frame and on the reduced mass $`m_n`$. The reaction rate $`N_k\sigma _{\gamma ,k}(E_C)V_C`$ with cross section (17) and for an absorber of nuclear density $`N_k`$ must be averaged over the nuclear spectra $`f_k(E_k)`$ and the neutron spectrum $`n(E_L)`$ (all spectra are normalized to one). The inverse nuclear burn-up time in an arbitrary point of the core is given by
$$\lambda _{\gamma ,k}(T)=N_k๐\stackrel{}{p}_k๐\stackrel{}{p}_Lf_k(E_k)n(E_L)\sigma _{\gamma ,k}(E_C)V_C.$$
(19)
At high temperatures the gas approximation is valid for heavy nuclei of the absorber. Changing to integration over the c.m. energy $`E_C`$ and the neutron energy $`E_L`$ and assuming a Maxwell nuclear spectrum, we get
$$\lambda _{\gamma ,k}(T)=N_kn(E_L)\sigma _{\gamma ,k}(E_C)V_CF(E_CE_L)๐E_L๐E_C,$$
(20)
$`F(E_CE_L)`$ is the transformation function from the c.m. to the lab system (for details see Ref. 46 ):
$$F(E_CE_L)=\frac{(A+1)}{2\sqrt{\pi ATE_L}}\left\{\mathrm{exp}\left[\frac{A}{T}\left(\sqrt{(1+\frac{1}{A})E_C}\sqrt{E_L}\right)^2\right]\mathrm{exp}\left[\frac{A}{T}\left(\sqrt{(1+\frac{1}{A})E_C}+\sqrt{E_L}\right)^2\right]\right\}$$
(21)
and $`A=M_A/m_n`$ is the mass of nucleus $`A`$ in units of the neutron mass.
Close to a resonance we can neglect the second term in Eq.(21) and evaluate the first term in integral (20) by the saddle-point method. As a result the integral (20) takes on the following form in the vicinity of a resonance:
$$\lambda _{\gamma ,k}(T)=N_k\frac{\pi }{2}\left(1+\frac{1}{A}\right)๐E_C\sigma _{\gamma ,k}(E_C)V_C๐E_Ln(E_L)\mathrm{\Gamma }\left[E_L\left(1+\frac{1}{A}\right)E_C\right],$$
(22)
where the Gaussian
$$\mathrm{\Gamma }\left[E_L\left(1+\frac{1}{A}\right)E_C\right]=\frac{1}{\sqrt{\pi }\mathrm{\Delta }_D}\mathrm{exp}\left\{\frac{\left[E_L\left(1+\frac{1}{A}\right)E_C\right]^2}{\mathrm{\Delta }_D^2}\right\}$$
(23)
is normalized to one and the Doppler width is equal to
$$\mathrm{\Delta }_D=\left[\frac{4E_LT}{A}\right]^{1/2}=\left[\frac{4E_CT}{A+1}\right]^{1/2}$$
(24)
The values of the Doppler widths for $`T=700`$K are shown in Table 9. Since all $`\mathrm{\Delta }_D\mathrm{\Gamma }_\gamma `$, function (23) can be replaced by $`\delta \left(E_LAE_C/(A+1)\right)`$ and integral (22) becomes
$$\lambda _{\gamma ,k}=N_k\frac{\pi }{2}\left(1+\frac{1}{A}\right)^2\sigma _{\gamma ,k}(E_C)V_Cn\left[\left(1+\frac{1}{A}\right)E_C\right]๐E_C.$$
(25)
The correction $`2/A`$ is of magnitude 1%. If the neutron spectrum is Maxwellian in the c.m. frame, then it is also Maxwellian (with reduced neutron mass) when the nuclear motion is taken into account. It can be shown that Eq.(25) is valid at all energies if the distribution of nuclei and neutrons is Maxwellian 47 .
Therefore it is not surprising that the authors of Ref. 9 did not notice any deviations from formula (25) in their numerical calculation that took account of the thermal motion of the target nuclei (Doppler effect). However, the situation is different if the neutron spectrum is not Maxwellian. In this case one must use formula (22) instead of the simple formula (25).
To average the capture cross section of samarium one normalizes the cross section, integrated over the neutron flux spectrum $`(n(E,T)v)`$, traditionally not by the integrated flux but by the product of the velocity $`v_n^0=2200`$ m/s and the integrated neutron density $`(n(E))`$ 47 ; 8b ; 9 :
$$\widehat{\sigma }_{\gamma ,k}(T)=\frac{\sigma _{\gamma ,k}(E_L)n(E_L)v_L๐E_L}{v_n^0n(E_L)๐E_L}.$$
(26)
If the cross section $`\sigma _{\gamma ,k}(E_L)`$ has a $`1/v_L`$ behaviour, then the integral of $`\widehat{\sigma }_{\gamma ,k}(T)`$ is constant. From formula (26) one has
$$\widehat{\sigma }_{\gamma ,k}(T)=\sqrt{\frac{4T}{\pi T_0}}\sigma _{\gamma ,k}(T),$$
(27)
where $`T_0=300`$ K $`=25.9`$ meV. Useful is also the relation 9
$$\sigma \mathrm{\Phi }=\widehat{\sigma }\widehat{\mathrm{\Phi }},\mathrm{where}\widehat{\mathrm{\Phi }}=\sqrt{\frac{\pi }{2}\frac{T_0}{T}}\mathrm{\Phi }.$$
(28)
We have evaluated the cross section $`\widehat{\sigma }_{\gamma ,Sm}(T)`$ of $`_{\mathrm{\hspace{0.25em}\hspace{0.25em}62}}^{149}\text{Sm}`$ without recourse to any approximations. For $`\lambda _{\gamma ,Sm}(T)/(N_{Sm}v_n^0)`$ we have
$$\widehat{\sigma }_{\gamma ,Sm}(T)=\frac{\sqrt{\pi }๐E_k๐E_Lf_{Sm}(E_k)\sigma _{\gamma ,Sm}(E_C)V_Cn(E_L)}{v_n^0๐E_Ln(E_L)}.$$
(29)
For the calculations we used the computer package MATHEMATICA 48 . In Fig.8 we show the values of $`\widehat{\sigma }_{\gamma ,Sm}(T,\mathrm{\Delta }E_r)`$ at six temperatures $`T=3001000`$K for a shift of the resonance position $`\mathrm{\Delta }E_r=\pm 0.2`$ eV. The curves have a maximum at negative shifts of the resonance; the maximum of the curves is higher at lower temperature $`T`$. At the point $`\mathrm{\Delta }E_r=0`$ and at $`T=293`$K the cross section calculated as the contribution of the closest resonance is equal to $`\sigma _{\gamma ,Sm}(293\text{K})=39.2`$ kb. The contribution of higher positive resonances is $`\sigma _{\gamma ,Sm}^+(293\text{K})=0.6`$ kb and negative ones is $`\sigma _{\gamma ,Sm}^{}(293\text{K})=0.3`$ kb 45 . The total cross section (as measured on a neutron beam) is $`\sigma _{\gamma ,Sm}^{tot}(293\text{K})=40.1`$ kb. At small energy shifts $`\mathrm{\Delta }E_r`$ $`\sigma _{\gamma ,Sm}^+`$ and $`\sigma _{\gamma ,Sm}^{}`$ practically do not change. In Refs.8a ; 8b ; 9 the total cross section $`40.1`$ kb has been used instead of the single resonance one $`39.2`$ kb.Therefore the curves in Refs. 8a ; 8b ; 9 are higher by 40.1 kb/39.2 kb, i.e. by 2.5%.
#### III.2.2 Taking account of the reactor spectrum
In Fig.LABEL:Fig9 we show the results of calculating $`\widehat{\sigma }_{\gamma ,Sm}(T_C,\mathrm{\Delta }E_r)`$ with the Maxwell spectrum replaced by the reactor spectrum $`n_R(E_L,T_C)`$. The central curve 2 is the result of the calculation using code MCNP4C for the fresh core with $`Y_{U2}(0)=49.4\%`$U in the dry ore and with $`\omega _{H_2O}^0=0.405`$ at $`T=725`$K. For comparison we also show the cross section averaged over the Maxwell spectrum at $`T=725`$K for the same composition of the core (curve 4). Curves 4 and 2 are significantly different, especially at negative $`\mathrm{\Delta }E_r`$: curve 2 lies distinctly lower. The maximum of curve 2 is 1.5 times lower than the maximum of curve 4. At lower temperatures this difference is even greater. Thus we conclude that we have proved a significant effect of the reactor spectrum on the cross section of $`_{\mathrm{\hspace{0.25em}\hspace{0.25em}62}}^{149}\text{Sm}`$.
In order to determine the dependence of $`\widehat{\sigma }_{\gamma ,Sm}(T_C,\mathrm{\Delta }E_r)`$ on the uncertainty in the initial active core composition, we have calculated the values for the two outermost curves of Fig.9. Curve 3 of this figure corresponds to an initial content of $`Y_{U1}(0)=38.4\%`$U in the ore, $`\omega _{H_2O}^0=0.355`$ and $`T_C=670`$K; curve 1 corresponds to an initial content of 49.4%, $`\omega _{H_2O}^0=0.455`$ at $`T=780`$K. Since the numerical constants are known only for values of $`T_C`$ which are multiples of 100, we have done the calculations for $`T_C=(600,700,800)`$K and interpolated to intermediate temperatures. The broadening of curve 2 on account of the scatter of temperatures is small. Experimental data of $`\widehat{\sigma }_{\gamma ,Sm}^{\mathrm{Exp}}(T)`$ for core 2 are presented in Ref. 8a (Table 10) (see also Refs. 52 ). The labeling of sample SC36-1418 indicates that the sample was taken from bore-hole SC36 at a depth of 14 m 18 cm. The mean value $`\overline{\widehat{\sigma }}_{\gamma ,Sm}^{\mathrm{Exp}}=(73.2\pm 9.4)`$ kb is shown in Fig.9. Curve 1 ($`T_C=670`$K) crosses the lower limit of $`\overline{\widehat{\sigma }}_{\gamma ,Sm}^{\mathrm{Exp}}=64`$ kb to the left of point $`\mathrm{\Delta }E_r^{(1)}=73`$ meV, and curve 3 ($`T_C=780`$K) to the right at $`\mathrm{\Delta }E_r^{(2)}=+62`$ meV. The possible shift of the resonance is therefore given by these limits:
$$73\mathrm{meV}\mathrm{\Delta }E_r62\mathrm{meV}.$$
(30)
#### III.2.3 Connection between $`\mathrm{\Delta }E_r`$ and $`\dot{\overline{\delta \alpha }}/\alpha `$
The shift $`\mathrm{\Delta }E_r`$ must be related to a variation of the fundamental constants, for instance to a shift of the electromagnetic constant $`\alpha =1/137.036`$. This has been done by Damour and Dyson 8a ; 8b . The change of the Coulomb energy contribution $`\mathrm{\Delta }H_C`$ to the energy of the level in the nuclear potential $`\mathrm{\Delta }E_C`$, that results from a change of $`\alpha `$, is given by
$$\frac{}{\alpha }\mathrm{\Delta }E_C=\frac{}{\alpha }\mathrm{\Delta }H_C.$$
(31)
In first perturbative approximation the dominant contribution to $`\mathrm{\Delta }E_C`$ is given by the isotopic effect (51 , p.568)
$$\mathrm{\Delta }E_C=E_R=\frac{2\pi }{3}\psi _e(0)^2Ze^2<R^2>.$$
(32)
Here $`\psi _e(0)`$ is the wave function of the $`s`$ wave electrons in the nucleons and $`R^2=(Ze)^1\rho R^2๐V`$, where $`\rho (R)`$ describes the proton charge distribution in the nucleus. Damour and Dyson have estimated the value of $`R^2`$ for the excited nucleus $`_{\mathrm{\hspace{0.25em}\hspace{0.25em}62}}^{150}`$Sm from the neighbouring isotopes. They found
$$=\mathrm{\Delta }E_R=(1.1\pm 0.1)\mathrm{MeV}.$$
(33)
Combining this value with the shift of the resonance $`\mathrm{\Delta }E_r`$ in formula (30), we get for $`\beta \mathrm{\Delta }E_r/`$
$$5.610^8<\overline{\delta \alpha }/\alpha <6.610^8.$$
(34)
Because of the negative value of $``$ the limits on $`\overline{\delta \alpha }/\alpha `$ change their places. For the past time ($`T_0`$) the product ($`T_0`$) is positive and hence the limits on $`\dot{\overline{\delta \alpha }}/\alpha \mathrm{\Delta }E_r(T_0)`$ are restored to their previous places. Note that traditionally $`\overline{\delta \alpha }/\alpha `$ is defined by $`\delta \alpha =(\alpha _{Oklo}\alpha _{now})/\alpha `$. This shift of $`\alpha `$ lies in a narrower range than in Ref. 8a ; 8b . Assuming a linear change of the e.m. constant during the time $`T_0`$, we get the following limit on the relative rate of change:
$$3.710^{17}\mathrm{year}^1<\dot{\overline{\delta \alpha }}/\alpha <3.110^{17}\mathrm{year}^1.$$
(35)
Thus, within the limits given by Eq. (35), the e.m. constant changes for the fresh reactor Oklo core with zero speed, i.e. it remains constant.
### III.3 Review of previous work
The work of Shlyakhter (1976, 1983) 6a ; 6b ; 6c . The authors of Refs. 6a ; 6b ; 6c ; 7 were the first to point out the possibility of using the data of the natural nuclear reactor Oklo to find the most precise limits on the rate of change of the fundamental constants. The most convenient data are those of the strong absorbers, e.g. of $`_{\mathrm{\hspace{0.25em}\hspace{0.25em}62}}^{149}\text{Sm}`$. For this isotope Shlyakhter calculated at $`T=300`$K the dependence of the change of the cross section on the resonance by an amount of $`\mathrm{\Delta }E_r`$: $`\sigma _{\gamma ,Sm}(T_C,\mathrm{\Delta }E_r)`$ (Fig.10) 6b .
He compared this curve with the experimental data available at the time: $`\overline{\sigma }_{\gamma ,Sm}^{\mathrm{Exp}}=(55\pm 8)`$ kb. A possible shift of the first resonance within two standard deviations (95% confidence level) was found to be
$$\delta E_{r2}^{\mathrm{Exp}}20\mathrm{meV}.$$
(36)
(In going from $`\overline{\sigma }_{\gamma ,Sm}(T,\mathrm{\Delta }E_r)`$ to $`\widehat{\sigma }_{\gamma ,Sm}(T,\mathrm{\Delta }E_r)`$, all values must be multiplied by 1.18 according to Eq. (27) and $`\overline{\widehat{\sigma }}_{\gamma ,Sm}^{\mathrm{Exp}}=(65\pm 9.5)`$ kb, but this does not affect the value of $`\delta E_{r2}^{\mathrm{Exp}}`$). To estimate $``$, Shlyakhter used data on the compressibility of the nucleus; he found $`=2`$ MeV 6b . With a linear dependence of the change of $`\alpha `$ with time for $`T_0=210^9`$ years the limit on the rate of change of $`\alpha `$ is
$$\dot{\overline{\delta \alpha }}/\alpha 0.510^{17}\mathrm{year}^1.$$
(37)
Using the present, more accurate value of $`1.1`$ MeV, we get
$$\dot{\overline{\delta \alpha }}/\alpha 110^{17}\mathrm{year}^1.$$
(38)
It should emphasized again that this limit was found for only one temperature: $`T=300`$K.
The work of Petrov (1977) 7 . In this paper the resonance shift $`\delta E_r`$ was estimated from the shifts of the widths of a few strong absorbers. The resonances of strong absorbers lie close to a zero energy of the neutron, and the resonance energy is of the capture width: $`\delta E_r\mathrm{\Gamma }_\gamma 0.1`$. The capture cross section of these absorbers for thermal neutrons changes sharply when the resonance is shifted by an amount of the order of $`\mathrm{\Gamma }_\gamma /2`$. The analysis of experimental data for $`_{\mathrm{\hspace{0.25em}\hspace{0.25em}62}}^{149}\text{Sm}`$ and $`_{\mathrm{\hspace{0.25em}\hspace{0.25em}63}}^{151}`$Eu, taking account of a threefold standard deviation and the uncertainly of the core temperature, shows that the shift $`\delta E_r`$ of the resonance since the activity of the Oklo reactor does not exceed $`\pm `$0.05 eV 7 . The results of the measurement of the concentration of rare earth elements with respect to $`_{\mathrm{\hspace{0.25em}\hspace{0.25em}60}}^{143}`$Nd (the second branch of the mass distribution of the fission fragments) in one of the Oklo samples are shown in Fig.7. A more conservative estimate in Ref. 7 is $`\overline{\delta E_r}50`$ meV, i.e. 2.5 times higher than Shlyakhterโs estimate. Using the modern value $`1.1`$ MeV, we find for $`\overline{\delta \alpha }/\alpha \overline{\delta E_r}/`$
$$\overline{\delta \alpha }/\alpha =\overline{\delta E_r}/4.510^8.$$
(39)
This is almost 5 times greater than Shlyakhterโs optimistic estimate. For the rate of change $`\dot{\overline{\delta \alpha }}/\alpha `$ we get
$$\dot{\overline{\delta \alpha }}/\alpha 2.510^{17}\mathrm{year}^1.$$
(40)
This is less by a factor of 2 than the limit found 20 years later by Damour and Dyson 8a ; 8b . The reason of this discrepancy is the use of only one temperature ($`T_C=300`$K). Although the dependence of $`\mathrm{\Delta }E_r`$ on $`T_C`$ was noted in Ref. 7 , no calculations of the effect of the temperature were carried out.
The work of Damour and Dyson (DD) (1996) 8a ; 8b . The dependence $`\mathrm{\Delta }E_r(T_C)`$ was analysed 20 years later in the paper DD 8a ; 8b . They have repeated the analysis of Shlyakhter and came to the conclusion that it was correct. DD also updated Shlyakhterโs data in three directions:
(i) They employed a large amount of experimental data (see Table 10).
(ii) They have taken account of the great uncertainty of the reactor temperature, $`T_C=(4501000)^0`$C (Fig. 11). As a result they made a conservative estimate of the mean shift of the resonance
$$120\mathrm{meV}\overline{\mathrm{\Delta }}E_r90\mathrm{meV}.$$
(41)
The range of the shift $`\mathrm{\Delta }E_{r1}\mathrm{\Delta }E_{r2}=210`$ meV is 1.5 times greater than the range of the shift in our paper.
(iii) They have calculated the value $`=(1.1\pm 0.1)`$ MeV, but used $`=1`$ MeV. For $`\overline{\delta \alpha }/\alpha `$ DD found
$$9.010^8\overline{\delta \alpha }/\alpha 1210^8.$$
(42)
This leads to the following limits on the rate of change $`\dot{\overline{\delta \alpha }}/\alpha `$:
$$6.710^{17}\mathrm{year}^1\dot{\overline{\delta \alpha }}/\alpha 5.010^{17}\mathrm{year}^1.$$
(43)
Since $`_{\mathrm{\hspace{0.25em}\hspace{0.25em}62}}^{149}\text{Sm}`$ burns up 100 times faster than $`_{\mathrm{\hspace{0.25em}\hspace{0.25em}92}}^{235}\text{U}`$, therefore the only $`_{\mathrm{\hspace{0.25em}\hspace{0.25em}62}}^{149}\text{Sm}`$ found in the stopped reactor is that which was produced immediately before the end of the cycle. As a consequence DD emphasized that one must know the detailed distribution of the nuclear reaction products of the end of the cycle to make a detailed analysis.
The work of Fujii et al. (2000, 2002) 9 ; 53 . The experimental data on the measurement of $`\widehat{\sigma }_{\gamma ,k}`$ of strong absorbers were presented in the papers of Fujii et al. Of five experimental points, four are from core RZ10 (Table 11) and one from another core, RZ13 , and therefore we have omitted it. Core RZ10 lies at a depth of about 150 m from the surface of the quarry. As we do not have any detailed data on the size and composition of core RZ10 , we shall assume them to be similar to those of core RZ2 . For $`_{\mathrm{\hspace{0.25em}\hspace{0.25em}62}}^{149}\text{Sm}`$ the mean value of the four points of Table 11 is
$$\overline{\widehat{\sigma }}_{\gamma ,Sm}^{\mathrm{Exp}}=(90.7\pm 8.2)\mathrm{kb}.$$
(44)
This value is noticeably greater than for core RZ2 \[(72.3$`\pm `$9.4) kb (see subsection 9). The error bars of both values do not even touch and their difference remains significant. For increased reliability the number of measurements for core RZ10 should be increased. Possibly this core has finished its cycle at a lower temperature.
In Fig.12 we show the dependence of $`\widehat{\sigma }_{\gamma ,Sm}(T_C,\mathrm{\Delta }E_r)`$ on the shift of the resonance for a Maxwell distribution 9 ; 52 . The authors have estimated (on the grounds of indirect considerations) the uncertainty in $`T_C=(180400)^0\mathrm{C}=(453673)`$K. They also show the experimental data of formula (44). The intersection of the limiting curves with the lower limit $`\widehat{\sigma }_{\gamma ,Sm}^{\mathrm{Exp}}=82.5`$ kb yield the following possible shift of $`\mathrm{\Delta }E_r`$:
$$105\mathrm{meV}<\mathrm{\Delta }E_r<+20\mathrm{meV}.$$
(45)
In Fig.11 we also show for comparison the curves $`\widehat{\sigma }_{\gamma ,Sm}(T_C,\mathrm{\Delta }E_r)`$ for the reactor spectrum of the fresh core at $`T_C=400^0`$C, $`Y_{U1}(0)=38.4\%`$U in the ore, $`\omega _{H_2O}^0=0.355`$ and $`P=100`$ MPa. The possible shift of $`\mathrm{\Delta }E_r`$ for the experimental data of formula (44) lie in a narrower interval than in the paper DD 8a ; 8b :
$$120\mathrm{meV}\mathrm{\Delta }E_r20\mathrm{meV}.$$
(46)
From Eq.(45) and using $`=1.1`$ MeV we get
$$1.810^8\overline{\delta \alpha /\alpha }9.510^8$$
(47)
and
$$5.310^{17}\mathrm{year}^1\dot{\overline{\delta \alpha /\alpha }}1.010^{17}\mathrm{year}^1.$$
(48)
Thus in this case too we do not find with certainty a nonzero deviation of the change of $`\alpha `$.
The work of Lamoreaux and Torgerson (LT) (2004) 53 ; 54 . These authors noted (two years after Ref. 12 ) that the reactor spectrum contains in addition to the Maxwell tail also the spectrum of the moderated neutrons (Fermi spectrum). They averaged $`\overline{\widehat{\sigma }}_{\gamma ,Sm}(T_C,\mathrm{\Delta }E_r)`$ at $`T_C=600`$K over this spectrum and, comparing this curve with the experimental data of Fujii et al. 53 , found a shift $`\mathrm{\Delta }E_r=(45_{15}^{+7})`$ meV. Thus they found a shift of the cross section of $`_{\mathrm{\hspace{0.25em}\hspace{0.25em}62}}^{149}\text{Sm}`$. Let us consider the LT model in more detail in order to understand this result. The age of the Oklo reactor is $`T_0=210^9`$ year. The relative content of $`_{\mathrm{\hspace{0.25em}\hspace{0.25em}92}}^{235}\text{U}`$ is $`\xi _5(0)=3.7\%`$ (and not 3.1%). The ratio of H to U is $`f_H=N_H/N_U=3`$. The cross section of the burning up admixtures (e.g. lithium) per atom of U is $`\beta _U=_iN_i\sigma _a^i/N_U=2`$ b. The ratio of the thermal neutron capture cross section to the slowing down cross section of epithermal neutrons is $`\mathrm{\Delta }=\mathrm{\Sigma }_a(V_T/(\mathrm{\Sigma }_S/2A)=2`$. The temperature of the core is $`T_C=600`$K$`=327^0`$C.
The following comments are appropriate concerning these parameters. The age of the reactor is $`1.810^9`$ year and not $`210^9`$ year. The value of $`\xi _5(0)=3.7\%`$ yields cross sections $`\mathrm{\Sigma }_{5,a}(0)`$ and $`\mathrm{\Sigma }_{5,f}(0)`$ which are too large by 1.2 times. The ratio $`f_H=N_H/N_U=3`$ holds approximately for $`Y_{U3}(0)=59.6\%`$ and $`\omega _{H_2O}^0=0.355`$ (see Table 3). For the condition $`f_H=3`$ to hold exactly one must reduce $`\omega _{H_2O}^0=0.355`$ to $`\omega _{H_2O}^0=0.323`$. The concentration of uranium nuclei can be kept in the calculations at its former value: $`N_U=0.720510^2`$ U/cm$``$b. The absorption cross section of $`_3^6`$Li nuclei per uranium nucleus, $`\beta =N_{Li}\sigma _{a,Li}/N_U=2`$ b at $`T=300`$K and normal pressure results in the following concentration of lithium nuclei: $`N_{Li}=2.08810^4`$ Li/cm$``$b.
The ratio of the capture cross section of thermal neutrons to the scattering cross section of epithermal neutrons is too large in the LT paper. At $`T=300`$K, LT use for the capture cross section the value $`_k\sigma _{a,k}N_k/N_U=31.1`$ b/U. The scattering cross section of a free hydrogen nucleus is 20.5 b 45 (for bound hydrogen it is greater). At $`f_H=3`$ the value of $`\mathrm{\Delta }`$ is $`\mathrm{\Delta }=231.1/320.5=1.01`$ and not 2. Even though, in repeating the calculation of LT we have used the value $`\mathrm{\Delta }=2`$. The curve $`\widehat{\sigma }_{\gamma ,Sm}(T_C,\mathrm{\Delta }E_r)`$ is shown in Fig. 13 (curve 1). The value of $`\widehat{\sigma }_{\gamma ,Sm}^{\mathrm{Exp}}(T)=(90.7\pm 8.2)`$ kb are taken from Table 11. Recall that Table 2 contains only four experimental points instead of five. This results in a change of $`\widehat{\sigma }_{\gamma ,Sm}^{\mathrm{Exp}(\mathrm{T})}`$ and its error as compared to LT . Curve 1 was obtained by interpolation of curves 2 for $`\mathrm{\Delta }=1`$ and for $`\mathrm{\Delta }=2`$ to the value $`\mathrm{\Delta }=2\sqrt{300\mathrm{K}/600\mathrm{K}}=\sqrt{2}`$. We have found a negative value for the energy $`\mathrm{\Delta }E_{r2}=24`$ meV, closer to zero than the value $`\mathrm{\Delta }E_{r2}=38`$ meV in LT . From Fig.13 one can clearly see a specific feature of the result of LT . It is enough to take $`\mathrm{\Delta }>1.41`$ for the curve not to intersect the error corridor, and for $`\mathrm{\Delta }<1.41`$ the shift $`\mathrm{\Delta }E_{r2}`$ is strongly reduced.
At a stiffness parameter $`\mathrm{\Delta }1`$ the spectrum is distorted on account of large absorption by the strong absorbers at small energies. Under these conditions the spectrum cannot be considered to be Maxwellian. It can be found only by direct Monte Carlo calculation. We have carried out the calculation of the reactor spectrum and of the distribution for a composition of the core as described in items 3โ5, using the code MCU REA. In the LT paper no absolute value of $`N_U(0)`$ was given, and we choose $`N_U=0.720510^2`$ U/cm$``$b. The calculations were done both without taking account of the power effect (PE) (curve 4) and with taking account of the PE (curve 5). Both calculations cross the error corridor at $`\mathrm{\Delta }E_{r2}^R>0`$ ($`\mathrm{\Delta }E_{r2}^R=4`$ meV). This means that in this case $`\dot{\overline{\delta \alpha }}/\alpha =0`$ as well.
In the LT model the neutron balance is maintained on account of a compensation of the fuel burn-up by the burn-up of strong absorbers. This is far from reality. Since strong absorbers burn up faster than $`_{\mathrm{\hspace{0.25em}\hspace{0.25em}92}}^{235}\text{U}`$, such a balance exists only at the beginning of the cycle. At the end of the cycle no strong absorber is left. The fast burning up strong absorber $`^{149}`$Sm that is still present to this day was formed only at the end of the cycle when the LT model does not work any more.
The results on a possible change of $`\alpha `$ based on the analysis of the cross section of $`_{\mathrm{\hspace{0.25em}\hspace{0.25em}62}}^{149}\text{Sm}`$ in the Oklo reactor are summarized in Table 12. For comparison we have included the cosmological results (Fig.14) 54a ; 54b and the results of laboratory measurements 55 . A review of a possible change of the fundamental constants (experiment and theoretical interpretation) was recently published by Uzan 56a ; 56b . All results show that there are no grounds for an assertion that the e.m. constant has changed in the distant past. However there is a possibility that this conclusion will be revised when the fuel burn-up is taken into account.
## IV Conclusions
We have built a complete computer model of the Oklo reactor core RZ2 . With the aid of present-day computational codes we have calculated in all detail the core parameters. The simulations were done for three fresh cores of different contents of uranium and water. We have also calculated the neutron flux and its spatial and energy distributions. For the three cores we have estimated the temperature and void effects in the reactor. As expected, the neutron reactor spectrum is significantly different from the ideal Maxwell distribution that had been used by other authors to determine the cross section of $`_{\mathrm{\hspace{0.25em}\hspace{0.25em}62}}^{149}\text{Sm}`$. The reactor cross section and the curves of its dependence on the shift of the resonance position $`\mathrm{\Delta }E_r`$ (as a result of a possible change of fundamental constants) differ appreciably from earlier results. The effect of an influence of the reactor spectrum on the cross section of $`_{\mathrm{\hspace{0.25em}\hspace{0.25em}62}}^{149}\text{Sm}`$ can be considered to be firmly established. We have studied the limits of the variation of this effect depending on the initial composition and the size of the core. The fresh bare core RZ2 is critical for $`T_C=(725\pm 55)`$K. At these temperatures the curves of $`\widehat{\sigma }_{\gamma ,Sm}(T_C,\mathrm{\Delta }E_r)`$ lie appreciably lower than for a Maxwell distribution. Possible values of $`\mathrm{\Delta }E_r`$ lie in the range of $`73`$ meV $`\mathrm{\Delta }E_r62`$ meV. These limits are 1.5 times more accurate than those of Dyson and Damour. For the rate of change of the e.m. constant we find $`3.710^{17}`$ year$`{}_{}{}^{1}\dot{\overline{\delta \alpha }}/\alpha +3.110^{17}`$ year<sup>-1</sup>. Within these limits we have $`\dot{\overline{\delta \alpha }}/\alpha =0`$. The analysis of all previous studies shows that none of them has reliably shown up a deviation from zero of the rate of change of the e.m. constant $`\alpha `$. Because of difficulties with the detailed calculation of the burning up in large reactors, which require accumulation of huge statistics, we have not determined the effect of the burn-up on the neutron spectrum and on the $`_{\mathrm{\hspace{0.25em}\hspace{0.25em}62}}^{149}\text{Sm}`$ cross section. Calculations of the influence of burn-up on the temperature of the active core and on the neutron spectrum are in progress.
###### Acknowledgements.
The authors express their thanks to V.A. Varshalovich and B.L. Ioffe for discussions, also their appreciation to N.N. Ponomarev-Stepnoi, E.A. Gomin, M.I. Gurevich, A.S. Kalugin and M.S. Yudkevich for making available codes MCU-REA and BURNUP. The authors consider it their pleasant duty to thank V.V. Kuzminov for making available the nuclear constants for code MCNP4C, J. Vallenius for consultations and W.B. von Schlippe for the translation. This work was done with the partial financial support of grant RFFI 02-02-16546-a.
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# The Microscopic Quantum Theory of Low Temperature Amorphous Solids
## I Introduction
During the past several decades, it has been gradually recognized in the condensed matter and materials science community that amorphous materials, while sharing many characteristics with the more common crystalline solids, represent a distinct solid state of matter. On the one hand, glasses exhibit rigidity and elastic response on humanly relevant time scales, thus qualifying them as solids for many practical purposes. In fact, until the relatively recent advent of systematic studies of the materialsโ response to mechanical and electromagnetic perturbation, as well as of their detailed microscopic structure, the only commonly known distinct attributes of amorphous substances had been their optical properties and the low magnitude and isotropic character of their thermal expansion. Those properties still undergird the main technological importance of amorphous materials. On the other hand, there are many ways in which glasses are fundamentally different from crystals. This is most noticeable in their properties at cryogenic temperatures.
We presently know very well that an amorphous solid is in reality a liquid caught locally in a small set of metastable free energy minima Xia and Wolynes (2000), each of which are separated from the much lower free energy crystalline arrangement by high barriers. Therefore, the glass transition, as manifest in the laboratory, is not strictly speaking a phase transition in the regular thermodynamic sense and is not accompanied by a symmetry change or appearance of a free energy singularity. In contrast, a liquid that was cooled below its melting point fast enough so as to avoid crystallization - i.e. has become supercooled - experiences a crossover to (highly viscous) activated transport. As the temperature is lowered further, the relaxation barriers grow in a very dramatic fashion thus confining the molecules in their metastable arrangements long enough to give the appearance of shear elasticity in the sample on the technologically relevant frequency scales. A quantitative understanding of the physics behind the glass transition has recently been achieved with the random first order transition (RFOT) theory of glasses Kirkpatrick *et al.* (1989); Xia and Wolynes (2000). This theory has provided a microscopic picture of molecular motions in supercooled liquids, such as first principle predictions of the length scales of these motions and the cooperativity lengths and the barrier heights of the activated transport. At any given time, a supercooled liquid is a mosaic of cooperatively rearranging regions, whose size becomes larger as the temperature is lowered. This article describes how the RFOT theory also provides the necessary microscopic input to understand the cryogenic anomalies observed in glasses.
In spite of the absence of periodicity, glasses exhibit, among other things, a specific volume, interatomic distances, coordination number and local elastic modulus comparable to those of crystals. Therefore it has been considered natural to consider amorphous lattices as nearly periodic with the disorder treated as a perturbation, often-times in form of defects, so such a study is not futile. This is indeed a sensible approach, as even the crystals themselves are rarely perfect, and many of their useful mechanical and other properties are determined by the existence and mobility of some sort of defects as well by interaction between those defects. Nevertheless, a number of low temperature phenomena in glasses have persistently evaded a microscopic model-free description along those lines. A more radical revision of the concept of an elementary excitation on top of a unique ground state is necessary Lubchenko and Wolynes (2001, 2003a); Lubchenko (2002).
Let us give a brief historical overview of some of the most outstanding issues in low temperature amorphous state physics. It was already noted in the 1960s that the thermal conductivity of amorphous solids is significantly lower than that of crystals. A low-temperature experimentalist using epoxy in his apparatus knew that its thermal conductivity at liquid helium temperatures went roughly as `constant`$`\times T^2`$, where the `constant` was practically the same for other amorphous substances as well Anderson (1999). Surprisingly, this had not particularly alarmed anyone, even though one would not ร priori expect low temperature properties of disordered solids to be different from crystals, as the appropriate thermal phonon length is much larger than the molecular scale which was presumed to characterize the relevant heterogeneity scale. It was not until Zeller and Pohl published their classic paper Zeller and Pohl (1971) that it became generally known that both the heat capacity and thermal conductivity of glasses were significantly different from those of crystals, and that these anomalies were correlated. The heat capacity turned out to be approximately linear in temperature and larger than the $`T^3`$ phononic contribution up to temperatures $`10`$ K. The challenge to the theorists was soon met by the so called Standard Tunneling Model (STM) Anderson *et al.* (1972); Phillips (1972), in which one assumes that due to a disordered pattern of molecular bonds in glasses, there are a number of defects in the lattice (something like โlooseโ atoms or โdangling bondsโ), which have two alternative positions in space separated by a sufficiently low tunneling barrier. At low temperatures, the dynamics of such a system is described well by a two-level system (TLS) hamiltonian. If one assumes that the spectral density of these TLSโs is flat, one recovers the linear heat capacity. One also finds that the inverse mean free path of a thermal phonon due to resonant scattering off the TLSโs is equal to $`l_{\text{mfp}}^1T`$, which implies thermal conductivity $`\kappa \frac{1}{3}_\omega C\text{ph}(\omega )l_{\text{mfp}}(\omega )c_sT^2`$. Here, $`C\text{ph}(\omega )`$ is the heat capacity of a phonon mode of frequency $`\omega `$ and $`c_s`$ is the speed of sound (one assumes here that heat is carried primarily by phonons, which was experimentally demonstrated explicitly four years later by Zaitlin and Anderson Zaitlin and Anderson (1975)). Note that the resonant character of phonon scattering implies that the scattering cross-section of low-frequency phonons would be independent of the scatterer size, but would scale with the phonon wavelength (squared) itself. Therefore no knowledge of scattererโs microscopic details are needed. Rather, only a single coupling parameter is needed to estimate the magnitude of scattering at low temperatures. The STM did prove to be very successful Phillips (1981), as it predicted, among other things, nonlinear sound absorption due to the saturation of the resonant absorption and the phonon echo, both of which were later observed Hunklinger *et al.* (1976); Golding and Graebner (1976). In spite of these successes, the microscopic nature of these defects had remained unknown, although there later appeared several indications in the literature that the tunneling centers are not single atom entities but rather involve motions within larger groups of atoms Mon and Ashcroft (1978); Guttman and Rahman (1986). On the experimental front, there had been a growing amount of evidence that the number of these additional excitations and their coupling to the phonons are correlated and also depend on $`T_g`$ Reynolds, Jr. (1979, 1980); Raychaudhuri and Pohl (1981), which culminated in the observation made by Freeman and Anderson Freeman and Anderson (1986), that the heat conductivities of all studied insulating glasses, if scaled by elastic constants, fall onto the same line in two regions, connected by a non-universal flat piece corresponding to the so called โplateauโ. In Fig.1 we show a facsimile of Fig.2 from Freeman and Anderson (1986) that demonstrates this heat conductivity universality.
The lower temperature straight line corresponds to the value $`150`$ of the ratio of the thermal phonon mean free path $`l_{\text{mfp}}`$ to the thermal Debye wave-length $`\lambda \mathrm{}c_s/k_BT`$. This region spans roughly 1.5 decades in temperature between several mK (lowest $`T`$ accessed so far for the heat conductivity measurements) and to 1-10 K, depending on the substance. The short linear region at higher temperatures (20-60 K) corresponds to $`l_{\text{mfp}}/\lambda 1`$, which actually implies complete phonon localization Graebner *et al.* (1986) according to the heuristic Ioffe-Riegel criterion. This implies, among other things, that one can no longer use kinetic theory expressions for heat transfer at these temperatures, as a diffusive mechanism must prevail <sup>1</sup><sup>1</sup>1This idea that the heat was transfered by a random walk was used early on by Einstein Einstein (1911) to calculate the thermal conductance of crystals but, of course, he obtained numbers much lower than those measured in the experiment. As we now know, crystals at low enough $`T`$ support well defined quasiparticles - the phonons - which happen to carry heat at these temperatures. Ironically, Einstein never tried his model on the amorphous solids, where it would be applicable in the $`l_{\text{mfp}}/\lambda 1`$ regime.. The intermediate region (โplateauโ) is usually observed between 1 and 30 K, and does not scale with the Debye temperature and speed of sound. The standard tunneling model of non-interacting two-level systems mentioned above is normally applied to the region where $`l_{\text{mfp}}/\lambda 150`$, that is generically below 1 K. The universality of $`l_{\text{mfp}}/\lambda `$ can be boiled down Phillips (1981) to the universality of the following combination of parameters: $`\overline{P}\frac{g^2}{\rho c_s^2}`$, where $`\overline{P}`$ is the spectral and spatial density of the TLSโs (empirically $`10^{45\pm 1}J^1m^3`$), $`g`$ is coupling to the elastic strain on the order of eV, $`\rho `$ is the mass density (for reference, $`\frac{g^2}{\rho c_s^2r^3}`$ would be the interaction strength between such TLSโs at distance $`r`$ from each other). Now, if the defects involved the motion of only a single atom, one would reasonably assume that the value of their spectral density and coupling to the lattice or their combination would be very strongly material dependent. Even though $`\overline{P}`$ and $`g^2`$ vary within almost two orders of magnitude (still surprisingly little), the combination $`\overline{P}\frac{g^2}{\rho c_s^2}`$ is constant within 50% for different materials ($`\rho `$ and $`c_s^2`$ vary considerably as well). It certainly takes a stretch of imagination to think that this is merely a coincidence, as pointed out in Leggett (1991). In 1988, Yu and Leggett proposed Yu and Leggett (1988) that the density of states of the TLS might itself be a result of dipole-dipole interactions between some original non-renormalized excitations. In short, this idea is motivated by the observation that for TLS coupled to the phonons with strength $`g`$, the coefficient at the dipole-dipole interaction term $`g^2/\rho c_s^2`$ has dimensions energy times volume. Therefore the interaction induced renormalized density of states $`\overline{P}`$ has to be the inverse of $`g^2/\rho c_s^2`$ with a coefficient, hence the universality of $`\overline{P}g^2/\rho c_s^2`$ for different materials. However, it so far has not proved possible to use their approach to justify the value of that coefficient to yield the experimental $`l_{\text{mfp}}/\lambda 150`$. This is surprising, since one expects such a simple dimensional argument to be very robust. (Several other studies of the universality Meissner and Spitzmann (1981); Coppersmith (1991) were undertaken at the time, that followed the paper by Freeman and Anderson Freeman and Anderson (1986).) There has been subsequent work applying a renormalization group style calculation to a system of interacting TLS Burin and Kagan (1996), but it seems from the results that renormalizations are relevant only at ultra-low temperatures ($`\mu `$Kโs and below) Neu *et al.* (1997). In spite of the difficulties in justifying the strong interaction scenario Caruzzo (1994); Lubchenko and Wolynes (2000), the works Yu and Leggett (1988); Leggett (1991) that first challenged the standard TLS paradigm remain the main conceptual motivation behind the present paper. One stresses however that the idea that the observed coupling constant, which is quite small, could be a result of some original โbareโ strong interaction, is consistent with the microscopic theory if we argue it is the molecular interactions behind the glass transition itself which become โrenormalizedโ in a somewhat unexpected fashion. This microscopic theory suggests Lubchenko and Wolynes (2001) that the phenomenological two-level systems are discrete energy levels representing resonantly accessible local degrees of freedom that exist in glasses due to the possibility of collective transitions between alternative structural configurations of compact regions encompassing roughly 200 molecular units. The theory of glassy ergodicity breaking shows the spectrum of these excitations is nearly flat and the density of states scales with the inverse glass transition temperature $`T_g`$, echoing the excitation spectrum of a random energy model (REM) with that glass transition temperature. Furthermore, the transitions are an alternative mode of motion that must be in equilibrium with phononic excitations at $`T_g`$. This equilibrium requirement makes one realize that TLS-phonon coupling $`g`$, $`T_g`$ and the materialโs elastic constants are intrinsically related. The universality of the $`l_{\text{mfp}}/\lambda `$ ratio is a consequence of this relationship reflecting the non-equilibrium character of the glassy state. The structural transitions, that become tunneling two-level systems at cryogenic temperatures, exist because a glassy sample, when it falls out of equilibrium, resides in a metastable configuration chosen from a very high density of states. The sample is broken up into a mosaic of dynamically cooperative regions. Alternatively speaking, the energy landscape is local in nature; that is rearrangements of compact regions will not change the structural state of the rest of the sample, but only deform the surrounding regions weakly and purely elastically. A (small) fraction of these rearrangements requires overcoming only a very low barrier and can therefore occur even down to sub-Kelvin temperatures. The tunneling occurs by consecutive molecular displacement within the cooperativity length established at $`T_g`$. The consecutive motion of atoms is conveniently visualized as a domain wall separating the two alternative local structural states, moving through the local region.
The thermal conductivity plateau has traditionally been considered by most workers a separate issue from the TLS. In addition to the rapidly growing magnitude of phonon scattering at the plateau, an excess of density of states is observed in the form of the so called โbumpโ in the heat capacity temperature dependence divided by $`T^3`$. The plateau is interesting from several perspectives. For one thing, it is non-universal if scaled by the elastic constants (say $`\omega _D`$ and $`c_s`$). It is, however, located between two universal regions and it is important to understand which other scales in the problem determine its location and shape. The excitations that give rise to the dramatically increased phonon absorption at the corresponding frequencies have been circumstantially associated with the excitations observed as the so called Boson Peak (BP), directly seen in the inelastic X-ray and neutron scattering experiments, also observed in the optical Brillouin and Raman scattering measurements. These experimental developments date well into 90-s and became possible, in the neutron spectroscopy case, due to the improved resolution in the neutronsโ velocity detection, combined with the ability to generate higher energy incident beams Foret *et al.* (1996). Similarly, meV resolution was needed to utilize the X-ray scattering technics to discern the small inelastic wings on the sides of the strong elastic peak Benassi *et al.* (1996). The term โBoson Peakโ comes from the fact that its intensity scales roughly according to the Bose-Einstein statistics. The extraction of the density of states from the spectra is unfortunately model dependent, and those models can be roughly divided Pilla *et al.* (2000) into the ones where the Boson peak signifies the energy scale on the edge of phonon localization, as promoted in Foret *et al.* (1996), and those following the other school of thought which asserts that these modes are propagating even well above the frequency of the BP, as supported by the interpretation in Pilla *et al.* (2000). As far as theoretical interpretation is concerned, it is our impression that most of theories of the Boson Peak, existing until recently, have postulated a sort of spatial heterogeneity in an otherwise perfectly elastic medium (see a partial list of references in Grigera *et al.* (2001)), with the notable exception of the soft-potential model (SPM) Karpov *et al.* (1983); Buchenau *et al.* (1992). It is, of course, always possible to recover the observed magnitude of the heat capacity excess at the BP temperatures by a particular choice of parameters. While a contribution of the lattice disorder to the density of states undoubtedly exists and can be very significant (see, for example, simulations of silicaโs heat capacity by Horbach at el. Horbach *et al.* (1999)), we must note that if amorphous lattices were purely harmonic, the phonon absorption at the BP frequencies would be of the Rayleigh type and should be significantly lower than observed in the experiment Anderson (1981); Joshi (1979). There must be internal resonances present in the bulk, that scatter phonons inelastically. Though phenomenologically introduced, this feature is present, for example, in the soft-potential model. An analysis of the higher temperature behavior of the tunneling transitions that give rise to the TLS at subKelvin energies was provided in the RFOT approach in Lubchenko and Wolynes (2003a). When these transitions occur at high enough temperature, the domain wall separating the two alternative states can have its surface vibrations thermally excited. The large degeneracy of these vibrational states, characteristic of a two dimensional membrane, that accompany the underlying structural transition, is sufficient to account for the enhancement of phonon scattering at the plateau, as compared to the TLS regime. Finally, the superposition of the domain wall vibrations on the underlying tunneling transition leads to an excess of density of states that reproduces well the bump in the heat capacity (these compound excitations we call โripplonsโ). We therefore arrive at a unified physical picture that allows a unified quantitative explanation of previously seemingly unrelated mysteries in the TLS regime and at the higher, plateau energies.
The paper is organized as follows: the first section outlines the basics of the RFOT theory and then proceeds in applying that theory to understanding the origin of the tunneling centers in amorphous solids. The spectrum of the two-level systems, their coupling to the phonons and the origin of the universality of phonon scattering are then discussed. Additionally, we show how details of the derived TLSโ tunneling amplitude distribution lead to a deviation of $`T`$ dependence of the heat capacity from a strict linear form. The second section explains how the high energy vibrational excitations (ripplons) of the tunneling interfaces gives rise to an excess of states which exhibits itself as the heat capacity bump and yields the rapidly rising phonon scattering at these higher energies. A short discussion of the relaxational absorption from these excitations is given and its frequency dependent part is derived. The contents of these first two sections are, for the most part, a detailed account of the calculations underlying two earlier brief letters Lubchenko and Wolynes (2001, 2003a) that have reported our explanation of the low temperature anomalies in glasses within a semiclassical approach. The third and fourth chapters are comprised of new results. There, we establish that, while not altering the main conclusions of the semiclassical picture, a purely quantum phenomenon of level mixing and repulsion has an observable effect on the density of states of the tunneling centers at low $`T`$. Finally, the interaction between tunneling centers, mediated by phonons, is estimated and this is argued to make a significant contribution to the negative thermal expansivity (and thus a negative Grรผneisen parameter) observed in many amorphous materials.
## II Overview of the Classical Theory of the Structural Glass Transition
From a physicistโs perspective, a theory of the glass transition describes what happens to a liquid when it is cooled down sufficiently but is not observed to crystallize. To a mathematician, this is a generalized problem of packing compact interacting objects of comparable size given a specific constraint on the density distribution (it is not periodic) and total energy of the system. A nearly complete conceptual, microscopic picture of the amorphous state has emerged in the course of the two last decades Stoessel and Wolynes (1984) Singh *et al.* (1985) Kirkpatrick and Wolynes (1987a, b) Kirkpatrick and Thirumalai (1987a, b); Kirkpatrick *et al.* (1989) Xia and Wolynes (2000, 2001a, 2001b) Lubchenko and Wolynes (2001, 2003a, 2003b, 2004). This framework has lead to a unified, quantitative understanding of many seemingly unrelated phenomena in supercooled liquids above and below the glass transition. The glasses we consider form at temperatures where quantum effects are small so classical statistical mechanics is used. We review such a classical glass transition in what follows.
First, we make several comments on the phenomenology of supercooled liquids. Strictly speaking, these are nonequilibrium systems: When cooled sufficiently slowly, most simple liquids will crystallize at a temperature just below the melting temperature $`T_m`$. Randomly atactic polymers become glassy but presumably never crystallize. The melting point is defined as the temperature at which the liquid and crystal free energies are equal. Cooling the liquid at least a bit below $`T_m`$ is necessary to create a free energy driving force so as to make the nucleation barrier finite and to allow the system to equilibrate. The crystal, once formed is different from the liquid in several ways, e.g. it scatters X-rays at precise angles and it is anisotropic. Crucially for us, a crystal supports transverse sound waves, at all frequencies (including $`\omega =0`$, thence the crystal retains its shape). In contrast, the supercooled liquid is a finite lifetime state since crystallization will eventually occur by nucleation. However, the growth of crystalline nuclei, inside the liquid, is subject to the slowing of all motions in liquids. Owing to this dramatic slowing of liquid motions upon lowering the temperature, one can supercool the liquid substantially below its melting point, which is the key to forming glasses. The extra nucleation barrier ensures there is adequate time to study the properties of the supecooled noncrystalline state. Local structures in supecooled liquids persist for some time, call it $`1/\omega _c`$. This time is longer than the time it takes to establish a Maxwell distribution of velocity, which is at most a few vibrational periods. Such an amorphous system will support transverse waves at frequencies $`\omega >\omega _c`$, just as a crystal would, but will in contrast exhibit a liquid like, equilibrium response to time dependent perturbations at frequencies $`\omega <\omega _c`$. As we have said, $`\omega _c`$ drops rapidly upon cooling. If one is intent on observing equilibrium response at some frequency range, one must prepare the sample by cooling it more slowly than $`\omega _c`$. Conversely, for any given cooling rate, no matter how slow, the liquid will fall out of equilibrium on all time scales and the sample will appear to be mechanically solid. We say the liquid has undergone the glass transition. (The corresponding $`\omega _c`$ usually ranges between $`10^2`$ and $`10^5`$ sec, depending on the experimenterโs patience.) The liquid just below the glass transition temperature $`T_g`$ is only subtly different from the liquid just above $`T_g`$. Structurally, first of all, the two are nearly identical. Even dynamically, both can flow, although the $`T`$-dependences of the corresponding transport coefficients are distinct in the two forms of the โequilibriumโ supercooled liquid and the nonequilibrium glassy state Lubchenko and Wolynes (2004). The residual dynamics below $`T_g`$ is referred to as โagingโ. Aging is at least as slow as the motions just above $`T_g`$, but can be much slower when the sample is studied well below $`T_g`$. This requires a greater amount of the experimenterโs patience in studying system properties than even needed for sample preparation. Finally, when the sample falls out of equilibrium at $`T_g`$, a jump in the heat capacity is measured by differential calorimetry, thus resembling, crudely, a phase transition.
The dramatic slowing down of molecular motions is explicitly seen in a vast area of different probes of liquid local structures. Slow motion is evident in viscosity, dielectric relaxation, frequency dependent ionic conductance, as well as in the speed of crystallization itself. In all cases, the temperature dependence of the generic relaxation time obeys to a reasonable, but not perfect approximation the empirical Vogel-Fulcher law:
$$\tau _{\text{rlxn}}e^{DT_0/(TT_0)}$$
(1)
For a review, see Angell *et al.* (2000); Bรถhmer *et al.* (1993). A specific example of a $`\tau (T)`$ dependence is shown in the l.h.s. panel of Fig.7. In the expression above, $`T_0`$ is a material dependent temperature at which the relaxation times would presumably diverge, if the experimenter had the patience to equilibrate the liquid at the corresponding temperatures. Needless to say, measurements of equilibrium dynamics near $`T_0`$ are essentially nonexistent. The coefficient $`D`$ is often called โfragilityโ, with larger values of $`D`$ corresponding to โstrongerโ substances, while smaller values are associated with โfragileโ liquids. This terminology apparently refers to the degree of covalent networking in the material Angell (1985), a qualitative trend later rationalized by a density-functional study of Hall and Wolynes (2003). Fragility appears to correlate with the Poisson ratio, at least for non-polymeric glasses Novikov and Sokolov (2004). At any rate, the value of coefficient $`D`$ is directly related to what glassblowers refer to as โshort glassesโ and โlong glassesโ, Pfaender (1996): (molten) glass can be worked or shaped in the range of viscosities $`10^410^9`$ Poise. If the corresponding temperature range is short, the glass is called โshortโ, and vice versa for the โlongโglass. The former and the latter obviously correspond to a small and large value of the parameter $`D`$ respectively.
The non-equilibrium character of a supercooled liquid is exhibited in the entropy of the liquid which is considerably larger at $`T_g`$ than that of the corresponding crystall at this temperature. This additional entropy corresponds to all the molecular translations, that would have otherwise frozen out at crystallization. In crystallization, this would appear as the latent heat of the liquid-to-crystal transition. In a supercooled liquid, the molecular structure is dense enough to define a lattice locally. Vibrations around lattice sites are small. The excess entropy associated with the locations of these lattice sites has traditionally been designated as the โconfigurationalโ entropy. This excess entropy, $`s_c`$, is temperature dependent. It refers to all possible liquid configurations that could be surveyed by the liquid if we wait long enough for molecular translations to occur. Experimentally, we determine the configurational entropy by relying on the third law of thermodynamics. Using the third law, we know the total entropy of the liquid at $`T_m`$ by integrating the crystalโs heat capacity (over $`T`$) and adding the entropy of melting. Now for the supercooled liquid, we integrate the heat capacity difference between the liquid and the crystal. To do this we, of course, assume the vibrational entropies of the ordered and aperiodic lattices are close. The heat capacity measured by differential calorimetry above the glass transition depends on the rate of the configurational and vibrational entropy decrease with temperature right above $`T_g`$. Below $`T_g`$ the structure of the liquid remains the same as of the moment of vitrification, apart from some (normally insignificant) aging. The vibrational entropy decreases as it did above $`T_g`$, but there is no component from configurational change. Thus one observes a non-zero heat capacity jump at $`T_g`$. Above $`T_g`$, the $`s_c`$ decreases and the density increases with lowering the temperature. This is expected because there are fewer ways to mutually arrange the molecules at higher densities. When extrapolated past $`T_g`$, as was done by Simon Simon (1937) and notably by Kauzmann in his review Kauzmann (1948), the configurational entropy vanishes at a temperature $`T_K`$, which is securely above the absolute zero. This suggests that only a non-extensive number of low energy aperiodic, liquid arrangements could be found at $`T_K`$ and the entropy of the liquid is thus equal to the corresponding crystal (correcting for differences in their vibrational spectrum). This phenomenon is sometimes referred to as the โentropy crisisโ, which, again, would presumably occur only under completely equilibrium cooling. Such an entropy crisis strictly occurs in several mean-field spin glass models with infinite interactions Gross and Mรฉzard (1984); Gross *et al.* (1985); Kirkpatrick and Wolynes (1987b). There are many sound arguments suggesting a strict singular vanishing of configurational entropy at $`T_K`$ is unlikely for real liquids Eastwood and Wolynes (2002); Stillinger (1988). Nevertheless, $`T_K`$ is a useful fiducial point for the analysis. None of the results of the present theory in the experimentally accessible regime depend on the configurational entropy truly vanishing at any point. As we shall see, the configurational entropy is macroscopic but decreases with temperature. $`s_c`$ is typically $`.8k_B`$ per movable unit at the conventional glass transition temperature corresponding to cooling rate of inverse hour and decreases at a rate proportional to $`\mathrm{\Delta }c_p/T_g`$. For simplicity, we will assume $`s_c`$ extrapolates so as to scale linearly with the proximity to the entropy crisis (see Richert and Angell (1998)): $`s_c=\mathrm{\Delta }c_p(TT_K)/T_K`$.
Before our formal discussion, let us make several qualitative statements about molecular transport above $`T_g`$. The motions of a supercooled liquid are much slower than the local vibrations. The potential felt by an individual molecule comforms to a local โcageโ. This local โcageโ is formed by the neighboring molecules, of course. In order to translate irreversibly a given molecule, as opposed to vibrating about the current position, the cage must be destroyed. In other words, a number of surrounding molecues must be translated as well. Upon lowering the temperature, the density increases and $`s_c`$ decreases, therefore fewer alternative states are available to any given group of molecules. Thus it is clear that conforming the liquid to an arbitrary translation of a given molecular unit will require readjusting the positions of more and more surrounding molecules at the same time. This leads to a larger cooperative region size, leading in turn to higher barriers for relaxation processes and higher viscosity. At a crude level, this picture underlies the arguments from Adam and Gibbs (1965), but those arguments fail to relate the size of the moving regions to the energy landscape itself. In contrast, the Random First Order Transition (RFOT) theory Kirkpatrick *et al.* (1989); Xia and Wolynes (2000) explicitly shows how these reconfigurational motions occur and thus establishes intrinsic connection between the kinetic properties and the thermodynamics of supercooled liquids. Our account is based on Lubchenko and Wolynes (2004) which also discusses the intrinsic connection between cooperative, activated motions in the supercooled liquid both above and the classical aging dynamics below the glass transition. These arguments also pave the way for understanding the quantum dynamics at cryogenic temperatures.
The main prerequisite of the RFOT theory is the existence of time scale separation between vibrational thermalization and equilibrating structural degrees of freedom that involve crossing saddle points on the free energy surface. This only occurs below a crossover temperature $`T_A`$ which is predicted by the theory itself. The existence of local trapping in cages is well established by experiment: there is a long plateau in the time dependent structure factor as measured by the inelastic neutron scattering Mezei (1991). In RFOT, such trapping was first established theoretically using a density functional theory (DFT) in Singh *et al.* (1985): This paper shows there are aperiodic free energy minima by computing the free energy of an aperiodic variational density distribution function: $`\rho (๐ซ)\rho (๐ซ,\{๐ซ_i\})=_i\left(\frac{\alpha }{\pi }\right)^{3/2}e^{\alpha (๐ซ๐ซ_i)^2}`$. The set of coordinates $`\{๐ซ_i\}`$ denotes a particular aperiodic lattice. The typical lattice spacing is $`a`$. A zero value for the parameter $`\alpha `$ would correspond to a completely delocalized, uniform liquid state, such as just below the liquid-vapor transition. $`\alpha \mathrm{}`$ would imply freezing into an infinitely rigid lattice. $`\alpha `$ can also be interpreted as the spring constant of an equivalent Einstein oscillator forcing each molecule to remain near its proper location in the aperiodic lattice. $`F(\alpha )`$ develops a metastable minimum, at non-zero $`\alpha =\alpha _00`$, only below some temperature $`T_A`$. This minimum has higher free energy than than the lowest minimum at $`\alpha =0`$ (see Fig.2).
In the mean field limit, the appearance of such minimum would lead a lattice stiffness and would represent a state with a divergent viscosity. This localization transition and the viscosity catastrophe of mode-coupling theories are essentially identical as was established in Kirkpatrick and Wolynes (1987a). A single such high lying free energy minimum would be thermodynamically irrelevant, but one must recall that this $`F(\alpha )`$ is computed for a single, particular aperiodic lattice, which is actually only one of many possibilities. Taking into acount the thermodynamically large number of alternative aperiodic packings increases the entropy of the (set of) localized, aperiodic state(s) and thus lowers the metastable free energy minima just the right amount to make them competitive with the mean-field uniform, delocalized state. The correspondence between the free energy difference in mean field theory and the configurational entropy was rigorously shown for the Potts Glass by Kirkpatrick and Wolynes (1987b) who argued such systems have similar symmetry properties to structural glasses. For structural glasses this correspondnce may also be shown more formally using a replica formalism Mรฉzard and Parisi (1999). The localization transition at $`T_A`$ is accompanied by a discontinuous change in the order parameter $`\alpha `$. This is why the transition is called โRandom First Orderโ. Although there is a discontinuity in $`\alpha `$, the actual structure in which the system freezes is chosen at random out of a multitude of possibilities (given by the configurational entropy) At the same time, such an ordered phase will persist only for finite times, therefore this is a true transition only for high-frequency motions, comparable at first to the vibrational time scale. This transition at $`T_A`$ only signifies a soft cross-over, as far as the whole dynamical range is concerned. We emphasize, there are many different โphasesโ below $`T_A`$, all of which are random packings. The number of random packings, thermally available to a region of size $`N`$, $`e^{s_cN}`$, decreases gradually with temperature. (This corresponds to gradual freezing out the translational degrees of freedom with lowering the temperature, as signified by the decreasing $`\omega _c`$.) Because the decrease is gradual, the random first order transition does not exhibit a latent heat. In a finite range system, different minima can interconnect by barrier crossing. (Such barriers would be infinite in mean field.) Even though the transition at $`T_A`$ is a crossover, the temperature $`T_A`$ itself is a useful parameter characterizing material properties.
The resulting time scale separation at and below $`T_A`$ has two important consequences. First, one may perform canonical averaging over the vibrations within a particular structural state. This gives a free energy of a particular structural state: $`\mathrm{\Phi }=ETS_{\text{vibr}}`$, where $`S_{\text{vibr}}`$ is the vibrational entropy. Note the vibrations are not necessarily harmonic. To define $`\mathrm{\Phi }`$, all that matters is that the local vibrations equilibrate much faster than the structural degrees of freedom. As a consequence, $`\mathrm{\Phi }`$ can be termed the bulk, microcanonical energy of a given structural state. To any value of this energy one may associate a bulk, microcanonical entropy $`S_c(\mathrm{\Phi })`$ counting states with similar contributions from energy and vibratrional entropy; both $`\mathrm{\Phi }`$ and $`S_c(\mathrm{\Phi })`$ scale linearly with the size of the system. One may thus to work with morphologically distinct, globally defined aperiodic phases without actually specifying their precise molecular constitution, so long as we know their spectrum, i.e. their number as a function of the microcanonical free energy. These statistics are directly measurable by calorimetry just as in our discussion of the Kauzmann paradox.
Having established the transitory existence of a global aperiodic structure, we may next enquire into how molecular motions allow the system to escape such a phase <sup>2</sup><sup>2</sup>2Of course, the issue of producing the aperiodic state in the laboratory would also involve estimating whether corresponding quenching rates can be experimentally achieved.. This occurs by replacing locally one part of the aperiodic packing by a different local packing. This will be an activated event. The RFOT theory allows one to compute the mean activation barrier and its distribution. Also, the theory determines critical region size and the spatial extent of the excitations corresponding to the cooperative rearrangement. The magnitude of an individual molecular displacement during the transition is determined by $`\alpha `$. To estimate the activation free energy, let us make the following construct. Considering a library of possible local aperiodic arrangements at a particular location, as illustrated in Fig.3.
This local library of states can be constructed based on the existence of the global library of states introduced earlier that we described by the energy variable $`\mathrm{\Phi }`$ and the corresponding entropy $`S_c(\mathrm{\Phi })`$ reflecting the spectrum. Clearly, the energy density $`\mathrm{exp}[S_c(\mathrm{\Phi })]`$ is extremely high and grows rapidly with $`\mathrm{\Phi }`$. We might perform a full survey of local states by mentally carving out a small region of size $`N`$, while freezing in place the lattice sites surrounding the region. One can then heat the local region and then allow that region to equilibrate. Unless the new local arrangement is exactly the same as the original one, its energy will likely be significantly higher: A local substitution statistically must cost free energy, stemming from a structural mismatch between two randomly chosen aperiodic packings of a given energy $`\mathrm{\Phi }`$. This mismatch energy corresponds to the usual surface energy, such as that between two different crystal forms or at a liquid-crystal interface. The free energy cost of locally replacing the initial phase (labelled as โinโ) by another phase, call it $`j`$, can therefore be written as
$$\varphi _j^{lib}(๐)\varphi _{\text{in}}^{lib}(๐)=\mathrm{\Phi }_j(N)\mathrm{\Phi }_{\text{in}}(N)+\mathrm{\Gamma }_{j,in},$$
(2)
where $`\mathrm{\Gamma }_{j,in}`$ is the mismatch energy and $`๐`$ is the location of the local region. As before, the capital $`\mathrm{\Phi }`$ denotes the bulk energy, corresponding to a distinct aperiodic packing, with the vibrational entropy already included. To compute the likelihood of such a local rearrangement, substitute for the specific surface energy $`\mathrm{\Gamma }_{j,in}`$ its average value which should scale with size: $`\gamma N^x`$. $`\gamma `$ depends on the material and on temperature. Naively, the usual surface energy scaling is $`N^{(d1)/d}`$, expected in $`d`$ dimensions. One can argue however that $`x`$ will actually turn out to equal 1/2. Such a surface tension renormalization was first conceived by Villain Villain (1985), in the context of the random field Ising model (RFIM). In RFIM, the Ising spins, in addition to their coupling, are subjected to a random static magnetic field obeying certain fluctuation statistics. A flat interface, or domain wall, between spin-up and spin-down domains will distort so as to conform to the local variation of the field. An RG argument incarnating this distortion on a hierarchy of length scales yields a scale dependent renormalization of the surface tension, giving a surface free energy exponent $`x=1/2`$ Villain (1985). The structure-structure interface in a supercooled liquid resembles the RFIM, owing to the fluctuations of local energies of the various aperiodic packings. The statistics of these fluctuating local energies require that $`\delta \mathrm{\Phi }(\delta N)\mathrm{\Phi }_0\sqrt{\delta N}`$, where $`\mathrm{\Phi }_0`$ is $`\delta N`$-independent, echoing the fluctuation statistics of the frozen random field of the RFIM. Thus, as Kirkpatrick *et al.* (1989) suggest, the originally thin flat interface will become diffuse yielding $`x=1/2`$. In the liquid case, a vivid interpretation of the surface energy renormalization is possible: Since the interface is distorted down to the smallest scale (allowed by the materialโs discretness), the region occupied by the now diffuse wall is neither of the two original structures it separates. Instead it may be interpreted as accommodating other structures. These intermediate packings interpolate structurally two randomly chosen, and thus ร priori energetically disagreeable packings. In other words, the original thin interface separating two given packings, is โwettedโ by other packings thus lowering the overall interface energy. As we shall see, real liquids have only modest size regions of rearrangement, so it is hard to argue about the exact value of the exponent. Nevertheless, we note two felicitous observations: With $`x=1/2`$, the usual scaling argument will give precisely a discontinuity in $`\mathrm{\Delta }c_p`$ at any ideal transition, to be seen at $`T_K`$. Also, while the RFIM itself remains the subject of discussion, Villainโs argument does give a length scale exponent agreeing with the majority of experiments and numerical studies Nattermann (1998); Belanger (1998).
The role of the interface mismatch energy in the reconfiguration process can be beneficially understood from a statistical point of view, as illustrated in Fig.3. It costs free energy to reconfigure a small number $`N`$ of molecules because considering a small region severely limits the number of available liquid configurations. The interface energy grows with $`N`$, however the available density of states, too, grows with $`N`$, both in terms of its absolute value and the distributionโs width. At some size $`N^{}`$, that will be computed shortly, all relevant liquid states become available. The rate of escape of a group of $`N`$ molecules to another structural state can be determined by a canonical type sum accounting for the multiplicity of the final states at energy $`\varphi _j`$:
$`k`$ $`=`$ $`\tau _{\text{micro}}^1{\displaystyle (d\varphi _j^{lib}/c_\varphi )e^{S_c(\mathrm{\Phi }_j)/k_B}e^{(\varphi _j^{lib}\varphi _{\text{in}}^{lib})/k_BT}}`$ (3)
$``$ $`\tau _{\text{micro}}^1e^{S_c(\mathrm{\Phi }_{\text{eq}})/k_B}e^{(\varphi _{\text{eq}}\varphi _{\text{in}}^{lib})/k_BT}.`$
In the second step, a steepest descent evaluation is made where $`\varphi _{\text{eq}}`$ maximizes the integrand. $`c_\varphi `$ is some constant of units energy that reflects the local curvatures of the energy landscape. The quantities $`\varphi _j^{lib}`$ and $`\mathrm{\Phi }_j`$ are related through Eq.(2). The time scale $`\tau _{\text{micro}}`$ is the time scale of a molecular scale non-activated process, typically of the order a picosecond. The value $`\varphi _{\text{eq}}`$ that maximizes the integrand above will be the internal (equilibrium) free energy characteristic of the system at the ambient (i.e. vibrational) temperature $`T`$. In other words, the greatest kinetically accessibility of a state, as embodied in the optimization in Eq.(3), implies that the state will be most frequently visited by the system, therefore it must be the equilibrium state. The integration in Eq.(3) is similar to a canonical sum; yet it is different in an important way: The summation in Eq.(3) is far more general than the usual expression for the partition function because when relaxation times are continuously distributed, one must explicitly weigh the contribution of a state (to the canonical sum) by its kinetic accessibility. The latter, in general, will depend on the spatial extent of the excitation corresponding to a transition between two states; in this regard, the integration variable $`\varphi _j`$ is, in a sense, a local microcanonical energy. Consequently, the energy $`\varphi _{\text{eq}}`$ corresponds to a canonical energy. Yet, $`\varphi _j`$ and $`\varphi _{\text{eq}}`$ would strictly become a microcanonical and canonical energy, in their conventional sense, only in the large $`N`$ limit, when the boundary effects are small. In contrast, the very thermodynamic relevance of the glassy state is due to the locality of the landscape and non-smallness of the surface term. Finally, since the bulk entropy $`S_c(\mathrm{\Phi }_{\text{eq}})`$ corresponds to the equilibrium energy, it will be given by the equilibrium configurational entropy $`S_c(T)`$, measured by calorimetry. Thus given $`\varphi _{\text{eq}}`$, one can compute the value of the typical escape rate to a structure where $`N`$ particles have moved. This gives:
$$k(N)=\tau _{\text{micro}}^1\mathrm{exp}\left\{S_c(N,T)\frac{\varphi _{\text{eq}}\varphi _{\text{in}}^{lib}}{k_BT}\right\}.$$
(4)
The number of particles that must be moved for complete equilibration is determined by the minimum of this expression over $`N`$. We thus determine an activation free energy profile
$`F^{}(N)`$ $`=`$ $`\varphi _{\text{eq}}\varphi _{\text{in}}^{lib}TS_c(N,T)`$ (5)
$`=`$ $`\mathrm{\Phi }_{\text{eq}}(N)\mathrm{\Phi }_{\text{in}}(N)+\gamma \sqrt{N}TS_c(N,T),`$
where we used Eq.(2) in the second equality. The maximum of the $`F(N)`$ curve defines the bottleneck location. This equation is suitable for finding the rate of structural rearrangement both in the equilibrated supercooled liquid (before it crystallizes!) and in the nonequilibrium glass, which ages below $`T_g`$.
Let us first consider equilibrium liquid rearrangements. In this case typically $`\mathrm{\Phi }_{\text{eq}}=\mathrm{\Phi }_{\text{in}}`$, apart from fluctuations. Thus one arrives at the following simple expression,
$$F(N)=\gamma \sqrt{N}Ts_cN,$$
(6)
where we have used the thermodynamic scaling of the configurational entropy, $`S_c(N)=s_cN`$. In the supercooled equilibrated liquid, molecular transport is driven by only the multiplicity of mutual molecular arrangements. For this reason, the reconfigurations following the activation profile from Eq.(6) have been called โentropic dropletsโ. The graph of the function in Eq.(6) is shown in Fig.4.
The transition state configuration will satisfy $`F/N=0`$, corresponding to an unstable saddle point of this free energy. This gives for fixed $`\gamma `$ a rearranging region size $`N^{}`$ that grows as $`s_c`$ diminishes: $`N^{}=(\gamma /2s_cT)^2`$. The resulting barrier also scales inversely proportionally to $`s_c`$:
$$F^{}=\frac{\gamma ^2}{4s_cT}.$$
(7)
An inverse scaling of the barrier with the configurational entropy was arrived at by Adam and Gibbs Adam and Gibbs (1965) in a different (and inequivalent) way. Notice if $`\gamma `$, as function of temperature, is smooth around $`T_K`$ and $`s_c`$ is described by the linear law $`s_c(TT_K)`$, the resulting activation barrier is exactly of the Vogel-Fulcher law form (1), which, as we have said, fits data well. Many arguments can lead to increasing relaxation times at low temperatures and with enough adjustable parameters, can fit data. What is different about the RFOT theory is that it establishes an intrinsic link between the rate law and the entropy crisis. In addition, if the entropy of the equilibrated fluid can be estimated, the density functional theory allows the vibrational entropy and thus, by substraction, the configurational entropy to be determined. Therefore $`T_K`$ can be estimated from the microscopic force laws. This has been done for simple soft spheres by Mezard and Parisi Mรฉzard and Parisi (1999), giving reasonable results. Hall and Wolynes Hall and Wolynes (2003) have also calculated $`T_0`$ and $`T_A`$ for a simplified model of a network fluid. Their study is consistent with known chemical trends for $`T_A`$ and $`T_K`$ as the network becomes more thoroughly crosslinked.
The idea of the configurational entropy itself driving liquid rearrangements still appears to generate some confusion. One possible reason for this is that $`s_c`$ is totally unambiguously defined only in the mean field limit. In the latter limit, rearrangements are infinite so dynamics driven by $`s_c`$ do not arise. This is a good place to emphasize that the RFOT theory is not mean-field! Only the local landscape, within an entropic droplet, is actually well described by a mean-field, Random Energy Model like approximation. We took advantage of this in extracting the energy spectrum of low energy structural excitations in a frozen glass Lubchenko and Wolynes (2001), as explained in detail in the following Section. We wish to point the reader to the recent elegant treatment of Bouchaud and Biroli (2004) re-analyzing the RFOT conclusions for rearrangements in an equilibrated fluid from the viewpoint of Derridaโs Random Energy Model (REM) Derrida (1981).
Now, calculations of $`T_A`$ and $`T_K`$ are plagued by the usual difficulties of liquid state structure theory and the accuracy of approximations some of which are hard to control. Still, even in the face of such approximations, such microscopic considerations lead us to expect a universal value of $`\gamma /T_g`$ at $`T_g`$ as we shall discuss below.
The RFOT theory allows the coefficient $`\gamma `$ in the mismatch energy to be estimated from a microscopic argument. It turns out to be proportional to $`T_K`$ and to depend logarithmically on the inverse square of the so called Lindemann ratio. Early in the 20th century, Lindemann argued that the thermal fluctuations of an atomโs position could only be a fraction of the lattice spacing $`a`$ in a solid, if the packing is to be mechanically stable Lindemann (1910). Since the threshold value of the vibrational amplitude of an atom in the lattice is finite, the transition in which the lattice disintegrates must be first order. For crystals, the Lindemann ratio of this threshold displacement $`d_L`$ to the lattice spacing is about $`1/10`$. For amorphous materials, the $`d_L/a`$ ratio can be obtained from the plateau in the self correlation functions measured by neutron scattering experiments Mezei (1991). Again, this ratio turns out to be approximately one-tenth (universally!). This number is reproduced in several microscopic calculations consistent with the RFOT theory, such as the self-consistent phonon theory and density functional theories Singh *et al.* (1985); Stoessel and Wolynes (1984), and dynamical mode coupling theory Bengtzelius *et al.* (1984); Kirkpatrick and Wolynes (1987a, b); Wolynes (1992), with modest quantitative variations. The meaning of $`\alpha 0.1`$ as a mechanical stability criterion has been also corroborated within the replica formalism Mรฉzard and Parisi (1999). In terms of the DFT calculation dicussed earlier, $`\alpha _L`$ corresponds with the metastable minimum that the free energy $`F(\alpha )`$ develops below the dynamical transition temperature $`T_A`$ (see Fig.2). It has a relatively weak temperature dependence. The logarithmic scaling of the surface tension coefficient with the Lindemann length follows from a detailed calculation by Xia and Wolynes (2000), but can be rationalized in a simple way: Below $`T_A`$, motions span only the length $`d_L`$, while in the liquids, they can move a distance $`a`$ before losing their identity with a neighboring molecule. The entropy of the โcagedโ fluid is less and thus the free energy cost of confining a molecule within length $`d_L`$, as opposed to $`a`$, can be assessed by recalling the free energy expression for an ideal monatomic gas: $`f=\frac{3}{2}k_BT\mathrm{ln}\left[\left(\frac{eV}{N}\right)^{2/3}\frac{mT}{2\pi \mathrm{}^2}\right]`$, written deliberately here so as to have a length scale squared in the logarithm.
$`\gamma `$ is proportional to $`T_K`$ and only logarithmically depends on a nearly universal quantity, the Lindemann ratio. If $`T_g`$ is near $`T_K`$, i.e. for slow quenches, $`\gamma /T_g`$ is thus nearly material independent and calculable: $`\gamma =\frac{3}{2}\sqrt{3\pi }k_BT_g\mathrm{ln}(\alpha _La^2/\pi e)`$. Quantifying the mismatch energy this specifically leads to many predictions about the dynamics near $`T_g`$, for a range of substances. First, the coefficient in the Vogel-Fulcher law $`D`$ is predicted to follow from the measured thermodynamics. Using the $`\gamma `$ value above, we find not only the VF dependence of the relaxation times on the temperature, $`e^{DT_0/(TT_0)}`$, but also a remarkably simple formula relating $`D`$ and the heat capacity jump: $`D=32.R/\mathrm{\Delta }c_p`$ Xia and Wolynes (2000). The coefficient $`32.`$ is nearly universal and, as we see, follows numerically from the microscopic theory since the universal value of the Lindemann ratio enters only logarithmically in the localization entropy cost. The numerical relation between $`D`$ and $`\mathrm{\Delta }c_p`$ from this simple explicit calculation is in rather remarkable agreement with experiment. In Fig.5, we plot the so called fragility index $`m`$, as computed from calorimetry and extracted from direct relaxation measurements. $`m`$ is proportional to the slope of the $`\mathrm{log}\tau `$ vs. $`1/T`$ relation at $`T_g`$ and thus is directly related to $`D`$ if the VF law is valid. ($`D`$ values in the literature are obtained from global fits of $`\mathrm{log}\tau `$ vs. $`1/T`$ and depend somewhat on the fitting procedure.)
Two other remarkable universalities emerge from the value of $`\gamma `$. First, at a reference laboratory time scale of 1 hr $`10^{17}\tau _0`$ we have a universal value of $`s_c0.8k_B`$. This implies $`s_c(T_g)/s_c(T_m)0.7`$, where $`s_c(T_m)`$ is, of course, also the fusion entropy. This relation is independent of question of what is the moving subunit. The relation holds very well. A second important universal feature emerges from the universal value of $`\gamma /T_g`$: the cooperative size at $`T_g`$ is nearly universal.
Let us now consider in greater detail the pattern of cooperative structural rearrangements in a supercooled liquid. These turn out to presage the existence of the residual degrees of freedom in a glass below $`T_g`$. Within a period of time shorter than the typical relaxation time $`\tau `$, the molecular motions within regions of size $`\xi ^3`$ will be highly correlated and, at the same time, approximately decoupled from the surrounding. That is, the liquid is broken up in to a (flickering) mosaic pattern of cooperative regions. This mosaic structure is directly manifested in the dynamical heterogeneity recently observed in supercooled liquids using single molecule experiments Russel and Israeloff (2000), nonlinear relaxation experiments Silescu (1999) and non-linear NMR experiments Tracht *et al.* (1998). (These experimental tools became available only a decade after the RFOT theory was first formulated.) The size of a typical mosaic cell is found from the thermodynamic condition $`F(N^{})=0`$. Unlike the regular nucleation of one distinct phase within another (as in crystal growth in the liquid), by crossing the barrier from Eq.(6) the local region arrives at a statistically similar but an alternative solution of the free energy functional, thus that solution still represents a typical liquid state! An informal analogy here is that distinct low energy dense local liquid packings are like the fingerprints of different individuals - different in detail, yet generically equivalent liquid states. Since we have agreed that $`F=0`$ is the liquid equilibrium free energy at this temperature (the crystalline state is assumed to be hidden behind a high enough crystal nucleation barrier), the condition $`F(N^{})=0`$ specifies the size of region to which an arbitrary liquid configuration is available. Therefore, a region of size $`N^{}`$ is able to reconfigure on the experimental time scale characterized by $`F^{}`$. In terms of physical length, $`F(N^{})=0`$ implies $`\xi N_{}^{}{}_{}{}^{1/3}a=a\left[\frac{8}{3\sqrt{3\pi }}\mathrm{ln}\left(\frac{\tau }{\tau _0}\right)/\mathrm{ln}\left(\frac{\alpha _La^2}{\pi e}\right)\right]^{2/3}5.8a`$ ($`N^{}190`$). The critical radius $`r^{}`$ at $`T_g`$ is a multiple of $`\xi `$. Droplets of size $`N>N^{}`$ are thermodynamically unstable and will break up into smaller droplets, in contrast to what prescribed by $`F(N)`$, if used naively beyond size $`N^{}`$. This is because $`N=0`$ and $`N=N^{}`$ represent thermodynamically equivalent states of the liquid in which every packing typical of the temperature $`T`$ is accessible to the liquid on the experimental time scale, as already mentioned. In view of this โsymmetryโ between points $`N=0`$ and $`N^{}`$, it may seem somewhat odd that $`F(N)`$ profile is not symmetric about $`N^{}`$. Droplet size $`N`$, as a one dimensional order parameter, is not a complete description. The profile $`F(N)`$ is a projection onto a single coordinate of a transition that must be described by $`e^{s_cN^{}}`$ order parameters - the effective number of distinct aperiodic packings explored by the liquid. At the point $`N^{}`$, the free energy functional actually has a minimum as a function of the (multi-component) order parameter. A more detailed discussion of this can be found in Ref.Lubchenko and Wolynes (2003b), where we compute the softening of the barrier $`F^{}`$ near $`T_A`$ due to order parameter magnitude fluctuations that are important near $`T_A`$.
We thus see that the length scale of the mosaic and number density of the mosaic domain walls is determined by the competition between the energy cost of a domain wall and the entropic advantage of using the large number of configurations. We emphasize again, the relative domain size $`\xi /a`$ depends only on the logarithms of the relaxation rate and the Lindemann ratio, nearly universal parameters themselves, and is therefore the same for different substances. This high temperature phenomenon of universality at $`T_g`$ has direct consequences for the universality of the ultra-low temperature glassy anomalies.
We have seen that the cooperative region, which represents a nominal dynamical unit of liquid, is of rather modest size, resulting in observable fluctuation effects. Xia and Wolynes computed the relaxation barrier distribution Xia and Wolynes (2001a). The configurational entropy must fluctuate, with the variance given by the usual expression: $`(\delta S_c)^2=C_p1/D`$ Landau and Lifshitz (1980). The barrier height for a particular region is directly related to the local density of states, and hence to the configurational entropy itself by Eq.(6), $`F^{}1/s_c`$. As a result, the barrier distribution width must correlate with the fragility. A gaussian approximation leads to a simple formula $`\delta F^{}/F^{}=1/2\sqrt{D}`$ Xia and Wolynes (2001a). There are also calculable deviations from gaussianity. The barrier distribution gives rise to non-exponentiality of relaxations. These are well fitted by a stretched exponential $`e^{(t/t_0)\beta }`$. The measured $`\beta `$ correlates with the fragility, in good agreement with the theory, see Fig.6.
We have so far presented a simplified picture of activated relaxation in liquids, which is more accurate at temperatures close to $`T_K`$, and thus sufficiently lower than $`T_A`$ \- the temperature at which activated processes emerge. The transition at $`T_A`$ where metastable minima emerge, along with a mosaic structure with intermediate tense regions, i.e. domain walls, is in many respects similar to a spinodal for an ordinary first order transition, except that the number of alternative phases is very large ($`e^{s_cN}`$ for a region of size $`N`$). The proper treatment of this transition must include fluctuations of the order parameter and consequent softening of the droplet surface tension at temperatures close to $`T_A`$. As a result of this, closer to $`T_A`$ the structural relaxation barriers are lowered from what would be expected extrapolating from near $`T_K`$ \- this gives deviations from the VF law. The corresponding length scales $`r^{}`$ and $`\xi `$ also should be smaller than would be predicted by the โvanillaโ, $`T_K`$-asymptotic version of the RFOT theory. These barrier โsofteningโ effects were quantitatively estimated in Lubchenko and Wolynes (2003b). They demonstrated that softening effects do vary between different substances and are more pronounced for fragile liquids. As a result, the value of the configurational entropy at $`T_g`$, as predicted by the RFOT theory with softening varies somewhat, within a factor of two or so among different substances. This is in contrast to the universal $`s_c(T_g)=.82`$ of the vanilla RFOT. Nonetheless, the value of $`\xi `$ at $`T_g`$ is much less sensitive and seems to be always within 5% of the simple estimate above. This is shown in the r.h.s. panel of Fig.7.
Understanding of the softening effect has allowed us to compute the activation barrier for liquid rearrangements in the full temperature range, including the high $`T`$ part near $`T_A`$, where the barriers become low, and the transport is dominated by activationless, collisional phenomena. Consistent with this predicted softening, the $`T`$ dependence of relaxation times, $`\tau =\tau _{\text{micro}}e^{\gamma ^2/4T^2s_c(T)}`$, as predicted by the RFOT (see Eq.(6)), fits well the experimental dependences in the low frequency range, but underestimates the viscosity near boiling. After softening is included, one can compute the activation component of the molecular transport, with the temperature $`T_A`$ as a fitting parameter of the theory Lubchenko and Wolynes (2003b). Fitting the viscosity was performed using the following obvious constraints: (a) at low temperatures, the order parameter $`\alpha `$ fluctuations are negligible, the barriers are fully established and high, and the transport is thus fully activated; (b) near boiling, the barrier vanishes, and the viscosity (known to be around a centipoise for all liquids) gives the value of $`\tau _{\text{micro}}`$. The fit, shown in Fig.7, demonstrates that of the 16-17 orders of the total dynamical range, about three orders, on the low viscosity side, are dominated by collisions. The experimental and activation-only theoretical curve differ from each other above a temperature $`T_{\text{cr}}`$. The three order of magnitude time scale separation, arising internally in the theory, is indeed consistent with the prerequisite of the transport being fully activated at $`T_{\text{cr}}`$ and below. The discussion above indicates that samples quenched (sufficiently fast) from a temperature $`T>T_{\text{cr}}`$ may exhibit somewhat distinct detailed molecular motions, also implying quantitative deviations form the RFOT predictions. At any rate, these sample, being caught in a very high energy state, are expected to have small cooperative regions, and also be very brittle and in general mechanically unstable. Such rapid quenches would be extremely difficult to produce in a lab, because $`T_{\text{cr}}`$ corresponds to relaxation times of the order $`10^{8\mathrm{}9}`$ sec. On the other hand, it is these ultra fast quenches, that must be currently employed by simulations owing to the limitations of computer power. We speculate that the thin โamorphousโ films made by vapor deposition on a cold substrate also may sometimes correspond to such ultra-quenches. While one may expect a number of behaviors in the bulk that are qualitatively distinct from what we have discussed here, various surface effects are likely to be important too: For one thing, such films are thin, have a large free surface, and strongly interact with the substrate. Further, there is a good reason to believe these films undergo local cracking, and spontaneous crystallization Perry (2004).
The present article deals with phenomena in glasses at temperatures much much lower than the temperatures at which the samples form. If a sample, upon vitrification, is cooled significantly below $`T_g`$, its lattice remains practically the same as of the moment of freezing. Indeed, the typical reconfiguration barrier is at least $`\mathrm{ln}(10^{15})35k_BT_g`$, as already mentioned. If, on the other hand, the sample in maintained at some temperature $`T`$ close enough to $`T_g`$, exceedingly slow structural relaxations take place. These attempts of the sample to equilibrate to a structure characteristic of temperature $`T`$ can be detected. Achieving quantitative accuracy in such experiments is difficult. Consistent with the notion that the lattice, and the barrier distribution, freeze in at the glass transition, the relaxation below $`T_g`$, obeys approximately the following temperature dependence:
$$k_{\text{n.e.}}=k_0\mathrm{exp}\left\{x_{\text{NMT}}\frac{\mathrm{\Delta }E^{}}{k_BT}(1x_{\text{NMT}})\frac{\mathrm{\Delta }E^{}}{k_BT_g}\right\},$$
(8)
where $`E^{}`$ is the equilibrated apparent activation energy at $`T_g`$ and $`x_{\text{NMT}}`$ lies between 0 and 1. This equation is part of the Nayaranaswany-Moynihan-Tool (NMT) empirical description of aging Tool (1946); Narayanaswamy (1971); Moynihan *et al.* (1976). The difference in the apparent activation energy above and below $`T_g`$, as expressed by the parameter $`x_{\text{NMT}}`$, will depend on how fast the barrier itself was changing, with cooling, above $`T_g`$, under โequilibriumโ cooling conditions. Since the rate of that change depends on the fragility, $`m=\frac{1}{T_g}\frac{\mathrm{log}_{10}\tau }{(1/T)}|_{T_g}=\frac{\mathrm{\Delta }E^{}}{k_BT_g}\mathrm{log}_{10}e`$, one expects that $`x_{\text{NMT}}`$ and $`m`$ are correlated. The RFOT based theory of aging in Lubchenko and Wolynes (2004) analyzes structural rearrangements in a non-equibrium glassy sample by means of Eq.(5), where the initial state is not equilibrium, but instead corresponds to the structure frozen-in at $`T_g`$. The predicted correlation between $`x_{\text{NMT}}`$ and $`m`$ is very simple: $`m19/x_{\text{NMT}}`$, and is consistent with experiment, see Fig.8.
For some of the comparison of theory and experiment it is necessary to be specific about the molecular length scale $`a`$ (a very detailed discussion of this quantity can be found in Lubchenko and Wolynes (2003b)). The molecular scale denotes the lattice spacing between molecular units (or โbeadsโ) that act as idealized spherical objects at the ideal glass transition at $`T_K`$. The determination of $`a`$, though approximate, is rather unambiguous and can be done using the knowledge of chemistry to give values accurate within 15%. For example, the number of beads in a chain molecule, that interacts with the surrounding only weakly, is always close to the number of monomers. Highly networked substances, such as amorphous silica, present a more difficult case, because it is not clear how covalent the intermolecular bonds in these substances are. Since melting also involves freeing up molecules, with encreased entropy, an independent check on the soundness of a particular bead number assignment can be done by comparing the fusion entropy of the substance (if it exists in crystalline form) with the known entropy of fusion of a hard sphere liquid or Lennard-Jones liquid, equal to $`1.16k_B`$ and $`1.68k_B`$ respectively Hansen and McDonald (1976). Note, however, the knowledge of the absolute value of $`a`$ is not required for most of the numerical predictions the theory will make in the quantum regime.
We thus see that the RFOT theory provides a rather complete picture of vitrification and the microscopics of the molecular motions in glasses. The possibility of having a complete chart of allowed degrees of freedom is very important, because it puts strict limitations on the range of ร priori scenarios of structural excitations that can take place in amorphous lattices. This will be of great help in the assessment of the family of strong interaction hypotheses mentioned in the introduction. To summarize, the present theory should apply to all amorphous materials produced by routine quenching, with quantitative deviations expected when the sample is partially crystalline. The presence and amount of crystallinity can be checked independently by X-ray. It is also likely that other classes of disordered materials, such as disordered crystals, will exhibit many similar traits, but of less universal character.
## III The Intrinsic Excitations of Amorphous Solids
### III.1 The Origin of the Two Level Systems
In this section we discuss how phenomena near the glass transition temperature, described in the previous subsection, dictate the existence and character of the quantum excitations in glasses at liquid helium temperatures and below. As mentioned earlier, a dynamical pattern of cooperative regions forms in a supercooled liquid below $`T_A`$. Each cooperative region is defined by the existence of at least two distinct configurations mutually accessible within the time scale $`\tau `$, which chatacterizes the life-time of the local mosaic pattern. Conversely, a molecular transport event is made possible by rearranging molecules within the cooperative length scale. The mosaic pattern โflickersโ on the time scale $`\tau `$; this process slows down dramatically upon vitrification and, below $`T_g`$ is referred to as โagingโ, as it corresponds to (very slow) structural changes. At $`T_g`$, the existent pattern of transitions (with distributed energy changes and reconfiguration barriers) freezes in because each cell is now surrounded by a rigid lattice (this is because the rearrangements of the neighboring domains were uncorrelated at $`T_g`$). Each region of the material can now explore the phase space as prescribed by the environment at the time of freezing. Below $`T_g`$, the mosaic is spatially defined by the molecular motions that were not arrested at $`T_g`$, and is thus strictly speaking only dynamically detectable. It is true that the weaker walls will probably be the site of (unstable) instantaneous normal modes in the fluid state with imaginary frequencies. This dynamical correlation pattern does not necessarily imply any easily discernible spatial heterogeneity in the atomic locations. In fact, there has been no direct evidence for any static type of heterogeneity of the appropriate scale in glasses so far, which definitely contributes to the (underappreciated) mystery of glasses <sup>3</sup><sup>3</sup>3We note, however, that there have been instances of mistaking polycrystalline samples for truly amorphous ones.. But can the dynamical heterogeneity be seen directly? We will claim later that this is done for us by thermal phonons: the magnitude of scattering at the plateau can only be explained by presence of dynamical heterogeneities. The latter are signified by structural transitions that scatter the phonons inelastically. Apart from aging (which we will ignore in the rest of the work), a particular pattern of flipping regions, as frozen-in at $`T_g`$, will persist down to the lowest temperature. The apparent size of each cell in this mosaic of flippable regions will depend on the observation time. The longer this time is, the more structural relaxation degrees of freedom (from the high barrier tail of the barrier distribution) one should observe. Eventually, in fact, the glass should crystallize<sup>4</sup><sup>4</sup>4Note that there are, in principle, other ways to move molecules in a glass, in addition to the cooperative rearrangements: for example by creating defects such as vacancies (the corresponding barriers are prohibitively high, of course). In order to estimate the number of tunneling centers that are thermally active at low temperatures we will have to find the size of the regions that allow for a rearrangement accompanied by a small energy change and, at the same time with a low barrier. It may reasonably seem that typically such barriers for multiparticle events would be very high. Nevertheless, the lattice is arrested in a high energy state. We can thus foresee the possibility of stagewise barrier crossing (or tunneling) events, when the width of the barrier for each consecutive atomic movement is only a small fraction of a typical interatomic distance, thus rendering individual atomic movements nearly barrierless. This is as if one could define an instantaneous mode of nearly zero frequency, at each point along the tunneling trajectory. (Yet at no point is the motion harmonic per se!) The presence of such low frequency modes should be expected given the high number of configurational states available to the sample as the moment of freezing, as reflected in the high value of the configurational entropy at $`T_g`$. After all, the material is unstable, both globally and locally! (Note, the extent of bond deformation during an individual atomic movement is low - within the Lindemann length - actually affording a few โhardโ places along the tunneling trajectory, where the โinstantaneousโ frequencies are not necessarily low.) One may contrast the situation above with, say, tunneling of a substitutional impurity in a crystal, a system which is indeed near its true ground state. Such tunneling would not contribute to the very low $`T`$ thermal properties owing to a large barrier. Also, we note that multiparticle barrier crossing events have been seen in computer simulations of amorphous systems Guttman and Rahman (1986), anticipated theoretically Mon and Ashcroft (1978); Heuer and Silbey (1993), and recently inferred from simulations of dislocation motions in copper Vegge *et al.* (2001).
We summarize the discussion so far by noting that the preceding Section has demonstrated that the possible atomic motions in a supercooled liquids are either purely vibrational excitations or structural rearrangements. Any possible motions below $`T_g`$, in terms of the classical basis set must be a subset of the motions above $`T_g`$, although the dynamics of these events become quantum-mechanical at low enough temperatures. Even after the system is cooled to an arbitrarily low temperature, it remains essentially in the configuration in which it got stuck at the glass transition. The density of directly accessible states at that high energy configuration is rather high; the total density of states is, of course, exponentially larger, but inaccessible on realistic time scale without other regions of the glass rearranging. Since the typical rearrangements near $`T_g`$ span about a length $`\xi `$ across, we may make the following, preliminarily conclusion: The non-equilibrium character of the glass transition necessarily dictates the existence of intrinsic additional non-elastic degrees of freedom in a glass, tentatively one per region of roughly size $`\xi `$, in addition to the usual vibrations of a stable lattice. The universality of $`\xi `$, in a sense, is the main clue to the cryogenic universality that is observed. A schematic of a cooperative region is shown in Fig.9.
Note that showing the existence of low energy tunneling paths is really a mathematically problem of finding hyper-lines, connecting two points of particular latitude on a high-dimensional surface, that meander within a certain latitude range. Visualizing high-dimensional surfaces is prohibitively difficult, while the field of topology, at its present stage of development, is of little help. Yet, a completely general argument is not required here: We only need to consider a very small subset of all surfaces, such that they satisfy the (very severe) constraint on the liquid density distribution above $`T_g`$, namely such that conforms to an โequilibratedโ liquid at $`T_g`$. Because (and with the help) of this constraint, it is possible to put forth a formal argument showing that there are indeed enough low energy structural transitions in a frozen glass: This argument will follow (albeit in the reverse order, in a sense) the argument from the preceding Section, where we found the typical trajectory for rearrangement. The key point of the microcanonical-like library construction from Section II is that the distribution width of energies of a region increases with region size. A region is guaranteed to have a state at some low energy, call it $`E_{\text{GS}}(N)`$, as found by integration in Eq.(3). Past a certain critical size, this energy decreases as $`N`$ grows larger, giving rise to the existence of a resonant state at a large enough $`N`$. One must bear in mind, however, that $`E_{\text{GS}}(N)`$ reflects the typical freezing energy. It really gives an upper bound on the lowest energy level. The actual lowest energy state fluctuates and always lies below $`E_{\text{GS}}(N)`$, although most likely not much below. Here we will look in detail the statistics of these energy states below the typical reconfiguration profile, with the aim to find the probability of a low energy trajectory for reconfiguring a region size $`N`$.
We will make several preliminary, quite general notions that will guide us in constructing an adequate approximation for the local statistics of the energy landscape of a frozen lattice. First, we give a general argument of the density of frozen-in excitations, valid, as we will see shortly, in the limit of infinitely slow aging: Since the atomic arrangement does not change upon freezing, the classical density of states of a frozen glass is that of the supercooled liquid at $`T_g`$. Those states correspond to configurations in which the system could have frozen at $`T_g`$ and in principle can explore, provided they are thermally accessible and have a sufficiently low barrier separating them from a given configuration. Take a generic liquid state at $`T_g`$ as the reference state. Then the Boltzmann probability to switch to a conformation higher in energy by amount $`ฯต`$ is $`e^{ฯต/T_g}`$. That a configuration with that energy was one of the allowed configurations upon freezing means there must have been $`n(ฯต)=\frac{1}{T_g}e^{ฯต/T_g}`$ of them. The factor $`1/T_g`$ arises because the energy spectrum by construction is continuous, while the actual local spectrum is discrete and $`ฯต=0`$ gives the upper bound on the location of the actual ground state. The latter must be somewhere between $`0`$ and $`\mathrm{}`$: $`_{\mathrm{}}^0๐ฯตn(ฯต)=1`$. This argument, however, is silent as to what the spatial characteristics of such excitations or their time scale are.
This inverse โBoltzmannโ density of states has been computed explicitly in frustrated mean-field spin systems Mรฉzard *et al.* (1985), but is of more general nature. Indeed, such distributions arise universally when describing the statistics of the lowest energy state of a wide class of energy distributions Bouchaud and Mรฉzard (1997), including the random energy model Derrida (1981), that will be used later on. Kinetic considerations did not explicitly enter our heuristic derivation above (or, the mean-field estimates in Mรฉzard *et al.* (1985); Bouchaud and Mรฉzard (1997)). This is directly seen by differentiating $`\mathrm{log}n(ฯต)/ฯต=1/T_g`$. Clearly, $`n(ฯต)`$ is the microcanonical density of states corresponding to the translational (liquid-like) degrees of freedom, and the system is assumed to be completely ergodic within that set of states. This corresponds to an approximation where we consider all degrees of freedom which are faster than a given time scale as very fast, and everything slower than that chosen time scale is regarded to be much slower than can be detected in the experiment. By using this same density $`n(ฯต)`$, as it was at $`T_g`$, also at $`T<T_g`$, we formally express the fact that this subset of the total density of states no longer thermally equilibrates but stays put where it was at $`T_g`$ \- the subsystem of the translational degrees of freedom has undergone an entropy crisis, a glass transition. Everywhere in the discussion above, we have been ignoring the contribution of the purely vibrational excitations to the total free energy. We thus assume that the spectrum of those elastic excitations is independent of precisely where on the glassy landscape the liquid is.
We now give the argument, first laid out in Lubchenko and Wolynes (2001), that allows one to estimate the classical density of states and will also simultaneously yield the size of the region where the excitation takes place. First we address the question of how many structural states are available to a compact fragment of lattice of size $`N`$, regardless of the barrier that separates those alternative states from the initial ones. This corresponds to the assumption of time scale separation mentioned just above. Within this assumption, the low energy limit of the spectrum must obey $`e^{E/T_g}`$ so as to give a glass transition at $`T_g`$. Next, the spectrum, when integrated, must give $`e^{s_cN}`$ for the total number of states available to the region. Notice further that we expect the reconfiguring regions to be relatively small. The atomic motions within these small regions are directly coupled and so a mean-field, gaussian density of states, that only describes lowest order fluctuations around the mean, should be accurate enough. An energy density satisfying the requirements above actually corresponds to the well known Random Energy Model (REM) Derrida (1981), which also describes the pure state free energy in mean field frustrated spin models:
$$\mathrm{\Omega }_N(E)\mathrm{exp}\left\{s_cN\frac{[E(N\mathrm{\Delta }ฯต+\gamma \sqrt{N})]^2}{2\delta E^2N}\right\},$$
(9)
where $`\delta E^2`$ is the variance to be determined shortly. Here, the factor $`e^{s_cN}`$ gives the correct total number of states, the term $`\gamma \sqrt{N}`$ takes into account the interface energy cost of considering distinct atomic arrangements with the region. Note the fluctuations in the surface term are expected but are automatically included in the fluctuation of the microcanonical energy $`E`$ itself.The term $`N\mathrm{\Delta }ฯต`$ is a bulk energy necessary to account for the observed excess energy of the frozen structure relative to the energy of the ideal structure at $`T_K`$. It is easy to relate to measured quantities: To do this, recall that the system freezes in its โground stateโ, with energy $`E_{\text{GS}}`$, when its entropy becomes non-extensive:
$$\mathrm{\Omega }_N(E_{\text{GS}})=1.$$
(10)
We take the energy $`E_g`$ of the liquid state at $`T_g`$ as the reference energy. Next note that in the absence of the surface energy term $`\gamma \sqrt{N}`$, the lowest available energy state is that of the liquid at $`T_K`$: $`(E_KE_g)/N=_{T_K}^{T_g}๐T\mathrm{\Delta }c_p(T)\mathrm{\Delta }c_p(T_gT_K)T_gs_c`$. (The two latter equalities are accurate for $`T_g`$ close to $`T_K`$. The corrections would be observable Lubchenko and Wolynes (2004, 2003b), but small.) One immediately gets from Eqs.(9) and (10) that $`\mathrm{\Delta }ฯต=\sqrt{2\delta E^2s_c}T_gs_c`$ ($`\gamma =0`$ must be used in this estimate, but nowhere else!). Further, using the microcanonical $`\mathrm{ln}\mathrm{\Omega }_N(E)/E|_{E=E_{\text{GS}}}=1/T_g`$ fixes the value of the variance $`\delta E^2=2T_g^2s_c`$. The resultant density of states is proportional to $`e^{(EE_{\text{GS}})/T_g}`$ at $`T_g`$, as already shown above by a general argument. Now that we have determined the thermodynamical quantities entering Eq.(9), we can find how the excess energy of an alternate ground state depends on the size $`N`$:
$$E_{\text{GS}}(N)=\gamma \sqrt{N}T_gs_cN,$$
(11)
where $`E_{\text{GS}}`$ is defined by Eq.(10). Only low energy excitations will be thermally active at the lowest temperatures. Therefore, we are looking for excitations that are nearly isoenergetic with the reference state. This imposes an additional condition $`E_{\text{GS}}(N)=0`$ thus prescribing the minimal size $`N_0`$ of a region such that has $`\mathrm{\Omega }_N(0)1`$ for $`NN_0`$. A region of this size has at least one alternative structural state at the same energy. One obtains from Eq.(11) that $`N_0=(\gamma /T_gs_c)^2=N^{}`$, consistent with our previous argument that any region of size $`N^{}`$ has a spectral density of states equal to $`\frac{1}{T_g}e^{E/T_g}`$. Note that Eq.(11) echoes the free energy profile of droplet growth from Eq.(6), but unlike Eq.(6), it can be used for $`N>N^{}`$ as well. Eq.(11) explicitly shows that a droplet of size larger than $`N^{}`$ has an exponentially increasing number of available configurations corresponding to lattices typical of $`T_g`$, consistent with the instability of droplets larger than $`N^{}`$ at temeperatures above $`T_g`$ mentioned earlier.
The microcanonical argument above is basically a gedanken experiment in which we had the demon-like ability to browse through all possible atomic arrangements, given the total number of allowed states equal to $`e^{Ns_c}`$. The total sample is thus comprised of regions of the type considered in the argument (the interface energy has been taken into account by the term $`\gamma \sqrt{N}`$, this energy may be viewed as the penalty for considering the states of a given compact region as if this region were totally independent from the rest of the sample; c.f. our earlier comments on the locality of the liquidโs energy landscape). Therefore, if the rate of conversion between the alternative glassy states can be ignored, the argument immediately yields the density of residual excitations in a frozen glass: $`\frac{1}{N^{}a^3T_g}e^{ฯต/T_g}`$ ($`N^{}a^3\xi ^3`$, of course). However, even though each of these imaginary regions has an alternative resonant state, there is so far only an undetermined chance to reach it within any particular time. In fact, the typical classical barrier for the excitations available to the regions of size $`N^{}190`$ is $`F^{}39k_BT_g`$. Such a barrier would seem to exclude the possibility for tunneling for a typical domain of size $`N^{}`$. But to account for kinetics issues, we should repeat the argument for the critical size, but also simultaneously include the life-time of each considered configuration as a selection criterion. In other words, one should compute the combined distribution of the excitation energies, their spatial extent and the corresponding tunneling amplitudes. Later in the paper, we will discuss one source of correlation between the excitation energy and the tunneling amplitude owing to level repulsion effects. Nonetheless here, we present a simpler argument, given in Lubchenko and Wolynes (2001), in which the tunneling rate distribution is assumed to be independent from that of $`ฯต`$. Simpler yet, we will look for the density of regions that allow for a rearrangement with a zero-height barrier. As vindicated post factum, all of these simplifications can lead only to at most a 10% error in the resulting density of states.
Imagine the process of conversion to another state as a step-wise process where the โnucleusโ of this new state is increased by adding one atom at a time, as signified by the horisontal axis in Fig.10.
Such addition involves moving the atom a distance of the order of the Lindemann distance $`d_L`$. It follows then that the path connecting the two states is likely to encounter a high barrier of the order $`F^{}`$, which effectively disconnects those two states. However, the possible configurations through which one can pass and therefore the barrier heights are distributed and there is a chance even for a region of size $`N^{}`$ to have an arbitrarily low barrier. What would such a distribution be? We first have to decide whether the tunneling probability is a sum of contributions of many (interfering) paths or, whether it is dominated by a single path, which has the lowest barrier. The first scenario would be realized in a highly quantum glass, where Debye temperature rivals or exceeds the glass transition temperatures Schmalian and Wolynes (2000). Such a highly quantum glass could in fact melt due to quantum fluctuations. In our case, freezing is a completely classical process, which is signified by the fact that the barriers are proportional to a classical energy scale $`T_g`$. We now assume more specifically that the contribution of a tunneling path is proportional to $`e^{\pi V^{}/\mathrm{}\omega ^{}}`$, where $`\omega ^{}`$ is a quantum frequency scale, a multiple of $`\omega _D`$, and barrier $`V^{}`$ scales with $`T_g`$, as mentioned earlier. This would be an accurate assignment in the case of a parabolic barrier. The form of the tunneling amplitude $`e^{\pi V^{}/\mathrm{}\omega ^{}}`$ conforms to our expectation that the tunneling trajectory is dominated by a single path with the lowest barrier, as $`V^{}`$ and $`\mathrm{}\omega ^{}`$ are taken from distributions characterized by scales $`k_BT_g`$ and `const`$`\times \mathrm{}\omega _D`$ respectively, the former one generally much larger than the latter (the `const` self-consistently will turn out to be less than one in subsection IV.2). Since the energy profile along the tunneling trajectory has a complicated shape formed by many intermediate states separated by small intermediate barriers (see below), it is fair to say that the state of the system at the highest barrier corresponds to the highest energy intermediate state (the โtransitionโ state). The statistics of energy states have already been found earlier. We therefore use distribution (9) with only one difference: we must double the variance because the barrier height is actually the difference between two fluctuating quantities: the energy of the final (or initial) energy and the highest energy along the path. As a result,
$$\mathrm{\Omega }_N(V)\mathrm{exp}\left\{s_cN\frac{[V(T_gs_cN+\gamma \sqrt{N})]^2}{4\delta E^2N}\right\}.$$
(12)
Distribution (12) thus gives the typical value of the barrier for the (quantum) growth of a droplet. It is easy to see from (12) that the highest barrier corresponds to rearranging a region of size $`N^{}`$ 14 and is equal to
$$V_{max}=F^{}/(2\sqrt{2}1)26T_gs_c.$$
(13)
Since this is the hardest place to get through, we must take it as the transition state. Hence, the final distribution of (transition state) barriers is the density of pure states corresponding to Eq.(12) with $`N=N^{}`$ (similar to the $`e^{ฯต/T_g}`$, obtained above). Thus,
$$\mathrm{\Omega }(V^{})\mathrm{exp}\left\{\frac{V^{}V_{max}}{\sqrt{2}T_g}\right\}=\mathrm{exp}\left\{18s_c+\frac{V^{}}{\sqrt{2}T_g}\right\}.$$
(14)
As one can see, the probability to have a small barrier path is exponentially suppressed. Nevertheless, owing to the large value of the energy parameter in distribution (12) the fraction of zero barrier paths per mosaic cell $`e^{18s_c}310^7`$ is actually not prohibitively small. A region larger by only 18 molecules (less than a single layer) will have $`e^{18s_c}`$ more final states (and therefore paths) to go to. We therefore conclude, any region of size $`200`$ molecules will have an accessible alternative state with spectral density $`1/T_g`$. Finally, we stress a remarkable feature of the tunneling paths statistics in glass. Mark the very rapid - exponential - scaling of the number of paths leading out of a particular local structural state on the size of the respective region. This means that the final estimate of the density of structural transitions that have low enough barriers to be thermally relevant is rather insensitive to the details of correlation between the energy of the transition and its tunneling amplitude. Consequently, even a very simple estimate of this density, such as the one above, is very robust. Finally, note that the tunneling argument above is, again, a microcanonical argument, such as the one leading to Eq.(9), that also takes into account (in a rather crude manner) the mutual accessibility between alternative energy states.
As we will see later, the tunneling barriers, and hence the relaxation times of the tunneling centers, are distributed. This would lead to a time dependent heat capacity. Ignoring this complication for now, the classical, long time heat capacity is easy to estimate already (assuming it exists): Since our degrees of freedom span a volume $`\xi ^3`$ and their spectral density is $`1/T_g`$ at low energy, one obtains for the low $`T`$ heat capacity per unit volume: $`T/T_g\xi ^3`$, up to an insignificant coefficient. The coefficient at the linear heat capacity dependence is often denoted $`\overline{P}`$. For silica, $`T_g=`$1500 K and realistic $`\xi =20\AA `$ yields $`\overline{P}610^{45}`$m<sup>-3</sup>J<sup>-1</sup> in agreement with the experiment (we took $`a=3.5\AA `$ \- a length scale appropriate for a tetrahedron formed by four oxygens around a silicon atom. These tetrahedra appear to be moving units in a-Silica Trachenko *et al.* (2000)). The assumption of the existence of the long-time heat capacity is empirically justified (within logarithmic accuracy), but is also consistent with the present theory, see Subsection III.3 and Section V.
In conclusion, the main result of this Subsection is that the non-equilibrium nature of the glass transition results in the existence of residual motional degrees of freedom, a significant fraction of which remain thermally active down to the lowest temperatures. These degrees of freedom are collective highly anharmonic atomic motions within compact regions of size $`\xi ^3`$, determined mainly by the length scale of the entropic droplet mosaic determined at $`T_g`$. The energy scale in the spectrum of these excitations is set by the glass temperature $`T_g`$ itself. We now turn to the question of what determines the strength with which these entropic droplet excitations couple to the phonons. This will explain the universality in the heat conductivity at temperatures below $``$1 K.
### III.2 The Universality of Phonon Scattering
First of all, what do we mean by โphononsโ in amorphous materials? There is no periodicity, therefore one can only strictly speak of elastic strain, even if the structure is completely stable. In the latter case, low gradient strains $`\varphi `$ are still described by a simple bi-linear form:
$$H\text{ph}d^3๐ซ\frac{\rho c_s^2}{2}(\varphi )^2.$$
(15)
The $`\varphi `$ field is defined on a isotropic, translationally invariant flat metric, as in continuum mechanics Landau and Lifshitz (1986), and so a wave-vector $`k`$ is an operational concept. It is easily seen, by dimensional analysis, that strains arising specifically due to disorder will be of higher order in $`k`$ than the term in Eq.(15): The corresponding energy terms should scale with some positive power of the lattice inhomogeneity lengthscale(s), $`l_{\text{inhom}}^\zeta `$ ($`\zeta >0`$), so as to vanish upon โzooming outโ. The terms will subsequently go as $`k^2(l_{\text{inhom}}k)^\zeta `$. But, as we have already seen, the amorphous lattice is not stable, that is there are anharmonic transitions with arbitrarily small energy and barrier. Still, the regions encompassing the transitions are quite small, at most 6 lattice spacings across, which is much less than the thermal de Broglie wave-length at 1 K (about $`10^3a`$ in a-silica). The unstable regions interact with the strains of the otherwise stable lattice. We conceptualize this interaction by approximating the strain with pure phonons and computing the phonon mean-free path. The latter will turn out to be about 150 times longer than the phonon wave-length, so that the phonon approximation is internally consistent in that the phonons are indeed reasonably good quasi-particles, at the plateau temperatures and below. Finally, note in Eq.(15), we have used only one phonon polarization for simplicity (it will be usually obvious how to account approximately for all three acoustic phononic branches at the end of a calculation; using this โscalarโ version of the lattice dynamics, for the purposes of this paper, boils down to neglecting the difference between the longitudinal and transverse sound, except in the later discussion of the Grรผneisen parameter).
The structural transitions interact with the phonons because the energy of the transition changes in the presence of a strain: To see this exlicitly, consider the elastic energy within a droplet-sized region capable of undergoing a low energy transition, as relevant at low $`T`$. For the sake of argument, assume there is no sheer deformation; a similar argument applies to the transverse phonon branches. The stress energy is then $`_{V_\xi }d^3๐ซKu_{ii}^2/2`$, where $`K`$ is a compressibility constant on the order of $`\rho c_s^2`$, which we are allowed to treat as a constant, with the error contributing in a higher order in $`k`$, as already mentioned. $`u_{ii}`$ is the trace of the elastic tensor Landau and Lifshitz (1986), which has the same meaning as $`\mathrm{\Delta }\varphi `$. Since the transition energy is low, the lowest order, quadratic expansion suffices. This implies that all individual bonds within the region distort by a very small amount (already shown not to exceed $`d_L`$, even at the transitionโs bottle-neck). We have demonstrated that such regions do indeed exist and found their density in the previous section. We separate the total elastic tensor $`u_{ij}`$ into a contribution $`\varphi _{ij}`$ due to the elastic stress and $`d_{ij}`$ due to the tunneling displacement. The $`d_{ii}\varphi _{ii}`$ cross term represents the coupling between the transition and the strain. If the phonon wave-length is larger than $`\xi `$, $`\varphi _{ii}`$ is constant within the integration boundaries and can be taken out of the integral. One consequently arrives at the following energy difference for the defect states in the presence of a phonon: $`\rho c_s^2\varphi _{ii}_{V_\xi }d^3๐ซd_{ii}`$. Here, $`d_{ii}`$ corresponds to the transition induced displacement between two given structural states. Each tunneling center is a multilevel system. However, since we are presently interested in the coupling of the lowest energy transition to the strain, we consider here the set of $`d_{ii}`$ corresponding to the two lowest energy states of the region. (Higher energy transitions turn out to be intimately related to the lowest energy transition, and are discussed later, in Section IV.) We therefore conclude that a tunnelling transition, active at low temperatures, is linearly coupled to a lattice strain with the strength defined as $`g=\rho c_s^2_{V_\xi }d^3๐ซd_{ii}/2`$; the corresponding term in the Hamiltonian reads:
$$H_{int}๐ \varphi \sigma _z.$$
(16)
We present next two independent ways to estimate the coupling constant $`g`$. The first one is based on the realization that at the glass transition, purely phononic excitations and a frozen-in structural transition must coexist, that is they are are of marginal stability with respect to each other. On the one hand, a local region posed to harbor a structural transition below $`T_g`$, must not be โcrumpledโ by a passing phonon. On the other hand the energy of the transition can only be so high, as to be sustainable by the lattice stiffness. In other words, an atom will be part a of frozen-in transition, if that atom is roughly in equilibrium between the transition driving forces and the ambient lattice strain. This stability condition gives at the molecular scale $`a`$, by Eqs.(15) and (16): $`๐ \sigma _z=\rho c_s^2a^3\varphi `$. The lattice strain will be distributed throughout the material in the usual manner, subject, of course, to the equipartition that fixes the variance of the strain so as to conform to the thermally availabe energy. We take advantage of this by multiplying the equilibrium condition by $`\varphi `$ and noting that thermal averaging is also ensemble averaging. This yields that for an atom posed to be part of low energy structural transitions below $`T_g`$, it is generically true that
$$|๐ \varphi \sigma _z|\rho c_s^2(\varphi )^2a^3k_BT_g.$$
(17)
Noting that $`๐ \varphi \sigma _z|๐ \varphi |`$, one arrives at a simple relation:
$$g=\sqrt{\rho c_s^2a^3k_BT_g},$$
(18)
which is the main result of this Subsection. It is understood that this estimate is accurate up to a number of order one, and $`g`$โs are likely distributed. At any rate, we observe that the TLS-phonon coupling is the geometric average between the glass transition temperature and an energy parameter $`\rho a^3c_s^2`$ ($``$ 10<sup>1</sup>eV) related to the cohesion energy of the lattice (note the quantity $`\rho c_s^2`$ is a multiple of the Youngโs modulus). We point out the estimate above applies quantitatively to low barrier transitions only. The mechanical stability criterion is a zero frequency, static condition. However, it takes a finite time for a structural transition to respond to an external stress. In other words, a region harboring a slow transition will likely appear perfectly elastic to a high frequency phonon. We thus arrive at the conclusion that TLS-phonon coupling must be frequency dependent, however deviations from the result obtained above will enter in a higher order in $`\omega ,k`$.
Alternatively, one may attempt to estimate the integral over the derivative of the displacement field that entered in the expression for the coupling constant $`g=\rho c_s^2_{V_\xi }d^3๐ซd_{ii}/2`$. Since $`d_{ii}`$ is the divergence of a vector, the integral is reduced to that over a surface within the dropletโs boundary: $`_{S_\xi }๐๐ฌ๐(๐ซ)`$, where $`๐(๐ซ)`$ is the tunneling displacement itself, near the boundary. How near? The Gauss theorem applies so long as the field $`๐`$ is continuous. This field is roughly $`d_L`$ in magnitude close to the dropletโs center and is zero outside of the region. A function defined on a discrete lattice is expressly discontinuous. Requiring that one be able to cast a continuous tunneling displacement field on a discrete manifold, so that the field interpolates smoothly between $`d_L`$ in the middle and zero outside, imposes constraints on the $`๐`$ values at the dropletโs boundary. We argue this value should generically go as $`(a/\xi )d_L`$, up to a constant, in order to realize the interpolation and spread evenly the tensile field of the domain wall throughout the droplet. In other words, $`(a/\xi )d_L`$ is quite obviously the lower bound, while higher values statistically imply a higher stress, and higher transition energy. Now, since $`๐(r)`$ is randomly oriented, the integral over the dropletโs border is of the order $`a^2\sqrt{N_{}^{}{}_{}{}^{2/3}}(a/\xi )d_L`$, where $`N_{}^{}{}_{}{}^{2/3}`$ is the number of molecules at the boundary. The Lindemann distance at $`T_g`$ is equal to the magnitude of a thermal fluctuation, hence $`d_L/a|\varphi |`$ <sup>5</sup><sup>5</sup>5For reference, $`|\varphi |(T_g/\rho c_s^2a^3)^{1/2}`$ at $`T_g`$ is about 0.05 for SiO<sub>2</sub>, 0.06 for B<sub>2</sub>O<sub>3</sub> (oxide glasses), 0.03 for PS and PMMA (polymer glasses), in agreement with the Lindemann criterion. We stress the sensitivity to the value of the molecular size $`a`$, which is somewhat arbitrary. Here, we have not calculated the bead size based on chemistry, but instead used the values of the speed of sound, as employed in the scaling procedure of Freeman and Anderson Freeman and Anderson (1986). According to the definition of the Debye frequency, $`a=(c_s/\omega _D)(6\pi ^2)^{1/3}`$.. As a result, the coupling to the extended defects is still about $`g\sqrt{\rho c_s^2a^3T_g}`$, again within a factor of two or so.
Note that, when considering a particular value of lattice distortion $`\varphi `$ in the discussion above, we did not specify the wave-length(s) of the phonons that contributed to this distortion, therefore the estimate of $`g`$ in Eq.(18) is correct as long as the form of the interaction term (Eq.(16)) is adequate. This surely holds for long-wave phonons relevant at the TLS temperatures.
With the knowledge of $`g`$, we can estimate the inverse mean free path of a phonon with frequency $`\omega `$. As done originally within the TLS model, the quantum dynamics of the two lowest energies of each tunneling center are described by the Hamiltonian $`H_{TLS}=ฯต\sigma _z/2+\mathrm{\Delta }\sigma _x/2`$. This expression, together with Eqs.(15) and (16), form a complete (approximate) Hamiltonian of the TLS plus the lattice vibrations. The phonon inverse mean free path is then calculated in a standard fashion Jรคckle (1972); Phillips (1981):
$$l_{\text{mfp}}^1(\omega )=\pi \frac{\overline{P}g^2}{\rho c_s^3}\omega \mathrm{tanh}\left(\frac{\mathrm{}\omega }{2k_BT}\right).$$
(19)
This yields
$$\lambda _{dB}/l_{\text{mfp}}=\frac{2\pi ^2}{3}\mathrm{tanh}(1/2)\left(\frac{a}{\xi }\right)^3,$$
(20)
where factor $`1/3`$ comes from the averaging with respect to different orientations of the defects and we used $`\overline{P}1/T_g\xi ^3`$. It follows that $`l_{\text{mfp}}/\lambda _{dB}(\xi /a)^3200`$ up to a constant of order one. Hence, the analysis above explains the universality of the combination of parameters $`\overline{P}g^2/\rho c_s^2`$, and relates it to the geometrical factor $`(a/\xi )^310^2`$, which is the relative concentration of cooperative regions in a supercooled liquid, an almost universal number within the random first order glass transition theory Xia and Wolynes (2000), depending only logarithmically on the speed of quenching. Strictly speaking, our argument predicts the universality in $`l_{\text{mfp}}/\lambda `$ only within 10%-20% or so. This is a consequence of indeterminacy of $`ฯต`$ vs. $`\mathrm{\Delta }`$ correlation that may be system specific, or could be due to deviations of the $`\xi /a`$ ratio from the universal 5.8 at $`T_g`$. Since the latter ratio depends on the glass preparation time, the corresponding experimental study seems worthwhile.
Numerically, Eq.(20) yields $`l_{\text{mfp}}/\lambda 70`$, a factor of $`2`$ less than the empirical $`150`$ Freeman and Anderson (1986). We could not have expected much better accuracy from our estimates, that used no adjustable parameters. Although it may seem that we have slightly over-estimated the number of scatterers, the size of the error is too small to reliably support this suggestion. However, this is a good place to make a few comments on the distribution of the tunneling matrix elements $`\mathrm{\Delta }`$, which will also prove useful for the discussion of the phonon scattering at higher frequencies in Chapter IV. The estimate for the phonon mean free path in Eq.(19) is not terribly sensitive to the form of tunneling amplitude distribution Phillips (1981) (within reasonable limits). This is because the contribution of an individual TLS to the total scattering cross-section is proportional to $`\mathrm{\Delta }^2/E^2`$, where $`E\sqrt{\mathrm{\Delta }^2+ฯต^2}`$. Two common distributions have been used in the literature. One distribution simply assumed $`(\mathrm{\Delta }^2/E^2)1`$ in the absence of a more specific knowledge and a flat distribution of the total energy splitting $`E`$ (this is actually the original TLS model). The earlier Standard Tunneling Model (STM) Anderson *et al.* (1972); Phillips (1972), on the other hand, postulates $`P(\mathrm{\Delta })\frac{1}{\mathrm{\Delta }}`$ (approximately supported by our own conclusions too), which predicts a nearly flat $`E`$ distribution as well. In the end, both models differ only in that the TLS model has to postulate the average $`(\mathrm{\Delta }^2/E^2)`$ value when calculating the scattering cross-section (this number is absorbed into the TLS-phonon coupling constant $`g`$). On the other hand, the STM allows for somewhat more closed-form derivations, however it still has to introduce the cut-off values $`\mathrm{\Delta }_{min}`$ and $`\mathrm{\Delta }_{max}`$ as parameters (fortunately, many measured quantities depend on these parameters only logarithmically). (The distinction between the two models is described in detail in the front article by Phillips in Phillips (1981).) One point in favor of the STM is that it necessarily predicts time dependence of the specific heat. While a time dependence has been observed, its specific functional form has not been unambiguously established in the experiment (see Pohl (1981); Hunklinger and Raychaudhuri (1986)). We also mention, for completeness, there is a different way to parametrize the two-level system motions within the more general, soft-potential model Karpov *et al.* (1983); Buchenau *et al.* (1992). At any rate, we see that while a two-level systemโs contribution to the total phonon scattering depends on the value of its $`\mathrm{\Delta }E`$ ratio, the precise form of the $`\mathrm{\Delta }`$ distribution will change the answer quantitatively, but not qualitatively. We note, that in the context of the present calculation it is preferable to consider the simpler, TLS setup that does not specify the $`\mathrm{\Delta }`$ distribution, because our argument has so far been only semi-classical. Indeed, so far the tunneling amplitudes have only interested us from the perspective of the volume density of allowed transitions. We saw that an indeterminacy of the density could not exceed a 10% or so due to a weak (logarithmic) dependence of that density on a specific $`\mathrm{\Delta }`$ distribution. Therefore, we are more confident in the numerical estimates using the TLS-model setup that does not require introducing additional parameters (such as $`\mathrm{\Delta }_{min}`$) explicitly. Nevertheless, the special role of $`\mathrm{\Delta }E`$ defects in scattering is worth noting. These defects have low classical energy splittings $`ฯต<\mathrm{\Delta }`$ and their dynamics are mostly determined by the quantum energy scale. These are the โfastโ, or โzero-barrierโ excitations discussed earlier in the literature Black and Halperin (1977); Geszti (1982), whose tunneling matrix element probably can not be directly estimated by WKB, but we can still guess that it scales with $`\omega _D`$. This suggests that using the same distribution function $`P(\mathrm{\Delta })`$ for all thermally defects may not be justified, as circumstantially supported by results of Black and Halperin Black and Halperin (1977) who found that the density of TLS derived from the heat capacity and conductivity measurements respectively are not exactly equal to each other. While this indeterminacy in the exact barrier distribution introduces only an error of order one in quantitative estimates Black and Halperin (1977) of the density of states and is not of special concern here, we note that the present theory, upon inclusion of the effects beyond the strict semi-classical picture, does in fact provide a mechanism for the excess of the โfastโ two-level systems, as will be explained in Section V.
Strong Interaction Scenarios. By deriving the density of states of structural transitions, and their coupling to the phonons based on the known properties of the amorphous lattice, we have constructively established the microscopics of glassy excitations in excess to the purely elastic excitations. It follows from the discussion that no other excitations are present in glasses at 1 K and below (see also the discussion on the exhaustive classification of excitations in amorphous lattices in Subsection IV.1). Importantly, the density of states (DOS) in excess of the phonons, is due to local motions. This is in contrast with Strong Interaction Scenarios (SIS) Yu and Leggett (1988); Leggett (1991); Burin and Kagan (1996); Coppersmith (1991) that posit that any local excitations (other than pure strain) would give rise to a โuniversalโ density of states. Such universal density of states arises in SIS as a consequence of long range, $`1/r^3`$, interaction mediated by the phonons, so that the actual observed DOS is a highly renormalized entity. The corresponding excitations are expressly non-local, possibly infinite in extent. The idea is very attractive because of its generality but remains a pure abstraction, until those bare excitations are constructively shown to exist in the first place. Additionally, even upon assuming some bare excitations are present, demonstrating the quantitative relationship between the effective density of states and the phonon coupling that conforms to the experimental $`\overline{P}g^2/\rho c_s^210^2`$ has so far proven elusive Leggett (1999); Caruzzo (1994); Lubchenko and Wolynes (2000). On the other hand, we have shown that local structural transitions, that interact with phonons with a particular strength, must indeed take place in amorphous solids. In order to determine where the current theory stands in relation to the SIS, one may inquire whether the phonon-mediated interaction leads to the emergence of some collective density of states. It should be immediately clear that no such additional, collective DOS appears at $`T_g`$, because the argument in the previous subsection has already included all the effects of the surrounding of a structural transition, which simply amounted to the thermal noise at $`T_g`$ delivered to the transition by the elastic waves. It, of course, does not matter what the phonon source is. What happens at low temperatures should be considered separately. The effects of interaction turn out to be small in the TLS regime and are discussed in detail in the final Section of this paper. Here, we give several qualitative estimates for the sake of completeness, both at high and low temperatures. The phonon-mediated interaction goes roughly as $`\frac{g^2}{\rho s_c^2}\frac{1}{r^3}`$, with a numerical factor less than one (see Section VI). Ignoring the factor, the interaction is expressed, with the help of Eq.(18), via the glass transition temperature according to $`k_BT_g\frac{a^3}{r^3}`$. Two neighbouring domains, a distance $`\xi `$ apart, would thus couple with strength $`J_{\text{neigh}}=T_g(a/\xi )^3`$. In order to assess the effects of interaction on the effective energies of individual transitions, or whether it even makes sense to talk about on-site energies after the interaction is turned on, one must compare the interaction strength to the width of the distribution of the on-site energies as derived in the absence of interaction, exactly the same way it is done in the context of Anderson localization. According to the previous Subsection the latter width, call it $`\mathrm{\Delta }E`$, is of the order $`T_g`$. The ratio $`J_{\text{neigh}}/\mathrm{\Delta }E(a/\xi )^3`$ is a small number, as expected. There will be no long range effects, due to resonant interactions, at high temperatures near $`T_g`$. At very low temperatures, only tunneling centers (TC) with energy splitting $`k_BT`$ or less are thermally active. While the relevant spread of the on-site energies $`E_Tk_BT`$ is now down by a factor $`T/T_g`$ compared to the glass transition temperature, the concentration of active TC is also down by the same factor, namely $`(T/T_g\xi ^3)`$, thus increasing the mutual separation between the regions of mobile transitions. The total dipole-dipole induced static field due to all those thermally active two-level systems at a given spot is simply $`\frac{g^2}{\rho c_s^2}(T/T_g\xi ^3)k_BT(a/\xi )^3`$, again much smaller than the relevant on-site energy range $`E_Tk_BT`$. The motions within the tunneling centers are quantum-mechanical at these low temperatures, and so one may consider possible effects of resonant interactions between distinct TCโs, as in the Burin-Kagan scenario Burin and Kagan (1996). These effects have been shown to become important only at ultra-cold temperatures of $`\mu `$K and below Neu *et al.* (1997), as already mentioned in the introduction.
Besides having explained the origin of the universality of combination $`\overline{P}\frac{g^2}{\rho c_s^2}`$ ubiquitous in the STM, we have also seen why the value of this parameter is different from $`1`$, suggested by the strong defect-defect interaction universality scenario. This value can be traced to the relative concentration of the domains $`(a/\xi )^3<0.01`$, as just mentioned. It is curious that the defect-phonon interaction in the long wave-length limit can be expressed as a surface integral. Besides supporting our picture of the residual excitations as motions of domain walls, it points at a connection with string theories, where the elementary particles exhibit internal structure at high enough energies, which is also true in our case. In fact, this internal structure is ultimately the cause of the phenomena observed in glasses at higher temperatures, namely the so called bump in the specific heat and the thermal conductivity plateau, which are dealt with in Chapter IV.
### III.3 Distribution of Barriers and the Time Dependence of the Heat Capacity
The STM postulated tunneling matrix element distribution $`P(\mathrm{\Delta })1/\mathrm{\Delta }`$ implies a weakly (logarithmically) time dependent heat capacity. This was pointed out early on by Anderson *et al.* (1972), while the first specific estimate appeared soon afterwards in Jรคckle (1972). The heat capacity did indeed turn out time-dependent, however its experimental measures are indirect, and so a detailed comparison with theory is difficult. Reviews on the subject can be found in Nittke *et al.* (1998); Pohl (1981). Here, we discuss the $`\mathrm{\Delta }`$ distribution dictated by the present theory, in the semi-classical limit, and evaluate the resulting time dependence of the specific heat. While this limit is adequate at long times, quantum effects are important at short times (this concerns the heat condictivity as well). The latter are discussed in Subsection V.1.
In the tunneling argument from Section III.1, we have suggested a WKB type expression for the tunneling amplitude:
$$\mathrm{\Delta }=\mathrm{\Delta }_0e^{\frac{\pi V^{}}{\mathrm{}\omega ^{}}},$$
(21)
which would be completely correct in the case of a parabolic barrier with frequency $`\omega ^{}`$ and height $`V^{}`$ and was motivated by the necessity to maintain the proper scaling with $`\mathrm{}`$ in the denominator of the exponent, given that the typical barrier height is determined by the classical landscape characteristics and should scale with $`T_g`$. According to Eq.(14), $`P(V^{})e^{\frac{V^{}}{\sqrt{2}T_g}}`$. It follows then that
$$P(\mathrm{\Delta })d\mathrm{\Delta }=A\left(\frac{\mathrm{\Delta }_0}{\mathrm{\Delta }}\right)^c\frac{d\mathrm{\Delta }}{\mathrm{\Delta }},$$
(22)
where $`c=\mathrm{}\omega ^{}/\sqrt{2}T_g`$ should be less then $`0.1`$ according to our estimates of $`\omega ^{}`$ (see section IV.2). $`A`$ is a constant, to be commented on very shortly. The distribution in Eq.(22) becomes $`P_{STM}(\mathrm{\Delta })d\mathrm{\Delta }d\mathrm{\Delta }/\mathrm{\Delta }`$ postulated in the standard tunneling model, if $`c0`$. As shown next, the non-zero $`c`$ gives rise to an anomalous exponent $`\alpha =c/2`$ in the heat capacity $`CT^{1+\alpha }`$ and a power law $`t^{c/2}`$ for the specific heat time dependence at long times, as opposed to a logarithmic one, predicted by the STM. While $`c0.1`$ implies $`\alpha 0.05`$, experimentally, $`\alpha `$ seems to vary from $`0.1`$ to $`0.5`$. This larger value is consistent with quantum mixing effects that go beyond the semiclassical analysis as we will discuss later.
Scaling $`\mathrm{}\omega ^{}`$ in the denominator of the tunneling exponent implies that $`\omega ^{}`$ must be a quantum energy scale and it is indeed shown in Section IV.2 that $`\omega ^{}`$ is proportional to the Debye frequency $`\omega _D`$. While the tunneling argument from Section III.1 only explicitly considered the statistics of the highest energy state along the tunneling trajectory, the expression in Eq.(21) actually does not use such a simplified picture but considers a finite vicinity of the barrier top. The conclusion of Section IV.2 that the barrier heights are distributed exponentially, such as in in Eq.(14), remains true in either case. The leads to a non-zero value of $`c`$, and here we explore what consequences this has for the low temperature heat capacity and conductivity. As follows from the discussions in Section III.1, constant $`A`$ in Eq.(22) is of order one.
Since the temperatures in question here are so low ($`1K`$ and below), we will ignore the energy dependence of $`n(ฯต)`$ in this section and take $`n(ฯต)=\overline{P}`$. In order to see the time dependence of the heat capacity we obtain the combined distribution of the TLS energy splittings $`E`$ and relaxation rates $`\tau ^1`$ \- $`P(E,\tau ^1)`$, much as was done in Jรคckle (1972), - and then count in only those TLS whose relaxation time $`\tau `$ is shorter than a particular experimental observation time $`t`$.
The (phonon irradiation induced) relaxation rate of a TLS is Jรคckle (1972):
$$\tau ^1\frac{3g^2\mathrm{\Delta }^2E}{2\pi \rho c_s^5}\mathrm{coth}(\beta E/2).$$
(23)
It follows from Eqs. (22) and (23) that
$$P(E,\tau ^1)=\overline{P}A\left(\frac{\tau }{\tau _{min}(E)}\right)^{c/2}\left(\frac{\mathrm{\Delta }_0}{E}\right)^c\frac{\tau }{2\sqrt{1\tau _{min}(E)/\tau }},$$
(24)
where
$$\tau _{min}^1(E)\frac{3g^2E^3}{2\pi \rho c_s^5}\mathrm{coth}(\beta E/2)$$
(25)
is the fastest relaxation rate of a TLS with energy splitting $`E`$, achieved at $`\mathrm{\Delta }=E`$, of course. As follows from Eq.(25), the rate scales roughly as the cube of temperature and is empirically of order an inverse millisecond at $`10^2K`$. The resultant sampleโs heat capacity per unit volume is then:
$$C(t)=_0^{\mathrm{}}๐E\left(\frac{\beta E}{2\mathrm{cosh}(\beta E/2)}\right)^2_{t^1}^{\tau _{min}^1}๐\tau ^1P(E,\tau ^1),$$
(26)
where $`[\beta E/2\mathrm{cosh}(\beta E/2)]^2`$ is the TLS heat capacity. With a change of variables, Eq.(26) reads:
$$C(t)=_0^{\mathrm{}}๐E\left(\frac{\beta E}{2\mathrm{cosh}(\beta E/2)}\right)^2_0^{\mathrm{log}(t/\tau _{min}(E))}๐z\frac{A}{2}\left(\frac{\mathrm{\Delta }_0}{E}\right)^c\frac{e^{\frac{c}{2}z}}{\sqrt{1e^z}}.$$
(27)
At long times the expression Eq.(27) yields a power law for both time and temperature dependence of the specific heat:
$$\underset{t\mathrm{}}{lim}C(t)t^{c/2}T^{1+c/2},$$
(28)
where, note, the temperature dependence also comes from the energy dependence of $`\tau _{min}(E)E^3`$ in Eq.(27). The value $`c0.1`$ implies the long time heat capacity should obey $`CT^{1+\alpha }`$ at low $`T`$, where $`\alpha 0.05`$, a smaller number than observed in amorphous materials. We must bear in mind that the issue of the exact form of the time dependence in the laboratory still appears to be unresolved, as there is no definite agreement between different experiments; for references, see Pohl (1981); Hunklinger and Raychaudhuri (1986); Nittke *et al.* (1998); Sahling *et al.* (2002). While there is no doubt that the specific heat is time-dependent, some experiments agree with the logarithmic time profile, as predicted by the STM, others give a lower or higher speed of variation with time. The present semiclassical prediction with $`c=0.1`$ would be hard to distinguish from a logarithmic law. Finally, even though a correction to the linear temperature heat capacity dependence with $`c=0.1`$ is most likely smaller than what is experimentally observed, the value of this correction is non-universal, consistent with empirical data.
The expression in Eq.(27) can be evaluated numerically for all values of $`t`$, and the results for three different waiting times are shown in Fig.11 for $`c=0.1`$. The value of $`\tau _{min}=2.0\mu sec`$ at $`E/T_D=5.710^4`$, derived from the present theory (also consistent with Goubau and Tait (1975)) was used. The results for $`t=10\mu sec`$ demonstrate that due to a lack of fast relaxing systems at low energies, short time specific heat measurements can exhibit an apparent gap in the TLS spectrum. Otherwise, it is evident that the power-law asymptotics from Eq.(28) describes well Eq.(27) at the temperatures of a typical experiment.
As clear from the discussion above, the long-time power law behavior of the heat capacity is determined by the โslowโ two-level systems corresponding to the higher barrier end of the tunneling amplitude distribution, argued to be of the form shown in Eq.(22). If one assumes that this distribution is valid for the zero-barrier tail of the $`\mathrm{log}(\mathrm{\Delta })`$ distribution as well, one would expect that the heat conductivity would scale as $`T^{2+c}`$ at the TLS temperatures, in contrast to an observed experimental sub-quadratic dependence $`T^{2\alpha ^{}}`$. As we shall see in Section V, other quantum effects are indeed present in the theory and we will discuss then how these contribute both to the deviation of the conductivity from the $`T^2`$ law and the way the heat capacity differs from the strict linear dependence, both contributions being in the direction observed in experiment. Finally, when there is significant time dependence of $`c_V`$, the kinematics of the thermal conductivity experiments are more complex and in need of attention. When the time-dependent effects are included, both phonons and two-level systems should ideally be treated by coupled kinetic equations. Such kinetic analysis, in the context of the time dependent heat capacity, has been conducted before by other workers Strehlow and Meissner (1999).
## IV The Plateau in Thermal Conductivity and the Boson Peak
In this section we continue to explore the consequences of the existence of the low temperature excitations in amorphous substances, which, as argued in Chapter III, are really resonances that arise from residual molecular motions otherwise representative of the molecular rearrangements in the material at the temperature of vitrification. We were able to see why these degrees of freedom should exist in glasses and explain their number density and the nearly flat energy spectrum, as well as the universal nature of phonon scattering off these excitations at low T ($`<1`$ K).
At higher temperatures ($`K_BT`$10<sup>-2.0</sup> to 10$`{}_{}{}^{0.5}\mathrm{}\omega _D`$), an apparently different kind of excitations begins to appear, leading to the so called bump in the heat capacity and plateau in the thermal conductivity, as was discussed in the Introduction. We argue in this chapter that the transitions between the mutually accessible frozen-in minima in the amorphous lattice that give rise to the two level system behavior at the lowest $`T`$ also explain the existence of the modes responsible for this โBoson Peakโ and the intense phonon scattering at the corresponding frequencies. This thus removes the need to invoke theoretically any additional mechanisms, although other contributions may well be present to some extent; we will try to assess this possibility in the following Section.
### IV.1 Introduction: Classification of Excitations in Glasses
While we believe to have mostly achieved a microscopic understanding of the excitations that are specific to the amorphous solids and are not present in other types of materials, this description is rather new and, naturally, there is certain lack of established language that could be efficiently used to characterize these excitations. In this section, we will introduce some terminology that will be used in the rest of the article. At the same time, we will provide a brief general analysis of what possible qualitatively distinct types of molecular motions can exist in glasses.
Any atomic motions that take place in a frozen glass, obviously may also be present in the liquid above $`T_g`$. For example, those high $`T`$ motions that correspond to shear attain stiffness (on realistic time scales) below freezing. The motions in the liquid, apart from pure volume change, corresponding to the longitudinal sound, are molecular translations, or, informally speaking, jumps. Above $`T_A`$, such jumps are not accompanied by a noticeable volume change and bond stretching, as no metastable structures form in the liquid at these temperatures. The barriers are therefore largely entropic. (It is nice to compare Feynmanโs discussion of the absence of energy barriers in superfluid He Feynman (1954) in this regard.) Below $`T_A`$, such hopping already involves moving a number of molecules from one local minimum of the free energy functional to another such minimum and thus requires structural rearrangement within a certain cooperative length owing to the formation of metastable local arrangements. Molecular translations do not conserve momentum, which subsequently must be provided by the rest of the bulk. We thus call these degrees of freedom, which are relics of the translational motions in the liquid above $`T_g`$, inelastic degrees of freedom. They are truly inelastic also in the macroscopic sense of the word, because the existence of alternative configurations in the solid bulk, which are also coupled to the phonons, ultimately leads to irreversible relaxation, if the sample is subjected to mechanical stress thus causing a shift in the thermal population of the alternative internal, structural states. This is the mechanism behind the so called bulk viscosity Landau and Lifshitz (1986) (incidentally, it also contributes to the so called relaxational phonon absorption, which we discuss in subsection IV.7). The switching from one energy minimum to another is accomplished by moving the domain wall - the interface between the two alternative configurations - through the local region. As mentioned earlier, this domain wall is something of an abstract entity, really a quasiparticle of a sort. Yet it has many ponderable attributes. For one thing, it has a mass (per unit area), which will be obtained in this section. It also has surface tension, therefore it can support surface vibrations, again, of a sort. Although these vibrations are realized through real atomic motions, it is more beneficial to think of them as vibrational modes of an imaginary membrane. In fact, as will be argued later, the oscillations of this membrane correspond to the indeterminacy in the exact boundary of the frozen-in domain that has more than one kinetically accessible internal state. Therefore, highly anharmonic atomic motions in the real space correspond to harmonic motions in the space where the domain walls are defined. This mental construction does the trick of enabling us to calculate the ripplon spectrum, as demonstrated in section IV.3. Now, since it was shown in Xia and Wolynes (2000), that the liquid degrees of freedom below $`T_A`$ consist of switching to alternative local energy minima; we can claim our assignment of different inelastic modes is exhaustive (but not unique, of course!). These are, again, translations and vibrations of the domain walls.
On the other hand, any purely elastic motions in the glassy lattice can be thought of as a sum of ordinary, affine, displacements and non-affine displacements (see e.g. Wittmer *et al.* (2001)). The affine component would be the only one present in a perfectly isotropic medium and would follow the stress pattern according to a Poisson equation (the situation with a non-isotropic crystal is conceptually the same). The non-affine displacements are a consequence of the absence of periodicity. They involve a small number of molecules and are characterized by a non-zero circulation of the displacement field. It is not clear at present whether the size of these non-affine โislandsโ could be inferred in present day computer simulations, since the amorphous structures that can be currently generated on modern computers still correspond to unrealistically rapid quenching rates. The resulting structure corresponds to a sample caught in a very high energy state with extraordinarily low barriers. As is clear by now from the random first order transition theory, such structures correspond to temperatures close to $`T_A`$ and will have very small cooperative regions approaching the molecular scale $`a`$.
We conclude this subsection by repeating ourselves that one important difference between the elastic and the inelastic modes is in how they absorb the phonons. While any static disorder can only provide Rayleigh scattering with a characteristic length scale equal to the size of the heterogeneity, the inelastic (resonant) absorptionโs cross-section scales as the square of the phonon wave-length, it thus will considerably dominate the Rayleigh mechanism for the longer wave-length phonons (absorption saturation in the TLSโs does not occur at the sound intensities typical of heat transfer).
### IV.2 The Multilevel Character of the Entropic Droplet Excitations
We hope to have convinced the reader by now that the tunneling centers in glasses are complicated objects that would have to be described using an enormously big Hilbert space, currently beyond our computational capacity. This multilevel character can be anticipated coming from the low temperature perspective in Lubchenko and Wolynes (2001). Indeed, if a defect has at least two alternative states between which it can tunnel, this system is at least as complex as a double well potential \- clearly a multilevel system, reducing to a TLS at the lowest temperatures. Deviations from a simple two-level behavior have been seen directly in single-molecule experiments Boiron *et al.* (1999). In order to predict the energies at which this multilevel behavior would be exhibited we first estimate the domain wall mass. Obviously, the total mass of all the atoms in the droplet is so large that the possibility of simultaneous tunneling of all atoms is completely excluded. The tunneling, we argue, occurs stagewise; each individual motion encounters a nearly flat potential, implying low frequency instantaneous modes.
In addition, the effective mass of the domain wall turns out to be low, also owing to the collective, barrierless character of the tunneling events. This is because moving a domain wall over a molecular distance $`a`$ involves displacing, at any one (imaginary!) time, individual atoms only a Lindemann length $`d_L`$. Suppose this occurs on the (imaginary) time scale $`\tau `$. The resulting kinetic energy is $`M_w(a/\tau )^2=N_wm(d_L/\tau )^2`$, where $`N_w(\xi /a)^2`$ is the number of molecules in the wall and $`m`$ is the molecular mass. Thus the mass of the wall $`M_w`$ is only $`m(\xi /a)^2(d_L/a)^2`$. Using $`(\xi /a)5.8`$ and $`(d_L/a)^20.01`$ gives $`M_wm/3`$. This implies the mass of the wall per atom is very small - about a hundredth of a molecular mass, consistent with the simulations of certain barrierless dislocation motions in copper Vegge *et al.* (2001). Using $`(d_L/a)^2k_BT_g/\rho c_s^2`$, derived earlier, one can express the wallโs mass through the material constants as $`M_w(\xi /a)^2k_BT_g/c_s^2`$. The wall mass estimate above, inspired by the Feynmanโs argument on the effective mass in liquid helium Feynman (1953), is entirely analogous to the well known estimate of the soliton mass in polyacetylene, see e.g. Heeger *et al.* (1988). In the latter, the soliton moves a large distance, while individual atoms undergo only small displacements leading to a low soliton mass.
We can now use the typical value of the barrier curvature from our tunneling argument in section III.1 (see Fig.10) to estimate the typical frequency $`\omega ^{}`$ of motion at the tunneling barrier top. We now express the barrier profile $`V(N)`$ as a function of the dropletโs radius $`ra(3N/4\pi )^{1/3}`$ and obtain
$$\omega ^{}=^2V/r^2/M_w1.6(a/\xi )\omega _D.$$
(29)
According to the quantum transition state theory Wolynes (1981), and ignoring damping, at a temperature $`T^{}\mathrm{}\omega ^{}/2\pi k_B(a/\xi )T_D/2\pi `$, the wall motion will typically be classically activated. This temperature lies within the plateau in thermal conductivity Freeman and Anderson (1986). This estimate will be lowered if damping, which becomes considerable also at these temperatures, is included in the treatment. Indeed, as shown later in this section, interaction with phonons results in the usual phenomena of frequency shift and level broadening in an internal resonance. Also, activated motion necessarily implies that the system is multi-level. While a complete characterization of all the states does not seem realistic at present, we can extract at least the spectrum of their important subset, namely those that correspond to the vibrational excitations of the mosaic, whose spectral/spatial density will turn out to be sufficiently high to account for the existence of the Boson Peak.
### IV.3 The Vibrational Spectrum of the Domain Wall Mosaic and the Boson Peak
At low temperatures the two-level system excitations involve tunneling of the mosaic cells typically containing $`N^{}200`$ atoms. The tunneling path involves stagewise motion of the wall separating the distinct alternative configurations through the cell untill a near resonant state is found. At higher temperatures, other final states are possible since the exact number and identity of the atoms that tunnel can vary. These new configurations typically will be like the near resonant level but will also move a few atoms at the boundary, i.e. at the interface to another domain. This is schematically shown in Fig.12.
Alternatively, due to the quantum mechanical uncertainty of the exact location of the domain wall, its shape is intrinsically subject to fluctuations (these are zero-point vibrations of the domain wall). It is thus not surprising that the ripplonโs frequencies turn out to be proportional to $`\omega _D`$, the basic quantum energy scale in the system. These fluctuations of the domain boundary shape can be visualized as domain wall surface modes (โripplonsโ). A detailed calculation of the ripplon spectrum would require a considerable knowledge of the mosaicโs geometry. At each temperature below $`T_A`$ the domain wall foam is an equilibrium structure made up of flat patches of no tension (remember the renormalized $`\sigma (r)r^{1/2}`$; however fluctuations will give rise to finite curvature and tension). To approximate the spectrum we notice that the ripples of wave-length larger than the size of a patch will typically sense a roughly spherical surface of radius $`R=\xi (3/4\pi )^{1/3}`$. The surface tension of the mosaic has been calculated from the classical microscopic theory and is given by $`\sigma (R)=\frac{3}{4}(k_BT_g/a^2)\mathrm{log}((a/d_L)^2/\pi e)(a/R)^{1/2}`$ Xia and Wolynes (2000), where $`d_L/a`$ is the universal Lindemann ratio. It could appear that such tension could collapse an individual fragment of the mosaic but this tension is, of course, compensated by stretching the frozen-in outside walls. We approximate the effect of this compensation by an isotropic positive pressure of a ghost (i.e. vanishing density) gas on the inside.
The eigen-frequency spectrum of the surface modes of a hollow sphere with gas inside is well known (see e.g. Morse and Feshbach (1953), as well as our appendix A). If we pretend for a moment that the surface tension coefficient $`\sigma `$ is curvature independent, the possible values of the eigen-frequency $`\omega `$ are found by solving the following equation:
$$\mathrm{cot}[\alpha _l(\omega R/c_g)]=\left(\frac{\rho _W}{\rho _gR}\right)\frac{(l1)(l+2)}{(\omega R)^2}\left(\frac{\sigma }{\rho _gR}\right),$$
(30)
where $`\rho _W`$ is the membraneโs mass per unit area, $`\rho _g`$ and $`c_g`$ are gasโ mass density and sound speed respectively. As stated earlier, Eq.(30) only applies for $`l2`$. Finally, function $`\mathrm{cot}[\alpha _l(z)]\left(l+z\frac{j_{l+1}(z)}{j_l(z)}\right)^1`$, where $`j_l(z)`$ is the spherical Bessel function of $`l`$-th order, does exhibit behavior similar to that of the regular trigonometric cotangent for arguments of the order unity and larger, going however to $`1/l`$ as $`z0`$. Its graph for $`l=2`$ is shown in Fig.13.
An inspection shows that for each $`l`$, the smallest solution of Eq.(30) gives the frequency of the proper eigen-mode of the shell itself (shifted due to the presence of the gas inside), whereas the rest of the solutions represent the standing acoustic waves in the gas. This is especially clear in the $`\rho _g0`$ limit, when the lowest frequency does not even depend on the gasโ sound speed, whereas the rest of the solutions are obviously determined by the inverse time it takes the sound in the gas to traverse the sphere.
Since we are interested only in the wallโs proper modes in the limit $`\rho _g0`$, we get unambiguously for the frequency of an $`l`$-th harmonic:
$$\omega _l^2=(l1)(l+2)\left(\frac{\sigma }{\rho _WR^2}\right);(l2).$$
(31)
Accounting for the unusual $`r`$ dependence of the surface tension $`\sigma (r)r^{1/2}`$ modifies the standard result from Eq.(31) by a factor of $`9/8`$. The reason is, the peculiar surface energy dependence $`F_{surf}(R)=4\pi R^2\sigma =4\pi \sigma _0R^{3/2}a^{1/2}`$ calls for the following dependence of pressure on the curvature: $`p=\frac{1}{4\pi R^2}\frac{F_{surf}(R)}{R}=\frac{3}{2}\frac{\sigma }{R}`$ (as compared to the regular $`p=2\frac{\sigma }{R}`$). The eigen-frequencies, in their turn, are determined by calculating the (frequency dependent) excess pressure due to a variation in curvature. Since now $`pR^{3/2}`$, varying $`p`$ with respect to $`R`$ brings down another factor of $`3/2`$, thus giving $`9/4`$ instead of the $`2`$ of the curvature independent case. Hence the (barely significant, but curious) correction factor of $`9/8`$ used in Lubchenko and Wolynes (2003a). Since we have been assuming that the amplitude is infinitesimally small, this factor is the only consequence of having a curvature dependent $`\sigma `$, which should have made the membrane oscillations even more non-linear (as compared to $`\sigma =`$ const case) in the case of finite displacements. Pinpointing this effect, however, is clearly beyond the accuracy attempted by the present model. Finally, one finds a spectrum with
$$\omega _l^2=\frac{9}{8}\frac{\sigma }{\rho _WR^2}(l1)(l+2);(l2),$$
(32)
where each $`l`$-th mode of a sphere is $`(2l+1)`$-fold degenerate. Using $`\rho _W=(d_L/a)^2\rho a`$, obtained earlier in the chapter and $`T_g\rho c_s^2a^3(d_L/a)^2`$ (section III.2), one finds
$`\omega _l`$ $`1.34`$ $`\omega _D(a/\xi )^{5/4}\sqrt{(l1)(l+2)/4}`$ (33)
$``$ $`0.15\omega _D\sqrt{(l1)(l+2)/4}.`$
Because of the universality of the $`(a/\xi )`$ ratio Lubchenko and Wolynes (2001), $`\omega _l`$ is a multiple of the Debye frequency. Apart from the barely significant $`(a/\xi )^{1/4}`$ factor, again, due to the $`R`$ dependent $`\sigma `$, the ubiquitous scaling $`\omega _l(a/\xi )\omega _D`$ stresses yet another time the significance of the scale $`\xi `$. Such a scale has been previously empirically deduced by interpreting inelastic scattering experiments but has been usually ascribed to the static heterogeneity length scale, in contrast with the dynamical nature of the mosaic in the present theory. We note, again, that this โstatic heterogeneityโ has never been unambiguously seen in X-ray diffraction. Owing to the materialโs discreteness, one does not expect harmonics of higher than $`\pi \left(\frac{3}{4\pi }N^{}\right)^{1/3}[(Ra/2)/R]\mathrm{9..10}`$th order, a relatively large number, which justifies the tacitly assumed continuum approximation. The lowest allowed ripplon mode is $`l=2`$ (corresponding frequency is $`1`$THz for silica, in remarkable agreement with the inelastic neutron scatering data Wischnewski *et al.* (1998)).
The requirement $`l2`$ can be understood from the symmetry considerations. $`l=1`$, the case of no restoring force, corresponds to a domain translation. Within our picture, this mode corresponds to the tunneling transition itself. The โtranslationโ of the defects center of mass violates momentum conservation and must be thus accompanied by absorbing a phonon. Such resonant processes couple linearly to the lattice strain and contribute the most to the phonon absorption at the low temperatures, dominated by one-phonon processes. $`l=0`$, on the other hand, corresponds to a uniform dilation of the shell. This mode is formally related to the domain growth at $`T>T_g`$, and is described by the theory in Xia and Wolynes (2000). It is thus possible, in principle, to interpret our formalism as a multi-pole expansion of the interaction of the domain with the rest of the sample. Harmonics with $`l2`$ correspond to pure shape modulations of the membrane.
The existence of the domain wall vibrations explains and allows us to visualize, at least in part, the multilevel character of the tunneling centers as exhibited at temperatures above the TLS regime. Curiously, the existence of TLSโs, even though displayed at the lowest $`T`$, is basically of classical origin due to the non-equilibrium nature of the glassy state. Yet the ripplons, even though seen at higher $`T`$, are mostly due to quantum effects and would not be predicted by a strictly semi-classical theory, in which $`\mathrm{}0`$. A schematic of the resultant droplet energy levels is shown in Fig.14.
The arrangement of the combined internal (configurational) and ripplonic density of states, as depicted in Fig.14, has the following motivation behind it. We include the possibility of distorting the domain wall during the tunneling transition by providing a set of vibrational states on top of the alternative internal state. The arrangement of the energy states as depicted in Fig.14 insures that only thermally active tunneling centers have mobile (and thus vibrationally excitable boundaries). The atomic motions at the inactive defectsโ sites (i.e. that cannot tunnel or cross the barrier) would be indistinguishable from the regular elastic lattice vibrations. Importantly, direct transitions to the ripplonic state can occur from one of the two lowest - โTLSโ - energy states of a tunneling center. This inherent assymetry between the two structural states of a tunneling center actually reflects the thermodynamical inequivalence of the two states at the glass transition temperature. While one of the states represents the local structure in (meta-stable) equilibrium with the current liquid arrangement around it, the other state is a configuration that must only be regarded as one of the structures along the many escape routes from the current equilibrium local state. At $`T_g`$, most of those escape routes become too costly energetically.
This is a good place to remind the reader that existing explanations of the large density of states at the BP energies have to do either with purely harmonic excitations of disordered, but perfectly stable lattice (see Introduction), or by generalizing the low energy inelastic, two-level degrees of freedom to multilevel systems, as was done e.g. by the soft potential model (SPM) Karpov *et al.* (1983); Buchenau *et al.* (1992). Such generalizations imply a connection between the anomalies seen in the TLS regime and at these higher energies. Such a connection is strongly suggested by experiment, most prominently by the strength of phonon scattering. The latter is inelastic at the BP energies, as it was at the TLS energies. We stress, the rate of increase of the ripplonic density of states is much much higher than that empirically assumed in the purely empirical SPM. Again, there is virtually no freedom to adjust the numbers in our theory.
In order to compute the heat capacity of the ripplons on top of the structural transitions we will need to consider the (classical) density of the inelastic states in more detail than in the previous section. The density of states $`n(ฯต)=\frac{1}{T_g}e^{ฯต/T_g}`$ was derived earlier taking as the reference state the generic global liquid state corresponding to the (high-energy) configuration frozen-in at $`T_g`$. It turns out that only transitions to states with $`ฯต<0`$ (relative to the liquid state!) will contribute to the TLS density of states. Indeed, as we have shown, the size of the region that permits a low-barrier rearrangement must be slightly (by $``$18 molecules) larger than the generic cooperative size at $`T_g`$. On the other hand, we know from the RFOT theory that larger cooperative regions correspond to lower energy liquid structures. Therefore one of the two alternative states must be lower in energy than the generic liquid state at $`T_g`$. As a result, the negative $`ฯต`$โs correspond to some of the very numerous but mostly unavailable lower lying energy states, now accessed by tunneling. Now, if each of those true local ground states is taken as the reference one, the spectral density will be now $`n(ฯต)=\frac{1}{T_g}e^{ฯต/T_g}`$ ($`ฯต>0`$). We consequently can let $`ฯต`$ from Fig.14 take both positive and negative values by writing
$$n(ฯต)=\frac{1}{T_g}e^{|ฯต|/T_g}.$$
(34)
We can now calculate each domainโs partition function by including all possible ways to excite the system:
$$Z_ฯต=1+\underset{\{n_{lm}\}}{}e^{\beta (ฯต+_{lm}n_{lm}\omega _{lm})}=1+e^{\beta ฯต}\underset{l}{}Z_l^{2l+1},$$
(35)
where $`Z_l1/(1e^{\beta \omega _l})`$ is the partition function of an $`l`$th order ripplon mode and we used $`m=l\mathrm{..1}`$. Here we assume each ripplon is a harmonic oscillator. Note that since the โharmonicโ excitations of frequency $`\omega _l`$ are on top of another (structural) excitation, we must consider the issue of the zero-point energy of these โharmonicโ excitations, that is no longer a matter of simply choosing a convenient reference energy. Note that this zero-point energy is actually several orders of magnitude higher than the subKelvin energies that are sufficient to excite some of the local structural transitions. And indeed, the energy that comprises the ripplonsโ ground state energy is not extracted from the thermal fluctuations of the medium, but, one may say, is simply โconvertedโ from the zero-point energy of local elastic vibrations of the lattice. At the site of a โslowโ (or, thermally inactive) structural transition, domain wall vibrations are indistinguishable from the regular lattice phonons, as already mentioned.
The specific heat corresponding to the partition function in Eq.(35) is found by computing $`c_ฯต=\beta ^2\frac{^2logZ_ฯต}{^2\beta }`$:
$$c_ฯต=\frac{\left[\beta ฯต+_l(2l+1)\frac{\beta \omega _l}{e^{\beta \omega _l}1}\right]^2}{\left[2\mathrm{cosh}\frac{\beta ฯต+_l(2l+1)\mathrm{log}(1e^{\beta \omega _l})}{2}\right]^2}+\frac{_l(2l+1)\left(\frac{\beta \omega _l}{2\mathrm{sinh}\frac{\beta \omega _l}{2}}\right)^2}{e^{\beta ฯต+_l(2l+1)\mathrm{log}(1e^{\beta \omega _l})}+1}.$$
(36)
Expression (36) clearly becomes the TLS specific heat $`c_{TLS}=\left(\frac{\beta ฯต}{2\mathrm{cosh}(\beta ฯต/2)}\right)^2`$ for $`T\omega _l`$.
In order to obtain the amorphous heat capacity per domain, we (numerically) average $`c_ฯต`$ with respect to $`n(ฯต)`$; the result is shown in figure 15 with the thin solid line.
### IV.4 The Density of Scatterers and the Plateau
In order to estimate the phonon scattering strength and thus the heat conductivity, we need to know the effective scattering density of states, the transition amplitudes and the coupling of these transitions to the phonons.
Any transition in the domain accompanied by a change in its internal state is coupled to the gradient of the elastic field with energy $`g\rho c_s^2๐๐ฌ๐(๐ซ)`$, where $`๐(๐ซ)`$ is the molecular displacement at the droplet edge due to the transition (see section III.2). An additional modulation in the domain wall shape due to the current vibrational state cancels out due to the high symmetry ($`l2`$), as easily seen when computing the angular part of the surface integral. We therefore conclude that any transitions between groups marked with solid lines in Fig.14 are coupled to the phonons with the same strength as the underlying (TLS-like) transition. (Notice this also implies inelastic scattering off those transitions!) Incidentally, no selection rules apply for the change in the ripplon quantum numbers, being essentially a consequence of strong anharmonicity of the total transition.
We do not possess detailed information on the transition amplitudes, however they should be on the order of the transition frequencies themselves, just as is the case for those two level systems that are primarily responsible for the phonon absorption at the lower $`T`$ which also have their transition amplitudes comparable to the total energy splitting. The argument is thus essentially the same as proposed earlier in section III.1. It should be noted, however, that the Hilbert spaces corresponding to the quantum in nature ripplons and the classical inelastic states are quite distinct (although overlapping); it thus should not be surprising that the matrix element beween superpositions of these spaces is on the order of the energy differences themselves. In what follows, we circumvent to an extent the question of what the precise distribution of the tunneling amplitudes of the TLS+ripplon transitions is and simply calculate the enhancement of the bare TLS induced scattering due to the presence of the ripplons. This is suggested by an earlier notion that the structural transitions in glasses couple to the phonons with the same strength even if accompanied by exciting vibrational modes of the mosaic.
We now calculate the density of the phonon scattering states. Since we have effectively isolated the transition amplitide issue, the fact of equally strong coupling of all transitions to the lattice means that the scattering density should directly follow from the partition function of a domain via the inverse Laplace transform. We will not proceed this way for purely technical reasons. In addition, we will separate the cases of positive and negative $`ฯต`$ (see Fig.14), corresponding to absorption from ground and excited states respectively.
The phonon-ripplon interaction exhibits itself most explicitly through the phonon scattering, which becomes so strong by the end of the plateau as to cause complete phonon localization. This interaction also results in other observable consequences, such as dispersion (or frequency shift) of the ripplon frequencies, as well as rendering the resonances finite width. Furthermore, we will argue, this interaction suffices to account for the non-universality of the plateau. First, however, we consider a simpler situation, where we assume the ripplon spectrum itself is unaffected by coupling to the phonons.
### IV.5 Phonon scattering off frictionless ripplons
If $`ฯต>0`$, the phonon absorbing transition occurs from the ground state. The total number of ways to admit energy $`\omega `$ into the system is
$$\rho (\omega )=_0^{\mathrm{}}๐ฯตn(ฯต)\underset{\{n_{lm}\}}{}\delta (\omega [ฯต+\underset{lm}{}n_{lm}\omega _{lm}])=1/T_g\underset{\{n_{lm}\}}{}\theta (\omega \underset{lm}{}n_{lm}\omega _{lm})e^{\beta _g(\omega _{lm}n_{lm}\omega _{lm})},$$
(37)
where we sum over all occupation numbers of the ripplons with quantum numbers $`l,m`$ ($`m=l..l`$). Using an integral representation of the step function $`\theta `$, this can be rewritten as
$$\rho (\omega )=\frac{1}{T_g}\underset{\{n_{lm}\}}{}\underset{ฯต_10^+}{lim}_{\mathrm{}}^{\mathrm{}}\frac{dk}{2\pi (ik+ฯต_1)}e^{ik(\omega _{lm}n_{lm}\omega _{lm})}e^{\beta _g(\omega _{lm}n_{lm}\omega _{lm})}.$$
(38)
The integral in Eq.(38) will be taken by the steepest descent method (SDM). The reason why we do not apply an analogous technique directly to the $`\delta `$-function in Eq.(37) is not only because we want to get rid of the $`ฯต`$ integration, but also because that the SDM proved more forgiving in terms of accuracy when used to approximate the $`\theta `$-function, rather than the $`\delta `$-function.
For each $`k`$ on the real axis, the sum over the occupation numbers $`n_{lm}`$ diverges, so each integral should be taken before the summation. However, in the vicinity of the point that will turn out to be the saddle point $`k_0`$ ($`\mathrm{}k_0<\beta _g`$) all the sums are finite, so we reverse the order of summation and integration. The integration contour is shifted as shown in see Fig.16.
Hence, the saddle-point approximation yields for the value of the integral in Eq.(38) ($`\kappa ik`$):
$$\rho (\omega )=\frac{1}{T_g}\frac{1}{\sqrt{2\pi |f^{\prime \prime }(\kappa _0)|}}\frac{1}{\kappa _0}\mathrm{exp}\{(\kappa _0\beta _g)\omega \underset{lm}{}\mathrm{log}[1e^{(\kappa _0\beta _g)\omega _{lm}}]\},$$
(39)
where the saddle point $`\kappa _0`$ is determined from
$$\omega =\underset{lm}{}\frac{\omega _{lm}}{e^{(\kappa _0\beta _g)\omega _{lm}}1}+\frac{1}{\kappa _0}$$
(40)
and the curvature at the saddle point is equal to
$$|f^{\prime \prime }(\kappa _0)|=\underset{lm}{}\frac{\omega _{lm}^2}{4\mathrm{sinh}^2[(\kappa _0\beta _g)\omega _{lm}/2]}+\frac{1}{\kappa _0^2}.$$
(41)
As is clear from (40), the approximation amounts to finding the effective temperature so as to populate the ripplonic states to match the excitation energy $`\omega `$. The expression for the curvature (41) appropriately involves the corresponding heat capacity of the excitations.
The $`\omega 0`$ and the barely relevant $`\omega \mathrm{}`$ asymptotics are easily found. As luck has it, the $`\omega 0`$ limit of Eq.(39), apart from $`1/T_g`$ factor, gives $`\mathrm{exp}(1)/\sqrt{2\pi }1.08`$, only 8% away from the correct $`1`$. The $`\omega \mathrm{}`$ yields, on the other hand, $`\rho (\omega )_{lm}(\omega /\omega _{lm})\omega ^{96}`$, as expected ($`_{l=2}^9(2l+1)=96`$). The SDM is thus reasonably accurate in this case, which could be at least somewhat evaluated by computing the value of the fourth order term under the exponent at โone sigmaโ distance from the extremal action point. This turns out to be satisfactorily small, as demonstrated in Fig.17, along with the density of states itself as a function of $`\omega `$.
When estimating absorption from the ground state, we totally ignore the depletion of ground state population at finite temperatures, when the system spends some time in an excited state. This is fine because by the relevant temperatures, the excited state absorption dominates anyway (see Fig.14 and note that $`\omega _l|ฯต|<\omega _l+|ฯต|`$). This case, i.e. $`ฯต<0`$, is somewhat less straightforward. Let us calculate
$$N_E(\omega )_0^E๐ฯตn(ฯต)\underset{\{n_{lm}\}}{}\delta (\omega [\underset{lm}{}n_{lm}\omega _{lm}ฯต]).$$
(42)
This expression gives the cumulative density of absorbing states between energies $`0`$ and $`E`$ (note the change of sign in front of $`ฯต`$). This expression can be used to estimate the total excited state absorption by computing
$$\rho _{exc}(\omega ,T)_0^{\mathrm{}}๐Ef(E,T)\frac{N_E(\omega )}{E},$$
(43)
where $`f(E,T)2/(e^{\beta E}+1)`$ gives the appropriate Boltzmann weights. The factor of $`2`$ is used in order to calibrate the excited state absorption relative to the ground state case: $`f(0,T)=1`$. We now have two $`\theta `$-functions and consequently two integrations. The SDM value for $`N_E(\omega )`$ is given by
$$N_E(\omega )=\frac{1}{T_g}\frac{1}{2\pi |\text{โDetโ}|^{1/2}}\frac{1}{\lambda _0\mu _0}\mathrm{exp}\{(\beta _g+\lambda _0\mu _0)\omega +\lambda _0E\underset{lm}{}\mathrm{log}[1e^{(\beta _g+\lambda _0\mu _0)\omega _{lm}}]\}.$$
(44)
The corresponding saddle points are determined from
$$\omega +E=\underset{lm}{}\frac{\omega _{lm}}{e^{(\beta _g+\lambda _0\mu _0)\omega _{lm}}1}+\frac{1}{\lambda _0}$$
(45)
and
$$\omega =\underset{lm}{}\frac{\omega _{lm}}{e^{(\beta _g+\lambda _0\mu _0)\omega _{lm}}1}\frac{1}{\mu _0}.$$
(46)
Here,
$$|\text{โDetโ}|\underset{lm}{}\frac{\omega _{lm}^2}{4\mathrm{sinh}^2[(\beta _g+\lambda _0\mu _0)\omega _{lm}/2]}\left(\frac{1}{\lambda _0^2}+\frac{1}{\mu _0^2}\right)+\frac{1}{\lambda _0^2\mu _0^2}$$
(47)
is the determinant of the curvature tensor in the direction (i.e. 2D subset) of the fastest descent in the 4-dimensional (complex) $`\lambda ,\mu `$ space. The steepest descent approximation turns out to perform well, except at very low frequencies ($`\omega <10^2\omega _D`$). However, even though it overestimates the answer, it is still very small compared to the $`\rho (\omega )`$ calculated earlier at these frequencies, much as the complete result would be. The appropriate graph is shown in Fig.18.
An accurate calculation of the heat conductivity requires solving a kinetic equation for the phonons coupled with the multilevel systems, which would account for thermal saturation effects etc. We encountered one example of such saturation in the expression (19) for the scattering strength by a two-level system, where the factor of $`\mathrm{tanh}(\beta \omega /2)`$ reflected the difference between thermal populations of the two states. Neglecting these effects should lead to an error of the order unity for the thermal frequencies. Within this single relaxation time approximation for each phonon frequency, the Fermi golden rule yields for the scattering rate of a phonon with $`\mathrm{}\omega k_BT`$:
$$\tau _\omega ^1\omega \frac{\pi g^2}{\rho c_s^2}[\rho (\omega )+\rho _{exc}(\omega ,T)].$$
(48)
The heat conductivity then equals $`\kappa =\frac{1}{3}_\omega l_{\text{mfp}}(\omega )C_\omega c_s`$. The mean free path cannot be less than the phononโs wave-length $`\lambda `$ (which occurs at the Ioffe-Riegel condition). Since our theory does not cover the phonon localization regime we account for multiple scattering effects by simply putting $`l_{\text{mfp}}=c_s\tau _\omega +\lambda `$. At high $`T`$, the heat is not carried by โballisticโ phonons, but rather is transfered by a random walk from site to site, as originally anticipated by Einstein Einstein (1911) for homogeneous isotropic solids. The resultant heat conductivity is shown in Fig.19
We postpone further discussion of the results above until we include the effects of coupling of the resonant transitions to the phonons on the transitionsโ spectrum.
### IV.6 The effects of friction and dispersion
A transition linearly coupled to the phonon field gradient will experience, from the perturbation theory perspective, a frequency shift and a drag force owing to phonon emission/absorption. Here we resort to the simplest way to model these effects by assuming that our degree of freedom behaves like a localized boson with frequency $`\omega _l`$. The corresponding Hamiltonian reads:
$$H=\omega _la^{}a+\underset{๐ค}{}\omega _kb_๐ค^{}b_๐ค+\underset{๐ค}{}\frac{(\mathrm{๐ ๐ค})}{\sqrt{2\omega _kV\rho }}(a^{}b_๐ค+b_๐ค^{}a).$$
(49)
The ensuing equations of motion are
$`\dot{a}=i\left[\omega _la+{\displaystyle \underset{๐ค}{}}{\displaystyle \frac{(\mathrm{๐ ๐ค})}{\sqrt{2\omega _kV\rho }}}b_๐ค\right],`$
$`\dot{b}_๐ค=i\left[\omega _kb+{\displaystyle \frac{(\mathrm{๐ ๐ค})}{\sqrt{2\omega _kV\rho }}}a\right].`$ (50)
We next introduce the following (retarded) Greenโs functions $`A(t)i\theta (t)[a(t),a^{}(0)]`$ and $`B(t)i\theta (t)[b(t),a^{}(0)]`$. The fourier transforms of these Greenโs functions will consequently obey
$`(\omega \omega _l)\stackrel{~}{A}={\displaystyle \frac{1}{2\pi }}+{\displaystyle \underset{๐ค}{}}{\displaystyle \frac{(\mathrm{๐ ๐ค})}{\sqrt{2\omega _kV\rho }}}\stackrel{~}{B}_๐ค`$
$`(\omega \omega _k)\stackrel{~}{B}_๐ค={\displaystyle \frac{(\mathrm{๐ ๐ค})}{\sqrt{2\omega _kV\rho }}}\stackrel{~}{A}.`$ (51)
From Eqns.(51), one determines the real and imaginary parts of the Greenโs functions self-consistently. We however can disregard the phononsโ dispersion and damping which introduces an error in a higher order, in so far as the shifted frequencies $`\omega _l`$โs are concerned. This yields
$$\stackrel{~}{A}=\frac{1}{2\pi }\left(\omega \omega _l\frac{1}{3}\frac{g^2}{4\pi ^2\rho c_s^2}\underset{ฯต_10^+}{lim}_0^{k_c}\frac{k^3dk}{\omega /c_s+iฯต_1k}\right)^1,$$
(52)
where $`k_c`$ is the cut-off wave-vector whose value will be discussed shortly (we have also replaced $`_๐คV\frac{d^3๐ค}{(2\pi )^3})`$. Eqn.(52) gives immediately for the inverse life-time of the internal resonance
$$\tau _{\omega _l}^1=\frac{g^2}{4\pi \rho c_s^2}\left(\frac{\omega }{c_s}\right)^3\frac{3\pi }{2\mathrm{}}T_g\left(\frac{\omega }{\omega _D}\right)^3,\omega \omega _c$$
(53)
and its frequency shift
$`\omega _l(\omega )`$ $`=`$ $`\omega _l{\displaystyle \frac{g^2}{4\pi ^2\rho c_s^2}}{\displaystyle _0^{\omega _c}}{\displaystyle \frac{d\omega ^{}(\omega ^{}/\omega _c)^3}{\omega ^{}\omega }}`$ (54)
$``$ $`\omega _l{\displaystyle \frac{3}{2\mathrm{}}}T_g\left({\displaystyle \frac{\omega _c}{\omega _D}}\right)^3{\displaystyle _0^{\omega _c}}{\displaystyle \frac{d\omega ^{}(\omega ^{}/\omega _c)^3}{\omega ^{}\omega }},`$
where the factor of $`1/3`$ has disappeared because we have accounted for the three phonon polarizations and also ignored the distinction between the longitudinal and transverse branches. The singularity in Eq.(54) at $`\omega \omega _c`$ is completely artificial, as the cut-off is not supposed to be sharp. In our numerical estimates, we use a cut-off smeared by $`\delta \omega _c=\omega _c/\sqrt{D}`$, where $`D`$ is the glassโ fragility (see Appendix A); this is however totally unimportant as the divergence is only logarithmic. According to Eq.(54), the frequency shift scales roughly with $`\omega _c^3`$ and is thus rather sensitive to its value. Due to the dispersion, the resonance in Eq.(52) is effectively broadened because the value of the integral in Eq.(54) is positive for sufficiently small $`\omega `$, but turns negative at a frequency which is a multiple of $`\omega _c`$.
We approximate the phonon coupling effects by replacing in our spectral sums in Eqs.(39-41), (44-47) the discrete summation over different ripplon harmonics by integration over โlorentzianโ profiles:
$$\underset{l}{}๐\omega \delta (\omega \omega _l)\underset{l}{}๐\omega \frac{\gamma _\omega /\pi }{[\omega \omega _l(\omega )]^2+\gamma _\omega ^2},$$
(55)
where $`\gamma _\omega \tau _\omega ^1`$ is a (frequency dependent) friction coefficient and $`\omega _l(\omega )`$ is the ripplon frequency shifted due to the dispersion effects. This approximation amounts to having the total inverse life-time of a transition involving more than one mode being the sum of the inverse life-times of the participating modes. This would be correct in the case of a frequency independent $`\gamma `$, but should be still adequate at the low $`T`$ end of the plateau, where the absorption is mostly due to single ripplon mode processes.
The value of the cut-off frequency $`\omega _c`$ is close to but larger than $`(a/\xi )\omega _D`$ (see Appendix B), as the phonons whose wave-length is shorter than $`\xi `$ cause an increasingly smaller effective gradient of the phonon field as sensed by a region of size $`\xi `$. These shorter wave-length phonons will still strongly interact with the droplets, however at this point we could only emulate that to some extent by increasing $`\omega _c`$. This also brings us back to the radiation life-timeโs frequency dependence. It is now clear that for $`\omega _l(\omega )>\omega _c`$, $`\gamma _\omega `$ will not follow the simple cubic dependence cited above, the latter being probably still a safe lower estimate. We will thus use the above expression as it makes little difference computationally in the region of such intense damping. At the corresponding temperatures, the scattering is probably better formally described by the stochastic resonance Gamaitoni *et al.* (1998) methodology anyway.
We are now ready to discuss the non-universality of the plateau. It is evident from Eqs.(53)-(54) that even though the absorbersโ frequencies are determined by the quantum energy scale $`\omega _D`$, the overall effective frequency shifts scale with $`T_g`$. The ratio $`T_g/\omega _D`$ seems to vary within the range of between 2 and 5 among different glasses, and the non-universality in this number could have a substantial effect subject to the value of $`\omega _c`$. As argued in Appendix B, a value of $`\omega _c<2.5(a/\xi )\omega _D`$ is justified. $`\omega _c=1.8(a/\xi )\omega _D`$ seems to yield the experimentally observed spread in the plateauโs position. Our results for three values of $`T_g/\omega _D`$ are shown in Fig.19. Since $`\omega _c`$ should be regarded as an adjustable parameter we can claim to possess only circumstantial evidence that the plateauโs non-universality is caused by the spread in the value of the ratio of the two main energy scales in the problem: the classical $`T_g`$ and the quantum $`\omega _D`$. On a speculative note, this phenomenon may be a sign of strong mixing (and thus level repulsion) between the ripplons and the phonons, as implicitly confirmed by a phonon localization transition at frequencies just above those at the plateau. Indeed, the self-energy of an internal resonance of dimensions $`\xi `$ coupled with strength $`g`$ to an elastic medium scales (within perturbation theory) as $`g^2/\rho c_s^2\xi ^3T_g`$. This can be viewed as lowering of an impurity band edge due to the interaction with the phonons, yet another way to express the existence of mixing between the resonant transitions and the elastic waves. Within our theory, the non-universality of the plateau is an internally consistent proof that the degrees of freedom causing the Boson Peak are inelastic ones, whose coupling with the phonons then must be equal to $`g`$ related to the value of $`T_g`$ through the stability requirement explained in Section III.
We now comment on the plateau slopes in Fig.19 being noticeably more negative than the experimental value. The explaination is, we did not solve the full kinetic equation for the interacting system, but used a simplistic single life-time approximation. We demonstrate this issue by briefly presenting a slightly different way to estimate $`\rho _{exc}(\omega ,T)`$ from Eq.(43). Here, we imagine we do not exactly know the thermal weight function $`f(E,T)`$ due to the lack of knowledge of the life-times in the multi-level system. On general grounds, however, this function should decrease rapidly for $`\omega >\alpha T`$, where the $`\alpha `$ is of order unity. This yields $`\rho _{exc}(\omega ,T)N_{\alpha T}(\omega )`$ (where $`N_E(\omega )`$ was defined in Eq.(42)). We show the result of this approximation for reasonable $`\alpha =1`$ and $`\omega _c=2(a/\xi )\omega _D`$.
Even though these curves resemble the experimental data better than in the previous figure, they do not really provide more material support for the theory than the earlier method. This discussion simply demonstrates that the basic estimates are robust enough to โsurviveโ different levels of treatment. Also, curiously, these curves reflect the experimental tendency that the higher $`T`$ plateaux seem to have a more negative slope as compared to the low $`T`$ ones (see Fig.1), which was less obvious in Fig.19.
Finally in this subsection, we return to the specific heat. The effects of the phonon coupling on the ripplon spectrum can be taken into account in the same fashion as in the conductivity case. Here, we replace the discrete summation in Eq.(36) by integration over the broadened resonances, as prescribed by Eq.(55). The bump, as shown in Fig.15, is also predicted to be non-universal depending on $`T_g/\omega _D`$. The predicted bump for $`T_g/\omega _D=2`$ seems to match the best the available data for a-SiO<sub>2</sub>, whereas the more appropriate $`T_g/\omega _D4`$ is about a factor of $`3`$ lower in temperature. It is somewhat unsatisfying that the plateauโs and the bumpโs position can not be thus both made to exactly match the experiment at the same time say by adjusting $`\omega _l`$, which is certainly allowed given the qualitative character of some of the estimates. However, since we had to employ an approximation when calculating the scattering density of states, the discrepancy does not warrant too much concern, in our opinion.
### IV.7 The Relaxational Absorption
In addition to the resonant absorption, an internal resonance will also provide a so called โrelaxationalโ scattering mechanism. Since a cross-over to the multilevel behavior of the tunneling centers leads to an increased resonant scattering, we must check whether the relaxation mechanism is enhanced as well. This latter mechanism arises because a passing phonon modifies the energy bias of a particular pair of internal states. This causes irreversible thermal equilibration processes within each pair, resulting in energy dissipation Jรคckle *et al.* (1976); Maynard (1975). This phenomenon is sometimes referred to as the bulk viscosity Landau and Lifshitz (1987). One important difference between the relaxational and resonant absorption is that the former does not saturate and can easily exceed the latter at low enough temperature and high enough sound intensity, which is what is usually observed in ultrasonic experiments unless special care is taken Hunklinger and Raychaudhuri (1986) (this saturation is not an issue in heat conductance, owing to the rather low sound intensities in these experiments). Applying the notion of the relaxational absorption to the two-level systems explained well the shape of the maximum in the temperature dependence of the sound speed at very low frequencies at $`1K`$ Hunklinger and Raychaudhuri (1986), which is one of the impressive achievements of the TLS model. In Hunklinger and Raychaudhuri (1986), the relation between the slopes of the logarithmic temperature profiles around the maximum was explained. At higher $`T`$, the logarithmic decrease in $`c_s`$ is followed by what has been viewed by others as a mysterious linear law Belessa (1978). At higher frequencies still, the logarithmic decrease is outweighed by the just mentioned linear $`T`$ dependence. We have argued that the increase in the density of the scattering states is due to thermal activation of the vibrational states of the domain walls, or matching of the thermal phonon frequency with that of a ripplon on a mobile domain wall. Does the existence of the vibrational modes modify the relaxational scattering as compared to a bare underlying two-level system? The answer is: not significantly, for the following reason. The magnitude of the dissipation due to the bulk viscosity depends on the number of local distinct molecular configurations, populated according to the Boltzmann statistics. A shift in this population results in relaxational dissipation. While having a domain wall excited may modify the energy scale in the Boltzmann distribution, which may produce some effect, it does not change the number of the intrinsic (โinelasticโ) glassy states, and thus will not on average enhance the relaxational scattering. This is to be compared to the resonant scattering, which depends on the degeneracy of the ripplon states and will thus intensify at higher $`T`$, subject to the degree of the ripplonโs linearity. While the relaxational mechanism thus seems to play only a minor role in the phonon absorption at the plateau temperatures, its effects are observable and can explain, as we will argue below, the temperature independent $`\mathrm{log}\omega `$ part in the sound speed variation as measured in Belessa (1978). According to Jรคckle *et al.* (1976), the variation in the speed of sound due to a collection of two-level systems is
$$\frac{\delta c_s}{c_s}|_\omega =\frac{g^2}{2\rho c_s^2}\left(\frac{ฯต_i}{E_i}\right)^2\frac{\beta }{\mathrm{cosh}^2\beta E_i}\frac{1}{1+\omega ^2\tau _i^2},$$
(56)
where
$$\tau _i\frac{3g^2\mathrm{\Delta }_i^2E_i}{2\pi c_s^5}\mathrm{coth}(\beta E_i/2)$$
(57)
is the radiative life-time of the $`i`$th TLS Jรคckle (1972) (see also Eq.(53)), and the double angular brackets denote averaging with respect to $`E_i`$, $`\mathrm{\Delta }_i`$ and $`\tau _i`$. While it would seem that detailed information on the relevant parametersโ distribution is necessary to use Eq.(56), some qualitative conclusions can be made on general grounds. First, for small $`\omega `$ the average is dominated by the long life-time systems, i.e. those with $`\mathrm{\Delta }E`$ and thus $`ฯตE`$. As a result, the averaging over these systems is not very sensitive to the possible correlation between $`E_i`$ and $`\tau _i`$, and thus the summation over the two-level system (nearly flat!) spectral density $`๐ฯต(1/T_g\xi ^3)`$ introduces, within order unity, only a numerical factor proportional to $`T/T_g`$ (and eliminates the explicit temperature dependence). As just argued, the $`(ฯต/E)^2`$ factor should only give a correction factor of order unity, and we are left with averaging expression $`1/(1+\omega ^2\tau _i^2)`$ with respect to the life-time distribution. At low frequencies $`\omega `$, this averaging will be dominated by the TLS with the long life-times. Quite generally, for large $`\tau `$, $`P(\tau )d\tau d\tau /\tau `$ because $`\tau ^1`$ scales algebraically with $`\mathrm{\Delta }`$, and the distribution of $`\mathrm{log}\mathrm{\Delta }`$ is flat (at least for small $`\mathrm{\Delta }`$), or, almost flat, up to a weak power law, as argued earlier. More specifically, for a two-level system coupled linearly to the elastic strain, $`\tau ^1\mathrm{\Delta }^2E`$ (Eq.57). Therefore at each $`E`$ (which is incidentally only weakly dependent on $`\mathrm{\Delta }`$ in the relevant long life-time case $`\mathrm{\Delta }/E1`$), obviously $`d(\mathrm{log}\mathrm{\Delta })=constd(\mathrm{log}\tau )=const`$. Thus the averaging w.r.t. $`\tau `$ will produce a term of the order $`(g^2/\rho c_s^2)(1/T_g\xi ^3)\mathrm{log}\omega `$, which is of the right order of magnitude (and sign!). Since the dimensionless factor in front of the $`\mathrm{log}\omega `$ term has been shown to be universal ($`(a/\xi )^3`$), the present theory predicts that it should not vary significantly among the insulating glasses; in fact, according to our argument, it is proportional to the coefficient $`\alpha `$ at the logarithmic temperature dependence of the sound speed variance in the TLS regime, a rather universal quantity indeed Leggett (1991). We stress however that the just predicted TLS-like property should be observed in the plateau regime. A deviation would be a sign of more than two inelastic states playing a role in the transition. We finally mention that the lower limit in the integral over the life-time distribution should produce a $`\mathrm{log}T`$ term, which would be however masked by the stronger linear dependence.
## V Quantum Effects beyond the strict Semi-Classical Picture
### V.1 Quantum Mixing of a Tunneling Center and the Black-Halperin Paradox
The preceding sections have shown that structural transitions, accompanied at high enough temperatures by vibrational excitations of the mosaic, account for the most conspicuous departures of the low temperature behavior of glasses from the prescriptions of a standard harmonic lattice theory - namely the existence of multilevel intrinsic resonances in a amorphous sample made by quenching a supecooled liquid. At the lowest temperatures these resonances behave for the most part as if they were two-level systems, while at higher $`T`$ the density of states of these intrinsic excitations grows considerably and leads to the Boson Peak phenomena. While we have computed the density of states accessible by tunneling even at the lowest temperatures, we have assumed, within a semi-classical approach, that having a small tunneling barrier between alternative local structural states does not affect significantly the corresponding spectrum $`n(ฯต)`$ of the lowest energy transitions from its classically defined value. Likewise, we have assumed that the vibrational spectrum of moving domain walls is unaffected by the presence of tunneling, that would in principle mix those vibrations quantum mechanically. Clearly, the transitions that are active at low $`T`$ must have some significant (even if small) overlap between the wave-functions corresponding to the alternative structural states. This overlap would lead to the familiar effects of repulsion between the semi-classically determined energy levels. This could be described as partial quantum melting of some tunneling centers, but it is probably better to use term โquantum mixingโ.
In this section we estimate the magnitude of these quantum mixing effects. Even though the strictly semiclassical theory agrees well with experiment as is, making such estimates that go beyond it is useful for two distinct reasons. First, we must check to what extent the semi-classical picture, tacitly assumed earlier, is a consistent zeroth order approximation to a more complete treatment. Second, it is important to ask whether the expected corrections to the strict semiclassical theory lead to observable consequences. In what follows, we provide approximate arguments that indeed such corrections are discernible and may even potentially answer some long-standing puzzles in this field.
Quantum Mixing. As the starting point in the discusion, we consider a simplified version of the diagram of a tunneling centerโs energy states from Fig.14 with $`ฯต<0`$, as shown on the left hand side of Fig.21.
We remind the reader that the $`ฯต<0`$ situation, explicitly depicted in Fig.21, implies lower transition energies than when the semiclassical energy difference $`ฯต>0`$ and thus dominates the low temperature onset of the Boson Peak and the plateau.
Accounting for tunneling, the low energy portion of the Hamiltonian that corresponds to Fig.21 is as follows:
$$H_i=\left(\begin{array}{cccc}0& \mathrm{\Delta }/2& 0& 0\\ \mathrm{\Delta }/2& |ฯต|& \mathrm{\Delta }_{i_1}& \mathrm{\Delta }_{i_2}\\ 0& \mathrm{\Delta }_{i_1}& \mathrm{}\omega _{i_1}& 0\\ 0& \mathrm{\Delta }_{i_2}& 0& \mathrm{}\omega _{i_2}\end{array}\right),$$
(58)
where the semi-classical values of the ripplonic energies are denoted as $`\mathrm{}\omega _i`$ and the transition amplitudes to those levels are $`\mathrm{\Delta }_i`$ respectively (only two lowest of those ripplonic states are shown in Eq.(58)). As argued in detail in Section IV, only one of the lowest two energy levels in a tunneling center (the top one in this case) is directly coupled to the higher, ripplonic, energy states. Obviously, virtual transitions to those high energy states will result in lowering the energy of the higher level. There are no direct transitions from the bottom state of energy $`0`$, as explained in the previous Section, and therefore its position is unaffected by the presence of the ripplons. Consequently, the effective energy splitting of the two-level system (with $`ฯต<0`$) will be lower than the classical value obtained earlier, and the smaller the original value of $`ฯต`$ was, the more pronounced the effect will be. In what follows we estimate the consquences of this effect on the apparent energy spectrum of the lower excitations, i.e. the empirical two-level systems. In the limit of inifinitely small tunneling amplitude $`\mathrm{\Delta }`$, the decrease in $`ฯต`$ could be estimated using a perturbative expansion:
$$|\stackrel{~}{ฯต}|=|ฯต|\underset{i}{}\frac{\mathrm{\Delta }_i^2}{\mathrm{}\omega _i|\stackrel{~}{ฯต}|}.$$
(59)
Here, $`\stackrel{~}{ฯต}`$ is the new value of the energy splitting, $`\omega _i`$โs are the ripplon frequencies and $`\mathrm{\Delta }_i`$โs are tunneling amplitudes of transitions that excite the corresponding vibrational mode of the domain wall. Those amplitudes will be discussed in due time; for now, we repeat, the expression above will be correct in the limit $`\mathrm{\Delta }_i/\mathrm{}\omega _i0`$. Finally, the renormalized value $`\stackrel{~}{ฯต}`$ was used in the denominator. While, according to Feenbergโs expansion Feenberg (1948), including $`\stackrel{~}{ฯต}`$ in the resolvent is actually more accurate, we do it here mostly for convenience.
Given that the semi-classical values of eigen-values $`\mathrm{}\omega _i`$ are known, the low energy portion of the energy level structure of the tunneling center, as shown in Fig.21, gives a quantitative idea of the eigen-energies of the full Hamiltonian only in the limit of a very small tunneling splitting $`\mathrm{\Delta }`$. In a complete treatment, all transition amplitudes must be included and the Hamiltonian diagonalized. In general, such diagonalization (and, in our case, the systemโs โquantizationโ) is difficult, however could still be conducted approximately in some cases of interest. Consider, for the sake of argument, the following situation, where $`\mathrm{\Delta }`$ is not necessarily smaller than $`ฯต`$ but $`_i\mathrm{\Delta }_i^2/\mathrm{}\omega _i`$ is. In this arrangement, the energy shift due to the higher lying states can be computed using perturbation theory and yields a โrenormalizedโ value of the classical energy difference that we have called $`\stackrel{~}{ฯต}`$. This procedure also modifies the tunneling amplitude $`\mathrm{\Delta }`$ of the underlying TLS by a multiplicative factor according to
$$\stackrel{~}{\mathrm{\Delta }}=\mathrm{\Delta }\left(1\frac{1}{2}\underset{i}{}\frac{\mathrm{\Delta }_i^2}{(\mathrm{}\omega _i)^2}\right).$$
(60)
Following this, the full energy splitting of the TLS tunneling transition is computed using $`E=\sqrt{\mathrm{\Delta }^2+\stackrel{~}{ฯต}^2}`$. The important feature of the argument is that $`\mathrm{\Delta }`$ (or $`\stackrel{~}{\mathrm{\Delta }}`$) is allowed to take arbitrarily large values relative to $`ฯต`$ and the ratio of the two parameters is not treated perturbatively. The lowering of $`ฯต`$ due to virtual transitions among the higher energy states changes somewhat the effective density of transition energies $`E`$ that directly enters into the heat capacity and conductivity calculations. While Eq.(59) is perturbative, it should accurately give finite effects in the mean-field limit of infinitely many transitions $`\mathrm{}\omega _i`$ coupled infinitely weakly to one of the two bottom states of the tunneling center. We will analyze the physical consequences assuming the accuracy of Eq.(59). We must, of course, bear in mind that while a tunneling center is nearly a meanfield entity, owing to the strong correlations, it is actually of finite, albeit molecularly large size. Let us plot $`|ฯต|`$ as a function of $`|\stackrel{~}{ฯต}|`$ (see Fig.22(a)).
Clearly, the smallest allowed value of the effective classical splitting $`|\stackrel{~}{ฯต}|`$ of zero corresponds to a finite value of $`|ฯต|`$ equal to
$$\mathrm{\Sigma }\underset{i}{}\frac{\mathrm{\Delta }_i^2}{\mathrm{}\omega _i}.$$
(61)
Therefore, smaller values of $`|ฯต|`$ do not correspond to physically realizable systems, according to Eq.(59). The overall excitation spectrum of the structural transitions with those (small) values of $`ฯต<\mathrm{\Sigma }`$ is strongly affected by the ease of tunneling. As a consequence, the $`ฯต`$ and $`\mathrm{\Delta }`$ distributions become correlated. The quantity $`\mathrm{\Delta }_i^2/\mathrm{}\omega _i`$ is of central importance for this section of the article, therefore we will discuss it now in some detail. The wave-functions of the highly quantum tunneling centers are heavily mixed combinations of the classical states corresponding to potential energy minima. Transitions between such states are strongly coupled to lattice vibrations and, among other things, would strongly scatter phonons. This result is expected on rather general grounds and was exploited earlier when we noted that the eigen-states of the classical potential energy are largely unrelated to the eigenstates of the vibrational modes of the domain wall, hence one expects that transition amplitudes of exciting a ripplon (which, loosely speaking, is a highly anharmonic combination of both structural and vibrational modes of the lattice) are expected to be comparable to the ripplon energy itself. If this is the case, then number $`\mathrm{\Sigma }`$ is actually very large relative to $`ฯต`$ and one, strictly speaking, should not use a perturbative expression, such as in Eq.(59) in order to assess the lowering of $`ฯต`$ due to quantum effects. We note, however, it is more likely that the seeming discrepancy simply stems from our โquantizationโ procedure being so far rather naive, so let us slow down a bit and attempt to outline briefly a more careful way to quantize the tunneling centersโs dynamics. First of all, recall that we actually know the values $`\omega _i`$ in a strongly quantum regime, because they were computed assuming a freely moving membrane. On the other hand, we know that in the classical limit domain wall vibrations are indistinguishible from the lattice vibrations. Remarkably, the vibrational eigen-frequencies of a โboxโ of dimensions $`\xi `$ \- $`\omega _D`$ times a number from $`a/\xi `$ to $`1`$ \- span roughly the same range of energies that the $`\omega _i`$โs do. Therefore, even though the quantization procedure will โreshuffleโ all the ripplonic states, it will not significantly shift their position as a whole. Next, since static lattice inhomogeneities scatter phonons only elastically, the coupling of the ripplonic excitations to the phonons must scale with a positive power of $`\mathrm{\Delta }`$ so as to vanish in the classical limit. Consequently, in the limit of small tunneling matrix element $`\mathrm{\Delta }`$, the transition amplitudes $`\mathrm{\Delta }_i`$โs must scale with a positive power of $`\mathrm{\Delta }`$. On the other hand, as mentioned many times, in the quantum regime, the $`\mathrm{\Delta }_i`$โs are of the order $`\mathrm{}\omega _i`$ and are not directly related to $`\mathrm{\Delta }`$. Therefore, a careful quantization procedure of introducing quantum tunneling in the system must combine and rationalize both of those seemingly conflicting notions. It will turn out that in the quantum regime, the $`\mathrm{\Delta }_i`$โs are related to the strength of the underlying tunneling transition $`\mathrm{\Delta }`$ only in a certain renormalized sense. We will now outline this renormalization/quantization procedure. This procedure imposes certain restrictions on how adding the possibility of tunneling to a local structural transition can be performed, so that the structure of the energy levels of the transition, that we know from general arguments, is preserved. Recall that โswitching onโ tunneling to the higher energy states $`ฯต+\mathrm{}\omega _i`$ not only lowered $`ฯต`$, but also made $`\mathrm{\Delta }`$ smaller by a factor of $`[1_i\mathrm{\Delta }_i^2/(\mathrm{}\omega _i)^2]`$. This expression is only valid if the sum is a small number, so that the whole correction factor is necessarily positive and only changes the effective magnitude of the matrix tunneling element, but not its sign. We must require that $`\mathrm{\Delta }`$ not change its sign in the course of the โquantizationโ, but only its absolute value, because the (ordinarily small) value $`\mathrm{\Delta }`$ only reflects the (ordinarily small) configurational overlap between two local structural states, while the sign (or complex phase, in general) bears no special meaning here because no particular spatial symmetry is involved in the problem. (Such spatial symmetry is important, for instance, when computing overlaps between eigenstates, or near-eigenstates of orbital momentum centered around close locations in space.) Thus, the final answer should depend only on $`|\mathrm{\Delta }|^2`$. This becomes especially evident after the following realization. Note that the expression in the brackets in Eq.(60) is a small coupling limit of what can be considered a Franck-Condon factor. The appearance of such a factor after the introduction of non-zero transition amplitudes is natural: the degrees of freedom that used to be classical and static, can now follow to some extent a selected motion in the system. The Franck-Condon (FC) factor is the overlap between the the initial and final wave-functions of these other (โripplonicโ) degrees of freedom, corresponding to the initial and final configuration of that selected motion. (A well known example of a FC factor arising in an analogous, dynamical fashion is the tunneling matrix element renormalization in the spin-boson problem Leggett *et al.* (1987).) Let us suppose, in a simplified manner, that the effective renormalization of all of the newly introduced tunneling amplitudes occurs in a similar fashion. This allows one to self-consistently close Eq.(60) and rewrite it for a representative amplitude $`\mathrm{\Delta }_Q`$:
$$\stackrel{~}{\mathrm{\Delta }}_Q=\mathrm{\Delta }_Q\left[1\frac{B}{2}\frac{\stackrel{~}{\mathrm{\Delta }}_Q^2}{(\mathrm{}\omega _D)^2}\right],$$
(62)
where replacing $`\mathrm{\Delta }_Q`$ by $`\stackrel{~}{\mathrm{\Delta }}_Q`$ inside the brackets preserves the approximationโs order. $`B`$ is a numerical constant, reflecting the sum over ripplon states with their vibrational frequencies $`\omega _i`$, and we have replaced $`\omega _i`$ by the Debye frequency $`\omega _D`$. The two must be related at the end of the renormalization, as we already know. Identifying the expression in brackets with a Franck-Condon factor reminds us that $`\mathrm{\Delta }_Q`$ (or $`\stackrel{~}{\mathrm{\Delta }}_Q`$) is not only a (generic) tunneling amplitude but also can be considered a coupling constant, therefore only its absolute value is physically relevant in the present context, not its sign. The smallness of $`\mathrm{\Delta }_Q`$ and $`\stackrel{~}{\mathrm{\Delta }}_Q`$ lets us recast Eq.(62) as
$$d\mathrm{log}\left(\frac{\stackrel{~}{\mathrm{\Delta }}_Q}{\mathrm{\Delta }_Q}\right)=\frac{B}{2}d\left[\frac{\stackrel{~}{\mathrm{\Delta }}_Q^2}{(\mathrm{}\omega _D)^2}\right],$$
(63)
where the reference state is $`\stackrel{~}{\mathrm{\Delta }}_Q=\mathrm{\Delta }_Q=0`$, $`lim_{\mathrm{\Delta }_Q0}(\stackrel{~}{\mathrm{\Delta }}_Q/\mathrm{\Delta }_Q)=1`$ (that is โno tunnelingโ $``$ โno renormalizationโ). Notice that the r.h.s. of Eq.(63) depends explicitly only on the effective tunneling element $`\stackrel{~}{\mathrm{\Delta }}_Q`$, but not on the original (tunable) perturbation strength $`\mathrm{\Delta }_Q`$. We therefore can use the differential relation in Eq.(63) to extend the perturbative construction from Eq.(62) into the region of arbitrarily large values of the bare coupling $`\mathrm{\Delta }_Q`$ by using the outcome of the previous (infinitesimal) change in $`\mathrm{\Delta }_Q`$ as the initial input in the subsequent increment $`d\mathrm{\Delta }_Q`$. Each of these increments is a small perturbation around a new, self-consistently determined, value of $`\stackrel{~}{\mathrm{\Delta }}_Q`$. One gets the following self-consistent equation for the effective tunneling amplitude as a result:
$$\stackrel{~}{\mathrm{\Delta }}_Q=\mathrm{\Delta }_Qe^{\frac{B}{2}\frac{\stackrel{~}{\mathrm{\Delta }}_Q^2}{(\mathrm{}\omega )^2}}.$$
(64)
Our renormalization procedure is internally consistent in that the physical value of the tunneling amplitude depends on the scaling variable - the bare coupling $`\mathrm{\Delta }_Q`$ \- only logarithmically. This bare coupling must scale with the only quantum scale in the problem - the Debye frequency, as pointed out yet in the first section.
In a more complete treatment, the $`\mathrm{\Delta }`$ renormalization would not be characterized by a single โFranck-Condonโ parameter, but by a distribution of Franck-Condon factors. Therefore, the exponential form in Eq.(64) might be possibly replaced by a different, perhaps a polynomial expression. In fact, one may think minimally of Eq.(64) as of one of the possible Padรฉ extensions of the perturbative formula (62). At any rate, such a Padรฉ approximant will retain the main feature of Eq.(64) in that the value of the observable tunneling matrix element $`\stackrel{~}{\mathrm{\Delta }}_Q`$ is bounded from above and depends strongly on the (semi-)classical energies $`\mathrm{}\omega _i`$. According to the discussion above, this restriction stems from a self-consistency condition, namely that the leading term in the exponent in Eq.(64) must scale with $`\stackrel{~}{\mathrm{\Delta }}_k^2`$, if this same number $`\stackrel{~}{\mathrm{\Delta }}_k`$ is on the l.h.s. in that equation. An important corollary of this is that the perturbative term inside the brackets of Eq.(60) must scale with $`\stackrel{~}{\mathrm{\Delta }}^2`$ itself, hence the perturbation correction $`\mathrm{\Sigma }`$ will scale with $`\stackrel{~}{\mathrm{\Delta }}^2`$ too (from now on, we will drop tildes from the symbols denoting the physical tunneling amplitudes, but retain them for the effective $`ฯต`$โs). That is, with our definition of $`B`$, it is roughly true that
$$\mathrm{\Sigma }=B\frac{\mathrm{\Delta }^2}{\mathrm{}\omega _D}.$$
(65)
We have thus demonstrated explicitly that the magnitude of quantum effects on the classical energy splitting $`ฯต`$ on a particular site should depend on the facility of tunneling at that same site. We have therefore established that the fact of quantum $`\mathrm{\Delta }_i`$ being close in value to a rather large energy scale $`\omega _i`$ is consistent with a relatively small value of the correction in Eq.(59) and its scaling with $`\mathrm{\Delta }^2`$. As another dividend from the argument, we obtain a ballpark estimate of the constant $`B`$. A structural transition that is thermally active at a plateau temeperature $`k_BT\mathrm{}\omega _i`$ and being an efficient phonon scatterer, will have $`\mathrm{\Delta }\mathrm{}\omega _i`$. Therefore, $`B=_i\omega _D/\omega _i`$, will be a number on the order of several hundreds, since the total number of the ripplonic modes (at the laboratory glass transition, at which $`\xi /a\mathrm{5..6}`$) is approximately a hundred, and $`\omega _i`$ is proportional to, but somewhat smaller than the Debye frequency $`\omega _D`$.
Eq.(65) implies that while the distributions of the clasical energy splittings $`ฯต`$ and the bare semiclassical tunneling amplitude $`\mathrm{\Delta }`$ may be uncorrelated, quantum corrections require that $`\stackrel{~}{ฯต}`$ and $`\mathrm{\Delta }`$ be correlated for systems with a sufficiently low barriers and which simultaneously have small energy difference between the initial and final structural state. Conversely, the independence approximation is valid when, roughly, $`|ฯต|>B\frac{\mathrm{\Delta }^2}{\mathrm{}\omega _D}`$. Since $`\mathrm{\Delta }`$ is proportional to $`\mathrm{}\omega _D`$, this criterion is a formal restatement of an earlier comment that the theory is strictly valid in the classical limit. (Note that there is also a (much stronger) $`\mathrm{}`$ dependence in the exponent of the tunneling element $`\mathrm{\Delta }`$ (see Eq.(21)). While, obviously, only a negligible fraction of the total number of structural rearrangements in the liquid at $`T_g`$ would not satisfy the classicality criterion, these particularly facile transitions do actually comprise a significant portion of those transitions that are thermally active at cryogenic temperatures. We will now indicate what the observable consequences of this deviation for the strict semiclassical limit are. In order to do this, let us discuss first the difference between the strongly quantum and the bare โclassicalโ structural transitions.
According to Eq.(59), for all transitions, whose diagonal energy difference would be $`|ฯต|<\mathrm{\Sigma }`$ in the classical limit, the effective diagonal splitting $`\stackrel{~}{ฯต}`$ is actually zero, meaning that the full energy splitting $`E`$ is entirely comprised of the originally off-diagonal energy scale $`\mathrm{\Delta }`$. This implies that the energy eigen-states of such highly mixed tunneling centers are heavy superpositions of the original classical structural states and would not be easily interpreted in terms of the atomic coordinates of the potential minima alone, but must include the kinetic energy term as well. This is directly related to the well known ambiguity in separating the energy of such systems into potential and kinetic components even at conditions that are entirely classical, such as at a high temperature. Of course, in such cases free energy formulations must be employed that allow one to count the number of configurational states unambiguously, while using โinherentโ structures based on potential energy stationary points alone is of limited utility. The strongly quantum case can be loosely understood by transcribing the complex multiparticle rearrangements onto a single collective โreactionโ coordinate (as in the soft potential model Galperin *et al.* (1991)). In fact, this analogy to a single coordinate soft potential model is quite loose because of the much higher density of states of the ripplons (that give rise to the Boson Peak and correspond to the vibrations of the membrane) compared with the density of states of the soft potential model, which is one dimensional so that only one coordinate is vibrationally excited. Nevertheless, following this analogy, consider a two-well potential (with very steep outer walls) with a barrier high enough so that the physical coordinate eigen-states corresponding to the particle being in the left or the right well are unambiguously definable and the diagonal component of the transitionโs energy is equal to the difference in the potential energy of the two well with high accuracy. Imagine next lowering the barrier. In the limit of zero barrier the system is simply a particle in a square box, whose energy scale is determined by the quantum energy scale in the problem - that is the particleโs kinetic energy alone. This analogy reminds us that just like the transition from a largely classical to quantum behavior in a double well potential, the transition at $`|ฯต|=\mathrm{\Sigma }`$ is not sharp (note, however, that unlike in a one dimensional soft-potential model, the density of excited states of a tunneling center is very high thus possibly leading to a sharper cross-over). Put another way, this โphase transitionโ clearly corresponds to term-crossing and therefore would be gradual in a finite system. From a mean-field perspective, the transition at $`ฯต=\mathrm{\Sigma }`$ resembles a de-localization phase transition (see e.g. Abou-Chacra *et al.* (1973)) which we may think of as quantum depinning of the domain wall. Alternatively, one could say that the local structure of classical energy levels melts out locally in that the energy variation on the mostly classical landscape (determined by $`T_g`$) happens locally to be smaller than the confinement kinetic energy of the domain wall motion. Of course, this is occuring for only small parts of an otherwise rigid matrix. Again, since the system is finite, one expects a soft cross-over rather than a sharp transition when such โmeltingโ occurs. Both ways of interpreting the quantum mixing/melting described above are consistent with our view of the tunneling process leading to the expression (21) for the tunneling amplitude $`\mathrm{\Delta }`$. The action exponent in Eq.(21) scales as the height of the barrier relative to the under-barrier frequency. The former quantity, while distributed, scales with the classical energy scale in the problem - $`T_g`$, while the latter is proportional to the Debye frequency (and, most likely, is somewhat distributed too). The quantum limit of large $`\mathrm{}\omega _D`$ corresponds to a narrow barrier and a short tunneling path. This would imply the relative unimportance of the classical energy landscape modulation during the tunneling process. Finally, in order to avoid ambiguity, we stress that the structural transitions of both types of tunneling centers, that we have called โclassicalโ (in that the wave functions are well localized near minima and are well defined structurally, i.e. in position space) and โquantumโ (i.e. in a superposition of structural states), at low temperature occur in a purely quantum mechanical fashion, that is by tunneling.
We now show that the presence of a somewhat distinct class of such low barrier, or โfastโ, two-level systems, whose effective diagonal splitting is zero, leads to additional phonon scattering in comparison with the strictly semiclassical analysis, which neglects the renormalization from quantum mixing effects. This additional scattering at low energy is consistent with the apparent subquadratic temperature dependence of the heat conductivity in the TLS regime. The mixing also leads to a super-linear addition to the heat capacity at subKelvin temperatures. These highly quantum tunneling centers in strongly mixed superpositions of structural states, therefore, give a mechanism to resolve a quantitative deviation from the standard tunneling model, which was brought up by Black and Halperin Black and Halperin (1977) in 1977. They noted that the short time heat capacity of a-SiO<sub>2</sub> is larger than would be predicted by the logarithmic dependence obtained in the STM, if one uses the TLS parameters extracted from ultra-sonic measurements. The quantitative mismatch appears to be as if there were two kinds of two-level systems: one set obeying the distribution postulated in STM, and another set of โfastโ tunneling centers responsible for the short time value of the heat capacity. We can see our analysis of mode mixing leading to the existence of a finite number of two-level systems with $`\stackrel{~}{ฯต}`$ very nearly $`0`$, as suggested by Eq.(59) is quite consistent with this empirical notion<sup>6</sup><sup>6</sup>6We must stress however that the Black-Halperin analysis has been conducted only for a single substance, namely amorphous silica, and systematic studies on other materials should be done. The discovered numerical inconsistency may well turn out to be within the deviations of the heat capacity and conductivity from the strict linear and quadratic laws repsectively. Finally, a controllable kinetic treatment of a time-dependent experiment would be necessary..
To see this more explicitly we note that Eq.(59) allows one to formulate the effects of quantum mode mixing as a change in the apparent distribution of the diagonal energy splitting. Whatever the old distribution of classical energy difference $`n(ฯต)`$, the new distribution of the effective classical component of the transition energy can be found using $`n(\stackrel{~}{ฯต})|d\stackrel{~}{ฯต}|=n(ฯต)|dฯต|`$. For $`\stackrel{~}{ฯต}`$โs not too close to $`\mathrm{}\omega _i`$ (case $`\stackrel{~}{ฯต}\mathrm{}\omega _i`$ will be discussed later), which is appropriate in the TLS regime, the function $`|ฯต/\stackrel{~}{ฯต}|`$ that describes the relative probability distribution of the two quantities, is given by
$$\left|\frac{ฯต}{\stackrel{~}{ฯต}}\right|=\mathrm{\Sigma }\delta (\stackrel{~}{ฯต})+1,$$
(66)
where the $`\delta `$-function is positioned to the right of the origin: $`_0^{0^+}๐ฯต\delta (ฯต)=1`$ (see also Fig.22b). Consequently, the distribution of the effective diagonal splitting is:
$$n_\mathrm{\Delta }(\stackrel{~}{ฯต})=\frac{1}{T_g\xi ^3}\left[B\frac{\mathrm{\Delta }^2}{\mathrm{}\omega _D}\delta (\stackrel{~}{ฯต})+e^{|\stackrel{~}{ฯต}|/T_g}\right].$$
(67)
The coefficient of the $`\delta `$-function reflects the โpile-upโ of the two-level systems that would have had a value of $`|ฯต|<\mathrm{\Sigma }`$ were it not for quantum effects. These fast two level systems will contribute to short time value of the heat capacity in glasses. The precise distribution in Eq.(67) was only derived within perturbation theory and so is expected to provide only a crude description of the interplay of clasical and quantum effects in forming low barrier TLS. Quantitative discrepancies from the simple perturbative distribution may be expected owing to the finite size of a tunneling mosaic cell, as mentioned earlier, and the finite life-times of each energy state due to phonon emission. These effects would also smoothen the local quantum melting transition as $`\stackrel{~}{ฯต}0`$. While various improvements of the functional form of $`n(\stackrel{~}{ฯต})`$ might be suggested, it seems unwarranted, at present, to use any more complicated expressions for this function. Thus, to see the main consequences of the quantum mixing effect, we will proceed with the perturbative expression. Assuming a particular value of the coefficient $`B`$ allows one to derive the contribution of the fast two-level systems to the heat capacity and scattering of the thermal phonons. Before we start, let us note that since we now have to deal with a specific coupled distribution of $`ฯต`$ and $`\mathrm{\Delta }`$, the generic two-level system model that only specifies the distribution of the total splitting $`E`$ is not sufficient. We must use the full tunneling model where the tunneling elements $`\mathrm{\Delta }`$โs are distributed according to Eq.(22). The exact value of constant $`A`$ in equation (22) depends (weakly!) on the (possibly $`ฯต`$-dependent) cut-off value of the $`P(\mathrm{\Delta })`$ distribution. Both the heat capacity and the phonon scattering strength depend on the coefficient $`A`$, therefore it is possible to check the relative contribution of the โquantumโ centers to both of those quantities, regardless of $`A`$โs value. The $`n(ฯต,\mathrm{\Delta })`$ distribution obtained in this way is now a product of the $`P(\mathrm{\Delta })`$ distribution from Eq.(22) and the density of states from Eq.(67). The new normalization coefficient $`A_1`$ is found from the requirement that $`๐ฯต๐\mathrm{\Delta }n(ฯต,\mathrm{\Delta })=1/\xi ^3`$. This gives $`A_1=\left[\frac{B}{T_g\mathrm{}\omega _D}\mathrm{\Delta }^{2c}+\frac{1}{c}\left(\frac{1}{\mathrm{\Delta }_{min}^c}\frac{1}{\mathrm{\Delta }_{max}^c}\right)\right]^1`$). In order to compute the life-time of a phonon of energy $`E`$, one averages the Golden Rule scattering rate $`\frac{\pi g^2\mathrm{\Delta }^2}{\rho c_s^2E}\mathrm{tanh}\frac{\beta E}{2}`$ with respect to $`n(ฯต,\mathrm{\Delta })`$, subject to the resonance condition $`E=\sqrt{ฯต^2+\mathrm{\Delta }^2}`$ Anderson *et al.* (1972); Jรคckle (1972); Phillips (1981). This yields two contributions to the decay rate:
$`\tau _E^1`$ $`=`$ $`{\displaystyle \frac{\pi }{3}}A_1\left({\displaystyle \frac{a}{\xi }}\right)^3E\left({\displaystyle \frac{\mathrm{\Delta }_{max}}{E}}\right)^c`$ (68)
$`\times `$ $`\left[{\displaystyle \frac{BE}{\mathrm{}\omega _D}}+{\displaystyle _{\mathrm{\Delta }_{min}/E}^1}๐x{\displaystyle \frac{x^{1c}}{\sqrt{1x^2}}}\right].`$
The first term in the square brackets is the contribution owing to the fast, or highly quantum, two-level systems. Note that this term scales faster with $`E`$ than the other term. Provided the magnitude of this first term is comparable to the other term, the fast modes will somewhat modify the overall scaling of the heat conductivity $`\kappa `$. Without the first term, $`\kappa `$ scales superquadratically according to $`T^{2+c}`$ (recall that the heat conductivity is inversely proportional to the scattering rate from Eq.(68)). If we use a numerical value of $`B`$ of the order 100, this leads to a subquadratic $`T`$ dependence of $`\kappa `$: Experimentally, $`\kappa (T)`$ scales like $`T^{1.9\pm .1}`$ as extracted from a decade and a half of data (see Fig.1). Without the fast TLS, one, again, would have $`\kappa T^{2+c}`$. Using the theoretical approximation for $`c`$, this differs from the empirically observed value at least by a factor of $`(10^{1.5})^{c+.1}2`$ at $`T10^2T_D`$. Obviously, this is a very crude estimate because, first, we do not know how far down in temperature the power law scaling of $`\kappa `$ goes; second, our correction, while going in the right direction, summed with the older result, is not strictly a power law. Since the integral in the square brackets of Eq.(68) varies between 1 and $`\pi /2`$ for $`0<c<1`$ ($`\mathrm{\Delta }_{min}/E1`$, surely at $`E10^2T_D`$), we conclude that the first term must be between $`10^0`$ and $`10^1`$ in order to make a sizable contribution to the phonon scattering and modify its functional form. Since $`E10^2T_D`$, this shows that $`B`$ indeed must be of the order of several hundreds, consistent with our expectations based on the number of vibrational modes in the Boson Peak.
Does this mixing induced correction to the density of states with the value of $`B`$ around a hundred make an appreciable contribution to the time-dependent heat capacity? Following the calculation from subsection III.3, but now using the new distribution $`n(ฯต,\mathrm{\Delta })`$, one finds:
$$C(t)=\frac{A_1}{T_g\xi ^3}_0^{\mathrm{}}๐E\left(\frac{\beta E}{2\mathrm{cosh}(\beta E/2)}\right)^2\left(\frac{\mathrm{\Delta }_{max}}{E}\right)^c\left[\frac{BE}{\mathrm{}\omega _D}\theta (t\tau _{min})+_0^{\mathrm{log}(t/\tau _{min}(E))}๐z\frac{e^{\frac{c}{2}z}}{2\sqrt{1e^z}}\right],$$
(69)
where $`\theta (t)`$ is the usual step-function and $`\tau _{min}`$ is the fastest possible relaxation time of a TLS with the total energy splitting $`E`$, defined in Eq.(25). Again, the first term in the square brackets gives the contribution of the โfastโ TLS. Using the same numbers as given in subsection III.3, it is straightforward to show that the second, regular, term is of the order a hundred at temperatures $`T10^2T_D`$ when measured on the time scale of minutes. At the same time, the first term is at most of order ten. Note that at the shortest times $`t\tau _{min}`$, when the regular two-level system only begin to contribute to the heat capacity, the theory with quantum corrections says the actual heat capacity is finite and is at the most one tenth of the long-time value. At the same time, the fast tunneling centers do not seem to contribute significantly to the long-time heat capacity. We note however that the result obtained $`c(T)T^{1+c/2}`$ with $`c=.1`$ gives a somewhat slower rise with temperature than seen in experiment. The quantum correction again goes in the right direction of increasing the rate of the heat capacity growth with temperature relative to the $`T^{1+c/2}`$ law.
We have established that effects beyond the strict semi-classical analysis give rise to a subset of tunneling centers that undergo faster tunneling than the rest. Nevertheless, there are some quantitative issues in the heat capacity magnitude that remain to be understood, namely that the computed contribution of the โfastโ centers seems somewhat lower than what is necessary to explain the deviation of the experimental $`T`$ dependence from the supelinear dependence $`T^{1+c/2}`$ predicted by the present (approximate) argument. It is posssible that ultimately a broader view of the time-dependence of the heat capacity needs to be taken. Since, in fact, the system will clearly be aging by tunneling at those low temperatures, the notion of fixed frozen-in โdefectsโ may no longer be adequate - essentially interactions between defects play a role. โAgingโ by definition implies irreversible structural changes. More work on understanding the long time evolution of the tunneling centers is necessary.
We have concentrated on the quantum corrections to the low lying tunneling states with low barriers. Quantum mixing applies to the higher energy states too. Energy shifts and quantum melting occur within sub-bands of the ripplonic states of order $`l`$ and respective degeneracy $`(2l+1)`$, thus mixing these states. As tunneling can take place on a given time scale and the vibrationally excited levels become observable, their apparent energies can not be degenerate because the levels are coupled through those same tunneling transitions. The magnitude of energy level repulsion from the quantum mixing can be assessed qualitatively. In the limit of weak coupling, the deviation of a ripplonic frequency from its classical value scales $`_i\mathrm{\Delta }_i^2/\mathrm{}\omega _i`$. The width of the ripplonic band of order $`l`$ is probably limited from above by the tunneling amplitude $`\mathrm{\Delta }_i`$ itself. Does this band broadening affect our previous results on the Boson Peak phenomena? Not very much. Since the observables depend mostly on the number of new excitations and the number of the ripplonic modes is not changed by these mixing effects, the essential core of our conclusions from Section IV remains intact. Nevertheless, some quantitative modifications are to be expected. For example, the lowest ripplonic energies may be lowered to the extent so as to cause a cross-over to a multi-level behavior in some of the internal resonances, thus possibly modifying the derived magnitude of the heat capacity and phonon scattering at sub-plateau temperatures. This effect will further contribute to the phonon interaction induced broadening of the ripplonic transitions, as estimated in Section IV.
### V.2 Mosaic Stiffening and Temperature Evolution of the Boson Peak
Eq.(59) raises another interesting point. According to that equation, the values of both the bare and the effective classical energy bias of a transition - $`ฯต`$ and $`\stackrel{~}{ฯต}`$ respectively - are limited from above by the lowest ripplon frequency ($`\omega _2`$). (Note that this is only realized in the $`ฯต<0`$ case, discussed in this section.) This is unimportant at low temperatures. But what happens at higher $`T`$, near this limit? Unlike in the low energy situation just discussed, one simply cannot ignore here that all the energy states have a rather short life-time. Therefore the singularity in Eq.(59) does not occur, but will be rounded. This observation does not completely answer the question that one should have asked in the first place on general grounds alone: what happens to the structure of the energy spectrum of a tunneling center, when the energy of the transition becomes comparable to a vibrational eigen-frequency of the domain wall<sup>7</sup><sup>7</sup>7We remind the reader that the tunneling transition energy could be also thought of as an eigen-energy of the wallโs motion, but of a lower, $`l=1`$ order, associated with the translational motion of the shellโs center of mass?
When attempting to answer this question, a general multi-level perspective on each tunneling center is somewhat easier to use than the very mechanical view of the wallโs excitations that we have mostly employed so far, in which the ripplonic energy states are obtained by quantizing vibrations of a freely moving classical mambrane. The โsingularityโ at $`|\stackrel{~}{ฯต}|\mathrm{}\omega _i`$ is actually a term-crossing phenomenon that, again, would not take place in the strict classical limit. Let us go back to our argument on the density of states, but consider a case when $`ฯต`$ is larger than a ripplonic frequency. As mentioned many times already, vibrational excitations of a domain wall can be defined meaningfully only when a structural transition takes place in a given region of the material. The energy of the transition must be the lowest excited state of a mosaic cell. On the other hand, the values of the ripplon frequencies are determined by a (fixed) surface tension coefficient and the wallโs mass density. They have fixed values. The necessary conclusion from this is that the tunneling centers will not have ripplons whose frequency is lower than the transition frequency. We provide a cartoon illustrating this idea in Fig.23.
We see the quantum mixing reduces the number of the lower frequency vibrational modes. The mosaic appears stiffer than expected. This effect may contribute to the temperature evolution of the Boson Peak as observed in inelastic scattering experiments. Wischnewski et al. find Wischnewski *et al.* (1998) that at temperatures between 51 K (numerically close to silicaโs $`\mathrm{}\omega _2`$) and above the glass transition, the left hand side of the Boson Peak decreases in size as the temperature was raised. At the same time, the high frequency side remained relatively unchanged. Note that, as temperature is raised, the total area of the peak in Fig.2 of Ref.Wischnewski *et al.* (1998) does not increase. In this temperature range mosaic cell motion loses oscillating character and becomes a rather featureless activated relaxation process.
To summarize this section, we have seen that the possibility of quantum tunneling between structurally close states in glass does have a predictable effect on the spectrum and must be taken into account when computing the density of low (and not so low) energy structural excitations in these materials. At the same time, the main conclusions of the original semi-classical argument remain valid: each structural transition may be thought of as a rearrangement of about 200 molecules accompanied by distortion of the domain wall that separates the two alternative local atomic arragements.
## VI The Negative Grรผneisen Parameter: an Elastic Casimir Effect?
With the exception of the plateauโs position and the quantum mixing effects, we have so far dealt with those anomalies in low temperature glasses that are more or less universal. These universal patterns are of particular interest because they cannot be easily blamed on chemical peculiarities of each substance. Indeed, given the flatness of the low energy excitation spectrum in glasses, the apparent universal ratio $`l_{\text{mfp}}/\lambda 150`$ is the dimensionless quantity that seems to express the general, intrinsic character of those low energy excitations, as arising from the non-equilibrium nature of the glass transition. The number $`150`$ reflects the size of nearly independent fragments into which a supecooled liquid is broken up at the laboratory glass transition. Yet, there is another dimensionless quantity, namely the Grรผneisen parameter $`\gamma `$, that also reflects the necessity of going beyond a harmonic picture for amorphous solids. This parameter is always a positive number of order one for simple cubic crystals (at low enough $`T`$), but varies wildly among amorphous materials Ackerman *et al.* (1984) (see also a discussion in Leggett (1991)). $`\gamma `$ in glasses has been reported to be as large as several tens and often negative in sign! A negative $`\gamma `$ implies a negative thermal expansion coefficient $`\frac{1}{V}\left(\frac{V}{T}\right)_p`$ (the linear expansion coefficient $`\alpha \frac{1}{L}\left(\frac{L}{T}\right)_p=\frac{1}{3}\frac{1}{V}\left(\frac{V}{T}\right)_p`$ is a commonly used quantity chracterizing anharmonicity too). Contraction with heating is observed in some crystals at not too low temperatures, owing to the details of the anharmonic couplings in a specific substance that may result in the negativity of the Grรผneisen parameter of a lattice mode of a finite frequency (see e.g. Wei *et al.* (1994)). Thermal contraction along a single direction in anisotropic materials is even more common. Nonetheless, as the temperature is lowered, the thermal expansion coefficient in an insulating crystal eventually becomes positive and approaches the cubic $`T`$ dependence predicted by standard thermodynamics. In contrast, an isotropic negative thermal expansivity is observed in many amorphous substances even at the lowest temperatures. In addition, the expansivity is not cubic in $`T`$. The most widely known example of a substance with a negative $`\alpha `$ is rubber. Rubber owes this property to the largely entropic nature of its elasticity. Here, we will see that a distinct mechanism of thermal contraction in glasses in the TLS temperature range arises, which is a direct consequence of the existence of the spatially extended tunneling centers that give rise to the universal phenomena considered earlier.
As shown above, the excitation spectrum of the tunneling centers may be represented as a combination of the two lowest energy levels, corresponding to the structural transition and a set of higher energy states involving vibrations of the moving domain wall. By the exchange of phonons, these local (quantum) fluctuations in the elastic stress will be attracted to each other much like in the Van der Waals interaction between neutral molecules. The elastic Casimir effect seems a more appropriate name for this phenomenon, since the moving domain walls are not point-like but, instead, resemble fluctuating membranes. While we do not claim this attraction is solely responsible for the negative expansion coefficient, it turns out to provide a large contribution to the thermal contraction in glasses. We will see how this effect arising from interaction of amorphous state excitations depends on the material constants and the preparation speed of the glass is derived and, therefore, is not universal!
We note first that not all amorphous substances actually exhibit a negative $`\alpha `$ in the experimentally probed temperature range. In such cases, it is likely that the contraction coming from those interactions in these materials is simply weaker than the regular, anharmonic lattice thermal expansion. Other contributions to the Grรผneisen parameter will be discussed below as well.
Coupling the motion of the mosaic cell (TLS and Boson Peak) to phonons is necesssary to explain thermal conductivity, therefore the interaction effects discussed below follow from our identification of the origin of amorphous state excitations. The emission of a phonon followed by its absorption by another cell will give an effective interaction, in the same way that photon exchange leads to inter-particle interactions in QED. The longest range coupling between local degrees of freedom coupled linearly to the elastic stress has the form of a dipole-dipole interaction. Since the structural transitions are of finite size, the dipole assumption is only approximate for the closer centers. For the time being, we take for granted that there is no first order, static, interaction between the vibrating domain walls, which, if non-zero, could be ร priori of either sign. The next, second order interaction is always negative in sign and is proportional to $`_{ij}\frac{1}{r_{ij}^6}\left(1\frac{\delta V}{V}\right)^21+2\frac{\delta V}{V}`$. This favors a sampleโs contraction ($`V`$ is the volume). This attractive force, which will be temperature dependent, is balanced by the regular temperature independent elastic energy of the lattice: $`F_{elast}/V=\frac{K}{2}\left(\frac{\delta V}{V}\right)^2`$. Calculating the equilibrium volume from this balance allows us to estimate the thermal expansion coefficient $`\alpha `$. More specifically, the simplest Hamiltonian describing two local resonances that interact off-diagonally is $`H=\frac{\omega _i}{2}\sigma _x^i+\frac{\omega _j}{2}\sigma _x^j+J_{ij}\sigma _z^i\sigma _z^j`$, where $`\omega _i`$ and $`\omega _j`$ would be the frequencies of ripplons on sites $`i`$ and $`j`$ and
$$J_{ij}\frac{3}{4\pi \rho c_s^2}\frac{(๐ _i๐ _j)3(๐ _i๐ซ_{ij})(๐ _j๐ซ_{ij})/r_{ij}^2}{r_{ij}^3}$$
(70)
is the dipole-dipole interaction following from Eqs.(15) and (16). (Having the interaction be off-diagonal automatically removes the first order term in $`J_{ij}`$.) The factor $`3`$ accounts in our usual simplistic way for all three acoustic phonon branches. This ignores a distinction between the longitudinal and transverse speed of sound. This simplification is however accurate enough for our purposes. Since $`g\sqrt{\rho c_s^2a^3k_BT_g}`$, the $`J_{ij}`$โs turn out to scale in a very simple way with the glass transition temperature and the molecular size $`a`$, giving $`J_{ij}k_BT_g\left(\frac{a}{r}\right)^3`$.
Since only mobile domain walls give rise to local dynamic heterogeneities, one may conclude intuitively that only the sites of thermally active structural transitions can contribute to $`\alpha `$. Therefore one expects that as temperature in increased, more tunneling centers will contribute to the Van der Waals attraction thus leading to negative expansivity. As already mentioned, the excitations of a tunneling centers are conveniently subdivided into a low energy TLS-like pair of states, and higher energy, โripplonicโ excitations corresponding to distortions of an active centerโs domain wall. Hence we may view the total Van der Waals attraction as having three somewhat distinct contributions: โTLS-TLSโ, โripplon-ripplonโ and โTLS-ripplonโ attractions. In this section, we focus on the relatively low, subplateau temperature regime, for reasons that will be explained later. At these low temperatures, transitions to the ripplonic states are only virtual, whereas the TLS structural may well be thermally active. This, in addition to the differences in the respective spectra of these excitations, will lead to some difference in the dependence of the mutual interactions between those excitations on temperature and other parameters. In order to assess the magnitude of those interactions let us consider the following, very simple, three-level Hamiltonian that is designed to model a transition of energy $`ฯต_i`$ between two different structures that may also be accompanied by a wall vibration of frequency $`\omega _i`$:
$$H_i=\left(\begin{array}{ccc}0& 0& 0\\ 0& ฯต_i& 0\\ 0& 0& ฯต_i+\omega _i\end{array}\right)$$
(71)
Note that, even though, for simplicityโs sake, we use the semiclassical energy $`ฯต`$ in the Hamiltonian above, the latter is meant as (the lowest energy portion of) the full, diagonalized Hamiltonian with quantum corrections included. This corresponds to the plain two-level system formalism that does not specify a distribution of the tunneling matrix element $`\mathrm{\Delta }`$. Also, in comparison with the general case of Eq.(35), we only include an excitation by a single quantum of a single ripplon. The latter simplification is obviously justified in the lowest perturbation order, where all pairs of excitations contribute to the total in an additive fashion. Considering only single-quantum excitations is a low temperature approximation, made mostly to avoid adopting extra modelling assumptions necessary to embody the mixed spin/boson statistics on each site. This simplification is nevertheless adequate, as will become clear later in the discussion.
Since the contributions of the three constituents of the Van der Waals attraction are additive, one can consider each contribution separately. This indeed proves to be convenient not only because all the contributions exhibit distinct scaling with the parameters, but each contribution comes to dominate the expansivity at somewhat distinct temperatures. We consider first the ripplon-ripplon attraction. This contribution appears to dominate the most studied region around 1 K. The off-diagonal (flip-flop) interaction between the ripplons has the form:
$$H_{ij}^{int}=J_{ij}|2_i3_j3_i2_j|+H.C.,$$
(72)
where the rows and columns in the unperturbed Hamiltonian from Eq.(71) are numbered in the conventional way from the upper left corner. The โripplon-ripplonโ case appears the simplest of the three because here, the issue of how many tunneling centers contribute to the effect is more or less separate from the strength of the interaction. The former is (qualitatively) determined by the number of thermally active two-level systems, that scales roughly with the heat capacity. The latter is nothing but the ground state lowering of a pair of resonances after interaction is switched on, which scales as $`J_{ij}^2/(\omega _i+\omega _j)`$ and is $`T`$-independent at these low temperatures. This contribution to the negative thermal expansion is therefore expected to be roughly quadratic in temperature (this corresponds to linear expansivity), which is similar, if not somewhat slower than observed in amorphous silica around 1 degree K.
Calculating the correction to the systemโs free energy in the lowest order in $`J_{ij}`$, that corresponds to the interaction term from Eq.(72) is entirely straightforward and yields:
$`\delta F_{rr}`$ $`=`$ $`{\displaystyle \underset{ij}{}}J_{ij}^2{\displaystyle \frac{e^{\beta (ฯต_i+ฯต_j)}(1+e^{\beta \omega _i})(1+e^{\beta \omega _j})}{Z_iZ_j}}`$ (73)
$`\times `$ $`{\displaystyle \frac{\omega _i\mathrm{tanh}(\beta \omega _i/2)\omega _j\mathrm{tanh}(\beta \omega _j/2)}{\omega _i^2\omega _j^2}},`$
Here, $`Z_i1+e^{\beta ฯต_i}+e^{\beta (ฯต_i+\omega _i)}`$ is the unperturbed on-site partition function, corresponding to Eq.(71). Here, subscript โ$`rr`$โ signifies the โripplon-ripplonโ contribution.
At low - subplateau - temperatures $`T<\omega _i`$, that we are primarily interested in here, the expression above reduces to the following Van der Waals energy:
$$\delta F_{rr}=\underset{ij}{}\underset{l_1l_2}{}\frac{J_{ij}^2}{\omega _{l_1}^i+\omega _{l_2}^j}\frac{1}{(1+e^{\beta ฯต_i})(1+e^{\beta ฯต_j})},$$
(74)
where we have explicitly written out summation over distinct ripplon harmonics $`l_1`$ and $`l_2`$ at sites $`i`$ and $`j`$.
A few intermediate calculations are needed to compute the sum in Eq.(74). First, averaging of $`J_{ij}^2`$ with respect to different mutual orientations of $`๐ _i`$, $`๐ _j`$ and $`๐ซ_{ij}`$ yields an effective isotropic attractive interaction $`\frac{2}{3}\left(\frac{3}{4\pi }\right)^2T_g^2\left(\frac{a}{r}\right)^6`$. Second, the sum over all harmonics amounts to $`_{l_1,l_2=2}^{l_{max}}\frac{(2l_1+1)(2l_2+1)}{\omega _{l_1}+\omega _{l_2}}`$, where $`\omega _l`$ is found using the dispersion relation from Eq.(33). Here we assume that $`\omega _i`$โs are not correlated with $`J_{ij}`$ and $`ฯต_i`$. As we already know from the discussion in the previous section, the latter assumption is adequate for values $`ฯต`$ smaller than the Boson Peak frequency. Now, recall that $`l_{max}`$ actually depends on the dropletโs perimeter, thus introducing an additional (cubic!) scaling with $`\xi /a`$. In the end, the sum over the $`l`$โs is equal, within sufficient accuracy, to $`1.5\omega _D^1(3/4\pi )\pi ^3(\xi /a)^{5/4}(\xi /a)^3`$. Finally, assuming $`J_{ij}`$โs and $`ฯต`$โs to be uncorrelated enables one to present the double sum over $`ฯต_i`$ as a product of two identical sums: $`\left(_i(1+\beta ฯต_i)^1\right)^2`$. Each sum is the effective concentration of thermally active tunneling centers: $`k_B(\mathrm{ln}2)\frac{T}{T_g\xi ^3}`$ as computed by integrating $`1/(1+e^{\beta ฯต})`$ with the density of states from Eq.(34). Note that here we use the simple $`1/T_g\xi ^3`$ expression for density of the tunneling transitions, in keeping with the assumption $`Eฯต_i`$ of the plain two-level system model adopted in this section. This is reasonable, given the qualitative character of this calculation. Finally, the summation over the ripplon sites can now be reduced to an integration with the lower limit equal to $`\xi (3/4\pi )^{1/3}`$.
As a result of the previous discussion, one recovers the following expression for the energy gain (per volume) due to a volume change $`\delta V`$: $`\delta F_{rr}/V1.5\left(\mathrm{ln}2\right)^2\pi ^2\frac{k_BT^2}{\xi ^3T_D}\left(\frac{a}{\xi }\right)^{7/4}\left(\frac{\delta V}{V}\right)`$. This works against the regular elastic energy $`\delta F_{elast}/V=\frac{K}{2}\left(\frac{\delta V}{V}\right)^2`$, introduced earlier. The equilibrium relative change $`\delta V/V`$ as a function of $`T`$ is obtained by setting $`F/V=0`$. Differentiating the equilibrium value of $`\delta V`$ with respect to temperature yields the following estimate for the thermal (volume) expansion coefficient:
$$\frac{1}{V}\left(\frac{V}{T}\right)_p3.0\left(\mathrm{ln}2\right)^2\pi ^2\frac{1}{K}\frac{k_BT}{\xi ^3T_D}\left(\frac{a}{\xi }\right)^{7/4}.$$
(75)
This can already be used to estimate the magnitude of the ripplon-ripplon contribution to the โCasimirโ effect numerically. One can do it in several ways. The simplest thing to do that does not require knowing $`K`$, is simply to use Eq.(75) to calculate the Grรผneisen parameter $`\gamma `$ itself according to $`\gamma =(p/T)_V/c_V`$ Kittel (1956), also using $`(p/T)_V=(p/V)_T/(T/V)_p`$. This yields a temperature independent Grรผneisen parameter:
$$\gamma _{rr}3.0\left(\mathrm{ln}2\right)^2\pi ^2\frac{T_g}{T_D}\left(\frac{a}{\xi }\right)^{7/4}.$$
(76)
Using $`(\xi /a)^3200`$ and silicaโs $`T_g/T_D1500/350`$ one obtains $`\gamma 3.`$, within an order of magnitude of what is observed in amorphous silica at low temperatures (that experimental number varies between $`5`$ and $`20`$ among different kinds of silica at $`1`$ K and seems to grow larger with lowering the temperature, see Fig.3 from Ackerman *et al.* (1984)). We will argue shortly that this growth may be explained by other contributions to the attraction between local resonances.
We can also directly compare the conribution in Eq.(75) to the linear thermal expansion coefficient $`\alpha =\frac{1}{3V}\left(\frac{V}{T}\right)_p`$ for silica as measured in Ackerman *et al.* (1984). According to the Fig.2 from Ackerman *et al.* (1984), the $`\alpha `$ of silica is linear (possibly slightly sub-linear) in temperature and equals $`1.010^9K^1`$ at 1K. The compressibility $`K`$ was obtained in Ackerman *et al.* (1984) from measured speed of sound and density. For internal consistency, we use the scalar elasticity to estimate $`K`$ in this way. Summing up three single polarization phonon Hamiltonians from Eq.(15) yields $`K\rho c_s^2/3`$ (remember, $`\mathrm{\Delta }V/V=3\mathrm{\Delta }\varphi `$). Using silicaโs constants, given in Fig.15 and the earlier obtained $`\xi =20\AA `$ and recalling that $`\delta l/l=\delta V/3V`$, Eq.(75) gives linear expansion coefficient $`\alpha 0.410^9K^1`$ at 1K, indeed strongly suggesting that attraction between the tunneling centers is a significant contributor to the negativity of the expansion coefficient. The numbers just obtained are also a convenient benchmark in assessing other contributions to the negative thermal expansivity.
Next, we estimate the magnitude of the attraction between virtual transition and the direct, lowest energy transitions on different sites. The corresponding coupling term - $`J_{ij}|2_i2_j3_i1_j|+H.C.`$ \- leads to the following contribution to the free energy in the lowest order:
$`\delta F_{rT}`$ $`=`$ $`{\displaystyle \underset{ij}{}}J_{ij}^2{\displaystyle \frac{e^{\beta ฯต_i}(1+e^{\beta \omega _i})}{Z_i}}`$ (77)
$`\times `$ $`{\displaystyle \frac{\omega _i\mathrm{tanh}(\beta \omega _i/2)ฯต_j\mathrm{tanh}(\beta ฯต_j/2)}{\omega _i^2ฯต_j^2}},`$
At subplateau temperatures, when $`\beta \omega _i1`$, $`\mathrm{tanh}(\beta \omega _i/2)`$ can replaced by unity. Furthemore, the summation with respect to $`ฯต_j`$ is no longer cut off by the temperature and the respective integral (weighted by $`n(ฯต)=\frac{1}{T_g}e^{|ฯต|/T_g}`$) picks up most of its value at $`ฯตT`$. Therefore $`\mathrm{tanh}(\beta ฯต_j/2)`$ may be replaced by unity as well. (Actually, both of those replacements must be made simultaneously lest the sum becomes potentially ill-behaved when $`\omega _iฯต`$.) As a result, the expression in Eq.(77) simplifies:
$$\delta F_{rT}=\underset{ij}{}J_{ij}^2\frac{1}{1+e^{\beta ฯต_i}}\frac{1}{\omega _i+ฯต_j}.$$
(78)
The $`ฯต_j`$ integral is related to an exponential integral $`E_1`$ and yields in the two lowest orders: $`(\mathrm{ln}(T_g/\omega _i)\gamma _E)`$, where $`\gamma _E=0.577\mathrm{}`$ is the Euler constant. As in the previous calculation, we regard $`ฯต_i`$, $`\omega _i`$ and $`J_{ij}`$ as uncorrelated. The summation over $`\omega _i`$ can be approximately represented as a continuous integral between 0 and $`l_{max}`$ and leads to a quantity that scales as the area of the domain wall with a logarithmic correction. The final result is $`\delta F_{rT}/V=0.5\frac{T}{\xi ^3}(a/\xi )^4\mathrm{ln}\left[2.0\frac{T_g}{\omega _D}\left(\frac{\xi }{a}\right)^{1/4}\right](1+2\frac{\delta V}{V})`$. Up to a logarithmic correction, the expression is independent of the energy parameters in the problem and thus must scale linearly in $`T`$. Note that we have written out the full expression of $`\delta F_{rT}/V`$ that includes the bigger, $`\delta V`$ independent term โ$`1`$โ, for the following reason: This larger negative term is linear in temperature, which apparently would lead to a non-zero (positive) entropy at $`T=0`$. This observation signals a breakdown of a perturbative picture of largely non-interacting two-level systems. For the sake of argument, let us estimate at what temperature this breakdown occurs we compare the magnitude of the $`\delta F_{rT}/V`$ term, assuming it is correct, to the free energy of non-interacting two-level systems per unit volume: $`\frac{dฯต}{T_g\xi ^3}e^{ฯต/T_g}[T\mathrm{ln}(1+e^{\beta ฯต})]`$, where we have appropriately chosen $`E=0`$ as the reference energy. The latter expression is equal to $`(\pi ^2/12)T^2/T_g\xi ^3`$ and becomes smaller (in absolute value) than the $`\delta F_{rT}/V`$ term at temperatures below $`10^3T_g`$. This temperature is actually less, but still within an order of magnitude from the lower end of the plateau, which is well within the empirical validity of the non-interacting two-level systems regime. Let us recall, however, that a perturbative expansion is an asymptotic one and therefore always overestimates the magnitude of a correction (we suspect that most of the error comes from the low $`ฯต`$ two-level systems). Therefore, a more accurate estimate would probably yield a break-down temperature lower in value than the estimate above. There is a reason to believe the โbreak-downโ temperature is just at the edge of the lowest temperatures routinely accessed in the experiments. This is suggested by several experiments such as on internal friction where deviations from the standard non-interacting two-level system picture have been seen (see, for example, a recent review by Pohl et al. Pohl *et al.* (2002)). In general, the effect of interaction between two-level systems could exhibit itself under several guises. One of those is an apparent gap in the excitation spectrum of the effective individual TLS. Such effects may have in fact been observed Thompson *et al.* (2000); Lasjaunas *et al.* (1978). The estimates above show thise effects are more likely to be observed in substances with a higher glass transition temperature, such amorphous silica, or, germania (GeO<sub>2</sub>). Note, however, that the effects of interaction on the apparent TLS spectrum must be separated from quantum effects of level repulsion on each sites, that we have considered in Section V. At any rate, the volume expansion coefficient, corresponding to the computed value of the ripplon-TLS term, is approximately equal to
$$\frac{1}{V}\left(\frac{V}{T}\right)_p1.0\frac{1}{\xi ^3K}\left(\frac{a}{\xi }\right)^4\mathrm{ln}\left[2.0\frac{T_g}{\omega _D}\left(\frac{\xi }{a}\right)^{1/4}\right].$$
(79)
Substituting the numerical values for a-SiO<sub>2</sub> in Eqs.(79) and (75) shows that at 1 K, the ratio of the ripplon-TLS contribution to the ripplon-ripplon term is about 1.2 - that is they contribute comparably to the โcontractionโ free energy at this temperature. However, since the ripplon-TLS $`\alpha `$ is temperature independent, it will dominate at subKelvin temperatures. The Grรผneisen parameterโs value corresponding to Eq.(79) is
$$\gamma _{rT}1.0\frac{T_g}{T}\left(\frac{a}{\xi }\right)^4\mathrm{ln}\left[2.0\frac{T_g}{\omega _D}\left(\frac{\xi }{a}\right)^{1/4}\right].$$
(80)
The ripplon-TLS term, as estimated here, therefore seems somewhat larger relative to the ripplon-ripplon term than seen in experiment, consistent with our earlier notion that it is somewhat overestimated. Still, qualitatively our estimates are consistent with the observed tendency of $`\gamma `$ to increase in magnitude, when the temperature is lowered. We point out that the results obtained above disregard possible effects of a specific distribution of $`\mathrm{\Delta }`$ that will influence the precise value of the coupling between phonons and tunneling centers.
Note that the heat capacity like expression reflecting the number of thermally active sites $`i`$ enters into the expressions from Eqs.(79) and (80) in a linear fashion. Therefore, in contrast to Eq.(76), the temperature dependence of expression (80) is expected to be largely independent of the exact $`T`$-scaling of the heat capacity. Therefore, according to Eq.(80), the Grรผneisen parameter should eventually scale as $`1/T`$ at low enough temperatures in all substances (however, unrealistically long observation times may be required to verify this prediction; see the discussion at the end of this section). And again, the apparent density of states of the tunneling centers may be modified at those low temperatures due to interaction effects (such as the Burin-Kagan Burin and Kagan (1996) effect).
According to Eqs.(76) and (80), the dimensionless contribution of the attractive forces between the tunneling centers can be expressed in a simple manner through the $`T_g/T_D`$ and $`T_g/T`$ ratios, as well as the relative size of the mosaic. Note that the effects of varying the quenching speed of the liquid on the number in Eqs.(76) and (80) add up. For instance, making the quenching faster will increase $`T_g`$ and decrease $`\xi `$. The ripplon-TLS term is especially convenient with regard to testing our results, because it is nearly insensitive to changes in the Debye temperature potentially induced by altering the speed of glass preparation.
Finally, we show that the second order coupling between direct tunneling transitions is subdominant to the already computed quantities. Consider an interaction of the form $`J_{ij}|1_i2_j2_i1_j|+H.C.`$. If one repeats simple-mindedly the steps leading to Eq.(73), one obtains the following simple expression for the free energy correction due to interaction between the underlying structural transitions:
$$\delta F_{TT}=\underset{ij}{}J_{ij}^2\frac{ฯต_i\mathrm{tanh}(\beta ฯต_i/2)ฯต_j\mathrm{tanh}(\beta ฯต_j/2)}{ฯต_i^2ฯต_j^2}.$$
(81)
Assuming, again, that $`J_{ij}`$ and $`ฯต_i`$โs are uncorrelated, the $`ฯต`$ summation can be performed via averaging with respect to the distribution from Eq.(34). One can show that the low temperature expansion of the expression above yields, within two leading terms, $`\delta F_{TT}/V=(2T_g/3\xi ^3)\left(\frac{a}{\xi }\right)^6[1+(\pi T/T_g)^2\mathrm{ln}(T_g/T)/3]`$. The $`T`$-independent term in itself is curious in that it is a contribution to the โvacuum energyโ of the lattice that is of purely glassy origin and is entirely due to the locality of the free energy landscape of a liquid. Indeed, as attested by its scaling with $`T_g/\xi ^3`$, this โvacuum energyโ contribution would disappear at the ideal glass transition at which the whole space is occupied by a non-extensive number of distinct aperiodic solutions of the free energy functional. However, this constant term will have no effect on the thermal expansion. The lowest order $`T`$-dependent term - $`T^2\mathrm{ln}T`$ \- actually has a slightly stronger temperature dependence than the ripplon-ripplon contribution, however the latter is larger by at least three orders of magnitude, mostly owing to the large number of ripplon modes. Apropos, we would like to stress again that the presence of vibrational modes of the (extended) mosaic walls is essential to the existence of the negative thermal expansivity effect that we just estimated. Therefore, while the present theory predicts that many (and most conspicuous) effects that distinguish amorphous lattices from crystals should be described well by a set of non-interacting two-level-like entities at cryogenic temperatures, the intrinsic multilevel character of the structural transitions, that follows from the present theory, in glasses exhibits itself even at these low energies in higher order perturbation theory.
To complete the discussion of the second order interaction between tunneling centers we note that the corresponding contribution to the heat capacity in the leading low $`T`$ term comes from the โripplon-TLSโ term and scales as $`T^{1+2\alpha }`$, where $`\alpha `$ is the anomalous exponent of the specific law. Within the approximation adopted in this section, $`\alpha =0`$. However it is easily seen that the magnitude of the interaction induced specific heat is down from the two-level system value by a factor of $`10.(a/\xi )^5(d_L/a)^210^{\mathrm{4..5}}`$ and therefore may be safely neglected.
We have so far considered the second order part of the induced interactions (square in $`J_{ij}^2`$, but forth order in $`g`$). There could be also, ร priori, lower order contributions - first order in $`g`$, and first order in $`J_{ij}`$. First, let us consider the term linear in $`J_{ij}`$, which also has to do with interaction, mediated by the phonons. If non-zero, it could be of either sign. In our case, it is identically zero for the following reason. It is known Yu and Leggett (1988); Neu *et al.* (1997), that the apparent TLSโs are only weakly interacting (one could also infer this implicitly from the smallness of the second order term that we have already estimated. The first order term, if non-zero, is comparable to the second order in a mean-field disordered system. The dipole-dipole $`1/r^3`$ interaction is long range and is indeed well described by the mean field). But we are dealing here with a non-polarized state, for which the first order term, linear in the average on-site magnetization, vanishes. In any case, even if the system were in a โferromagneticโ state, the first order term would still be only very weakly temperature dependent and thus would not contribute to the thermal contraction. Whether to consider such first order term non-zero or not is, to some degree, a matter of choice. If non-zero, it must be simply thought of as the effective Weiss-like field that is part of molecular field at each site. That field implies a hard gap of the order $`T_g`$ and indeed is negligible at low $`T`$. Yet, at low enough temperatures - microKelvins or so Neu *et al.* (1997), the phonon-mediated first order interaction between the tunneling centers may become important and one can no longer use the bare frozen-in values of the on-site TLS energies, but those determined by the interaction. In this regime an independent two-level system picture breaks down and more complicated renormalized excitations may begin to play a role Burin and Kagan (1996).
On the other hand, the other possible contribution to $`\alpha `$, a term linear in $`g`$ does not have to do with interactions between the anharmonic amorphous solid excitations but is due to the direct coupling of the tunneling centers with the phonons. This direct TLS-phonon interaction has so far been the main suspect Phillips (1973); Papoular (1972); Galperin *et al.* (1985) behind the anomalous thermal expansion properties of the glassses. This mechanism requires however the existence of a correlation Phillips (1973) (in our notation) between the on-site values $`g`$ and $`ฯต`$, or else between $`\mathrm{\Delta }`$ and $`\mathrm{\Delta }/\varphi _{ii}`$. In other words, the value of either classical or quantum splitting of a two-level system must be correlated with the way its energy changes when elastic stress is applied locally. The $`\mathrm{\Delta }`$ with $`\mathrm{\Delta }/\varphi _{ii}`$ correlation has been argued to make a small contribution relative to the $`g`$ versus $`ฯต`$ correlation because of the smallness of the value of the $`\mathrm{\Delta }`$โs for the majority of the thermally active TLS Phillips (1973). On the other hand, a correlation between $`g`$ and $`ฯต`$ could produce, in principle, both a negative or positive Grรผneisen parameter and therefore could explain, by itself, the observed variety of expansion anomalies in the low $`T`$ glasses. However, the degree of correlation between $`g`$ and $`ฯต`$ and its temperature dependence is not really known and has to be parametrized. The soft-potential model offers enough richness in behavior to accomodate two possible contributions - one dilating and the other contracting - to the sampleโs volume. In fact, Galperin et al. Galperin *et al.* (1985) suggest that those two types of the TLS may well be the two types of the tunneling centers that were postulated early on by Black and Halperin Black and Halperin (1977) in order to resolve the apparent discrepancy between the value of the TLS density $`\overline{P}`$ as deduced from the phonon scattering experiments and the equilibrium and time dependent heat capacity measurements. This, of course, could be checked experimentally by comparing the degree of the discrepancy in $`\overline{P}`$ and the sign of the thermal expansion coefficient in different substances. (We have shown in the previous section how the Black-Halperin paradox is, at least partially, explained by quantum corrections to the semi-classical landscape picture of structural transitions in glass.) With regard to the linear in $`g`$ effect, we suggest here a modification to the original argument of Phillips from Phillips (1973). According to Phillips (note some notational differences), $`|\gamma |=12g\alpha _0\mathrm{ln}2/\pi ^2k_BT`$ (he also assumed a linear heat capacity linear in $`T`$). Here, $`|\alpha _0|1`$ is an (unknown) coefficient that reflects the degree of correlation between $`g`$ and $`ฯต`$: $`ฯต_ig_i(ฯต_i)=\alpha _0ฯต_ig`$. $`\alpha _0=\pm 1`$ means complete correlation and $`\alpha _0=0`$ means no correlation. Now, due to symmetry, $`\alpha _0`$ must be odd power in $`ฯต`$, the dominant term being therefore linear (see the form of $`\alpha _0(ฯต)`$, somewhat cryptically mentioned as a remark of B. Halperin, at the end of Phillipsโ article). We must note, that although we have pretended, within our one-polarization phonon theory, that $`g`$ is a vector quantity, it is in reality a tensor, if the phonons are treated properly. The off-diagonal terms, corresponding to interaction with shear, will indeed be uncorrelated with $`ฯต`$ due to symmetry. However, the trace of the tensor, corresponding to coupling of the TLS to a uniform volume change could be, in principle, correlated with the energy of the transition. For example, it may happen that when the sample is locally dilated, the structural transitions in that region will require less energy to occur. At present, we do not have an argument in favor of or against such a correlation. Note, however that at the glass transition temperature, when the current arrangement of the defects freezes in, most structural transitions involve a thermal energy around $`T_g`$. On the other hand, the energy spliting $`ฯต`$ of the tunneling centers relevant at the cryogenic temperatures is significantly smaller. Informally speaking, relative to the thermal energy scale at $`T_g`$ all two-level systems with low splitting will feel the same to the phonons. Therefore, qualitatively, the correlation factor $`\alpha _0`$ should be at least a factor of $`ฯต/T_g`$ down from the largest value of one. Note, this coincides with the form $`\alpha _0(ฯต)ฯต`$, suggested by Halperin. Therefore, the contribution of TLS-phonon coupling to the thermal expansivity of Phillips (who left the issue of the degree of correlation open at the time) should be multiplied by a factor of $`T/T_g`$. This takes into account, in a very naive way, both the symmetry and our knowledge of the energy scales relevant at the moment of the tunneling centers formation. This modifies Phillipsโ result to yield
$$|\gamma |<12g\mathrm{ln}2/\pi ^2k_BT_g=\frac{12\mathrm{ln}2}{\pi ^2}\sqrt{\frac{\rho c_s^2a^3}{T_g}}=\frac{12\mathrm{ln}2}{\pi ^2}\frac{a}{d_L}.$$
(82)
The temperature independence of this contribution to the Grรผneisen constant is the main difference between Eq.(82) and the original calculation of Phillips. The numerical value of the expression should be nearly the same for all substances and is about 8. This suggests that the direct coupling to phonons is a potential contributor to the elastic Casimir effect at temperatures around 1 K. Remember, however, the sign of the expression in Eq.(82) is unknown and its numerical value of $`10^1`$ only provides an estimate from the above.
From the qualitative analysis in this section, we conclude, tentatively, that there are several contributions of comparable magnitude to the thermal expansion at low temperatures. Higher order effects may also be present. In this case, it may be more straightforward to estimate the interaction between ripplons as extended membranes without using a multipole expansion, as indeed is done when computing the regular Casimir force between extended plates.
The qualitative treatment above of the second-order interaction between the ripplons on different sites can be extended to higher temperatures as well. It is easily seen from Eq.(73) that an excitation of energy $`\omega _l`$ will contribute only $`\beta J_{ij}^2`$ at temperatures comparable to $`\omega _l`$ and above. Therefore one might expect that at the temperatures near the end of the plateau the ripplonic transitions become thermally saturated and this attractive mechanism becomes increasingly less important. The expression in Eq.(77), in contrast, is subject to thermal saturation to a lesser degree. Still, we have seen that its scaling with temperature is subdominant to the ripplon-ripplon term at temperatures above 1 K. Finally, we remind the reader about the effect of mosaic stiffening explained in the previous sections. This should also diminish the attraction between the tunneling centers, owing to a smaller number of resonant modes at the sites of centers thermally active at these higher temperatures. On the other hand, the usual anharmonic effects also become more significant at a higher $`T`$ leading to a turn-over in the temperatrure dependence of $`\alpha `$, as circumstantially supported by the old data on several materials cited in Ref.Krause and Kurkjian (1968). However, in order to assess this โcross-overโ temperature, one needs to know the magnitude of the regular thermal expansion due to the non-linearities of the lattice. This is something that would be extremely difficult to measure independently, because even a crystal with the same stoichiometry as the respective glass, is not guaranteed to have the same non-linearity. Direct computer simulation estimates of the Grรผneisen parameter, on the other hand, may be problematic due to the current difficulty of generating amorphous structures corresponding to realistic quenching rates. This is the main reason we have confined ourselves here to sub-plateau temperatures.
Finally, we note again that even at the low temperatures we have been discussing, not all glasses have been shown to exhibit a negative $`\alpha `$. According to our theory, however, the โCasimirโ contribution to $`\alpha `$ is negative and sub-linear in $`T`$, whereas the regular non-linear expansion coefficient is positive but only cubic in temperature. Therefore, there should be a (perhaps very low) temperature at which the Casimir force should dominate. Data for many substances, although still positive at the achieved degree of cooling, do extrapolate to negative values of $`\alpha `$ at finite temperatures. This is not the case, however, for all substances Ackerman *et al.* (1984). Even excluding the possibility of error in these difficult experiments, this is not necessarily inconsistent with our theory for the following reasons. As the temperature is lowered, it takes a long time (proportional to $`T^3`$) for the tunneling transitions to occur and appear thermally activated. For these same reasons, like the amorphous heat capacity, the direct interaction effect is time dependent at low temperatures. It may therefore take an excessively long time to actually observe the effects, discussed in this section, at very low temperatures, thus making it difficult to see a sign change in $`\alpha `$ for lattices with relatively large anharmonicity. Incidentally, this analysis predicts that the response of the length of an amorphous sample to a temperature change at sub-plateau temperatures must be time dependent (such time-dependence, acompanied by heat release, has been observed in polycrystalline NbTi Escher *et al.* (2000)). Since the interaction effect is quadratic in concentration, one expects qualitatively that the relative rate of the expansionโs time dependence should be twice that of the specific heat.
## VII Conclusions
In summary, this work elucidates the origin of the thermal phenomena observed in the amorphous materials at temperatures $`T_D/3`$ and below, down to the so far reached milliKelvins. The nature of these phenomena can be boiled down to the existence of excitations other than elastic strains of a stable lattice. The peculiarity of these excitations is exhibited most conspicuously in the following phenomena: The specific heat obeys a nearly linear dependence on the temperature at the lowest $`T`$, greatly exceeding the Debye contribution. At the same temperatures, the heat conductivity is nearly quadratic in $`T`$ and is universal if scaled in terms of the elastic constants. At higher temperatures ($`T_D/30`$), the density of these mysterious excitations grows considerably leading to enhanced phonon scattering and thus a plateau in the temperature dependence of the heat conductance. This increase in the density of states is also directly observed as the so called Boson Peak in the heat capacity data, as well as inelastic scattering experiments.
We have argued that the origin of these excitations is a necessary consequence of the non-equilibrium nature of the structural glass transition. This transition, not strictly being a phase transition at all in a regular equilibrium sense, occurs if the barriers for molecular motions in a supercooled liquid become so high as to prohibit any macroscopic shape changes in the material on the scales of hours and longer Xia and Wolynes (2000). The origin of these high barriers lies in a cooperative character of the molecular motions, which involve around $`200`$ molecules at the glass transition temperature. Unlike regular crystals, where the correlation between the molecular motions is rather long range, thus leading to the emergence of translational symmetry below solidification, the motions within the cooperative regions in a supercooled liquid, or entropic droplets, are only weakly correlated with their surrounding. In the language of the energy landscape paradigm, a crystal is a (possibly non-unique) ground state of the sample (thus the long-range correlation!), whereas a glass is caught in a high energy state, not being able to reach the true ground state for kinetic reasons. The respective dense energy spectrum at these energies exhibits itself in the existence of alternative mutually accessible conformational states of regions, or domains, of about $`200`$ molecules in size. It was argued that quantum transitions between these alternative states are the additional excitations observed in glasses at low temperatures. The knowledge of the spectral and spatial density of these excitations allowed us to estimate from first principles the magnitude of the observed linear specific heat. The relevant energy scale here is the glass transition temperature $`T_g`$ itself.
Stability requirements for the existence of these alternative conformational states at $`T_g`$ allowed us also to estimate the strength of their coupling to the regular lattice vibrations, which is determined by $`T_g`$, the material mass density and the speed of sound. This enabled us to understand the universality of the phonon scattering at the low temperatures.
The novelty of this picture is that we have established rather generally a multiparticle character of the tunneling events. This is counter-intuitive because, naively, the larger the number of particles involved in a tunneling event, the larger the tunneling mass is, and the harder the tunneling becomes. This is indeed the case for systems like disordered crystals or crystals with substitutional impurities, where the tunneling mass is that of an atom, and the barrier heights are determined by the energy of stretching a chemical bond by a molecular distance; this virtually excludes the possibility of tunneling. The existence of structural rearrangements in a macroscopically rigid system is a sign of the system being in a high energy state in which the available phase space is potentially macroscopically large. However, a decrease in this density of states for glass transitions occuring at a slower pace of quenching would result in the necessity to engage a larger number of atoms in these structural rearrangements. Transitions between the internal states of a domain involve only a very minor length change of each individual bond and atomic displacements not exceeding the Lindemann length, which is of the order one-tenth of the atomic length scale. It is not particularly beneficial to picture the tunneling events as individual atomic motions but rather as the motion of an interface between the alternative states of the domain. This domain wall is a quasi-particle of a sort, which has a low mass indeed: per molecule in the domain, it is only about one-hundredth of the atomic mass. The contributes to the ease of the tunneling events that are thermally relevant at cryogenic temperatures: These events are subject to only very mild potential variations and are possible, again, because the lattice is frozen-in in a high energy state.
The spatially extended character of the domain wall excitations along with their strongly anharmonic nature explains also higher temperature phenomena, such as the Boson peak and the plateau in the heat conductivity. By using our knowledge of the surface tension and the mass density of the domain wall we were able to calculate the energy spectrum of vibrational excitations of the active domain walls, or ripplons. This spectrum is in good agreement with the observed frequency of the Boson peak. The ripplonic excitations accompany the transitions between the domainโs internal states and thus are strongly coupled to the phonons. This has enabled us to understand the experimentally observed rapid drop in the phonon mean free path at the plateau temperatures. In addition, we have investigated the effects of phonon coupling on the spectrum of the ripplons. These spectral shifts scale with $`T_g`$ and seem to be the cause of the non-universal position of the plateau.
We have carrried out an analysis of the multi-level structure of the tunneling centers that goes beyond a semi-classical picture of the formation of those centers at the glass transition, that was primarily employed in this work. These effects exhibit themselves in a deviation of the heat capacity and conductivity from the nearly linear and quadratic laws respectively, that are predicted by the semi-classical theory.
A Van der Waals attraction between the domain walls undergoing tunneling motions was argued to contribute to the puzzling negative expansivity, observed in a number of low $`T`$ glasses.
Finally, we note that the conclusions of this work strictly apply only to glasses made by quenching a supercooled liquid. One may ask, nevertherless, to what extent the present results are pertinent to other types of disordered solids, such as โamorphousโ films made chemically or by vapor depositions, or, say, disordered crystals. Indeed, phenomena, reminiscent of real glasses, such as an excess density of states, are observed in many types of disordered materials, although they do not appear to be as universal as in true glasses (see, for example Pohl *et al.* (2002)). In this regard, we note that most of the phenomena discussed in the present work should indeed take place in other types of aperiodic structures. What makes quenched glasses special is the intrinsic character of their additional degrees of freedom that stems from the non-equilibrium nature of the glass transition. Since the characteristics of this transition (while not being a transition in a strict thermodynamic sense!) are nearly universal from substance to substance, many low (and not so low) temperature properties of all those substances can be understood within a unified approach.
## Acknowledgments
We thank J.Schmalian, A.Leggett and A.C.Anderson for helpful discussions. This work was supported by NSF grant CHE 0317017.
## Appendix A Rayleigh Scattering of the Phonons due to the Elastic Component of Ripplon-Phonon Interaction
In this Appendix, we present an argument on the strength of the phonon scattering due to the direct coupling with the ripplons via lattice distortions, but not due to the inelastic momentum absorbing transition in which the internal state of the domain changes. We thus consider phonon scattering processes which do obey selection rules and couple to the lattice strain only in the second and higher order. This scattering is of the Rayleigh type (and higher order) and occurs off the domain walls as localized modes. Importantly, we will use only derived quantities and no adjustable parameters in this estimate. We show here that, indeed, this absorption mechanism is not significant compared to the resonant scattering by the inelastic transitions between the internal states of a thermally active domain.
First, it proves handy to rederive the ripplon spectrum from Eq.(32) in the less general case $`\rho _g=0`$ (but non-zero pressure!). As argued in Section IV, the droplet wall is at equilibrium pressure $`p=\frac{3}{2}\frac{\sigma }{R}=\frac{3}{2}\frac{\sigma _0a^{1/2}}{R^{3/2}}`$. If the surface is distorted locally by $`\mathrm{\Omega }`$, this results in an extra force on this portion of the wall due to a changed curvature Morse and Feshbach (1953). The second Newtonโs law (as applied per unit area) yields then:
$$\frac{9}{8}\frac{\sigma }{R^2}\left[2+\frac{1}{\mathrm{sin}\theta }\frac{}{\theta }\left(\mathrm{sin}\theta \frac{\mathrm{\Omega }}{\theta }\right)+\frac{1}{\mathrm{sin}^2\theta }\frac{^2\mathrm{\Omega }}{\varphi ^2}\right]=\rho _W\frac{^2\mathrm{\Omega }}{^2t},$$
(83)
where $`\theta `$ and $`\varphi `$ are the usual polar and asimuth angular coordinates on the surface and we took into account the $`r`$ dependence of pressure. The equation above can be solved by a linear combination of the eigen-functions of angular momentum in 3D:
$$\chi \underset{lm}{}\mathrm{\Omega }_{lm}(t)Y_{lm}(\theta ,\varphi ),$$
(84)
$`Y_{lm}(\theta ,\varphi )`$ are the spherical Laplace functions ($`m=l\mathrm{..1}`$). Substituting a harmonic of $`l`$-th order in Eq.(83) yields the equation for $`\omega _l`$ derived in text as Eq.(32). We will absorb the $`9/8`$ factor into the definition of $`\sigma `$ in the rest of the Appendix.
A (fake) potential energy, yielding the equation of motion (83), is (c.f. the discussion of surface waves on a spherical liquid droplet in Landau and Lifshitz (1987)):
$$f_{surf}=\sigma ๐\varphi d(\mathrm{cos}\theta )\left\{(R+\mathrm{\Omega })^2+\frac{1}{2}\left[\left(\frac{\mathrm{\Omega }}{\theta }\right)^2+\frac{1}{\mathrm{sin}^2\theta }\left(\frac{\mathrm{\Omega }}{\varphi }\right)^2\right]\right\}.$$
(85)
Although varying Eq.(85) w.r.t. $`\mathrm{\Omega }`$ does produce the Eq.(83), note that it differs (by a factor of $`9/8`$!) from the original surface energy $`\sigma 4\pi r^2`$. The resulting error is sufficiently small for our purposes, however this subtlety may be worth thinking about as this could reveal an extra friction mechanism due to the wetting phenomenon and surface tension renormalization mentioned in our discussion of the random first order transition in Section II.
While the domain wall positions are not strictly tied to the atomic locations, they are tied to the lattice as a continuum and follow the lattice distortions. Let us employ our usual โscalarโ phonons descibed by Hamiltonian
$$H\text{ph}=d^3๐ซ\left[\frac{\pi ^2}{2\rho }+\frac{\rho c_s^2(\psi )^2}{2}\right],$$
(86)
where $`[\psi (๐ซ_1),\pi (๐ซ_2)]=i\mathrm{}\delta (๐ซ_1๐ซ_2)`$. The surface energy due to the presence of both $`\mathrm{\Omega }`$ and $`\psi `$ is:
$$H_{surf}=\sigma ๐\varphi d(\mathrm{cos}\theta )\left\{(R+[\psi \psi (r_i)]+\mathrm{\Omega })^2+\frac{1}{2}\left[\left(\frac{(\psi +\mathrm{\Omega })}{\theta }\right)^2+\frac{1}{\mathrm{sin}^2\theta }\left(\frac{(\psi +\mathrm{\Omega })}{\varphi }\right)^2\right]\right\},$$
(87)
where $`\psi `$ is taken on the sphere of radius $`R`$ with the center located at $`๐ซ_i`$. The potential energy in Eq.(87) thus provides an explicit form of phonon-ripplon interaction due to the liquid free energy functional solutions being imbedded in the real space.
If we expand the value of the displacement field $`\varphi `$ in terms of spherical harmonics according to $`\psi _{lm}๐\varphi d(\mathrm{cos}\theta )\psi (r=R)Y_{lm}^{}(\varphi ,\theta )`$, it is then possible to write down equations of motion for the ($`l,m`$)-components of both ripplon and phonon displacements:
$$\frac{^2\mathrm{\Omega }_{lm}}{t^2}+\omega _l^2\left(\mathrm{\Omega }_{lm}+\psi _{lm}\right)=0.$$
(88)
The equation of motion for the phonon field can be obtained e.g. from $`\ddot{\psi }=i[H\text{ph}+H_{surf},\pi /\rho ]`$ to yield:
$`\ddot{\psi }c_s^2\mathrm{\Delta }\psi `$ $`=`$ $`{\displaystyle \frac{\sigma }{\rho }}{\displaystyle _0^{2\pi }}d\varphi ^{}{\displaystyle _1^1}d(\mathrm{cos}\theta ^{}){\displaystyle }dr^{}\delta (r^{}R)\{2(R+[\psi (๐ซ^{})\psi (๐ซ_i)]){\displaystyle \frac{1}{\mathrm{sin}\theta ^{}}}\left(\mathrm{sin}\theta ^{}{\displaystyle \frac{\psi }{\theta ^{}}}\right)`$ (89)
$``$ $`{\displaystyle \frac{1}{\mathrm{sin}^2\theta ^{}}}\left({\displaystyle \frac{^2\psi }{\varphi ^2}}\right)+{\displaystyle \underset{lm}{}}\mathrm{\Omega }_{lm}[2+l(l+1)]Y_{lm}(\theta ^{},\varphi ^{})\}\delta (๐ซ๐ซ^{}).`$
The terms with $`\psi `$ on the r.h.s. serve only to modify the local elastic constants, and therefore give rise to the regular Rayleigh scattering, so we will ignore them from now on.
Equations (88-89) can be used to write down equations of motion for the retarded Greenโs functions, which are preferable due to their convenient analytical properties (see Zubarev (1960) for our conventions). We are interested in the systemโs response to โpluckingโ the latice at site $`๐ซ=\mathrm{๐}`$ at time zero, hence the choice of the Greenโs function corresponding to an operator $`X`$: $`i\theta (tt^{})[X(t),\psi (๐ซ=\mathrm{๐},t^{}=0)]`$. Eqs.(88-89), if rewritten for the corresponding Greenโs functions, will preserve except there will be an additional term $`\frac{1}{\rho }\delta (t)\delta ^3(๐ซ)`$, corresponding to the โpluckingโ event, in the r.h.s. of Eq.(89) (note also a change in units). Thus obtained equations are possible to rewrite in the Fourier space:
$$\omega ^2\stackrel{~}{\mathrm{\Omega }}_{lm}^i+\omega _l^2\left[\stackrel{~}{\mathrm{\Omega }}_{lm}^i+\stackrel{~}{\psi }_{lm}^i\right]=0$$
(90)
and
$$\omega ^2\stackrel{~}{\psi }_๐ค+c_s^2k^2\stackrel{~}{\psi }_๐ค=\underset{i}{}\frac{\sigma }{\rho }\underset{lm}{}\stackrel{~}{\mathrm{\Omega }}_{lm}^i[2+l(l+1)]\frac{e^{i\mathrm{๐ค๐ซ}_๐ข}}{2\pi ^2}Y_{lm}(๐ค/k)i^lj_l(kR)\frac{1}{(2\pi )^4\rho },$$
(91)
where $`\stackrel{~}{\psi }_{lm}^id^3๐ค\stackrel{~}{\psi }_๐คe^{i\mathrm{๐ค๐ซ}_i}(4\pi )i^lj_l(kR)Y_{lm}^{}(๐ค/k)`$ and we used the expansion of a plane wave in terms of the spherical harmonics: $`e^{i\mathrm{๐ค๐ซ}}=4\pi _{l=0}^{\mathrm{}}_{m=l}^li^lj_l(kr)Y_{lm}^{}(๐ค/k)Y_{lm}(๐ซ/r)`$. Here, $`j_l(x)\sqrt{\pi /2x}J_{l+1/2}(x)`$ is the spherical Bessel function, which scales as $`x^l`$ for small $`x`$, hence we see that the ripplonsโ coupling with the phonons is quadratic or higher order in $`๐ค`$ as the second harmonic is the lowest order term allowed. Modes $`l=0`$ and $`l=1`$ have the meaning of the dropletโs growth and translation respectively, as was discussed in Section IV.3. These modes are not covered by this Sectionโs formalism. Even though the theory as a whole could be thought of as a multipole expansion of a molecular cluster interacting with the rest of the lattice, the modes of different orders end up being described by different theories.
The system of Eqs. (90) and (91) can now be used to determine the sound dissipation due to the interaction with the ripplons. Since the system is infinite and has a continuous spectrum, all excitations will have finite life-times, which can be, in principle, obtained self-consistently by using e.g. the Feenbergโs perturbative expansion Feenberg (1948); Abou-Chacra *et al.* (1973) (one in the end arrives at Greenโs functions that are well behaved at infinity, as implied in the thus greatly simplified derivation). We do not have to do this self-consistent self-energy determination as long as we are interested in the lowest order estimate, as justified in the end by the smallness of the obtained value of the perturbation. Substituting Eq.(90) into Eq.(91) yields
$`\omega ^2\stackrel{~}{\psi }_๐ค+c_s^2k^2\stackrel{~}{\psi }_๐ค`$ $`=`$ $`(4\pi )^2{\displaystyle \frac{\sigma }{\rho }}{\displaystyle \underset{i}{}}{\displaystyle \underset{lm}{}}{\displaystyle \frac{[2+l(l+1)]\omega _l^2}{\omega _l^2\omega ^2}}`$ (92)
$`\times `$ $`{\displaystyle \frac{d^3๐ค_1}{(2\pi )^3}e^{i(๐ค_1๐ค)๐ซ_๐ข}(1)^lj_l(k_1R)j_l(kR)Y_{lm}(๐ค/k)Y_{lm}^{}(๐ค_1/k_1)\stackrel{~}{\psi }_{๐ค_1}}.`$
Since the spatial locations $`๐ซ_i`$ of active droplets are not correlated <sup>8</sup><sup>8</sup>8This is not strictly true - they, of course, can not be on top of each other., we can replace the summation over the droplets by a continuous integral, assuming at the same time that the ripplon frequency corresponding to $`\omega _l`$ varies from droplet to droplet within a (normalized) distribution $`๐ซ_l(\omega )`$ centered around $`\omega _l`$ and having a characteristic width $`\delta \omega _l`$, whose value will be discussed shortly. There is no reason to believe that the frequency and location of the tunneling centers are correlated, therefore one obtains
$`\omega ^2\stackrel{~}{\psi }_๐ค+c_s^2k^2\stackrel{~}{\psi }_๐ค`$ $`=`$ $`n{\displaystyle \frac{\sigma }{\rho }}{\displaystyle \underset{l}{}}{\displaystyle ๐\omega ^{}๐ซ_l(\omega ^{})\frac{4\pi [2+l(l+1)](2l+1)\omega ^2}{\omega ^2(\omega +iฯต)^2}j_l^2(kR)\stackrel{~}{\psi }_๐ค},`$ (93)
where $`n`$ is the concentration of the active domain walls to be estimated shortly and we have displaced $`\omega `$ by $`ฯต`$ into the upper half-plane because we are looking for the retarded Greenโs function. Also, in order to derive Eq.(93), we have used the summation theorem for the spherical functions $`P_l(\mathrm{๐ง๐ง}^{})=\frac{4\pi }{2l+1}_{m=1}^lY_{lm}^{}(๐ง^{})Y_{lm}(๐ง)`$, as well as $`P_l(1)=(1)^l`$, where $`P_l`$ is the Legendre polynomial. If we ignore the real part of the r.h.s. of Eq.(93), responsible only for the dispersion, the poles of the resultant phonon Greenโs function are found by solving $`\omega ^2c_s^2k^2+i\mathrm{\hspace{0.17em}2}\omega \tau _\omega ^1=0`$, where $`\tau _\omega ^1`$ clearly has the meaning of the inverse life-time of a phonon of frequency $`\omega `$ and is given by
$$\tau _\omega ^1=n\frac{\sigma }{\rho }\underset{l=2}{\overset{9}{}}\pi ^2[2+l(l+1)](2l+1)j_l^2(kR)๐ซ_l(\omega ),$$
(94)
where we have ignored the contribution of the peaks centered around $`(\omega _l)`$. We remind the reader that $`l_{max}9`$ is dictated by the finite size of a droplet.
One can find the value of $`\delta \omega _l`$ from an argument identical to the one used in Xia and Wolynes (2001a) to obtain the width of the distribution of the barriers for the droplet growth free energy profile. At the glass transition, a liquid breaks up into dynamically cooperative regions, so that a translation of one atom involves moving about $`200`$ atoms around it, which involves overcoming a large (on average) barrier. This barrierโs height is determined, together with the domain surface tension coefficient, by the configurational entropy density, which in its turn reflects the number of metastable states available to a particular volume of liquid at this temperature. Even though a good description of freezing is achieved by assuming that this number of available states does not strongly depend on where exactly on the free energy surface a particular molecular cluster is Xia and Wolynes (2000), it should vary from domain to domain. The size of the variation can be estimated from the known magnitude of the entropy fluctuations at constant energy, so that the ratio of the variance to the mean is related to the jump in the heat capacity at $`T_g`$ and subsequently turns out to be $`1/2\sqrt{D}`$ Xia and Wolynes (2001a), where $`D`$ is the liquidโs fragility, entering the Vogel-Fulcher law for relaxation times in a supercooled liquid $`\tau _{relaxation}e^{\frac{DT_K}{TT_K}}`$. We conclude then that the lower bound on the fluctuations of the ripplon frequency $`\omega _l`$ is given by $`\delta \omega _l\omega _l/2\sqrt{D}`$.
Lastly, in order to use Eq.(94) to compute the phonon absorption due to this particular mechanism, we need to estimate the density of the active domain walls. It will suffice for our purposes here to consider as active the defects that contribute to the specific heat, that is, roughly, $`n\frac{1}{\xi ^3}T/T_g`$. A more accurate estimate would be similar to the one we made when calculating the bump in the heat capacity in Section IV.3.
We are now ready to give a numerical estimate of the expression in Eq.(94). We will compute here the contribution of the $`l=2`$ term in the plateau region. It is convenient to represent $`kR`$ from Eq.(94) as $`kR\frac{\omega }{0.4(a/\xi )\omega _D}`$. For the reference, $`(a/\xi )T_D0.2T_D`$ is at the high temperature end of the plateau, whereas its middle is about an order of magnitude lower depending on the substance (see $`\kappa `$ vs. $`T/T_D`$ plot in Fig.1). We can now use our usual expressions connecting $`\sigma ,T_g,\omega _D,c_s,\rho ,a`$ etc to obtain a numerical estimate of Eq.(94) at the plateau frequencies $`\omega _{plateau}10^{1.5}\omega _D`$. Even if one favorably assumes that $`\omega _2\omega _{plateau}`$ (it is somewhat larger according to Section IV.4), one still gets $`l_{\text{mfp}}/\lambda >10^4`$ at the plateau frequency, whereas the resonant absorption by the TLS would give $`l_{\text{mfp}}/\lambda 10^2`$. The amplitude of this type of absorption is small due to the weakness of direct coupling to the ripplons for the processes not accompanied by a change in the domainโs internal state.
## Appendix B Frequency Cutoff in the Interaction Between the Tunneling Centers and the Linear Strain
As argued in Section III.2, the coupling of the tunneling transition to a phonon can be found from an additional energy cost of moving the molecules within the domain in the presence of a strain and is given by an integral over the dropletโs volume (we consider only longitudinal strain for simplicity):
$$g=\rho c_s^2_Vd^3๐ซ(\stackrel{}{\varphi })(๐),$$
(95)
where $`\rho c_s^2`$ is basically the elastic modulus, $`\stackrel{}{\varphi }`$ and $`๐`$ are elastic and inelastic components of the atomic displacements respectively. If the phononโs wave-length is much larger than $`\xi `$, the elastic component is constant throughout the integration region and the integral reduces to one over the dropletโs surface and thus the $`g`$ estimate obtained in text. Otherwise, one obtains:
$$g=\rho c_s^2\left\{_S๐๐๐(\stackrel{}{\varphi })_Vd^3๐ซ(๐)(\stackrel{}{\varphi })\right\}.$$
(96)
The volume integral will give a higher order term in $`k`$, so for now, we focus on the surface integral. The displacement due to the phonon is conveniently expanded it terms of the spherical waves: $`e^{i\mathrm{๐ค๐ซ}}=4\pi _{l=0}^{\mathrm{}}_{m=l}^li^lj_l(kr)Y_{lm}^{}(๐ค/k)Y_{lm}(๐ซ/r)`$. Since it is the first derivative with respect to $`๐ซ`$ that we are interested in, we only need the $`l=1`$ term from this expansion. The angular part contributes only to the overall constant, but it is the spherical function $`j_1(kr)`$ that sets the cut-off value of the wave-vector, above which the phonons do not produce significant linear uniform stress on the domain. In Fig.24, we plot the derivative $`j_1(x)/x`$ (or, rather, we plot the square of it, which enters into all the final expressions).
We see that it is not unreasonable to assume that only the phonons with $`kR<6`$ will exert an appreciable linear strain on the domain. $`kR=6`$ translates into $`\omega _c2.5(a/\xi )\omega _D`$.
While we are at it, we estimate the interaction of the domain with the higher order strain, at least due to the term (95), in the frequency region of interest. The next order term in the $`k`$ expansion in the surface integral from Eq.(96) has the same structure, but is scaled down from the linear term by a factor of $`kR`$. At the plateau frequencies $`\omega _D/30`$, $`kR<0.5`$ as immediately follows from the previous paragraph. While this is not a large number, it is not very small either. Therefore, this interaction term is of potential importance.
The volume integral in Eq.(96) produces a quadratic term, which is roughly equal to
$`(\stackrel{}{\varphi })_Vd^3๐ซ(๐๐ค)`$. We then proceed in a completely identical fashion to our earlier estimate of $`g`$. Assuming the diplacements within the droplet are random, one gets for the integral $`\frac{1}{4}\sqrt{N^{}}a^3d_Lk`$, where factor of $`1/4`$ comes about because the displacement is assumed do decrease from $`d_L`$ in the center of the droplet to zero at the edge Lubchenko and Wolynes (2001). This yields that this term becomes comparable to the linear one at frequencies $`\omega \omega _D\sqrt{(a/\xi )}\mathrm{\hspace{0.17em}4}/(6\pi ^2)^{1/2}0.4\omega _D`$ \- well beyond the high $`T`$ end of the plateu.
We must note, there are other sources of non-linearity in the system, such as the intrinsic anharmonicity of the molecular interactions present also in the corresponding crystals. While these issues are of potential importance to other problems, such as the Grรผneisen parameter, expression (95) only considers the lowest order harmonic interactions and thus does not account for this non-linear effect. We must note that if this non-linearity is significant, it could contribute to the non-univrsality of the plateau, in addition to the variation in $`T_g/\omega _D`$ ratio. It would be thus helpful to conduct an experiment comparing the thermal expansion of different glasses and see whether there is any correlation with the plateauโs location.
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# GRB 050509b: the elusive optical/nIR/mm afterglow of a short-duration GRBBased on observations taken with the 1.2 m Mercator telescope at Observatorio del Roque de los Muchachos, with the 1.5 m telescope at Observatorio de Sierra Nevada, with the 2.2 m and 3.5 m telescopes at the Centro Astronรณmico Hispano Alemรกn (CAHA) at Calar Alto (operated jointly by the Max-Planck Institut fรผr Astronomie and the Instituto de Astrofรญsica de Andalucรญa (CSIC)) and with the 6.0 m telescope at the Special Astrophysical Observatory in Russia.
## 1 Introduction
Gamma Ray Bursts (GRBs) are flashes of high energy ($``$1 keVโ10 GeV) photons (Fishman and Meegan 1995), occurring at cosmological distances. The detection of counterparts at other wavelengths for long duration GRBs, revealed their cosmological origin (see van Paradijs et al. 2000 for a review) and it is accepted nowadays that they are associated with the death of massive stars. On the other hand, the nature of short duration GRBs, a class that comprises about 25% of all events (Mazets et al. 1981; Kouveliotou et al. 1993), still remains a puzzle. No counterparts have been discovered so far, in spite of intense efforts in order to detect the optical, infrared and radio counterparts to several short, hard bursts (Kehoe et al. 2001; Gorosabel et al. 2002; Hurley et al. 2002; Klotz et al. 2003). A possible optical transient, related to GRB 000313 was proposed by Castro-Tirado et al. (2002) but a firm conclusion could not be established.
The Swift satellite (Gehrels et al. 2004) offers unique capabilities for the detection of GRBs thanks to its high sensitivity and imaging capabilities at $`\gamma `$-rays, X-rays and optical wavelengths. The short GRB 050509b was discovered by Swift/BAT detector on 9 May 2005. The burst started at 04:00:19.23 UT and lasted for $`40`$ ms, putting it in the โshort-durationโ class of GRBs. It had a fluence of (0.95 $`\pm `$ 0.25)$`\times 10^8`$ erg cm<sup>-2</sup> in the 15-150 keV range (Gehrels et al. 2005). The prompt dissemination (13.7 s) of the GRB position (Hurkett et al. 2005) enabled prompt responses of automated and robotic telescopes on ground, like ROTSE-III (Rykoff et al. 2005), RAPTOR (Wozniak et al. 2005) and BOOTES-1 (shown in this paper), although no prompt afterglow emission was detected. By the time when Swift slewed and started data acquisition (about 60 s after the event onset), a fading X-ray emission was detected by the Swift/XRT; this can be considered as the first clear detection of an afterglow in a short duration GRB (Kennea et al. 2005). This triggered a multiwavelength campaign at many observatories aimed at detecting the afterglow at other wavelengths, as in the case of the long duration GRB class. Here we report the results of the multiwavelength observations carried out by our group, from millimetre (mm) wavelengths to the optical bands.
## 2 Observations and data reduction
The BOOTES-1 very wide field camera, located at Estaciรณn de Sondeos Atmosfรฉricos (INTA-CEDEA) in Huelva (Spain), observed the region of the sky containing the Swift/BAT error box of GRB 050509b as part of its routine observing schedule (Castro-Tirado et al. 2004). A 30 s exposure started at 04:00:00 UT (19 s prior to the onset of the 40 ms short burst), with the following frame starting at 04:01:00 UT. A limiting (unfiltered, airmass 4.0) magnitude of 6.0 is derived for any prompt optical flash arising from this event.
ToO observations in the optical were triggered starting 0.53 hr after the event at the 1.2 m Mercator telescope (+ MEROPE CCD camera) at Observatorio del Roque de los Muchachos in La Palma (Spain). Subsequently, ToO observations were made at the 1.5 m telescope at Observatorio de Sierra Nevada in Granada (Spain) and at the 6.0 m BTA telescope (+ SCORPIO) at the Special Astrophysical Observatory (SAO-RAS) in Nizhnij-Arkhyz (Russia), and at the 2.2 m telescope (+BUSCA) at Calar Alto (Spain). Near infrared (nIR) observations were obtained at the 3.5 m telescope (+ OMEGA2000) at Calar Alto as part of the ALHAMBRA<sup>1</sup><sup>1</sup>1http://alhambra.iaa.es:8080 back-up programme. Table LABEL:tabla1 displays the observing log.
Additionally, mm observations were obtained at the Plateau de Bure Interferometer (PdBI) as part of our ToO programme. The PdBI observed the source on May 11 and May 13 (6 antennas compact D configuration) and May 16 (5 antennas compact D configuration). The data reduction was done with the standard CLIC and MAPPING software distributed by the Grenoble GILDAS group<sup>2</sup><sup>2</sup>2http://www.iram.es/IRAMFR/GS/gildas/gildas.html; the flux calibration is relative to the carbon star MWC349. Table LABEL:tabla2 displays the observing log.
In order to subtract in all of our optical and nIR images the contribution of the bright elliptical galaxy 9$`\stackrel{}{.}5`$ away from the XRT error box, we modeled it using the ELLIPSE routine under IRAF<sup>3</sup><sup>3</sup>3IRAF is distributed by the NOAO, which are operated by USRA, under cooperative agreement with the US NSF.. We used the residual to perform further analysis in the optical and nIR images. The optical field was calibrated using the field photometry provided by Henden (2005). The nIR images were calibrated using the 2MASS Catalogue.
## 3 Results and discussion
The main observational result is the lack of any variable optical/nIR/mm counterpart in our images, within the refined Swift/XRT error box (Gehrels et al. 2005), in spite of the intensive searches, in agreement with the upper limits reported by ROTSE-III and RAPTOR and Swift/UVOT (Gehrels et al. 2005) for the prompt optical emission and by PAIRITEL (Bloom et al. 2005) for the prompt nIR emission. Similarly, the deep upper limits reported at the Keck (Bloom et al. 2005; Cenko et al. 2005) and at the VLT (Hjorth et al. 2005a; Covino et al. 2005) are in agreement with our conclusions drawn from the deep 6.0BTA observations.
### 3.1 No optical/nIR/mm afterglow at all ?
The fact that the detected X-ray afterglow for GRB 050509b (Gehrels et al. 2005) is the faintest of all afterglows detected by the Swift/XRT so far may indicate that the density of the surrounding medium where the progenitor took place is much lower than the typical value of $``$ 1 cm<sup>-3</sup> derived for several long-duration GRBs. Assuming that the X-ray afterglow is caused by synchrotron emission, as in the case of the long-duration family, one should also expect some contribution at nIR and optical wavelengths, but a simple S<sub>X-ray</sub>/S<sub>optical</sub> scaling would predict prompt optical fluxes $``$ 10<sup>2</sup> (i.e. $``$ 5 magnitudes) fainter than the optical afterglows observed so far for the long-duration events. Moreover, the combined Swift/XRT and Chandra X-ray Observatory (CXO) observations (Patel et al. 2005) imply that the decay exponent $`\alpha `$ is 1.1 (Gehrels et al. 2005), i.e., an optical afterglow might have rapidly decayed in brightness with a similar power-law decay index for $`\nu _{\mathrm{opt}}`$ $`<`$ $`\nu _X`$ $`<`$ $`\nu _c`$ (Sari et al. 1998).
### 3.2 Is GRB 050509b at $`z`$ = 0.225 ?
GRB 050509b is located in the direction of the NSC J123610+285901 cluster of galaxies (Gal et al. 2003) at $`z`$ = 0.225 (see also Bloom et al. 2005). In fact, during the prompt search for the X-ray afterglow detected by Swift/XRT, CXO has detected the diffuse X-ray emission from the intracluster gas rather than the point source itself (Patel et al. 2005).
If this would be the case, taking into account the proximity (9$`\stackrel{}{.}5`$) of the X-ray error box to the elliptical galaxy 2MASX J12361286+2858580 (Fig. 1), the relationship to it cannot be discarded. It would be only $``$ 33 kpc in projection (3.563 kpc/<sup>โฒโฒ</sup>), as the angular size distance $`D_A`$ = 735 Mpc, considering a Hubble constant of H<sub>0</sub> = 71 km s<sup>-1</sup> Mpc<sup>-1</sup>, a matter density $`\mathrm{\Omega }_m`$ = 0.3, and a cosmological constant $`\mathrm{\Omega }_\mathrm{\Lambda }`$ = 0.7. Assuming this, the progenitor would have been originated in the halo of the elliptical galaxy and could favour a neutron-star merger origin (Goodman et al. 1987; Eichler et al. 1989; Narayan et al. 1992), a physical scenario that can explain a short-duration burst like GRB 050509b. We note that past searches for correlations between clusters of galaxies and GRBs did not reveal positive results (Hurley et al. 1997; Gorosabel & Castro-Tirado 1997).
Radio emission in (or close to) the center of this galaxy has been detected at the WSRT (van der Horst et al. 2005) although no emission lines are seen in the optical spectrum (Bloom et al. 2005). The restframe colours (and therefore the associated K-corrections) have been obtained based on the HyperZ code (Bolzonella et al. 2000). Fitting synthetic spectral energy distributions (SED) templates to the $`B`$ band magnitude (19.73 $`\pm `$ 0.08), derived from the Bok telescope image (Engelbracht & Eisenstein 2005) and from our own data and taking the nIR magnitudes from the 2MASS, produces $`UB`$ = $``$0.50 $`\pm `$ 0.20, $`BV`$ = 1.08 $`\pm `$ 0.20, $`VR`$ = 0.31 $`\pm `$ 0.20 and $`RI`$ =0.69 $`\pm `$ 0.20, (correcting all for the Galactic reddening E($`BV`$) = 0.019; Schlegel et al. 1998). At $`z`$ = 0.225, this is a rather luminous galaxy, with M<sub>B</sub> = $``$20.6 and M<sub>R</sub> = $``$22.2. According to HyperZ, the $`BRIJHK`$ band SED of the neighbour galaxy at $`z`$ = 0.225 favours ($`\chi ^2`$/d.o.f = 2.4) a moderately extincted ($`A_\mathrm{v}0.4`$ mag) galaxy harbouring an evolved dominant stellar population (age $`360`$ Myr). Figure 2 shows the SED of the elliptical galaxy.
In this scenario, an intriguing possibility arises if the event would be the result of a stellar collapse, similarly to the long duration GRBs. At such a redshift, an underlying type Ib/c SN similar to SN 1998bw/GRB 980425 (Galama et al. 1998) should have peaked at $`R`$ $``$ 21, about 20 days since the burst onset. An underlying Type Ia SN is also expected if the event is the result of the gravitational collapse of a C/O white dwarf into a neutron star (Dar & De Rรบjula 2004). Our optical limits in the R-band 18.5 days after the event onset imply that the peak flux of any underlying supernova should have been $``$ 3 magnitudes fainter than the one observed for the type Ib/c SN 1998bw/GRB 980425 (Galama et al. 1998), and 2.3 magnitudes fainter than a typical type Ia SN (Filippenko 1997 and references therein), in agreement with the VLT results (Hjorth et al. 2005a).
### 3.3 Is GRB 050509b at high redshift ?
It is also plausible that GRB 050509b has been originated at a redshift considerably higher than 0.225, i.e. it could lie in the S1 galaxy or in one of the fainter sources (S2-S6) detected within the XRT error box (Fig. 1) at unknown redshift (see Bloom et al. 2005) or it occurred in a much more distant object, at very high redshift. In the former case, extinction might have played a considerable role in order to hide optical variability in the first hours/days following the event. However, the lack of detection of a nIR transient in the observations presented here disfavours this argument. On the contrary, if GRB 050509b arose from a high-redshift host galaxy, it would have easily been beyond the limit of the optical/nIR telescopes, because of the Lyman $`\alpha `$ blanketing affecting the optical band.
## 4 Conclusions
We have shown multiwavelength observations of the short duration gamma-ray burst detected by Swift (GRB 050509b) between 0 s and $``$ 18.8 days after the event. No optical/nIR/mm afterglow emission has been detected, in spite of the reported Xโray afterglow detection by Swift few minutes after the event, confirming the elusiveness of the afterglow of the short duration events. The spectral energy distribution of the neighbouring, potential host galaxy, favours a system harbouring an evolved dominant stellar population (age $`360`$ Myr), unlike most long duration GRB host galaxies observed so far, i.e. thus giving support to a compact binary merger origin. Any underlying supernova that could be associated with this particular event should have been at least 3 magnitudes fainter than SN 1998bw and 2.3 magnitudes fainter than a type Ia SN.
GRB 050509b is the second short duration GRB which is detected by Swift/BAT after GRB 050202 (which occurred too close to the Sun and could not be properly followed-up), and the first one localised with high accuracy by Swift/XRT. In spite of CGRO/BATSE detecting about 1/4 of events belonging to the short duration class, Swift has only detected 2 (out of $``$ 40 events), most likely due to its softer threshold energy. However, thanks to its extraordinary repointing capabilities, the accurate localisations for future events and the corresponding multiwavelength follow-up, will shed more light on the origin of the short-duration GRBs.
###### Acknowledgements.
We acknowledge the comments from the anonymous referee, and C.W. Engelbracht and D.J. Eisenstein for making available their B-band image. This work is based partly on observations carried out with the IRAM Plateau de Bure Interferometer, supported by INSU/CNRS (France), MPG (Germany) and IGN (Spain); and on data products from the Two Micron All Sky Survey (2MASS), which is a joint project of the Univ. of Massachusetts and the IPAC/CalTech, funded by NASA and the NSF. The work was partially supported by the grant A3003206 of the Grant Agency of the Academy of Sciences of the Czech Republic, by the RFBR grant 04-02-16300a and by the Program of the Presidium of RAN entitled โNon-stationary Phenomena in Astronomy 2005โ and by the Spanish MEC programmes AYA2004-01515 and ESP2002-04124-C03-01.
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# 1 Introduction
## 1 Introduction
Recently, a domain wall solution to type IIB supergravity was discovered which explicitly breaks supersymmetry while retaining stability . This solution, dubbed โJanusโ after the Roman god with two faces, includes a dilaton and a 5-form. The dilaton develops a kink profile in this solution, approaching different boundary values on either side of the domain wall (hence, two faces). Non-perturbative stability was proven for Janus-type solutions in AdS<sub>d</sub> in . This extends to IIB with the mild assumption that KK modes on the $`S_5`$ factor of the geometry do not destabilize the solution. In , a gauge theory was proposed and investigated as an AdS/CFT dual of this solution. This gauge theory is ordinary $`๐ฉ=4`$ Yang-Mills with different gauge couplings on either side of a defect (or interface) that the boundary theory inherits from the bulk domain wall. It was found that in the field theory, partial SUSY could be restored by the addition of certain counter-terms localized to the defect at the cost of breaking the R-symmetry at least as far as SU(3).
We now propose a supergravity dual of this field theory, which we dub Super Janus. The uplift formulas of allow one to obtain solutions of IIB SUGRA in 10D from solutions of $`๐ฉ=8`$ gauged SUGRA in 5D. We restrict ourselves to the subsector of this theory which is invariant under the SU(3) R-symmetry of the field theory. Under this restriction we find 4 scalars. This subsector is well described by gauged $`๐ฉ=2`$ SUGRA with one hypermultiplet, where there is well-known formalism for solving the equations of motion for supersymmetric domain walls .
## 2 Review of Janus
### 2.1 Janus itself
The Janus ansatz is an AdS<sub>4</sub>-sliced domain wall supported by a dilaton and 5-form:
$`ds_{10}^2`$ $`=`$ $`e^{2U(\mu )}\left(g_{ij}dx^idx^j+d\mu ^2\right)+ds_{S^5}^2`$ (2.1)
$`\varphi `$ $`=`$ $`\varphi (\mu )`$ (2.2)
$`F_5`$ $`=`$ $`4\left(e^{5U(\mu )}d\mu \omega _{AdS_4}+\omega _{S^5}\right),`$ (2.3)
where $`g_{ij}`$ is the metric of AdS<sub>4</sub> with unit scale in an arbitrary slicing. With this ansatz, the IIB SUGRA equations of motion can be simplified to the following two equations:
$`\varphi ^{}(\mu )`$ $`=`$ $`ce^{3U(\mu )}`$ (2.4)
$`U^{}(\mu )`$ $`=`$ $`\sqrt{e^{2U}1+{\displaystyle \frac{c^2}{24}}e^{6U}},`$ (2.5)
where the constant, $`c`$, can be thought of as an (arbitrary) integration constant arising from the dilatonโs second order equation of motion:
$$_\mu \left(\sqrt{g}g^{\mu \nu }_\nu \varphi \right)=0.$$
(2.6)
The largest root (in terms of $`U`$) of the equation $`U^{}(\mu )=0`$ then determines the value of the warp factor on the domain wall, $`U_0`$, and the range of the angular coordinate $`\mu `$ is determined by integrating equation (2.5):
$$\mu _0=_{U(0)}^{\mathrm{}}๐U\sqrt{e^{2U}1+\frac{c^2}{24}e^{6U}}.$$
(2.7)
Note that undeformed $`AdS_5`$ has $`\mu _0=\pi /2`$ while Janus has $`\mu _0>\pi /2`$. There is a critical value of $`c`$ above which the geometry becomes singular:
$$c_{cr}=\frac{9}{4\sqrt{2}}.$$
(2.8)
For $`cc_{cr}`$, the zeros (in $`U`$) of equation (2.5) become complex. When this happens, it becomes impossible for the warp factor to have a turning point at the wall. Instead, the warp factor must be allowed to grow arbitrarily negative as you approach the wall, and the geometry becomes nakedly singular. This is discussed in greater detail in .
An alternative interpretation is that one is free to pick any negative value for $`U_0`$ (within a range we will see shortly). The constant $`c`$ is then determined by the vanishing of equation (2.5):
$$c=\sqrt{24e^{6U_0}(1e^{2U_0})},$$
(2.9)
and the value of $`\mu _0`$ is determined as before. From this perspective, as we decrease $`U_0`$ from zero, we find that $`c`$ increases rapidly from zero to a maximum of $`c_{cr}`$ at $`U_0=\mathrm{ln}\left(\frac{2}{\sqrt{3}}\right)`$. We must stop there, as we are interested in the *largest* root of equation (2.5). It would not do for the warp factor to have two turning points, so we must restrict our choice to $`U_0(\mathrm{ln}\left(\frac{2}{\sqrt{3}}\right),0)`$. While choosing the value of $`c`$ directly is more convenient for the study of Janus and its dual, when studying Super Janus it becomes more convenient to specify the behavior of the warp factor on the domain wall (for numerical integration purposes). We will find Super Janus possesses an even narrower critical range for the values of $`U_0`$. The nature of this critical range of parameters in Super Janus appears to be somewhat different from the critical range we find in Janus. It is also interesting to note that there is no obvious interpretation of this criticality in the dual of either theory. In a gauge theory with jumping coupling, there is no obvious reason why the absolute amount of the jump in coupling strength should cause a breakdown of any of the salient features required for AdS/CFT duality. We postpone deeper study of this for the future.
### 2.2 The Gauge theory dual of Janus
The boundary theory dual to Janus was constructed and examined in detail in . The essential point is that the Lagrangian is simply that of $`๐ฉ=4`$ Yang-Mills with a jump in the gauge coupling at the defect inherited from the bulk. Spatial dependence of the gauge coupling means integration by parts is no longer trivial, and consistency requires that one write the scalar field kinetic term as $`\frac{1}{2}X^ID^\mu D_\mu X^I`$, where $`D_\mu X^I=_\mu X^I+i[A_\mu ,X^I]`$. If we denote by $`^{}`$ the Lagrangian for $`๐ฉ=4`$ Yang-Mills with this modified scalar kinetic term normalized with the gauge coupling appearing only as $`1/g^2`$, we may write the action for the gauge theory dual of Janus as follows:
$$S_{Janus}=d^4x\left(1\gamma \epsilon (x_3)\right)^{}(x),$$
(2.10)
where $`\epsilon `$ is an odd step function, $`^{}`$ is understood to have the average gauge coupling ($`g^2=\overline{g}^2=\frac{1}{2}(g_+^2+g_{}^2`$), and
$$\gamma =\frac{g_+^2g_{}^2}{g_+^2+g_{}^2}.$$
(2.11)
The $`\pm `$ subscripts refer to the sign of $`x_3`$, i.e. on which side of the defect one is located.
### 2.3 Restoring SUSY to the field theory dual of Janus
In this subsection we briefly review the construction of partial interface SUSY from appendix A of . This is the field theory for which we will then construct a gravitational dual. In this construction was used as a proof that it was not possible for partial supersymmetry to creep back into the boundary gauge theory dual of the explicitly non-supersymmetric Janus solution. The construction does, however, produce a perfectly good, if somewhat peculiar, gauge theory.
The first step is to consider the $`๐ฉ=4`$ fields as being made up of one $`๐ฉ=1`$ vector multiplet and 3 $`๐ฉ=1`$ chiral multiplets. One then considers how the SUSY variations of the $`๐ฉ=1`$ Lagrangians are modified by the inclusion of a Janus-style spatially varying coupling constant (varying only in the $`x_3`$ direction which we will now denote as $`z`$). In practice, one treats the coupling as a continuous function and takes the limit as it approaches a step function. The chiral Lagrangian we consider is
$$=_\mu \varphi ^{}^\mu \varphi \frac{i}{2}\overline{\psi }\mathrm{\Gamma }^\mu _\mu \psi +F^{}F+W^{}F\frac{i}{2}W^{\prime \prime }\overline{\psi }P_+\psi ,$$
(2.12)
where $`P_\pm =(1\pm \mathrm{\Gamma }^5)`$, and all dependence on the spatially varying coupling constant is assumed to be in the superpotential, $`W`$. When considering the SUSY variation of this Lagrangian, one finds that there are now terms which cannot be written as a total derivative:
$$\delta i\sqrt{2}_zg\overline{\epsilon }\left(P_+\mathrm{\Gamma }^z\frac{\delta W_{}^{}{}_{}{}^{}}{\delta g}+P_{}\mathrm{\Gamma }^z\frac{\delta W^{}}{\delta g}\right)\psi .$$
(2.13)
Now suppose that we preserve 2 supercharges in the general spirit of supersymmetric defect conformal field theories . Using the projector condition $`\mathrm{\Pi }\epsilon =\epsilon `$, where
$$\mathrm{\Pi }=\frac{1+i\mathrm{\Gamma }^5\mathrm{\Gamma }^z}{2},$$
(2.14)
the offending terms in $`\delta `$ (2.13) can be rewritten as
$$\sqrt{2}_zg\overline{\epsilon }\left(P_{}\frac{\delta W^{}}{\delta g}+P_+\frac{\delta W^{}}{\delta g}\right)\psi .$$
(2.15)
This expression happens to be a SUSY variation of $`2_zg\mathrm{Im}\frac{\delta W}{\delta g}`$. Thus, we may restore partial susy to the chiral sector by subtracting this term from the Lagrangian, effectively adding a counterterm. This counterterm takes the form of a delta function in the limit where $`g`$ becomes a true step function, so we say the counterterm is localized on the defect. Interestingly, this prescription fails if one normalizes the Lagrangian such that the coupling appears as an overall $`1/g^2`$ or if one writes the scalar kinetic term as $`\varphi ^{}\mathrm{}\varphi `$.
For the vector multiplet, the analogous prescription fails for the โoppositeโ normalization where we normalize to put $`g`$ in the numerator. We are forced to normalize the vector Lagrangian as follows:
$$=\frac{1}{4g^2}F_{\mu \nu }^aF^{a\mu \nu }\frac{i}{2g^2}\overline{\lambda }^a\mathrm{\Gamma }^\mu D_\mu \lambda ^a+\frac{1}{2g^2}D^aD^a.$$
(2.16)
The offending terms in the SUSY variation of this Lagrangian may be written as a supersymmetric variation of $`_z\left(\frac{1}{4g^2}\right)\overline{\lambda }^a\mathrm{\Gamma }^5\lambda ^a`$, after again imposing the projector condition $`\mathrm{\Pi }\epsilon =\epsilon `$, with $`\mathrm{\Pi }`$ given in (2.14).
Attempting to reassemble these Lagrangians into an $`๐ฉ=4`$ multiplet one finds no further obstruction to supersymmetry from gauge covariant derivatives and Yukawa interactions. However, the counterterms treat the multiplets very differently, and, thus, the $`๐ฉ=1`$ gaugino can no longer be mixed with the fermions from the chiral multiplet. Therefore, the maximal R-symmetry such a theory may possess is SU(3). In the next section we construct the gravity dual of the theory with this maximal R-symmetry. Note that the boundary gauge theory only has two supercharges but still possesses defect conformal symmetry, so the SUSY preserved by the dual gravity theory will be $`๐ฉ=1`$.
## 3 $`๐ฉ=8`$ gauged SUGRA in 5D
The $`๐ฉ=8`$ gauged SUGRA theory has 42 scalars which comprise an $`E_{6(6)}/USp(8)`$ coset. These 42 scalars can be organized in terms of the SO(6) R-symmetry of $`๐ฉ=4`$ Yang-Mills by looking at the symmetric traceless component of various operators. Scalar masses provide $`6\times 6(20^{}+1)_s+15_a`$ which gives us 20 traceless scalars. Fermion masses provide $`4\times 410_s+6_a`$, and $`\overline{4}\times \overline{4}\overline{10}_s+\overline{6}_a`$, which gives us a 10 and $`\overline{10}`$, for 20 more scalars. The axion and dilaton round out our set of 42. Next, we truncate this to the SU(3) invariant subsector. This sector of the theory has been extensively studied in . The breaking pattern of our scalar reps as SO(6) breaks to SU(3) is
$`20^{}`$ $``$ $`6+\overline{6}+8`$ (3.17)
$`10`$ $``$ $`1+6+\overline{3}`$ (3.18)
$`\overline{10}`$ $``$ $`1+\overline{6}+3,`$ (3.19)
giving us an additional 2 scalars which are singlets under our R-symmetry. Thus, we find a total of 4 SU(3) invariant scalars.
These 4 scalar fields live on an $`\mathrm{SU}(2,1)/\mathrm{SU}(2)\times \mathrm{U}(1)`$ coset, and thus they can be naturally assembled into the scalar sector of the universal hypermultiplet of $`๐ฉ=2`$, 5D SUGRA. For this theory, there exists plentiful machinery for solving the BPS equations . Additionally, the work of showed that with only hypermultiplets present, a simple consistency condition guarantees that the BPS equations will solve the equations of motion. We present and solve numerically a supersymmetric domain wall ansatz which is smooth and displays the correct qualitative features to be a gravity dual of the supersymmetric version of the Janus boundary field theory, namely one scalar field develops a kink profile, and the warp factor turns around at the domain wall.
### 3.1 $`๐ฉ`$=2 Flow Equations
The scalar manifold of the $`๐ฉ=2`$ theory with a single hypermultiplet has 8 isometries to gauge. Which isometries are gauged determine a triplet of SU(2) Killing pre-potentials which in turn determine the scalar potential and the flow equations for the scalars. The most general form of the 8 Killing vectors and the 8 corresponding pre-potentials are given in . Our theory really lives in IIB in 10D, so we are forced to choose our gauging very carefully in order to match the fixed points of the scalar potential. From the 10D perspective, we expect a 2D plane of critical points where the 10D dilaton and axion can take any value, but the other scalars are fixed. Thus we are forced to pick a gauging for the $`๐ฉ=2`$ theory that will give us a plane of fixed points. As in the $`\mathrm{SU}(2)\times \mathrm{U}(1)`$ invariant subsector studied in , this requirement uniquely fixes the gauging (up to a global symmetry transformation). Additionally, we may identify the scalars parameterizing this plane as the 5D dilaton and axion. Unfortunately, the uplift formulas generally entangle the 5D dilaton/axion with the metric so that it is not straightforward to read off 10D behavior from 5D or to unambiguously identify 10D fields with 5D counterparts, even though in principle tells us how to uplift. Though complicated, the existence of an uplift is guaranteed by , which details an $`๐ฉ=2`$ truncation of the FGPW flow in $`๐ฉ=8`$ 5D SUGRA . The field content and gauge structure of our model can be embedded in that truncation.
We parameterize the universal hypermultiplet with the 4 scalars $`V,\sigma ,r,\alpha `$. This choice can be obtained from the standard parametrization used in by redefining their $`\theta ,\tau `$ fields as $`\theta =r\mathrm{sin}\alpha ,\tau =r\mathrm{cos}\alpha `$. The metric of the scalar manifold is
$$ds^2=\frac{1}{2V^2}dV^2+\frac{1}{2V^2}d\sigma ^2\frac{2r^2}{V^2}d\sigma d\alpha +\frac{2}{V}dr^2+\frac{2r^2}{V}\left(1+\frac{r^2}{V}\right)d\alpha ^2.$$
(3.20)
In the language of , the Killing vectors of the quaternionic manifold are reorganized into generators of SU(2,1). Constants $`\alpha _i`$ denote coefficients of SU(2) generators, and $`\beta `$ denotes the coefficient of the compact U(1) (we do not consider non-compact gaugings). With only a hypermultiplet, the gauged isometry must live in this $`\mathrm{SU}(2)\times \mathrm{U}(1)`$ subgroup. Our model differs from the โToy model with only a hypermultipletโ of only in the slicing of the domain wall, so the gauge and potential structure must be the same. In order to match the fixed point structure of the potential as discussed above, we gauge the isometry corresponding to the choice $`\beta =\alpha _3,\alpha _1=\alpha _2=0`$. If one chooses $`\beta |\stackrel{}{\alpha }|`$, then the fixed point structure consists of an isolated critical point, UV in nature if $`|\beta |<|\stackrel{}{\alpha }|`$ and IR if $`|\beta |>|\stackrel{}{\alpha }|`$. The $`\beta =\alpha _3`$ gauge choice corresponds to the constant shift of our angular scalar field $`\alpha \alpha +c`$. The critical point of the superpotential occurs at $`r=0`$ with $`V`$ and $`\sigma `$ free to take on any value. Thus, we identify $`V`$ and $`\sigma `$ as linear combinations of the 5D dilaton and axion. Since $`\sigma `$ never appears in the scalar potential, either at the fixed point or away from it, we can specify $`\sigma `$ as the axion. Furthermore, turning on $`r`$ and $`\alpha `$ corresponds to turning on 3-form flux dual to the SU(3) singlet operators coming from the $`10,\overline{10}`$, i.e. gaugino mass terms (modulus dual to $`r`$, phase dual to $`\alpha `$).
The pre-potential is written in terms of an SU(2) phase, $`Q^s`$, and a superpotential, $`W`$, in the following way:
$$P^r=\sqrt{\frac{3}{2}}WQ^r.$$
(3.21)
Our gauge choice gives us the superpotential
$$W=\left(1+\frac{r^2}{V}\right),$$
(3.22)
and the SU(2) phase
$$Q^s=\frac{1}{V+r^2}(2r\sqrt{V}\mathrm{sin}\alpha ,2r\sqrt{V}\mathrm{cos}\alpha ,Vr^2).$$
(3.23)
With this gauging and parameterization, the full scalar potential has a very simple form:
$$๐ฑ=6+\frac{3r^4}{V^2}\frac{3r^2}{V}.$$
(3.24)
Our AdS-sliced domain wall ansatz is:
$`ds^2`$ $`=`$ $`e^{2U(z)}ds_{AdS4}^2+dz^2,`$ (3.25)
$`V`$ $`=`$ $`V(z)`$ (3.26)
$`\sigma `$ $`=`$ $`\sigma (z)`$ (3.27)
$`r`$ $`=`$ $`r(z)`$ (3.28)
$`\alpha `$ $`=`$ $`\alpha (z).`$ (3.29)
We can now use the machinery of to calculate the flow equations for the warp factor and the hyper-scalars from the vanishing of the fermionic supersymmetry variations. The function,
$$\gamma =\sqrt{1\frac{\lambda ^2e^{2U}}{g^2W^2}},$$
(3.30)
is often used to simply express the resulting equations. Throughout this paper we use conventions such that lower-case Greek indices ($`\mu ,\nu `$) refer to bulk space-time, analogous lower-case latin indices (m,n) refer to space-time coordinates along the domain wall, lower-case Latin indices from later in the alphabet (r,sโฆ) refer to SU(2) structure, and upper-case Latin indices refer to hyperscalars or the associated quaternionic geometry. Quantities expressed as an SU(2) triplet (e.g. $`P^r`$) can be re-expressed as a $`2\times 2`$ matrix in the usual way:
$$P_i^ji(\sigma _r)_i^jP^r.$$
(3.31)
For a BPS solution we require the fermionic supersymmetry variations to vanish. In the following expressions, $`๐_\mu `$ is the total covariant derivative (including both gravity and gauge structure), $`f_X^{iA}`$ is the quaternionic vielbein, and $`๐ฉ^{iA}=\frac{\sqrt{6}}{4}f_X^{iA}K^X(q)`$, where $`K^X(q)`$ is the gauged Killing vector. The fermionic SUSY variations are then
$`\delta \psi _{\mu i}`$ $`=`$ $`๐_\mu \epsilon _i{\displaystyle \frac{i}{\sqrt{6}}}g\gamma _\mu P_i^j\epsilon _j,`$ (3.32)
$`\delta \zeta ^A`$ $`=`$ $`{\displaystyle \frac{i}{2}}f_X^{iA}\gamma ^\mu (_\mu q^X)\epsilon _ig๐ฉ^{iA}\epsilon _i.`$ (3.33)
We use the residual supersymmetry projector of :
$$i\gamma _5\epsilon _i=[A(r)Q_i^j+B(r)M_i^j]\epsilon _j,$$
(3.34)
where $`M^r`$ is an SU(2) phase which depends on the scalar fields and is orthogonal to $`Q^r`$ (i.e. $`Q^rM^r=0`$). This is the most general ansatz for residual supersymmetry consistent with an AdS-sliced domain wall. It was found in that the quantities $`A(r),B(r)`$ are constrained up to signs by consistency and integrability constraints. There is an additional consistency constraint on the choice of $`M^r`$ that is derived in :
$$[\theta ,_z\theta ]=\sqrt{\frac{2}{3}}g[\theta ,P],$$
(3.35)
and it should be noted that the conventions of set $`g=\sqrt{3/2}`$. Writing this matrix equation in terms of SU(2) triplets we find
$$\theta ^r_z\theta ^s\epsilon ^{rst}\sigma ^t=\sqrt{\frac{2}{3}}g\theta ^rP^s\epsilon ^{rst}\sigma ^t.$$
(3.36)
Once this consistency condition is satisfied, the BPS equations will tell us the evolution of the scalar fields and the geometry.
Integrability of the gravitino variation condition along the wall ($`\delta \psi _m=0`$) and transverse to the wall ($`\delta \psi _5=0`$) give two different expressions involving $`U^{}`$. These may be solved for $`U^{}`$ and the previously unknown function $`A`$:
$`U^{}`$ $`=`$ $`\pm \gamma (z)|gW|`$ (3.37)
$`A`$ $`=`$ $`\gamma (z),`$ (3.38)
where the sign choice in these two equations is correlated. Note that consistency of the projector equation (3.52) determines the magnitude but does not fix the sign of the function $`B`$:
$$B=\pm \sqrt{1\gamma ^2},$$
(3.39)
where this upper/lower sign choice is independent of that in $`A`$.
As long as we satisfy equation (3.36), vanishing of the hyperini variations will now give us the flow equations of the hyper scalars, according to the formulas of (equivalent expressions can be found in using somewhat different language, and a nice summary and dictionary between the two languages can be found in ). We first write down the general form of the flow equations with an arbitrary projector, $`M^r`$. We will later make a specific choice and prove its consistency for the Super Janus gauging.
Let $`q^X`$ denote the $`X^{\mathrm{th}}`$ hyper scalar, $`g`$ the gauge coupling, $`R_{XY}^t`$ the SU(2) curvature (see for the curvature formulas). The general flow equations are
$$q^X=3g\left(Ag^{XY}+2B\epsilon ^{rst}M^rQ^sR^{tXY}\right)_YW.$$
(3.40)
The signs of the $`A`$ and $`B`$ functions can be determined by demanding that the warp factor has a turning point at the domain wall (matching to Janus) and that the Killing spinor equation be continuous across the domain wall. The turning point condition gives
$`\gamma (0)`$ $`=`$ $`0`$ (3.41)
$`A`$ $`=`$ $`\mathrm{sgn}(z)\gamma .`$ (3.42)
Continuity of the Killing spinor tells us that $`B`$ does not flip sign at the wall; thus,
$$B=\sqrt{1\gamma ^2}.$$
(3.43)
This procedure is analogous to that followed in for the original Janus solution. In those papers, geodesic completeness of the space-time demanded that the warp factor be analytically continued as an even function of $`z`$ to the other side of the wall. (N.B: In those papers, the warp factor was denoted by $`A(z)`$ rather than $`U(z)`$). We believe that this symmetry of the original Janus solution was somehow accidental, as Super Janus is slightly asymmetric, and there is no trace of this symmetry in the dual of either theory.
Before fixing the projector phase, $`M^r`$, the complete set of flow equations for our parameterization is given by the following:
$`U^{}(z)`$ $`=`$ $`\pm gW\gamma =\pm \sqrt{g^2\left(1+{\displaystyle \frac{r^2}{V}}\right)^2\lambda ^2e^{2U}}`$ (3.44)
$`V^{}(z)`$ $`=`$ $`6gr[\pm r\gamma {\displaystyle \frac{1}{V+r^2}}(\sqrt{V}\sqrt{1\gamma ^2}(2M^3r\sqrt{V}+`$ (3.46)
$`(M^2\mathrm{cos}\alpha +M^1\mathrm{sin}\alpha )(r^2V)))]`$
$`\sigma ^{}(z)`$ $`=`$ $`{\displaystyle \frac{3gr\sqrt{(1\gamma ^2)}}{\sqrt{V}(r^2+V)}}(M^2(3r^4+5r^2V+2V^2)\mathrm{sin}\alpha `$ (3.48)
$`+2\mathrm{cos}\alpha (M^1(r^2V)^2+M^3r^3\sqrt{V}\mathrm{sin}\alpha ))`$
$`r^{}(z)`$ $`=`$ $`3gr(\gamma {\displaystyle \frac{r\sqrt{1\gamma ^2}}{(r^2+V)\sqrt{V}}}(2M^3r\sqrt{V}`$ (3.50)
$`+(M^2\mathrm{cos}\alpha +M^1\mathrm{sin}\alpha )(r^2V)))`$
$`\alpha ^{}(z)`$ $`=`$ $`{\displaystyle \frac{3gr\sqrt{(1\gamma ^2)}}{2\sqrt{V}(r^2+V)}}\left(\left(2M^1\mathrm{cos}\alpha 3M^2\mathrm{sin}\alpha \right)(r^2V)+M^3r\sqrt{V}\mathrm{sin}2\alpha \right).`$ (3.51)
The sign choices come from the projector function $`A`$, thus, we must choose the upper sign for $`z>0`$ and the lower sign for $`z<0`$. We again emphasize that this is consistent and smooth because Janus-like solutions require $`U^{}\gamma 0`$ as $`z0`$, and all terms that flip sign at the domain wall are linear in $`\gamma `$ and thus continuous.
We follow the strategy of , first we pick $`\alpha (0)=0`$ as part of our initial conditions, then an easy guess for our projector phase is
$$M^s=(0,Q^3,Q^2)=\frac{1}{V+r^2}(0,r^2+V,2r\sqrt{V}).$$
(3.52)
Now we must check equation (3.36) using $`\theta ^r=\gamma Q^r+\sqrt{1\gamma ^2}M^r`$. Since the Pauli matrices are linearly independent, we may drop them and regard (3.36) as three separate equations. Then, because $`\theta ^1=P^1=0`$, only the $`t=1`$ component is non-trivial. The action of the covariant derivative on $`\theta ^r`$ may be written as
$$_z\theta ^r=q_{}^{X}{}_{}{}^{}_X\theta ^r+U^{}_U\theta ^r.$$
(3.53)
With $`\alpha (0)=0`$ and the projector (3.52), we find $`\alpha ^{}=\sigma ^{}=0,`$ and
$`V^{}(z)`$ $`=`$ $`6g\left(\pm r^2\gamma +r\sqrt{V}\sqrt{1\gamma ^2}\right)`$ (3.54)
$`r^{}(z)`$ $`=`$ $`3g\left(r\gamma +{\displaystyle \frac{r^2}{\sqrt{V}}}\sqrt{1\gamma ^2}\right).`$ (3.55)
This greatly simplifies the task of checking consistency:
$`\theta ^2_z\theta ^3\theta ^3_z\theta ^2=\left(1+{\displaystyle \frac{r^2}{V}}\right)\sqrt{1\gamma ^2}`$ (3.56)
$`\theta ^2P^3\theta ^3P^2={\displaystyle \frac{1}{g}}\sqrt{{\displaystyle \frac{3}{2}}}\left(1+{\displaystyle \frac{r^2}{V}}\right)\sqrt{1\gamma ^2},`$ (3.57)
so we see that (3.36) is satisfied. Since $`M^r`$ has three components, and there are 3 independent consistency equations, this completely specifies the projector phase.
We now have a system of three coupled, non-linear, ordinary differential equations for the warp factor and two running scalars (3.44, 3.54, 3.55). We will solve these numerically in the next subsection.
If one were to choose purely the upper sign in equations (3.44, 3.54, 3.55), then one obtains the solution of , which is nakedly singular at the domain wall. Super Janus avoids this by the requirement that the warp factor have a turning point at the domain wall, which we enforce through a careful choice of initial conditions. We believe the initial conditions chosen in were such that they guaranteed a curvature singularity. Indeed, we will find that only a narrow range of parameter space allows a turning point.
### 3.2 Numerics
We enforce the turning point condition with a suitable choice of initial condition for $`V`$. Setting equation (3.44) to zero at $`z=0`$, we obtain
$$V(0)=\frac{gr^2(0)}{\pm \lambda e^{U_0}g}$$
(3.59)
The scalar field $`V`$ must be strictly positive<sup>1</sup><sup>1</sup>1Strictly speaking, $`V`$ must be strictly of one sign. In many theories, it plays the role of the volume of a Calabi-Yau manifold, so it is conventional to choose it positive. See for more details., so we must choose the plus sign, as there is no positive value of $`\lambda e^{U(0)}`$ that gives positive $`V`$ when the minus sign is chosen. With the plus sign choice, we must still require $`\lambda e^{U_0}>g`$ for positivity of $`V`$. An additional constraint on the initial condition for the warp factor comes from the requirement that the turning point be a minimum: $`U^{\prime \prime }>0.`$ We may easily calculate $`U^{\prime \prime }`$ using the BPS equations to replace $`r^{}`$, $`V^{}`$, and $`U^{}`$ as needed. The result is a surprisingly simple expression:
$$U^{\prime \prime }(z)=\lambda ^2e^{2U(z)}\frac{6g^2r^4}{V^2}\frac{6g^2r^2}{V}.$$
(3.60)
Imposing the boundary condition for $`V`$ gives us
$$U^{\prime \prime }(0)=\lambda e^{U_0}\left(6g5\lambda e^{U_0}\right)$$
(3.61)
Positivity of this expression requires $`\lambda e^{U(0)}<\frac{6}{5}g`$. Both constraints together limit us to a very narrow band:
$$1<\lambda e^{U(0)}<\frac{6}{5}g.$$
(3.62)
We may now numerically integrate equations (3.44, 3.54, 3.55). As long as our initial conditions are within the critical range, we find the scalar $`V`$ develops a profile very much like the 5D dilaton in Janus. By adjusting $`r(0)`$ (arbitrarily) and $`\lambda e^{U_0}`$ (within criticality), we may adjust the average value of $`V`$ as well as the split between the two different asymptotic values of $`V`$.
We now plot the numerical results for a typical choice of parameters: $`\lambda e^{U_0}=1.02,r_0=0.25`$.
### 3.3 Asymptotic behavior
Using our knowledge from the numerics of the asymptotic behavior of the fields $`V`$ and $`r`$ and the warp factor $`U`$, we can use the flow equations (3.54,3.55) to determine the subleading behavior for matching against the field theory dual. The asymptotic flow equations are
$`U^{}(z)`$ $`\underset{z\pm \mathrm{}}{}`$ $`\pm g`$ (3.63)
$`r^{}(z)`$ $`\underset{z\pm \mathrm{}}{}`$ $`3gr`$ (3.64)
$`V^{}(z)`$ $`\underset{z\pm \mathrm{}}{}`$ $`6g\left(\pm r^2+r\sqrt{V}{\displaystyle \frac{\lambda e^U}{1+\frac{r^2}{V}}}\right).`$ (3.65)
Setting $`g=1`$ as we did in the numerics, equations (3.63, 3.64) may be easily solved to yield
$`e^{U(z)}`$ $`\underset{z\pm \mathrm{}}{}`$ $`c_Ue^{\pm z}`$ (3.66)
$`r(z)`$ $`\underset{z\pm \mathrm{}}{}`$ $`c_re^{3z},`$ (3.67)
where we have left out a constant term in $`r(z)`$ because numerics show $`r\underset{z\pm \mathrm{}}{}0`$. Plugging these back into equation (3.65) and Taylor expanding in $`e^z`$ we find
$$V^{}6\left(c_r^2e^{6z}+c_rc_U\lambda \sqrt{V}e^{4z}\right),$$
(3.68)
which gives us
$$VV_\pm c_r\frac{3\lambda c_U}{2}\sqrt{V_\pm }e^{4z}.$$
(3.69)
This behavior is what we expect for $`V`$ being dual to a dimension 4 operator with source term and $`r`$ to a dimension 3 operator with no source term, as expected if $`V`$ is the 5D dilaton and $`r`$ is related to 3-form flux (or gaugino mass terms). In order to read off the vevs of the corresponding operators in the dual gauge theory we need to determine the two integration constants $`c_U`$ and $`c_r`$, which only get fixed by the IR behavior of the solution and hence at the moment can only be determined from our numerical solution.
## 4 Conclusion
We have found a numerical solution for a supersymmetric domain wall in 5D, $`๐ฉ=2`$ gauged supergravity which is supported by two hyperscalars. The BPS equations allow us to determine analytically that the asymptotic behavior of these hyperscalars is appropriate for the duals of a dimension 4 and a dimension 3 operator. The uplift of this solution to type IIB supergravity in 10D is the Super Janus solution, dual to the gauge theory of with partial SUSY restoring counter-terms. The 10D domain wall solution will be supported by dilaton, 5-form flux, and 3-form flux dual to gaugino mass.
The existence of uplifts to 10D is guaranteed by the embedding in of the $`๐ฉ=2`$ theory with a single hypermultiplet into the FGPW solution of the $`๐ฉ=8`$ theory, where the uplift formulas of apply. Unfortunately, to realize the uplift one must understand how the hyperscalars of the $`๐ฉ=2`$ theory sit inside the 27-bein of the $`E_{6(6)}/USp(8)`$ coset of the $`๐ฉ=8`$ theory, which is far from straightforward. Solving this would immediately give the uplift of Super Janus. Another potential route is to solve the 10D equations of motion directly with some suitable ansatz deforming the AdS$`{}_{5}{}^{}\times S^5`$ geometry to something asymptotically AdS<sub>5</sub> crossed with an internal space possessing only SU(3) isometry.
It would also be interesting to understand the critical range of parameters in Janus and Super Janus from the perspective of the boundary gauge theories. There is no obvious reason for the gauge theory to care about the value of the jump in the coupling constant. Finding the importance of this critical range in the gauge theory would shed light on the Janus and Super Janus solutions.
## Acknowledgements
We wish to thank D. Freedman, M. Schnabl, J. Distler, R. Kallosh, M. Zagerman, and E. Silverstein for useful comments and suggestions. Special thanks to G. DallโAgata for comments on an earlier draft. Additionally, ABC would like to thank R. Van de Water and A. OโBannon for advice on numerics and helpful discussions. Both ABC and AK are supported in part by DOE contract #DE-FG03-96-ER40956.
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# Nevanlinna theory for the difference operator
## 1 Introduction
<sup>1</sup><sup>1</sup>footnotetext: *Mathematics Subject Classification 2000:* Primary 30D35; Secondary 39A10, 39A12.<sup>2</sup><sup>2</sup>footnotetext: The research reported in this paper was supported in part by the Leverhulme Trust and by the Finnish Academy (grant number 204819).
Nevanlinnaโs theory of value distribution is concerned with the density of points where a meromorphic function takes a certain value in the complex plane. One of the early results in this area is a theorem by Picard which states that a non-constant entire function can omit at most one value. Nevanlinna offered a deep generalization of Picardโs theorem in the form of his second main theorem , which implies the defect relation:
$$\underset{a}{}\left(\delta (a,f)+\theta (a,f)\right)2$$
(1.1)
where the sum is taken over all points in the extended complex plane, $`f`$ is a non-constant meromorphic function and the quantities $`\delta (a,f)`$ and $`\theta (a,f)`$ are called the deficiency and the index of multiplicity of the value $`a`$, respectively (see Section 2.1). The defect relation (1.1) yields, for instance, Picardโs theorem as an immediate corollary. It also implies that the maximum number of totally ramified values is at most four for any meromorphic function.
The appearance of the ramification index $`\theta (a,f)`$ in the defect relation (1.1) means that the density of high-multiplicity $`a`$-points is relatively low for most $`a`$. Similarly in this paper it is shown that $`a`$-points appearing in pairs with constant separation are rare for finite-order meromorphic functions, unless the function in question is periodic with the same period as the separation. For instance, if $`f`$ is of finite order and not periodic with period $`c`$, then
$$\underset{a}{}\left(\delta (a,f)+\pi _c(a,f)\right)2$$
(1.2)
where the sum is taken over all points in the extended complex plane, and $`\pi _c(a,f)`$ is a measure of those $`a`$-points of $`f`$ which appear in pairs separated by the constant $`c`$ (in other words, those points $`z_0`$ where $`f(z_0)=a=f(z_0+c)`$, see Section 2.1 for the exact definition.) The sharpness of inequality (1.2) is shown by giving an example of a finite-order meromorphic function, which is not periodic with period $`c`$, satisfying $`_a\pi _c(a,f)=2`$.
The defect relation (1.1) follows by an analysis of the behavior of the derivative $`ff^{}`$ in the ramification term of the second main theorem. In what follows, (1.2) is obtained by proving a version of the second main theorem where the derivative of $`f`$ is replaced by the exact difference $`f\mathrm{\Delta }f=f(z+c)f(z)`$ of a meromorphic function. In the remainder of this paper difference analogues of Picardโs theorem and Nevanlinnaโs theorem on functions sharing five values are given. In addition, the sharpness of the obtained results is discussed with the help of examples, and an application to difference equations is presented.
## 2 Nevanlinna theory for exact differences
Before going into details of value distribution of exact differences we must first give a precise answer to the following question: What is the difference analogue of a point with high multiplicity? By a formal discretisation of the derivative function $`f^{}(z)`$ we obtain
$$\frac{f(z+c)f(z)}{c}=:\frac{\mathrm{\Delta }_cf}{c},$$
(2.1)
where $`c`$. As noted in the introduction, those $`a`$-points of $`f`$ where the derivative vanishes, called ramified points, play a special role in Nevanlinna theory. The discretisation (2.1) of $`f^{}(z)`$ suggests that $`a`$-points appearing in pairs separated by a fixed constant $`c`$ may have similar importance with respect to the operator $`\mathrm{\Delta }_c`$. This indeed turns out to be the case as seen in the following sections.
### 2.1 Lemma on the exact difference
We first briefly recall some of the basic definitions of Nevanlinna theory. We refer to for a comprehensive description of the value distribution theory. The Nevanlinna deficiency is defined as
$$\delta (a,f):=\underset{r\mathrm{}}{lim\; inf}\frac{m(r,a)}{T(r,f)},$$
where $`a`$ is in the extended complex plane, $`m(r,a)`$ is the Nevanlinna proximity function and $`T(r,f)`$ is the characteristic function of $`f`$. The ramification index is
$$\theta (a,f):=\underset{r\mathrm{}}{lim\; inf}\frac{N(r,a)\overline{N}(r,a)}{T(r,f)},$$
where $`N(r,a)`$ is the counting function of the $`a`$-points of $`f`$, counting multiplicities, and $`\overline{N}(r,a)`$ the counting function ignoring multiplicities. The point $`a`$ is a totally ramified value of $`f`$ if all $`a`$-points of $`f`$ have multiplicity two or higher.
The following theorem is a recently obtained difference analogue of the lemma on the logarithmic derivative .
###### Theorem 2.1
Let $`f`$ be a non-constant meromorphic function of finite order, $`c`$ and $`\delta <1`$. Then
$$m(r,\frac{f(z+c)}{f(z)})=o\left(\frac{T(r,f)}{r^\delta }\right)$$
(2.2)
for all $`r`$ outside of a possible exceptional set $`E`$ with finite logarithmic measure $`_E\frac{dr}{r}<\mathrm{}`$.
In the original statement of Theorem 2.1 in the error term on the right side of (2.2) has $`T(r+|c|,f)`$ instead of $`T(r,f)`$. But by the following lemma, \[3, Lemma 2.1\], we have $`T(r+|c|,f)=(1+o(1))T(r,f)`$ for all $`r`$ outside of a set with finite logarithmic measure, whenever $`f`$ is of finite order.
###### Lemma 2.2
Let $`T:(0,+\mathrm{})(0,+\mathrm{})`$ be a non-decreasing continuous function, $`s>0`$, $`\alpha <1`$, and let $`F^+`$ be the set of all $`r`$ such that
$$T(r)\alpha T(r+s).$$
(2.3)
If the logarithmic measure of is $`F`$ infinite, that is, $`_F\frac{dt}{t}=\mathrm{}`$, then
$$\underset{r\mathrm{}}{lim\; sup}\frac{\mathrm{log}T(r)}{\mathrm{log}r}=\mathrm{}.$$
Let $`f(z)`$ be a non-constant meromorphic function of finite order, and let $`a(z)`$ be a finite-order periodic function with period $`c`$ such that $`f(z)a(z)`$. Denote
$$\mathrm{\Delta }_cf:=f(z+c)f(z),$$
and $`\mathrm{\Delta }_c^nf:=\mathrm{\Delta }_c^{n1}(\mathrm{\Delta }_cf)`$ for all $`n`$, $`n2`$. Then by applying Theorem 2.1 with the function $`f(z)a(z)`$, we have
$$\begin{array}{cc}\hfill m(r,\frac{\mathrm{\Delta }_cf}{fa})& =m(r,\frac{f(z+c)a(z+c)}{f(z)a(z)})+O(1)\hfill \\ & =o\left(\frac{T(r,fa)}{r^\delta }\right)+O(1)\hfill \end{array}$$
(2.4)
outside of a possible exceptional set with finite logarithmic measure. We denote by $`๐ฎ(f)`$ the set of all meromorphic functions $`g`$ such that $`T(r,g)=o(T(r,f))`$ for all $`r`$ outside of a set with finite logarithmic measure. Functions in the set $`๐ฎ(f)`$ are called small compared to $`f`$, or slowly moving with respect to $`f`$. Also, if $`g๐ฎ(w)`$ we denote $`T(r,g)=S(r,f)`$.
Since by (2.4)
$$m(r,\frac{\mathrm{\Delta }_cf}{fa})=S(r,fa)$$
(2.5)
we arrive at the following lemma by induction and using the fact that
$$T(r,f(z+1))(1+\epsilon )T(r+1,f(z))$$
for any $`\epsilon >0`$ when $`r`$ is large .
###### Lemma 2.3
Let $`c`$, $`n`$, and let $`f`$ be a meromorphic function of finite order. Then for all small periodic functions $`a๐ฎ(f)`$
$$m(r,\frac{\mathrm{\Delta }_c^nf}{fa})=S(r,f),$$
where the exceptional set associated with $`S(r,f)`$ is of at most finite logarithmic measure.
Finally, an identity due to Valiron and Mohonโko is needed in the following section. It states that if the function $`R(z,f)`$ is rational in $`f`$ and has small meromorphic coefficients, then
$$T(r,R(z,f))=\mathrm{deg}_f(R)T(r,f)+S(r,f).$$
(2.6)
For the proof see also .
### 2.2 Second main theorem
The lemma on the logarithmic derivative is one of the main components of the proof of the second main theorem of Nevanlinna theory. The following theorem is obtained by combining the standard method of proof for the second main theorem together with Theorem 2.1. As a result a version of the second main theorem is obtained where, instead of the usual ramification term, there is a certain quantity expressed in terms of paired points of the considered function $`f`$. Since periodic functions are the analogues of constants for exact differences, it is natural to consider slowly moving periodic functions as target functions of $`f`$.
###### Theorem 2.4
Let $`c`$, and let $`f`$ be a meromorphic function of finite order such that $`\mathrm{\Delta }_cf0`$. Let $`q2`$, and let $`a_1(z),\mathrm{},a_q(z)`$ be distinct meromorphic periodic functions with period $`c`$ such that $`a_k๐ฎ(f)`$ for all $`k=1,\mathrm{},q`$. Then
$$m(r,f)+\underset{k=1}{\overset{q}{}}m(r,\frac{1}{fa_k})2T(r,f)N_{pair}(r,f)+S(r,f)$$
where
$$N_{pair}(r,f):=2N(r,f)N(r,\mathrm{\Delta }_cf)+N(r,\frac{1}{\mathrm{\Delta }_cf})$$
and the exceptional set associated with $`S(r,f)`$ is of at most finite logarithmic measure.
Proof. By denoting
$$P(f):=\underset{k=1}{\overset{q}{}}\left(fa_k\right),$$
we have
$$\frac{1}{P(f)}=\underset{k=1}{\overset{q}{}}\frac{\alpha _k}{fa_k},$$
where $`\alpha _k๐ฎ(f)`$ are certain periodic functions with period $`c`$. Hence, by (2.5), we obtain
$$m(r,\frac{\mathrm{\Delta }_cf}{P(f)})\underset{k=1}{\overset{q}{}}m(r,\frac{\mathrm{\Delta }_cf}{fa_k})+S(r,f)=S(r,f),$$
and so
$$m(r,\frac{1}{P(f)})=m(r,\frac{\mathrm{\Delta }_cf}{P(f)}\frac{1}{\mathrm{\Delta }_cf})m(r,\frac{1}{\mathrm{\Delta }_cf})+S(r,f).$$
(2.7)
By combining the first main theorem, (2.7) and the Valiron-Moโhonko identity (2.6), we have
$$\begin{array}{cc}\hfill T(r,\mathrm{\Delta }_cf)& =m(r,\frac{1}{\mathrm{\Delta }_cf})+N(r,\frac{1}{\mathrm{\Delta }_cf})+O(1)\hfill \\ & m(r,\frac{1}{P(f)})+N(r,\frac{1}{\mathrm{\Delta }_cf})+S(r,f)\hfill \\ & =qT(r,f)\underset{k=1}{\overset{q}{}}N(r,\frac{1}{fa_k})+N(r,\frac{1}{\mathrm{\Delta }_cf})+S(r,f)\hfill \\ & =\underset{k=1}{\overset{q}{}}m(r,\frac{1}{fa_k})+N(r,\frac{1}{\mathrm{\Delta }_cf})+S(r,f).\hfill \end{array}$$
Thus, by (2.5),
$$\begin{array}{cc}\hfill m(r,f)+\underset{k=1}{\overset{q}{}}m(r,\frac{1}{fa_k})& T(r,f)+N(r,\mathrm{\Delta }_cf)+m(r,\mathrm{\Delta }_cf)\hfill \\ & N(r,\frac{1}{\mathrm{\Delta }_cf})N(r,f)+S(r,f)\hfill \\ & T(r,f)+N(r,\mathrm{\Delta }_cf)+m(r,f)\hfill \\ & N(r,\frac{1}{\mathrm{\Delta }_cf})N(r,f)+S(r,f)\hfill \\ & =2T(r,f)+N(r,\mathrm{\Delta }_cf)N(r,\frac{1}{\mathrm{\Delta }_cf})\hfill \\ & 2N(r,f)+S(r,f).\hfill \end{array}$$
$`\mathrm{}`$
Let us now analyze the assertion of Theorem 2.4 more closely. By Lemma 2.2 $`N(r+|c|,f)=(1+o(1))N(r,f)`$ for all $`r`$ outside of a set with finite logarithmic measure. Therefore,
$$\begin{array}{cc}\hfill N_{pair}(r,f)& N(r,f)N(r+|c|,f)+N(r,\frac{1}{\mathrm{\Delta }_cf})\hfill \\ & =N(r,\frac{1}{\mathrm{\Delta }_cf})+S(r,f)\hfill \end{array}$$
so clearly Theorem 2.4 is telling us something non-trivial about the value distribution of finite-order meromorphic functions. In order to better interpret the meaning of the pair term $`N_{pair}(r,f)`$ we introduce the counting function $`n_c(r,a)`$, $`a`$, which is the number of points $`z_0`$ where $`f(z_0)=a`$ and $`f(z_0+c)=a`$, counted according to the number of equal terms in the beginning of Taylor series expansions of $`f(z)`$ and $`f(z+c)`$ in a neighborhood of $`z_0`$. We call such points $`c`$-separated $`a`$-pairs of $`f`$ in the disc $`\{z:|z|r\}`$.
For instance, if $`f(z)=a`$ and $`f(z+c)=a`$ with multiplicities $`p`$ and $`q<p`$, respectively, then the $`q`$ first terms in the series expansions of $`f(z)`$ and $`f(z+c)`$ are identical, and so this point is counted $`q`$ times in $`n_c(r,a)`$. Similarly, if in a neighborhood of $`z_0`$
$$f(z)=a+c_1(zz_0)+c_2(zz_0)^2+\alpha (zz_0)^3+O\left((zz_0)^4\right)$$
and
$$f(z+c)=a+c_1(zz_0)+c_2(zz_0)^2+\beta (zz_0)^3+O\left((zz_0)^4\right)$$
where $`\alpha \beta `$, then the point $`z_0`$ is counted $`3`$ times in $`n_c(r,a)`$.
The integrated counting function is defined as follows:
$$N_c(r,a):=_0^r\frac{n_c(t,a)n_c(0,a)}{t}๐t+n_c(0,a)\mathrm{log}r.$$
Similarly,
$$N_c(r,\mathrm{}):=_0^r\frac{n_c(t,\mathrm{})n_c(0,\mathrm{})}{t}๐t+n_c(0,\mathrm{})\mathrm{log}r,$$
where $`n_c(r,\mathrm{})`$ is the number of $`c`$-separated pole pairs of $`f`$, which are exactly the $`c`$-separated $`0`$-pairs of $`1/f`$. This means that if $`f`$ has a pole with multiplicity $`p`$ at $`z_0`$ and another pole with multiplicity $`q`$ at $`z_0+c`$ then this pair is counted $`\mathrm{min}\{p,q\}+m`$ times in $`n_c(r,\mathrm{})`$, where $`m`$ is the number of equal terms in the beginning of the Laurent series expansions of $`f(z)`$ and $`f(z+c)`$ in a neighborhood of $`z_0`$. Of course, if $`pq`$ then $`m=0`$.
Note that $`n_c(r,a)`$ is finite for any finite $`r`$, provided that the given function $`f`$ is not periodic with period $`c`$. Otherwise there would be a point $`z_0`$ in a neighborhood of which the series expansions of $`f(z)`$ and $`f(z+c)`$ would be identical. But this means that $`f(z)f(z+c)`$ in the whole complex plane, which contradicts the assumption. However, it is possible that $`n_c(r,a)`$ is strictly greater than the counting function $`n(r,a)`$.
A natural difference analogue of $`\overline{N}(r,a)`$ is
$$\stackrel{~}{N}_c(r,a):=N(r,a)N_c(r,a)$$
which counts the number of those $`a`$-points (or poles) of $`f`$ which are not in $`c`$-separated pairs. We also use the notation $`N_c(r,\frac{1}{fa})`$ instead of $`N_c(r,a)`$ and $`N_c(r,f)`$ instead of $`N_c(r,\mathrm{})`$ when we want to emphasize the connection to the meromorphic function $`f`$. With this notation we may state the main result of this paper.
###### Theorem 2.5
Let $`c`$, and let $`f`$ be a meromorphic function of finite order such that $`\mathrm{\Delta }_cf0`$. Let $`q2`$, and let $`a_1(z),\mathrm{},a_q(z)`$ be distinct meromorphic periodic functions with period $`c`$ such that $`a_k๐ฎ(f)`$ for all $`k=1,\mathrm{},q`$. Then
$$(q1)T(r,f)\stackrel{~}{N}_c(r,f)+\underset{k=1}{\overset{q}{}}\stackrel{~}{N}_c(r,\frac{1}{fa_k})+S(r,f)$$
where the exceptional set associated with $`S(r,f)`$ is of at most finite logarithmic measure.
Before proving Theorem 2.5 we briefly discuss its implications. Analogously to the classical Nevanlinna theory, the counting function $`\stackrel{~}{N}_c(r,a)`$ satisfies $`\stackrel{~}{N}_c(r,a)=T(r,f)+S(r,f)`$ for all except at most countably many values $`a`$ (see \[5, pp. 43-44\] for a proof of this). However, unlike $`N(r,a)`$, the counting function $`\stackrel{~}{N}_c(r,a)`$ may, for some values $`a`$, be negative for all sufficiently large $`r`$. This fact has interesting consequences. By Theorem 2.5 any finite-order meromorphic function $`f`$ is either periodic with period $`c`$, or it can have at most one non-deficient value $`a`$ such that whenever $`f(z)=a`$ also $`f(z+c)=a`$ and the first two terms in the series expansions of $`f(z)`$ at $`z`$ and $`z+c`$ are identical. For instance, consider the function $`g(z):=\mathrm{}(z)+\mathrm{exp}(z)`$ where $`\mathrm{}(z)`$ is a Weierstrass elliptic function with a period $`c2\pi i`$. Then $`T(r,g)=N(r,g)+S(r,g)`$ and each pole of $`g`$ contributes $`2`$ to $`n(r,g)`$ but $`2`$ to $`\stackrel{~}{n}_c(r,g)`$. Therefore $`T(r,g)=\stackrel{~}{N}_c(r,g)+S(r,g)`$ and so $`\stackrel{~}{N}_c(r,a)=T(r,g)+S(r,g)`$ for all $`a`$ by Theorem 2.5.
Proof of Theorem 2.5. By Theorem 2.4 and the first main theorem, we obtain
$$\begin{array}{cc}\hfill (q1)T(r,f)& N(r,f)+\underset{k=1}{\overset{q}{}}N(r,\frac{1}{fa_k})N(r,\frac{1}{\mathrm{\Delta }_cf})\hfill \\ & +N(r,\mathrm{\Delta }_cf)2N(r,f)+S(r,f).\hfill \end{array}$$
(2.8)
We denote by $`N_0(r,f)`$ the counting function for those poles of $`f`$ having Laurent series expansions at $`z_0`$ and $`z_0+c`$ with identical principal parts, multiplicity counted according to the number of equal terms in the beginning of the analytic part of the series expansions. (For instance, if $`f(z)=c/(zz_0)^2+b/(zz_0)+a+\alpha (zz_0)+O((zz_0)^2)`$ and $`f(z+c)=c/(zz_0)^2+b/(zz_0)+a+\beta (zz_0)+O((zz_0)^2)`$ the pole at $`z_0`$ is counted once in $`N_0(r,f)`$ whenever $`\alpha \beta `$.) Since $`N(r,f)=N(r+|c|,f)+S(r,f)`$ by Lemma 2.2, inequality (2.8) takes the form
$$\begin{array}{cc}\hfill (q1)T(r,f)& N(r,f)+N_0(r,f)+\underset{k=1}{\overset{q}{}}N(r,\frac{1}{fa_k})N(r,\frac{1}{\mathrm{\Delta }_cf})\hfill \\ & +N(r,\mathrm{\Delta }_cf)2N(r+|c|,f)N_0(r,f)+S(r,f).\hfill \end{array}$$
(2.9)
The rest of the proof consists of estimates on different terms on the right side of (2.9). First, by the definition of a paired point, we have
$$N_0(r,f)+\underset{k=1}{\overset{q}{}}N_c(r,\frac{1}{fa_k})N(r,\frac{1}{\mathrm{\Delta }_cf})$$
for all $`r`$, and thus
$$N_0(r,f)+\underset{k=1}{\overset{q}{}}N(r,\frac{1}{fa_k})N(r,\frac{1}{\mathrm{\Delta }_cf})\underset{k=1}{\overset{q}{}}\stackrel{~}{N}_c(r,\frac{1}{fa_k}).$$
(2.10)
Second, assume that $`z_0`$ is such that $`f(z_0+kc)=\mathrm{}`$ for all $`k`$ with multiplicities $`p_k0`$. Here $`p_k=0`$ means that $`f(z_0+kc)`$ is finite. (Note that the case $`p_k=0`$ for all $`k0`$ is not ruled out.) Out of these points only finitely many are inside the disc $`\{z:|z|r+|c|\}`$ for any $`r>0`$. By redefining $`z_0`$ if necessary, we may assume that these points are $`z_0+jc`$, $`j=0,\mathrm{},K`$, where $`K`$ is a constant depending only on $`r`$. Then $`z_0+c,\mathrm{},z_0+(K1)c`$ are inside the disc with radius $`r`$ centered at the origin, and $`\mathrm{\Delta }_cf`$ has a pole with multiplicity $`\mathrm{max}\{p_j,p_{j+1}\}m_j^{}`$ at $`z_0+jc`$, where $`j=1,\mathrm{},K1`$ and $`m_j^{}`$ is the number of equal terms in the beginning of the principal parts of the Laurent series expansions of $`f(z)`$ and $`f(z+c)`$ at $`z_0+jc`$. If principal parts are completely identical, the number of equal terms in the beginning of the analytic parts of the series at $`z_0+jc`$ is denoted by $`m_j^{\prime \prime }`$, and moreover $`m_j:=m_j^{}+m_j^{\prime \prime }`$. Therefore the contribution to
$$n(r,\mathrm{\Delta }_cf)2n(r+|c|,f)n_0(r,f)$$
from the points $`z_0+jc`$, $`j=0,\mathrm{},K`$, is
$$\begin{array}{cc}\hfill \underset{j=1}{\overset{K1}{}}& \left(\mathrm{max}\{p_j,p_{j+1}\}m_j^{}\right)2\underset{j=0}{\overset{K}{}}p_j\underset{j=1}{\overset{K1}{}}m_j^{\prime \prime }\hfill \\ & =\underset{j=1}{\overset{K1}{}}\left(\mathrm{max}\{p_j,p_{j+1}\}m_j^{}m_j^{\prime \prime }\right)\hfill \\ & \left(p_0+\underset{j=0}{\overset{K1}{}}\left(\mathrm{max}\{p_j,p_{j+1}\}+\mathrm{min}\{p_j,p_{j+1}\}\right)+p_K\right)\hfill \\ & \underset{j=1}{\overset{K1}{}}\left(\mathrm{min}\{p_j,p_{j+1}\}+m_j\right).\hfill \end{array}$$
(2.11)
The quantity on the right side of (2.11) is by definition exactly the same as the contribution to $`n_c(r,f)`$ from the points $`z_0+jc`$, $`j=0,\mathrm{},K`$. Therefore, by summing over all poles of $`f`$, we obtain
$$N(r,f)+N(r,\mathrm{\Delta }_cf)2N(r+|c|,f)N_0(r,f)\stackrel{~}{N}_c(r,f).$$
(2.12)
The assertion follows by combining (2.9), (2.10) and (2.12). $`\mathrm{}`$
### 2.3 Defect relation and Picardโs theorem
Nevanlinnaโs second main theorem is a deep generalization of Picardโs theorem, and as such it has many important consequences for the value distribution of meromorphic functions. In this section we present difference analogues of a number of these results, including Picardโs theorem and Nevanlinnaโs theorems on the total deficiency sum and completely ramified values of a meromorphic function. All of the results in this section follow from Theorem 2.5.
A difference analogue of the index of multiplicity $`\theta (a,f)`$ is called the $`c`$-separated pair index, and it is defined as follows:
$$\pi _c(a,f):=\underset{r\mathrm{}}{lim\; inf}\frac{N_c(r,a)}{T(r,f)},$$
where $`a`$ is either a slowly moving periodic function with period $`c`$, or $`a=\mathrm{}`$. Similarly, we define
$$\mathrm{\Pi }_c(a,f):=1\underset{r\mathrm{}}{lim\; sup}\frac{\stackrel{~}{N}_c(r,a)}{T(r,f)},$$
which is an analogue of
$$\mathrm{\Theta }(a,f)=1\underset{r\mathrm{}}{lim\; sup}\frac{\overline{N}(r,a)}{T(r,f)}$$
in the usual value distribution theory.
The following corollary says that a non-periodic meromorphic function of finite order cannot have too many $`a`$-points which appear in pairs. It is a difference analogue of Nevanlinnaโs theorem on deficient values.
###### Corollary 2.6
Let $`c`$, and let $`f`$ be a meromorphic function of finite order such that $`\mathrm{\Delta }_cf0`$. Then $`\mathrm{\Pi }_c(a,f)=0`$ except for at most countably many meromorphic periodic functions $`a`$ with period $`c`$ such that $`a๐ฎ(f)`$, and
$$\underset{a}{}\left(\delta (a,f)+\pi _c(a,f)\right)\underset{a}{}\mathrm{\Pi }_c(a,f)2.$$
(2.13)
By the second main theorem it follows that $`\mathrm{\Theta }(a,f)=0`$ for all except at most countably many values $`a`$, see, for instance, \[5, pp. 43โ44\]. The same reasoning can be applied to prove that Theorem 2.5 implies Corollary 2.6.
Probably the most distinct difference between the classical Nevanlinna theory and its difference analogue is that, although $`0\mathrm{\Theta }(a,f)1`$ for all meromorphic functions $`f`$ and for all $`a`$ in the extended complex plane, the maximal deficiency sum
$$\underset{a}{}\mathrm{\Pi }_c(a,f)=2$$
may be attained by a single value $`a`$. For instance, the function $`g(z)=\mathrm{}(z)+\mathrm{exp}(z)`$, where $`\mathrm{}(z)`$ is a Weierstrass elliptic function with a period $`c2\pi i`$, satisfies $`\mathrm{\Pi }_c(\mathrm{},g)=2`$. In fact, by the definition of $`\mathrm{\Pi }_c(a,f)`$ alone, it is not even clear that $`\mathrm{\Pi }_c(a,f)`$ has an upper bound what so ever. The fact that $`\mathrm{\Pi }_c(a,f)2`$ for all $`a`$ follows by Corollary 2.6.
We say that $`a`$ is an exceptional paired value of $`f`$ with the separation $`c`$ if the following property holds for all except at most finitely many $`a`$-points of $`f`$: Whenever $`f(z)=a`$ then also $`f(z+c)=a`$ with the same or higher multiplicity. Clearly $`N(r,a)N_c(r,a)`$ for all exceptional paired values $`a`$ of $`f`$. Note also that by this definition all Picard exceptional values of $`f`$ are also exceptional paired values. The following corollary is an analogue of Picardโs theorem.
###### Corollary 2.7
If a finite-order meromorphic function $`f`$ has three exceptional paired values with the separation $`c`$, then $`f`$ is a periodic function with period $`c`$.
Corollary 2.7 implies that if a finite-order meromorphic function $`w`$ has two groups of three exceptional paired values with two different separations, say $`c_1`$ and $`c_2`$ independent over the reals, then either $`w`$ is a constant or $`w`$ is an elliptic function with periods $`c_1`$ and $`c_2`$ and therefore exactly of order $`2`$.
There is no hope of extending Corollary 2.7 (or Corollary 2.6) to include all infinite order meromorphic functions, since the function $`\mathrm{exp}(\mathrm{exp}(z))`$ has three exceptional paired values with the separation $`\mathrm{log}2`$: In addition to the Picard exceptional zeros and poles, the value $`1`$ is exceptionally paired, although non-deficient.
An example of a finite-order meromorphic function which has exactly two exceptional paired values with the separation $`2K`$ is given by the elliptic function $`\text{sn}(z,k)`$, where $`k(0,1)`$ is the elliptic modulus and $`K`$ is the complete elliptic integral. The function $`\text{sn}(z,k)`$ is periodic with the periods $`4K`$ and $`2iK^{}`$, and it attains the value zero at points $`2nK+2miK^{}`$ and has its poles at $`2nK+(2m+1)iK^{}`$, where $`n,m`$. The function $`\text{sn}(z,k)`$ has no deficient values, but it has the maximal four completely ramified values at $`\pm 1`$ and $`\pm 1/k`$. Therefore, the function $`g(z)=\text{sn}(z,k)`$ satisfies
$$\underset{a}{}\pi _{2K}(a,g)=2$$
and, moreover,
$$\underset{a}{}\left(\theta (a,g)+\pi _{2K}(a,g)\right)=4.$$
Analogously to complete ramification, we say that a point $`a`$ is completely paired with the separation $`c`$ if whenever $`f(z)=a`$ then either $`f(z+c)=a_j`$ or $`f(zc)=a_j`$, with the same multiplicity. Then a non-periodic meromorphic function of finite order can have at most four values which only appear in pairs.
###### Corollary 2.8
Let $`c`$, and let $`f`$ be a meromorphic function of finite order such that $`\mathrm{\Delta }_cf0`$. Then $`f`$ has at most four completely paired points with separation $`c`$.
Similarly, a non-periodic finite-order function $`f`$ can have at most three values $`a`$ which only appear such that for some $`z_0`$, $`f(z_0)a`$, $`f(z_0+jc)=a`$ with the same multiplicity for each $`j=1,2,3`$, and $`f(z_0+4c)a`$. We say that such values appear in lines of three. Similarly, a finite-order meromorphic function can have a maximum of two values which appear only in lines of four or more.
### 2.4 Functions sharing values
Another consequence of Nevanlinnaโs second main theorem is the five value theorem, which says that if two non-constant meromorphic functions share five values ignoring multiplicity then these functions must be identical. By considering periodic functions instead of constants, and by ignoring paired points instead of multiplicity, we obtain a difference analogue of the five value theorem.
We say that two meromorphic functions $`f`$ and $`g`$ share a point $`a`$, ignoring $`c`$-separated pairs, when $`f(z)=a`$ if and only if $`g(z)=a`$ with the same multiplicity, unless $`a`$ is a $`c`$-separated pair of $`f`$ or $`g`$. In short, all paired points are ignored when determining whether or not $`f`$ and $`g`$ share $`a`$. This also means that if $`f`$ has a paired $`a`$-point at $`z_0`$ and $`g`$ has a single $`a`$-point at the same location, this point is not shared by $`f`$ and $`g`$.
###### Theorem 2.9
Let $`c`$, and let $`f`$ and $`g`$ be meromorphic functions of finite order. If there are five distinct periodic functions $`a_k๐ฎ(f)`$ such that $`f`$ and $`g`$ share $`a_k`$, ignoring $`c`$-separated pairs, for all $`k=1,\mathrm{},5`$ then either $`f(z)g(z)`$ or both $`f`$ and $`g`$ are periodic with period $`c`$.
Proof. We follow the reasoning of the proof of the five value theorem . Suppose first that $`f`$ is periodic with period $`c`$. Then by definition all $`a`$-points of $`f`$ are paired. Since $`f`$ and $`g`$ share five points, ignoring pairs, $`g`$ has at least five exceptional paired values, and therefore it must also be periodic with period $`c`$ by Corollary 2.7.
Assume now that neither $`f`$ nor $`g`$ is periodic with period $`c`$, and that $`fg`$. Then by Theorem 2.5, for any $`\epsilon >0`$,
$$(4+\epsilon )T(r,f)\stackrel{~}{N}_c(r,f)+\underset{k=1}{\overset{5}{}}\stackrel{~}{N}_c(r,\frac{1}{fa_k})$$
(2.14)
and
$$(4+\epsilon )T(r,g)\stackrel{~}{N}_c(r,g)+\underset{k=1}{\overset{5}{}}\stackrel{~}{N}_c(r,\frac{1}{ga_k})$$
(2.15)
outside a set with finite logarithmic measure. Since
$$\stackrel{~}{N}_c(r,\frac{1}{ga_k})=\stackrel{~}{N}_c(r,\frac{1}{fa_k})$$
for all $`k=1,\mathrm{},5`$, inequalities (2.14) and (2.15) imply
$$\begin{array}{cc}\hfill T(r,\frac{1}{fg})& T(r,f)+T(r,g)+O(1)\hfill \\ & \frac{2}{3+\epsilon }\underset{k=1}{\overset{5}{}}\stackrel{~}{N}_c(r,\frac{1}{fa_k})\hfill \\ & \frac{2}{3+\epsilon }N(r,\frac{1}{fg})\hfill \\ & \frac{2}{3+\epsilon }T(r,\frac{1}{fg}).\hfill \end{array}$$
This is only possible when $`fg`$ is a constant, say $`g(z)=f(z)+k`$. But now, since $`f(z)`$ and $`f(z)+k`$ share five points out of which at most two can be either exceptionally paired or Picard exceptional, $`k=0`$, and the assertion follows. $`\mathrm{}`$
The elliptic functions $`\text{sn}z`$ and $`1/\text{sn}z`$ show that the number five cannot be replaced by four in Theorem 2.9. Namely, for both functions zero and infinity are exceptional paired values, and they share the points $`1`$ and $`1`$, counting multiplicities. Therefore, $`\text{sn}z`$ and $`1/\text{sn}z`$ share the points $`1`$, $`0`$, $`1`$ and $`\mathrm{}`$, ignoring pairs.
### 2.5 An application to difference equations
In this section we give an example of how to apply the obtained results to study meromorphic solutions of difference equations. We consider the equation
$$w(z+1)+w(z1)=\frac{a_2w(z)^2+a_0}{1w(z)^2}$$
(2.16)
where the right side is irreducible in $`w`$ and the coefficients $`a_j`$ are constants. Equation (2.16) is a subcase of a more general equation studied in where it was shown that the existence of one finite-order meromorphic solution is sufficient to reduce a large class of difference equations into a difference Painlevรฉ equation or into a linear difference equation, provided that the solution does not simultaneously satisfy a difference Riccati equation. Suppose that (2.16) has a finite-order meromorphic solution $`w(z)`$ and consider a Laurent series expansion of $`w`$ in a neighborhood of a point $`z_0`$ such that $`w(z_0)=\delta `$ with the multiplicity $`k1`$, where $`\delta :=\pm 1`$. Then $`w`$ has a pole of order at least $`k`$ at $`z_01`$ or $`z_0+1`$.
Consider first the case where $`w(z_0+1)=\mathrm{}`$ with the multiplicity $`k`$ and $`w(z_01)`$ is either finite or a pole with multiplicity strictly less than $`k`$. Then by iterating (2.16), we have
$$\begin{array}{cc}\hfill w(z+4n)& =\delta +\alpha (zz_0)^k+O\left((zz_0)^{k+1}\right)\hfill \\ \hfill w(z+2n+1)& =\frac{\left((1)^n(\frac{1}{4}n+\frac{1}{8})\frac{1}{8}\right)(a_0+a_2)}{\alpha \delta }(zz_0)^k+O\left((zz_0)^{1k}\right)\hfill \\ \hfill w(z+4n+2)& =a_2\delta +O\left((zz_0)\right)\hfill \end{array}$$
(2.17)
for all $`n\{0\}`$ and for all $`z`$ in a suitably small neighborhood of $`z_0`$, provided that $`a_20`$. Since we assumed the right side of (2.16) to be irreducible $`a_0+a_20`$ and so $`w(z+2n+1)=\mathrm{}`$ for all $`n\{0\}`$. The iteration in the case where $`w(z_0+1)`$ is finite, or a pole with low order, and $`w(z_01)=\mathrm{}`$ is symmetric with (2.17).
Suppose now that $`w(z_0)=\delta `$ and $`w(z_0\pm 1)=\mathrm{}`$ all with the same multiplicity $`k`$. Then, assuming $`c_1`$ and $`c_1`$ such that $`c_1c_10`$, we have
$$\begin{array}{cc}\hfill w(z+4n)& =\delta +\alpha (zz_0)^k+O\left((zz_0)^{k+1}\right)\hfill \\ \hfill w(z+4n+2)& =a_2\delta +O\left((zz_0)\right)\hfill \\ \hfill w(z+2n+1)& =c_{2n+1}(zz_0)^k+O\left((zz_0)^{1k}\right)\hfill \end{array}$$
(2.18)
for all $`n`$ as long as none of the constants $`c_{2n+1}`$ vanish. But if $`c_{2n_0+1}=0`$ for some $`n_0`$ then we are back in the situation (2.17) with the starting point $`z_0+2n_0+1`$ instead of $`z_01`$. Note also that a closer inspection of the iteration in (2.18) shows that
$$c_{k\pm 4}=c_k+\frac{a_2+a_0}{2\alpha }$$
(2.19)
for all $`k`$.
The final possibility is that $`w(z_0)=\delta `$ with the multiplicity $`k`$ and $`w(z_0\pm 1)=\mathrm{}`$ for both choices of the sign with the multiplicity strictly greater than $`k`$. But in this case it is immediately seen that $`w(z)`$ has a pole with the same order in $`z_0+2n+1`$ for all $`n`$.
We conclude that all poles, $`1`$-points and $`1`$-points of $`w`$ appear in lines where each point is separated from its neighbors by the constant $`4`$, with the possible exception of the endpoints of sequences of points appearing as a part of (2.17). In fact for our purposes it is sufficient to know that all poles and $`\delta `$-points of $`w`$ appear in groups of four or more, with $`4`$-separation. Assume that $`w`$ is not periodic with period four. Then by Theorem 2.5,
$$\begin{array}{cc}\hfill T(r,w)& \stackrel{~}{N}_4(r,\mathrm{})+\stackrel{~}{N}_4(r,1)+\stackrel{~}{N}_4(r,1)+S(r,w)\hfill \\ & \frac{1}{4}N(r,\mathrm{})+\frac{1}{4}N(r,1)+\frac{1}{4}N(r,1)+S(r,w)\hfill \\ & \frac{3}{4}T(r,w)+S(r,w),\hfill \end{array}$$
which is a contradiction. Therefore, either $`a_2=0`$, or $`w`$ is periodic with period $`4`$ or of infinite order.
Suppose finally that $`w`$ is periodic with period $`4`$. Then $`1`$ and $`1`$ are Picard exceptional values of $`w`$ by (2.17), (2.18) and (2.19). Therefore all poles of $`w`$ appear in lines where each pole is separated from its neighbors by the constant $`2`$, and so $`w`$ is periodic with period $`2`$. But then, by periodicity, $`w(z+1)`$, $`w(z1)`$ and $`w(z+1)+w(z1)`$ are infinite simultaneously. On the other hand, the right side of (2.16) is never infinite since the values $`\pm 1`$ are Picard exceptional. Hence also the value infinity is Picard exceptional for $`w`$, and therefore $`w`$ is a constant by Picardโs theorem. We conclude that if (2.16) has a non-constant meromorphic solution of finite order then $`a_2=0`$.
The existence of finite-order meromorphic solutions of (2.16) is guaranteed in the case $`a_2=0`$, $`a_00`$. Then (2.16) has solutions of the form
$$w(z)=\frac{\alpha \text{sn}(\mathrm{\Omega }z+C)+\beta }{\gamma \text{sn}(\mathrm{\Omega }z+C)+\delta }$$
(2.20)
where $`C`$ is arbitrary, and $`\alpha ,\beta ,\gamma ,\delta ,\mathrm{\Omega }`$ are certain constants depending on another free parameter. The meromorphic solutions (2.20) are of order $`2`$ and periodic, but not of period $`4`$.
## 3 Discussion
Nevanlinnaโs second main theorem implies that a non-constant meromorphic function cannot have too many points with high multiplicity. In this study a difference analogue of the second main theorem of Nevanlinna theory was given, which shows that a non-periodic finite-order meromorphic function cannot have many values which only appear in pairs, separated by a fixed constant. Then a number of results on the value distribution of finite-order meromorphic functions were derived by combining existing proof techniques from Nevanlinna theory together with the difference analogue of the second main theorem. These include analogues of Picardโs theorem, the theorem on the deficiency sum and the theorem on meromorphic functions sharing five values. Sharpness of these results was discussed with the help of examples. Also, an example of how to apply some of these results to study complex difference equations was given.
All concepts of Nevanlinna theory related to ramification have a natural difference analogue. For instance, constant functions are analogous to periodic functions, and a pole with multiplicity $`n>1`$ is analogous to a line of $`n`$ poles with the same multiplicity, each separated from its neighbors by a fixed constant. Similarly as a pole is counted only once in the counting function $`\overline{N}(r,f)`$ regardless of its multiplicity, only one pole from the above line of poles contributes to $`\stackrel{~}{N}_c(r,f)`$. However, some notions in the difference Nevanlinna theory seem to go, in a sense, further than their classical counterparts. If a line of points consists of poles with different multiplicities, the contribution from these poles to $`\stackrel{~}{N}_c(r,f)`$ is nevertheless strictly less than the contribution to $`N(r,f)`$. Therefore this situation is still exceptional in the sense of the difference deficiency relation (1.2). On the other hand, if all poles in the line have similar enough Laurent series expansions, then the contribution to $`\stackrel{~}{N}_c(r,f)`$ from these poles may be negative. This implies that the maximal value two in the difference deficiency relation (1.2) may be attained by one value $`a`$, which is impossible for the classical deficiencies (1.1).
## 4 Open problems
In addition to his ground-breaking results in the field of value distribution theory, Nevanlinna proposed a number of problems many of which have remained open until recently. In this section we briefly discuss two of them.
### 4.1 Inverse problem
The inverse problem for the deficiency relation is to find a meromorphic function $`f`$ which at prescribed points has certain non-zero deficiencies and ramification indices. This problem was proposed and partially solved by Nevanlinna himself, see , but the complete solution had to wait until 1977 when Drasin settled the issue by a clever use of quasi-conformal mappings. Later on Drasin established a related corollary by F. Nevanlinna, which states that if a meromorphic function $`f`$ has finite order $`\lambda `$ and $`_a\delta (a,f)=2`$ then $`2\lambda `$ is a natural number greater or equal to two. In the view of Corollary 2.6 it is natural to ask under what conditions it is possible to find a meromorphic function of finite order for which the pair index $`\pi (a,f)`$ and the deficiency $`\delta (a,f)`$ have certain non-zero values at prescribed points $`a`$?
### 4.2 Slowly moving targets
Another question proposed by Nevanlinna is whether or not the relation (1.1) remains valid if the sum is taken over all small functions with respect to $`f`$. Partial answer was given by Steinmetz and Osgood who showed that
$$\underset{a}{}\delta (a,f)2$$
where the sum is taken over distinct small functions with respect to $`f`$. A complete solution to this problem was given recently by Yamanoi who showed that (1.1) indeed remains valid if the sum is taken over the larger field small functions, rather than just constants. Similarly we propose that the (2.13) remains valid even if the sum is taken over the field $`๐ฎ(f)`$. It can be immediately seen, by a modification of the reasoning in \[5, p. 47\], that the assertion holds for at most three functions.
R. G. Halburd
Department of Mathematical Sciences, Loughborough University, Loughborough, Leicestershire, LE11 3TU, UK.
*E-mail address:* r.g.halburd@lboro.ac.uk
R. J. Korhonen
University of Joensuu, Department of Mathematics, P. O. Box 111, FIN-80101 Joensuu, FINLAND.
*E-mail address:* risto.korhonen@joensuu.fi
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# Issues with Core-Collapse Supernova Progenitor Models
## 1 Inroduction
The core-collapse supernova mechanism is still an unsolved problem. The failure of โstate-of-the-artโ one-dimensional core collapse simulations, utilizing multienergy-multiangle neutrino transport schemes and realistic opacities, and the plethora of evidence that the phenomenon is inherently multidimensional has inaugurated a new era of supernova modeling. Supernova codes are now beginning to couple multidimensional hydrodynamics to multidimensional neutrino transport, or at least neutrino transport along radial rays. Furthermore, spectacular advances are being made in the microphysics, particularly in the equation of state and the neutrino opacities. Many of the issues and prospect of both the microphysics and macrophysics involved in realistic core collapse supernova modeling are discussed in this volume.
Here I discuss the input data to these core collapse simulations, namely the progenitor models. Supernova modelers (myself included) tend to take delivery of these progenitor models without being fully cognizant of the approximations that are made in order to evolve these models from the main sequence to the point of core-collapse. Phenomena such as convection, rotational instabilities, and mass loss involve huge ranges of spatial and temporal scales and/or uncertain physics and require some sort of parameterization. Some of the macrophysics is inherently multidimensional. The purpose here is gather together most of the salient approximations and parameterizations that are made in computing progenitor models, and to increase thereby the awareness of the supernova community to the many aspects of current progenitor models open to future revision. Unfortunately, without a viable model of the core-collapse supernova mechanism, it is difficult to assess the effect on the core-collapse supernova scenario of variations within the uncertainties of the progenitor models. Conventional wisdom would point to the mass of the precollapse iron core, the density profile in the adjacent silicon and oxygen layers, and its rate of rotation as being particularly important. The reader will find many details of recent progenitor models reviewed by Maeder & Meynet (2000) and Woosley et al. (2002).
## 2 Convection
Convection at the high Reynolds number characterizing flows inside stars is highly turbulent and chaotic characterized by eddies on a vast spectrum of scales. A great source of uncertainty in current stellar evolutionary calculations is how to model the thermal and compositional mixing at convective boundaries, and how to model the reactive flows that occur during late evolutionary phases when convective and nuclear time scales become comparable. First-principled numerical simulations of turbulence with the necessary resolution are not yet practicable, and needless to say it has been impossible to couple a first-principled calculation of turbulence with a stellar evolution code. Almost all stellar evolution codes model convection with some variant of โmixing length theoryโ (MLT) (Bรถhm-Vitense, 1958) which attempts to capture the effects of convection by an essentially one parameter diffusion process. The convective diffusivity is taken to be $`K_{\mathrm{conv}}=\frac{1}{3}v_{\mathrm{conv}}\mathrm{}`$ where $`v_{\mathrm{conv}}`$ is the mean velocity of a typical convective eddy as it traverses a mean free path, $`\mathrm{}`$. The mean velocity is computed from the buoyancy of the eddy and Newtonโs laws, and the mean free path, $`\mathrm{}`$, referred to as the mixing length, is the free parameter of the theory. It is typically taken to be some fraction of the pressure scale height. A number of uncertainties attend this attempt to model convection and an attempt will be made to describe these below.
Some sort of convective motions will occur if a fluid is unstably stratified, that is, if a displaced fluid element finds itself subjected to a buoyancy force tending to amplify the displacement. Whether a fluid element will be unstable, and if so the mean velocity that it will acquire, will depend on the assumptions made as to how the fluid element is displaced. If it is displaced adiabatically (the typical assumption) and at constant composition, then the resulting convection if it occurs is referred to as Ledoux convection. If it moves adiabatically but maintains the same composition as the background, then the resulting convection is referred to as Schwarzschild convection. If the composition gradient is zero the criterion for the two is the same. As the background composition gradient in a star, when nonzero, almost always goes from heavier to lighter elements as a function of radius (e.g., in the wake of a retreating convective region), the composition gradient tends to be stabilizing. Thus Ledoux convection is more restrictive, in the sense that a fluid can be unstable to Schwarzschild convection but stable to Ledoux convection.
Regions unstable to Schwarzschild convection but stable to Ledoux convection can be doubly diffusive unstable (Kato, 1966), a phenomenon usually referred to as semiconvection, although this term has been used to refer to a multitude of sins. A fluid element perturbed outward under these conditions will find itself hotter than the background and therefore tend to continue the displacement, but will be stabilized by its heavier composition. Thermal diffusion, if faster than compositional diffusion, will tend to thermally equilibrate the fluid element with the background while leaving it with a compositional difference tending to drive it back. What can result is an oscillation of the fluid element with growing amplitude. It is unclear how to mix the material under these conditions. Two dimensional numerical simulations (Merryfield, 1995) suggest a complicated situation. Large-amplitude standing waves which break and mix over a distance of the order of a wavelength will arise if the instability is strongly driven. If the instability is weakly driven short waves arise initially and then organize themselves into longer waves which occasionally overturn and mix, and ultimately come to resemble horizontally propagating solitary waves. It is not clear how to connect the results of these simulations, which were unable to reach steady state, with the statistical steady state that presumably develops over the evolutionary time scales of stars. Extreme assumptions among stellar evolution modelers are that semiconvective mixing is fast and the use the Schwarzschild criterion for convection is therefore appropriate, or that it is slow and the use of the Ledoux criterion is therefore appropriate.
Semiconvection originally referred to another ambiguity that arises when a hydrogen burning core moves outward in mass as its helium content grows (Schwarzschild & Hรคrm, 1958). This happens in some massive star models as the pressure in the convective core becomes more dominated in time by radiation and convective instability is more easily achieved. A chemical discontinuity arises at the convective core boundary. If electron scattering dominates the opacity, as is the case for massive stars, then the opacity increases across the convective core boundary and a problem arises as to the placement of this boundary. As the boundary is approached from the inside the radiative gradient becomes equal to the adiabatic gradient. But just outside the boundary the opacity increase implies that the radiative gradient must exceed the adiabatic gradient. Hence the boundary should be moved outward. Doing so removes the composition gradient, however, and hence removes the need to move the boundary outward in the first place. This poses a dilemma, and a number of schemes have been proposed for dealing with it. These have been summarized by Stothers (1970).
A related mixing ambiguity, also referred to as semiconvection, can happen in stars with expanding helium burning cores (Schwarzschild & Hรคrm, 1969; Paczyลski, 1970; Castellani et al., 1971b, a; Robertson, 1971; Robertson & Faulkner, 1972). In this case the electron scattering is the same just inside and just outside the core, but the carbon rich mixture inside the core has a higher free-free opacity. This by itself does not prevent the boundary of the convective core from being located unambiguously, as curves a to c in Fig. 1 illustrate. The convective boundary occurs where the radiative gradient becomes equal to the adiabatic gradient, and curve b has correctly located this boundary. As the helium burning core grows, however, a point is reached where the opacity develops a minimum inside the convective core and then increases outward to the core boundary. In this case the attempt to find the convective core boundary leads to the possibilies illustrated by curves d to f. If curve d is chosen to locate the core boundary, the region between i and j is not convective, contradicting the choice. If curves e or f are chosen, the material at the edge of the core will be unstable to further convection since $`_{\mathrm{rad}}>_{\mathrm{rad}}`$ there. What is frequently done is to assume that curve e, modified by the horizontal segment connecting points m and n represents the correct choice. This is achieved by assuming that the requisite โsemiconvectiveโ compositional mixing takes place between points m and n causing $`_{\mathrm{rad}}=_{\mathrm{rad}}`$ there.
Another problem with the MLT parameterization of convection is overshooting, which refers to the tendency of convective eddies to penetrate the radiative layers surrounding a convective zone and hence induce a mixing of a region larger (in mass) than classically allowed by the strict adoption of the Schwarzschild or Ledoux criterion. This is a consequence of the fact that while the acceleration of the convective motions cease at the boundary of a convective region, there velocities do not. Thus convective overshooting may be present at the border of any convective region, and is not confined to any particular evolutionary phase. The effect of overshooting is to increase the mass in the convective region that is mixed which, in turn, can have a number of consequences for stellar ages, nucleosynthesis, and presupernova structure. Unfortunately, is not a natural outcome of MLT, due to the local nature of the theory (Renzini, 1987).
The radial extent, $`\mathrm{}_{\mathrm{OV}}`$, of the thermal and chemical mixing from the formal convective core boundary is typically parameterized by an ad hoc formula such as $`\mathrm{}_{\mathrm{OV}}=\alpha _{\mathrm{OV}}\mathrm{min}(r,H_\mathrm{p})`$, where $`H_\mathrm{p}`$ is the pressure scale height, $`r`$ is the core radius, the distance from the core edge to the surface, or some other such scale that naturally limits the extent of overshooting, and $`\alpha _{\mathrm{OV}}`$ is the parameter of the theory, typically below unity. The extent of convective overshooting will be a function of the Pรฉclet number (e.g., Zahn, 1999), which is the ratio of the convective to the radiative diffusivity. For large Pรฉclet numbers at the border of a convective region (typical of convective regions well below the stellar surface), the convective eddies exchange little heat with the background and therefore establish a nearly adiabatic gradient beyond the unstable region. They are therefore decelerated by the stable stratification. For small Pรฉclet numbers, however, radiation diffusion will tend to thermally equilibrate the convective eddies with the background as they penetrate beyond the formal convective boundary which will weaken their deceleration. In this case little heat is transported, but chemicals and momentum can be transported an appreciable distance. (Technically, the former (large Pรฉclet number) case is referred to as convective penetration, the latter is referred to as convective overshooting (Zahn, 1991).)
Some observational constraints suggesting a value of $`0.2`$ for $`\alpha _{\mathrm{OV}}`$ are provided by the size of gaps (blue loops) in open star cluster color-magnitude diagrams (Maeder & Mermilliod, 1981; Stothers & Chin, 1991a, b; Stothers, 1991; Nordstrรถm et al., 1997), the asteroseismology of $`\eta `$ Bootis (Di Mauro et al., 2003), accurate stellar dimensions derived from well-detached double-lined eclipsing binaries (Gimรฉnez et al., 2004) (which suggests a somewhat larger value of $`\alpha _{\mathrm{OV}}`$ for massive stars). Beyond this the value of $`\alpha _{\mathrm{OV}}`$ must be inferred from numerical simulations or guessed at.
A number of numerical simulations investigating the nature of turbulent compressible convection and convective overshooting have been performed. These include two dimensional simulations (Hurlburt et al., 1986, 1994; Dintrans et al., 2003), three dimensional simulations (Cattaneo et al., 1991; Muthsam et al., 1995; Singh et al., 1998; Stein & Nordlund, 1998; Brummell et al., 2002), three dimensional simulations with rotation (Brummell et al., 1996; Browning et al., 2004), three dimensional simulations with ionization (Rast et al., 1993; Rast & Toomre, 1993a, b), and two dimensional simulations (Bazan & Arnett, 1994, 1998). and they reveal a rather complicated picture. Depending on the density contrast, upward-moving flows are typically broader and slower moving than downward-moving flows (a trend seen in the neutrino driven convecting regions in post collapse stellar cores). Ionization regions can exagerate this trend. The downward flows traverse multiple scale heights and penetrate the stable layers below by a significant fraction of the local pressure scale height. Because of the low filling factor of the plumes, however, they do not establish an adiabatic gradient there. Convective overshooting can excite gravity waves which leads to further mixing. The use of MLT during shell oxygen burning and later nuclear burning stages is particularly problematic, as nuclear burning timescales at the base of the convecting region and convective timescales across the convective region become comparable. The simulations show inhomogeneities in the composition and strong fluctuations in space and time unlike the smooth, steady flow presupposed by one-dimensional stellar evolutionary calculations with MLT.
In conclusion we observe that MLT is a phenomenological parameterization of convection which is applied to a variety of convective phenomena in a physically motivated but crude way. Different prescriptions for MLT convection can lead to substantial differences in the interior structure of massive stars in their late evolutionary phases. We note just one example. The nonrotating models computed by Hirschi et al. (2004), who used the Schwarzschild criterion for convection with overshooting, and with convective diffusion beyond He burning, develop considerably larger Si core masses than the models computed by Rauscher et al. (2002), who used the Ledoux criterion for convection with semiconvection and without overshooting, and convective diffusion beyond He burning.
### 2.1 <sup>12</sup>C($`\alpha ,\gamma `$)<sup>16</sup>O Reaction Rate and Convection
The structure and explosive yields of massive stars depends on the mass fraction, $`X_\mathrm{C}`$, of <sup>12</sup>C left after He burning, and this depends both on the combined effects of the <sup>12</sup>C($`\alpha ,\gamma `$)<sup>16</sup>O reaction rate and the treatment of overshooting and semiconvection which governs the growth of the helium burning core (Weaver & Woosley, 1993; Thielemann et al., 1996; Imbriani et al., 2001). The triple-$`\alpha `$ reaction and the <sup>12</sup>C($`\alpha ,\gamma `$)<sup>16</sup>O reaction compete with each other, and the ratio of the two rates determines directly the ratio of <sup>12</sup>C and <sup>16</sup>O produced during core helium burning. The mixing of fresh He fuel into the He-burning core at late stages and high temperatures, when a core without growth by semiconvection and overshooting of convection eddies would have already ceased to process any He fuel, probes the <sup>12</sup>C rate at higher temperature with the effect of turning much of the remaining <sup>12</sup>C into <sup>16</sup>O. However, despite years of effort, the <sup>12</sup>C and <sup>16</sup>O cross section is still unknown to within a factor of a few (Buchmann et al., 1996; Angulo et al., 1999). Furthermore, as discussed above the treatment of convection in stellar evolutionary codes is by means of MLT, which is rudimentary and phenomenological, and cannot address the question of convective overshooting.
The implication of the uncertainty in the value of $`X_\mathrm{C}`$ left after helium burning is that its value affects the later structure of the star principally through its effect during the time interval that elapses between the end of the central C burning and the beginning of the central Ne burning. During this time the CO core experiences a phase of gravitational contraction which is partially alleviated by the formation of one (or more) convective C shell episodes. These convective episodes stop for a while the outwardly advancing C-burning front while the reservior of fuel contained in the convective shell is consumed. During this time the C-burning front remains essentially fixed in mass and slows down the contraction of the region above the front. A larger value of $`X_\mathrm{C}`$ after core carbon burning allows a more effective support of the layers above the C-burning front during C-shell burning and hence the formation of a less steep mass-radius relation. These differences in the mass-radius relations that form before the Ne ignition remain through later core contractions until the final explosion.
The situation is shown schematically in Fig. 2, which is adapted from Fig. (12) of Imbriani et al. (2001). The lines denoted by $`X_\mathrm{C}^{\mathrm{high}}`$ and $`X_\mathrm{C}^{\mathrm{low}}`$ represent the mass-radius relations for a star at the onset of core collapse having a high and a low value of $`X_\mathrm{C}`$, respectively, after core He-burning. Once the explosion commences and the shock wave moves outward, it is radiation dominated, gains or loses only a small fraction of its energy to the matter, and therefore expands essentially adiabatically (Weaver & Woosley, 1980). The temperature behind the shock is therefore a function only of its radius and the explosion energy. At the same time, the matter which is subject to complete silicon burning, incomplete silicon burning, or oxygen burning, and whose final composition therefore depends only on its initial proton fraction, $`Y_\mathrm{e}`$, is determined only by the peak post shock temperatures, and therefore by the geometrical distance of the matter from the core center. Because of the mass-radius relation (Fig. 2), the mass of material undergoing incomplete explosive Si burning and explosive O burning, which produce the bulk of the elements from <sup>28</sup>Si to <sup>55</sup>Mn, is greater for small $`X_\mathrm{C}`$. On the other hand, the lighter elements from <sup>20</sup>Ne to <sup>27</sup>Al are produced in the C convective shell and partial destroyed by the shock, and their production therefore scales with $`X_\mathrm{C}`$. Ignoring the subtleties of many of the production chains, it is seen that a large $`X_\mathrm{C}`$ favors the production of elements at the lighter end of the <sup>20</sup>Ne to <sup>55</sup>Mn range, while the opposite is true of a small $`X_\mathrm{C}`$. A dramatic illustration of this trend is shown in Fig. 4 of Weaver & Woosley (1993), who tried to put limits on the <sup>12</sup>C($`\alpha ,\gamma `$)<sup>16</sup>O reaction rate by requiring that the final explosive yields to have a scaled solar relative abundance. (Arnett (1971) made an analogous attempt to fix the <sup>12</sup>C($`\alpha ,\gamma `$)<sup>16</sup>O reaction rate by using the observed <sup>12</sup>C to <sup>16</sup>O ratio.) It must be remembered that these attempts to fix the <sup>12</sup>C($`\alpha ,\gamma `$)<sup>16</sup>O reaction rate by using the results of stellar evolutionary calculations actually fix a combination of this rate and the particular MLT scheme used.
Further problems with MLT convection for the evolution of massive stars, in brief, are the fact that it fails to deal with the interaction between convection and rotation, and the generation and transport of magnetic fields (Zahn, 1999), and convective nuclear burning (Bazan & Arnett, 1994, 1998; Asida & Arnett, 2000). Attempts have been made to overcome the limitations of MLT by a first principles approach (i.e., direct numerical solutions of the fundamental equations) (e.g., Stein & Nordlund, 1998; Singh et al., 1998; Deupree, 1998; Asida & Arnett, 2000), or by more sophisticated convection models (e.g., Canuto & Mazzitelli, 1991, 1992), but these have not made there way into evolutionary calculations of massive stars to core collapse.
### 2.2 Weak Interactions
Weak interactions affect Y<sub>e</sub>, the proton fraction, and therfore play an important role in determining both the presupernova stellar structure and the nucleosynthesis. They affect the structure because, at all times, the pressure is mostly due to electrons - at first, nonrelativistic and nondegenerate, but later neither. They affect the nucleosynthesis because the synthesis of all nuclei except those with equal numbers of neutrons and protons is sensitive to Y<sub>e</sub>.
The weak interaction rates after oxygen burning are particularly difficult to calculate as a large number of excited states with uncertain properties become populated so that their decay must be dealt with statistically. Early attempts in this direction were made by Hansen (1968); Mazurek et al. (1974); Takahashi et al. (1973). However, it was (Fuller et al., 1980, 1982b, 1982a; Fuller, 1982; Fuller et al., 1985) who recognized the key role played by the Gamow-Tellar resonance and noted that measured decay rates exploited only a small fraction of the available strength. More recently new shell-model calculations of the distribution of Gamow-Tellar strength have resulted in an improvedโand often reducedโestimate of its strength (Mart nez-Pinedo & Langanke, 1999; Langanke & Mart nez-Pinedo, 2000; Mart nez-Pinedo et al., 2000; Langanke & Mart nez-Pinedo, 2003). Inclusion in presupernova evolutionary models of these new rates for electron capture and beta dacay (Heger et al., 2001) lead to slightly higher central proton fractions and smaller outer core entropies at the time of core collapse, leading to slightly smaller iron core masses. Incorporation of the new rates in core collapse simulations (Langanke et al., 2003; Hix et al., 2003) lead to an increased importance of nuclear vs free proton electron capture and reduced initial mass behind the shock with lower densities, proton fractions, and entropies. However, the reduced electron capture in the outer layers slows their collapse and allows the shock to reach a slightly larger maximum radius. The collapsing core encounters a range of large and neutron rich nuclei whose beta strengths have not yet been calculated in detail, underscoring the need for more work in this area.
### 2.3 Rotation
Massive stars are observed to be rapid rotators, with equatorial velocities spanning the range $`100400`$ km $`s^1`$ (Fukuda, 1982; Halbedel, 1996; Penny, 1996; Howarth et al., 1997). As a consequence, a number of instabilities leading to composition mixing and angular momentum transport are predicted to occur within these stars as they evolve, leading to differences in the structure of supernova progenitors. Furthermore, the rotation rate of progenitor cores may play a role in the supernova mechanism and is dependent on the degree to which angular momentum transport has occurred during the course of prior evolution.
A number of observations point to the operation of rotationally induced mixing processes in massive stars. The ratio B/R, the number of blue to red supergiants, increases with the metalicity, Z, (e.g., Langer & Maeder, 1995; Maeder & Meynet, 2000), and this cannot be accounted for by mass loss or convection. For a number of reasons (Maeder & Meynet, 2001) rotation favors the development of the red supergiant structure. The increase of the B/R ratio with Z results from the increase of the mass loss rate with Z, and with it the loss of angular momentum of the star, rapidly reducing its rotation rate during the MS phase and thus reducing its propensity to become a red supergiant during later phases. Rotational mixing in the radiative envelopes of massive stars will modify their surface abundances. One would naively expect a depletion of an initial surface abundance of fragile nuclei, such as <sup>3</sup>He, <sup>6</sup>Li, <sup>7</sup>Li, <sup>9</sup>Be, <sup>10</sup>B, and <sup>11</sup>B, mixed down and destroyed by proton capture at higher interior temperatures. At the same time, hydrogen burning in massive stars is governed by the CNO cycle, and this has the effect of converting most of the initial <sup>12</sup>C and <sup>16</sup>O into <sup>14</sup>N. Rotational mixing to the surface of material in which the CNO cycle was operative should be manifested as a depletion of <sup>12</sup>C and <sup>16</sup>O and an enhancement of <sup>14</sup>N. These effects have been observed. For example, Proffitt & Quigley (2001); Venn et al. (2002) have observed boron depletions in B type stars in OB associations, consistent with the predictions of Fliegner et al. (1996) and the rotating models of Heger & Langer (2000). Some non-supergiant B stars show a moderate increase in N abundance (Gies & Lambert, 1992; Lennon et al., 1996).
Ideally, the evolution of rotating stars should be calculated multi-dimensionally, with the composition and angular momentum transport arising directly from the calculation itself. This program cannot be carried out with current computer resources. Rather, the equations of stellar structure are kept one-dimensional. Initially this was accomplished by replacing the usual spherical coordinates by new coordinates characterizing the equipotentials (which have cylindrical symmetry) (Kippenhahn et al., 1970). More recently, this is accomplished by making the assumption (Zahn, 1992) of highly anisotropic turbulence in radiative layers. In particular, turbulence generated by, say, shear in the presence of differential rotation is much stronger in the direction perpendicular to gravity (โhorizontal directionโ) than in the vertical direction, the latter being suppressed by the stable vertical stratification. If true, the strong horizontal turbulence makes the angular velocity $`\mathrm{\Omega }`$ and the composition nearly constant on isobaric surfaces, rather than cylinders, giving rise to a โshellularโ rotation law. In this case, the motion is not cylindrical. Nevertheless, a consistent 1-D scheme has been formulated (Meynet & Maeder, 1997, 2000).
The critical assumption of highly anisotropic turbulence in radiative stellar zones has indirect observational support, both in the fact that turbulent motions caused by shear stresses in the Earthโs atmosphere are highly anisotropic in those regions where the stratification is stable, and in the study of the solar tachocline (Spiegel & Zahn, 1992). (The tachocline is the transition zone between the rigid rotation in the radiative interior and the external convective zone, where rotation varies with latitude.) If the horizontal turbulence is intense, then the tachocline is very thin, and the latter is supported by helioseismological observations.
Keeping the equations of stellar structure one-dimensional allows stellar evolutionary calculations to be performed, but requires that various instabilities leading to angular momentum transport and the mixing of chemical elements, which play a major role in massive star evolution, be parameterized. Since the diffusion and advection of composition and angular momentum operate on $`\mathrm{\Omega }`$ and $`(r\mathrm{sin}\theta )^2\mathrm{\Omega }`$, respectively, their vertical transport rates are different, being much smaller for the composition. Gradients in composition ($`\mu `$-gradients) that develop during the evolution of the star tend to reduce the vertical transport rates, so the effect of these $`\mu `$-gradient effects must either also be parameterized (Heger et al., 2000) or attempt to incorporate them more consistently in the instability and mixing algorithms (Maeder & Zahn, 1998).
Rotation in convective zones is relatively easy to handle until oxygen shell and particularly silicon burning. Chemical homogeneity can be assumed and, if the high viscosity associated with turbulence tends to solid-body rotation, then rigid body rotation can also be assumed. An alternative (Endal & Sofia, 1976) is that convection preserves the angular momentum of the convective elements leading to equalization of the specific angular momentum. Supporting the tendency towards rigid body rotation over alternatives, however, is the observation that the solar convection zone deviates from solid body rotation by less than 5% (Antia et al., 1998). Complications in handling convective zones begin with oxygen shell burning. The times scales for convective mixing, nuclear burning and angular momentum transport become similar, and the feasibility of constructing self-consistent models with one-dimensional evolution equations becomes highly suspect.
At convective boundaries and in radiative zones a number of instabilities can lead to significant transport of composition and angular momentum (e.g., Endal & Sofia, 1978; Heger et al., 2000; Meynet & Maeder, 2000; Maeder & Meynet, 2000, and many others). These include the Eddington-Sweet circulation (von Zeipel, 1924; Eddington, 1925; Vogt, 1925) (a circulation that arises because a component of the radiation stress is directed along equipotential surfaces), shear instabilities (Spiegel & Zahn, 1970; Zahn, 1974) (dynamical: arising when the free energy in differentially rotating layers exceeds the work against restoring forces required to adiabatically overturn the fluid; secular: as above but allowing thermal diffusion in the overturning fluid to remove a stabilizing temperature gradient), the Solberg-Hรธiland instability (Tassoul, 2000) (analogous to the Ledoux criterion but including the angular momentum gradients), and the Goldreich-Schubert-Fricke instability (Goldreich & Schubert, 1967; Fricke, 1968) ((1) a secular analogue to the Solberg-Hรธiland stability criterion, and (2) a criterion for the generation of meridional flows for nonconservative rotation profiles).
During the main-sequence evolution of rapidly rotating massive stars, angular momentum transport in the outer radiative regions, principally by the Eddington-Sweet circulation, quickly establishes a steady state in which the diffusion of angular momentum is balanced by advection of angular momentum due to circulation (Zahn, 1992; Urpin et al., 1996; Talon et al., 1997). Neglecting angular momentum loss at the surface, this leads to a differential rotation in which the angular velocity at the convective core boundary is about a factor of 1.15 that at the surface.
Concerning the late evolutionary stages, which are of most interest to supernova modelers, there have been two recent stellar evolutionary calculations of rotating stars that have been carried out to these stages (without magnetic fields). These are by Heger et al. (2000) (HLW), who evolve to core collapse, and Hirschi et al. (2004) (HMM), who evolve through central oxygen or silicon burning. The two groups employed different numerical methods of incorporating the effects of rotation, and different parameterizations of convection.
HMM find that for stars with zero age main-sequence masses (M<sub>MS</sub>) $`<`$ 30 M rotation tends to increase the mass of the iron core prior to collapse. This trend is confirmed by HLW who confine their study to M$`{}_{\mathrm{MS}}{}^{}<25`$ M. The larger iron cores in the rotating models result from rotational mixing in prior evolutionary phases, mainly the H-burning phase where the mixing and nuclear timescales are comparable. The larger He cores resulting from the rotational mixing during H-burning cause them to have lower densities and higher temperature during He-core-burning, and this leads to lower C to O ratios at core He exhaustion. HLW, who parameterize the inhibiting effect of $`\mu `$-gradients on mixing, find on varying this parameter that the final iron core masses are a sensitive functions of this parameter. The more efficient the rotational mixing (or less strong are the inhibiting effects of the $`\mu `$-gradients), the greater the core masses. (This underscores the fact that the precollapse structure of rotating stars is highly dependent on approximate numerical treatments of complicated physics.) The larger precollapse iron core models resulting from rotation imply, of course, a smaller initial M<sub>MS</sub> that will lead to core collapse. HMM obtain significantly greater core masses then HLW, due to their different treatments of rotational mixing and convection, again underscoring the effect of different approximations. Both groups find that only convective process are rapid enough to notably redistribute angular momentum after core helium burning The effect is to leave each convective region with a constant angular velocity. The distribution of angular momenta in the final models is therefore characterized by rounded saw-tooth patterns, each saw tooth being the constant angular velocity imprint of a convective zone. For M<sub>MS</sub> $`>`$ 30 M, HMM find that rotationally enhanced mass loss causes the star to enter the Wolf-Rayet phase earlier and to therefore spend more time undergoing heavy mass loss. This erodes the star and results in smaller cores at the pre-supernova stage.
The two groups find that the angular momentum of precollapse cores tend to converge to a value that would imply a rotation rate upon collapse to a neutron star of about 1 ms or less, which is near breakup for even the slowest rotators. The corresponding angular momentum is about 100 times the angular momentum of the fastest rotating pulsars observed. The implication of this in unclear, as it is not yet established how rapidly newly formed pulsars rotate. For example, the fastest rotating young pulsar, PSR J053726910 (Marshall et al., 1998), has a period of 16 ms. But with its estimated age of $`5\times 10^3`$ yr an extrapolation to an initial rotation rate of $`1`$ ms, while not the only possibility, is not unreasonable. The Crab (and by implication others), on the other hand, may never have rotated near breakup (Trimble & Rees, 1970) unless this rotational energy ($`3\times 10^{52}`$ ergs) were radiated as gravitational waves. Otherwise this energy would surely have been seen in the optical and manifested now in the expansion velocity of the nebula.
### 2.4 Rotation and Mass Loss
Mass loss from the stellar surface (stellar winds) significantly affects the evolution of massive stars, particularly the Upper MS stars (Chiosi & Maeder, 1986; Maeder & Meynet, 1987) where an appreciable percentage of the mass of these star can be lost during their evolution. Type 1b and 1c supernovae probably arise from stars having suffered extensive mass loss. Mass loss also affects rotating stars, as these winds can transport large amounts of angular momentum out of the star. This is particularly true for equatorial mass loss by anisotropic stellar winds (Maeder, 1999). However, the uncertainties in these mass-loss rates, particularly for red supergiants (Lamers & Cassinelli, 1999) are considerable due both to the uncertain physics involved and the uncertainties in the observational data and their interpretation. (For example, mass loss rates from hot O and B stars have rencently been revised, generally downward (Nugis & Lamers, 2000; Vink et al., 2000, 2001; Bouret et al., 2004), owing to a improved treatments of clumping and multiple scatering.
Stellar evolutionary calculations of nonrotating stars with mass find that stars with solar metalicity initially more massive than $`3035`$ M converge to a hydrogen free star of roughly 5 M (Maeder, 1990; Schaller et al., 1992; Woosley et al., 1993; Meynet et al., 1994). This assumes the loss of the hydrogen envelope either during the main sequence phase (as luminous blue variables) or as red supergiants), and a mass loss rate for hydrogenless Wolf-Rayet stars being given by a positive power of the remaining mass (Langer, 1989b, a). The latter causes the convergence in mass. However, while the masses of these stars may converge, the thermal and chemical structures of these stars retain some memory of their former masses (Woosley et al., 1993). For example, a helium core undergoing carbon burning and trimming down from a larger mass will have a higher temperature and lower density than a constant mass helium core of the same final mass. It will therefore burn more of its helium to oxygen than its constant mass counterpart. Thus, despite similar final iron core masses, stars that have trimmed down from larger initial masses will have larger carbon-oxygen mantles with larger oxygen to carbon ratios and shallower density gradients in the outer regions than stars of the same final mass that have trimmed down from smaller initial masses. Nonrotating stars that do not succeed in loosing their hydrogen envelopes during hydrogen and core helium burning have internal structures that are little affected by the mass loss, as the interior evolution is largely decoupled from the surface. However, these stars may give rise to different supernova types (Type IIL versus Type IIP, for example) depending on the remaining mass of the hydrogen envelope.
Rotating stars undergoing mass loss can lose considerable angular momentum during evolution (Packet et al., 1980; Heger & Langer, 1998). The mass loss itself is affected by rotation through centrifugal forces, nonradial forces, and gravitational darkening (von Zeipelโs theorem). The mass loss rate is increased by rotation, but it is now appreciated that it could be either oblately (i.e., predominantly equatorial) or prolately (i.e., predominantly polar) asymmetrical (Owocki et al., 1998; Petrenz & Puls, 2000), with much less angular momentum being lost by the star for a given amount of mass lost in the latter case.
### 2.5 Rotation and Magnetic Fields
The growth of magnetic fields and their influence on the evolution of rotating stars has recently begun to be considered in stellar evolutionary calculations of stars to core collapse. The prime importance of magnetic fields here is in their potential for transferring angular momentum. As discussed in the previous Section, without magnetic fields the various hydrodynamic instabilities leading to angular momentum transport (as presently understood) are two weak to prevent the cores of supernova progenitors from ending up with about 100 times the angular momentum of the shortest period young pulsar. While this large angular momentum might be appropriate for the collapser model of gamma-ray bursts (MacFadyen et al., 2001), it is likely too high to account for pulsar rotation rates at birth. Magnetic effects might provide additional angular momentum transport during stellar evolution. A possible case in point is the Sun, whose near-uniform rotation of the radiative core and the small difference in rotation rate between the core and the convective envelope, the latter being continuously spun down by the solar wind torque, has been established through helioseismology (Corbard et al., 1997; Schou et al., 1998). On the other hand, it has been known for a long timed that hydrodynamic instabilities alone are incapable of accounting for this (e.g., Spruit et al., 1983).
In convective zones, angular momentum transport by turbulent viscosity is very efficient. Where the transport of angular momentum by magnetic fields may be critical is in radiative zones. Spruit (1999, 2002) has summarized much of the literature pertaining to this issue and described how a magnetic dynamo based on the an instability studied by Taylor (1973) and others and tapping the free energy available from differential rotation could operate in stably stratified zones. The basic picture is that differential rotation will wind up an initially weak field producing a predominantly toroidal (azimuthal) field. This field is subject to a number of instabilities, the m = 1 Taylor instability, exhibited in Fig. 3, being likely the most relevant. The poloidal component generated by this instability will be stretched into a strong toroidal component which will in turn be subjected to instabilities, forming a dynamo. An estimate of the equilibrium strength of the field generated by this โTaylor-Spruitโ dynamo is made by equating the growth timescale to the attenuation timescale by mangetic diffusivity. The radial component over the largest unstable lengths is chosen to evaluate these timescales as this determines the maximum saturation field. The result is obtained for the two limiting cases in which the stability is provided either by a thermal or $`\mu `$ gradient.
Comparisons of estimates of the angular momentum transport by the magnetic fields generated by the Taylor-Spruit dynamo with the angular momentum transport by hydrodynamic instabilities indicate that the former could be of major importance (Maeder & Meynet, 2004). A recent stellar evolutionary calculation (Heger et al., 2004) incorporating the estimates of magnetic torques produced by the Taylor-Spruit magnetic fields show a reduction by a factor of 10 in the final iron core angular momentum. It must be observed, however, that angular momentum transport by magnetic fields has the potential of being completely ineffective in a star or so effective as to lead to near uniform rotation throughout the entire star with the result of a much too slow rotation of the remnant (Spruit & Phinney, 1998). Thus, while the above results are encouraging, the fact that magnetic fields have an almost โjust soโ effectiveness must regarded as extremely preliminary.
### 2.6 Global Asymmetries
A potentially important effect in the late evolutionary stages of massive stars, which would require a multidimensional evolutionary code to follow, is the generation of overstable g-modes driven by shell nuclear burning (Goldreich et al., 1996). The idea is that nuclear burning rates in silicon and oxygen burning shells are extremely temperature sensitive. Consider an $`\mathrm{}=1`$ perturbation of the shell to the right, for example. The left-hand side of the shell will then be compressed and heated. Nuclear burning in the compressed region will be greatly enhanced and will generate a large local overpressure which will push the shell back to the left. If overstable, this g-mode will oscillate with growing amplitude with the very interesting possibility of generating a significant global asymmetry of the stellar core just prior to collapse.
### 2.7 Conclusions
The evolutionary calculations of massive stars to core collapse are extremely difficult, involving a variety of physics on multiple length and time scales and in multiple dimensions. Teams are incorporating more realistic physics into the stellar codes, but are still constrained to model inherently multidimensional phenomena in one dimension. This review has attempted to make supernova modelers aware of the many approximations and parameterizations that are perforce made in the course of evolving a star from the main-sequence to core collapse.
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# Magnetism via superconductivity in SF proximity structures
## I Introduction
Proximity effects are phenomena stipulated by a โpenetrationโ of an order parameter (of some state) from one material into another, which does not possess such type of the order itself, due to the materials being in contact. The leakage of superconducting correlations into a non-superconducting material is an example of the superconducting proximity effect. For a nonmagnetic normal metal (N) in contact with a superconductor (S), the proximity effect has been intensively studied and well understood many years ago . However, in the case of SN structures we deal with a single type of order - superconductivity. When a normal nonmagnetic metal is replaced by a ferromagnet (F), the physics of proximity effect is much more interesting and rich \[2-20\]. There are two competing states with different order parameters: superconductivity and ferromagnetism. Due to the difference in energy between spin-up and spin-down electrons and holes under the exchange field of a ferromagnet, a singlet Cooper pair, adiabatically injected from a superconductor into a ferromagnet, acquires a finite momentum $`\mathrm{\Delta }pH_e/\mathrm{}v_F`$ (here $`H_e=\mu _Bh_F`$ is an extra energy caused by the intrinsic magnetic field $`h_F`$ in ferromagnet; $`v_F`$ is the Fermi velocity, $`\mathrm{}`$ is the Planck constant, and $`\mu _B`$ is the Bohr magneton). As a result, proximity induced superconductivity of the F layer is spatially inhomogeneous and the order parameter contains nodes where the phase changes by $`\pi `$. Particularly, transport properties of tunnel SF structures in the $`\pi `$-phase state have turned out quite unusual. The phase shift of $`\pi `$ in the ground state of the junction is formally described by the negative critical current $`J_C`$ in the Josephson current-phase relation: $`J(\phi )=J_C\mathrm{sin}(\phi )`$. The $`\pi `$-phase state of an SFS weak link due to Cooper pair spatial oscillation was first predicted by Buzdin et al., . Experiments that have been performed by now on SFS weak links \[6-8\] and SIFS tunnel junctions directly prove the $`\pi `$-phase superconductivity (I denotes an insulator). Planar tunneling spectroscopy also reveals a $`\pi `$-phase shift in the order parameter, when superconducting correlations coexist with ferromagnetic order . The superconducting phase was also measured directly using SQUIDโs made of $`\pi `$-junctions.
There is another interesting case of a thin F layer, $`d_F<<\xi _F`$, being in contact with an S layer. As far as the thickness of the F layer, $`d_F`$, is much less than the corresponding superconducting coherence length, $`\xi _F`$, there is spin spiltting but there is no order parameter oscillation in the F layer. Surprisingly, but it was recently predicted Refs. that for SFIFS tunnel structures with very thin F layers one can, if there is a parallel orientation of the F layers magnetization, turn the junction into the $`\pi `$-phase state with the critical current inversion; if there is an antiparallel orientation of the F layers internal fields, one can even enhance the tunnel current. It was shown in Refs. , that physics behind the inversion and the enhancement of the supercurrent in this case differs from that proposed by Buzdin et al.
While proximity induced superconductivity of the F metal in SF hybrid structures has been intensively studied, much less attention has been paid to a modification of the electron spectrum of a superconductor in a region near the S/F interface due to a leakage of magnetic correlation into the superconductor. Some feature of the induced magnetism (e.g., the spin-splitting of the density of states) were found by numerical calculations in Refs. \[16-18\]. To our knowledge, only recently the question of S metal magnetization has been addressed in Refs. . (Here we do not touch S/FI systems, where FI stands for a ferromagnetic insulator (semiconductor). In such systems conduction electrons penetrate the magnetic layer on much smaller distances than in the case of metals and are totally reflected at the S/FI boundary. The S/FI boundary being magnetically active rotates spins of reflected electrons. This spin rotation occurs only as a result of a tunneling by the quasiparticle into the classically forbidden region of the boundary. Due to the spin rotation the exchange field is induced in a superconductor on a distance of order of superconducting coherence length $`\xi _S`$ near an S/FI surface \[21-23\]. However, in contrast to ferromagnetic metals, where the proximity effect is pronounced, this effect is drastically reduced in S/FI structures.)
The investigation of a โmagnetic proximityโ effect in SF nanostructures is the purpose of this report. To tackle the physics, we consider an interesting and practicable case of an SF structure of a massive superconducting and thin ferromagnetic layers. Using a quasiclassical theory of superconductivity for proximity coupled bilayer (Sec. II), we will show that for some limits the problem can be solved analytically. Two limits will be discussed here: (i) a weak and (ii) a strong proximity effect. Section III is the key one; here we describe the examples of the magnetic proximity effect manifestation. We show that due to induced magnetism of the S metal: (i) the superconducting phase jumps at the S/F interface; as well as, there are (ii) additional suppression of the order parameter near the S/F interface; (iii) the spin splitting of the quasiparticle density of states (DOS); (iv) the appearance of the local bands inside the energy gap; and, directly, in (v) induced equilibrium electronic magnetization of the S layer that spreads over distance of the order of the superconducting coherence length $`\xi _S`$. We also briefly discuss recent experiments. Summarizing the results in Conclusion we draw attention to the fact that in the general case, for proximity coupled SF hybrid structures both phenomena - induced superconductivity of the F metal and induced magnetism of the S metal - take place simultaneously and should be considered self-consistently.
## II Quasiclassical theory of superconductivity of SF bilayer
### II.1 Bilayer model
Let us consider proximity effects in the bilayer of a massive superconducting ($`d_S>>\xi _S`$) and a thin ferromagnetic ($`d_F<<\xi _F`$) metals, with arbitrary transparency of the S/F interface. Here $`\xi _S=(D_S/2\pi T_C)^{1/2}`$ and $`\xi _F=(D_F/2H_e)^{1/2}`$ stand for the superconducting coherence lengths, $`D_{S,F}`$ are the diffusion coefficients, $`d_{S,F}`$ are the thicknesses of the S and F layers. (Henceforth, we have taken the system of units with $`\mathrm{}=k_B=1`$.) We assume the โdirtyโ limit for both metals, i.e., $`\xi _{S,F}>>l_{S,F}`$ where $`l_{S,F}`$ are the electron mean free paths. It is also assumed that the superconducting critical temperature of the F material equals zero. All quantities are assumed to depend only on a single coordinate $`x`$ normal to the interface surface of the materials. We also expect that the F layer has a homogeneous (monodomain) magnetic structure with magnetization aligned parallel to the interface, so that there is no spontaneous magnetic flux penetrating into the S layer. Under these conditions, the only magnetic interaction which can affect the superconductor is the short-range exchange interaction between the superconducting quasiparticles and magnetic moments into the ferromagnet.
### II.2 Main equations
As is well know, the superconductivity of โdirtyโ metals is conveniently described by the quasiclassical Usadel equations for the normal, $`G_{\sigma \sigma ^{}}(x,\omega )`$ and $`\stackrel{~}{G}_{\sigma \sigma ^{^{}}}(x,\omega )`$, and anomalous, $`F_{\sigma \sigma ^{^{}}}(x,\omega )`$ and $`\stackrel{~}{F}_{\sigma \sigma ^{^{}}}(x,\omega )`$, Green functions, integrated over energy and averaged over the Fermi surface. (Green functions are defined in a standard way, see, e.g. Ref. ). It can be shown, that for singlet pairing and in the absence of spin-flip scattering, the whole system of Usadel equations decomposes into two equivalent subgroups, which go over to each other under interchange of the spin indices ($`\sigma =,`$) $``$ and reversal of the exchange field sign, $`H_eH_e`$
It is convenient to take into account the normalization of the Green function, $`G_F\stackrel{~}{G}_F+F_FF_F^+=1`$, explicitly and to introduce modified Usadel functions $`\mathrm{\Phi }_S`$, $`\mathrm{\Phi }_F`$, defined by the relations $`\mathrm{\Phi }_S=\omega F_S/G_S`$, $`\mathrm{\Phi }_F=\stackrel{~}{\omega }F_F/G_F`$, etc. Then we can recast the equations for the S layer in terms of these functions. We specialize the discussion to a geometry when all quantities depend on a single coordinate $`x`$, normal to the S/F interface. For the superconducting metal we have ($`x0`$):
$$\mathrm{\Phi }_S=\mathrm{\Delta }_S+\xi _S{}_{}{}^{2}\frac{\pi T_C}{\omega G_S}[G_S{}_{}{}^{2}\mathrm{\Phi }_{S}^{}]^{},\text{ }G_S=\frac{\omega }{(\omega ^2+\mathrm{\Phi }_S\stackrel{~}{\mathrm{\Phi }}_S)^{1/2}},$$
(1)
with the superconducting order parameter $`\mathrm{\Delta }_S(x)`$ determined by the self-consistency equation:
$$\mathrm{\Delta }_S\mathrm{ln}(T/T_C)+2\pi T\underset{\omega >0}{}[(\mathrm{\Delta }_S\mathrm{\Phi }_SG_S)/\omega ]=0,$$
(2)
Here the prime denotes differentiation with respect to a coordinate $`x`$, and in Eq. (2) the summation over frequencies is cut off by the Debye frequency $`\omega _D`$. For the F metal we have ($`d_Fx<0`$):
$$\mathrm{\Phi }_F=\xi ^2\frac{\pi T_C}{\stackrel{~}{\omega }G_F}[G_F{}_{}{}^{2}\mathrm{\Phi }_{F}^{}]^{},\text{ }G_F=\frac{\stackrel{~}{\omega }}{(\stackrel{~}{\omega }^2+\mathrm{\Phi }_F\stackrel{~}{\mathrm{\Phi }}_F)^{1/2}}$$
(3)
Here $`\stackrel{~}{\omega }=\omega +iH_e`$, and $`\omega \omega _n=\pi T(2n+1)`$, $`n=\pm 1,\pm 2,\pm 3,\mathrm{}`$ is Matsubara frequency. Assuming the symmetry of the system with respect to the rotation in spin space both in the F and in the S layers, we drop the spin indices, apart from the specified cases. We also put for the F metal a vanishing value of the bare superconducting order parameter $`\mathrm{\Delta }_F=0`$, while the pair amplitude $`F_F0`$ due to proximity with the superconductor.
The equations for the functions $`\stackrel{~}{\mathrm{\Phi }}_S`$ and $`\stackrel{~}{\mathrm{\Phi }}_F`$ have a form analogous to (1)โ(3); note that $`\stackrel{~}{\mathrm{\Phi }}(\omega ,H_e)=\mathrm{\Phi }^{}(\omega ,H_e)`$. In Eq. (3) we write our formulas for the F metal using the effective coherence length of normal nonmagnetic (N) metal with the diffusion coefficient $`D_F`$, $`\xi =\left(D_F/2\pi T_C\right)^{1/2}`$, instead of $`\xi _F=\left(D_F/2H_e\right)^{1/2}`$, to have a possibility to analyze both limits $`H_e0`$ (SN bilayer) and $`H_e>>\pi T_C`$. The relation on the ferromagnetic layer thickness one may read as $`d_F<<\mathrm{min}(\xi _F,\xi )`$.
### II.3 Boundary conditions
The Eqs. (1)-(3) should be supplemented with the boundary conditions in the bulk of the S metal and at the external surface of the F layer. Far from the S/F interface, $`x>>\xi _S`$, for the S layer we have the usual boundary conditions in the bulk of the S metal: $`\mathrm{\Phi }_S(\mathrm{})=\mathrm{\Delta }_S(\mathrm{})=\mathrm{\Delta }_0(T)`$, where $`\mathrm{\Delta }_0(T)`$ is the BCS value of the order parameter. At the external surface of the F metal $`\mathrm{\Phi }_F^/(d_F)=0`$. The relations at the S/F interface we obtain by generalizing the results of Kupriyanov and Lukichev for interface between two superconductors.
The first condition on the Usadel equations ensures continuity of the supercurrent flowing through the S/F boundary at any value of the interfacial transparency. Going over to the modified Usadel functions $`\mathrm{\Phi }_S`$ and $`\mathrm{\Phi }_F`$, the first boundary condition has the form:
$$\frac{1}{\stackrel{~}{\omega }}\gamma \xi G_F{}_{}{}^{2}\mathrm{\Phi }_{F}^{}|_{x=0}=\frac{1}{\omega }\xi _SG_S{}_{}{}^{2}\mathrm{\Phi }_{S}^{}|_{x=0},$$
(4)
Here $`\gamma =\rho _S\xi _S/\rho _F\xi `$ is the proximity effect parameter, which characterizes the intensity of superconducting correlations induced in the F layer, and vice versa, an intensity of magnetic correlation induced into the S layer; $`\rho _{S,F}`$ are the resistivities of the metals in the normal state.
The boundary condition (4) takes into account the effect of quasiparticle DOS of the metals in contact. The second relation takes into consideration the effects of a finite transparency (electrical quality) of the interface. For $`\mathrm{\Phi }(\omega ,x)`$ \- parametrization the second boundary condition becomes
$$\xi \gamma _{BF}G_F\mathrm{\Phi }_F^{}|_{x=0}=\stackrel{~}{\omega }G_S(\mathrm{\Phi }_S/\omega \mathrm{\Phi }_F/\stackrel{~}{\omega })|_{x=0},$$
(5)
where $`\gamma _{BF}`$ is the parameter that characterizes the effects of a finite transparency of the interface. For $`\gamma _{BF}=0`$, i.e., for a fully transparent boundary, condition (5) goes over to $`\mathrm{\Phi }_S/\omega =\mathrm{\Phi }_F/\stackrel{~}{\omega }`$. The expression for $`\gamma _{BF}`$ can be written through more convenient values: $`\gamma _{BF}=R_B/\rho _F\xi `$, where $`R_B`$ is the product of the S/F boundary resistance and its area .
The relations (4) and (5) generalize the proximity effect problem with an arbitrary interface transparency for the case of a normal metal with ferromagnetic order. The additional physical condition we assumed is that the exchange splitting of the momentum subbands, $`p_F^\pm =\sqrt{2m}\sqrt{E_F\pm H_e}`$, is substantially smaller than the Fermi energy $`E_F`$ ($`m`$ is the effective mass of an electron). For most magnetic materials the momentum renormalization is not so important as the frequency renormalization, because $`H_e>>\omega _n\mathrm{~}\pi T_C`$ while $`H_e<<E_F`$ and due to this the difference in the DOS and transparencies of the S/F interface for electrons with opposite spin orientations can be neglected.
According to the Green functions formalism, if the functions $`G_{S,F}(x,\omega )`$ and $`F_{S,F}(x,\omega )`$ are known, that is all we need to be able to describe, at least in principle, any superconducting and magnetic properties of the system. We draw attention to the feature important for further conclusions: due to superconductivity these is only a single space length - the respective superconducting coherence length, $`\xi _S`$ or $`\xi _F`$, - that encounters into the differential equations (1) and (3). So, due to superconductivity the coordinate dependences of both superconducting and magnetic properties of each layer have the same typical space scale.
### II.4 Analytical solutions
The proximity effect for an SF structure with a thin F metal, $`d_F<<(\xi _F,\xi )`$, can be reduced to consideration of the boundary value problem for the S layer . Indeed, the differential equation (3) can be solved by iteration with respect to the parameter $`d_F/\xi _F`$ ($`d_F/\xi `$). To a first approximation one can neglect the nongradient term and, taking into account that $`\mathrm{\Phi }_F^{^{}}(d_F)=0`$, we obtain $`\mathrm{\Phi }_F(x)=const`$. In the next approximation in $`d_F/\xi `$ we find, after linearizing Eq.(3),
$$\mathrm{\Phi }_F(\omega ,x)=\frac{\stackrel{~}{\omega }\mathrm{\Phi }_F(\omega ,0)}{\xi \pi T_CG_F}\left(x+d_F\right)$$
(6)
Here we have again taken into account the condition that $`\mathrm{\Phi }_F^{^{}}(d_F)=0`$. Determining $`\mathrm{\Phi }_F^{^{}}(0)`$ from Eq. (6) and substituting it into boundary conditions (4) and (5), we obtain the boundary condition for the function $`\mathrm{\Phi }_S(\omega ,x)`$. We have (here we restore the spin index):
$$\xi _SG_S\mathrm{\Phi }_S^{^{}}|_{x=0}=\gamma _M\stackrel{~}{\omega }_\sigma \mathrm{\Phi }_S\left[\pi T_C\left(1+\frac{2G_S\gamma _B\stackrel{~}{\omega }_\sigma }{\pi T_C}+\frac{(\gamma _B\stackrel{~}{\omega }_\sigma )^2}{(\pi T_C)^2}\right)^{1/2}\right]^1|_{x=0}$$
(7)
where $`\stackrel{~}{\omega }_\sigma \omega +i\sigma H_e`$. The unknown value of the function $`\mathrm{\Phi }_F(\omega ,x=0)`$ is defined by the relation:
$$\mathrm{\Phi }_F(\omega ,0)=G_S\mathrm{\Phi }_S\left[\omega \left(\frac{\gamma _B}{\pi T_C}+\frac{G_S}{\stackrel{~}{\omega }_\sigma }\right)\right]^1|_{x=0}$$
(8)
We introduce the effective boundary parameters, $`\gamma _M=\gamma d_F/\xi `$ and $`\gamma _B=\gamma _{BF}d_F/\xi `$, instead of $`\gamma `$ and $`\gamma _B`$. As a result, the problem of the proximity effect for a massive superconductor with a thin ferromagnet layer reduces to solving the equations (1) and (3) for a semi-infinite S layer with the boundary conditions (7) on the external side and BCS type on infinity. The spatial dependence of the function $`\mathrm{\Phi }_F(\omega ,x)`$ in the F layer can be neglected due to the mesoscopic thickness $`d_F<<(\xi ,\xi _F)`$ of the latter; Eq. (8) determines the value of $`\mathrm{\Phi }_F(\omega ,0)`$.
One can directly see, that via the boundary condition, Eq. (7), electronic spin โupโ and spin โdownโ subbands lost its equivalence in the S matel too. Spin discrimination means magnetism of a metal. The penetration of the magnetic correlation into the superconducting layer is governed by the proximity effect parameter $`\gamma _M`$, i.e., by the electron density of states on contacting metals. For high quasiparticleโs density in the F metal in comparison to that in the S counterpart (a large value of $`\gamma _M`$) the equilibrium diffusion of these quasiparticles into the superconductor leads to an effective leakage of magnetic order into the S layer and strong suppression of superconductivity near the S/F interface. In the opposite case, $`\gamma _M<<1`$, the influence of the F layer on properties of the S metal is weak; it even vanishes if $`\gamma _M0`$. Opposite is the behavior of the superconductivity on this parameter. So, to increase magnetic correlation near the S/F interface one should increase the parameter $`\gamma _M`$; in order to increase superconducting correlation one should decrease this parameter. Of course, the electric quality of the interface is also important for the penetration of magnetic and superconducting correlations from one metal into another.
There are three parameters which enter the model: $`\gamma _M`$ is the measure of the strength of proximity effect between the $`S`$ and $`F`$ metals, $`\gamma _B`$ describes the electrical quality of the SF boundary, and $`H_e`$ is the energy of the exchange correlation in the $`F`$ layer. In a general case, the problem needs self-consistent numerical solution. Here, to consider the new physics we are interested in, we will not discuss the quantitative calculations, but will use analytical ones obtained earlier in Refs. for two limits: (a) $`\gamma _M<<1`$, small strength of the proximity effect - low suppression of the order parameter in the S layer near the S/F boundary, and (b) $`\gamma _M>>1`$, strong suppression of the order parameter in the S layer near the S/F boundary. Note that the results obtained are applicable to any value of the S/F boundary transparency, as we made no specific assumption about $`\gamma _B`$ in the derivation below.
Weak proximity effect. If $`\gamma _M<<1`$, one can find an explicit expression for $`\mathrm{\Phi }_S(\omega ,x)`$ in the form
$$\mathrm{\Phi }_S(\omega ,x)=\mathrm{\Delta }_0\{1\gamma _M\beta \stackrel{~}{\omega }\frac{\mathrm{exp}(\beta x/\xi _S)}{\gamma _M\beta \stackrel{~}{\omega }+\omega A(\omega )}\}$$
(9)
where $`\beta =[(\omega ^2+\mathrm{\Delta }_0^2)/\pi T_C]^{1/2}`$ and $`A(\omega )=\left[1+\gamma _B\stackrel{~}{\omega }\left(\gamma _B\stackrel{~}{\omega }+2\omega /\beta ^2\right)/(\pi T_C)^2\right]^{1/2}`$. As one can expect, the magnetic correlation spreads into the S film over a distance of about $`\xi _S`$ and it can much exceed the distance of the superconducting correlation spreading into the F film $`\xi _F`$. If H$`{}_{e}{}^{}0`$ (i.e., $`\stackrel{~}{\omega }\omega `$) the result (9) restores with that for the SN bilayer in the limit in question (see, e.g., Ref. ). For the function $`\mathrm{\Phi }_F(\omega ,0)`$ we obtain
$$\mathrm{\Phi }_F(\omega ,0)=\mathrm{\Delta }_0\stackrel{~}{\omega }_S/(\gamma _B\stackrel{~}{\omega }\beta ^2+\omega )$$
Strong proximity effect. When $`\gamma _M>>1`$, the behavior of $`\mathrm{\Phi }_S(\omega ,x)`$ near the S/F boundary, $`0<x<<\xi _S`$, is given by
$$\mathrm{\Phi }_S(\omega ,0)=B(T)\{(\pi T_C+\gamma _B\stackrel{~}{\omega })/\gamma _M\stackrel{~}{\omega }\}$$
(10)
Here $`B(T)=2T_C[1(T/T_C)^2][7\zeta (3)]^{1/2}`$ (see Ref. ) and $`\zeta (3)1.2`$ is the Riemann $`\zeta `$ function. The function $`\mathrm{\Phi }_F(\omega ,0)`$ in this approximation is read
$$\mathrm{\Phi }_F(\omega ,0)=B(T)\pi T_C/\gamma _M\omega $$
It is seen that the proximity-induced superconductivity in the F layer is independent of the boundary transparency, but decreases with increase of $`\gamma _M`$. To obtain the results for larger distance, $`x\xi _S`$, the equations should be solved numerically by a self-consistent procedure. We will not discuss these results here.
## III Magnetic proximity effect manifestation
An important feature of the results obtained for the SF structure is that the modified Usadel function for the S layer $`\mathrm{\Phi }_S(\omega ,x)`$, Eqs. (9) and (10), directly depends on the exchange field of the F metal. That is the reason to speak about the exchange correlation that has been induced into the S layer due to superconductivity. In this section we discuss a few examples of such โmagnetic proximity effectโ manifestation.
### III.1 Phase variation at SF interface
Comparing the results for an SF structure with those for an SN bilayer, one can find a fundamental aspect, that leads to new physical consequences; namely, the $`\mathrm{\Phi }_S(\omega ,x)`$ is a complex function near the S/F interface. As a result, the additional โsuperconducting phase rotationโ (a phase jump on the S/F interface for our approximation of a thin ferromagnetic layer $`d_F<<\xi _F`$) occurs at the S/F interface. To illustrate this, let us take, for simplicity, a structure with favorable for magnetic effects interface parameters: $`\gamma _M>>1`$ and $`\gamma _B=0`$. Then, as follows from Eq. (10), the modified Usadel function at S/F interface $`\mathrm{\Phi }_S(\omega ,0)`$ can be written in the form
$$\mathrm{\Phi }_S(\omega ,0)=B(T)(\pi T_C/\gamma _M)\frac{\mathrm{exp}(i\theta )}{(\omega ^2+H_e^2)^{1/2}}\text{,}$$
(11)
with $`\theta =\mathrm{arctan}(H_e/\omega )`$. Taking into account that a typical value of $`\omega \mathrm{~}\pi T_C`$, one can see that in the limit $`H_e>>\pi T_C`$ the correlation function acquires an additional $`\pm \pi /2`$ phase shift in comparison with the similar function for the SN bilayer. For an SF multilayred system with strong enough ferromagnetism of the layers the phase shift can be summarized or subtracted, depending on mutual orientation of F layers magnetizations, leading to new effects in superonductivity of SF hybride structures. Namely, one can show, that proximity induced magnetizm of the S layers makes preferrable the $`\pi `$-phase superconductivity of the system for parallel directions of the exchange fields; for antiparallel magnetizations orientation and low temperature, the critical current can be even enhanced \[12-15\].
### III.2 Suppression of the superconducting order parameter by an exchange field
Another feature of the S/F boundary is that the gap $`\mathrm{\Delta }_S(x)`$ is suppressed near the interface more strongly than in the SN case. This is not surprising, since one would expect that induced ferromagnetism suppresses the superconducting order parameter at some distance into the S layer in excess of that for nonmagnetic normal layer. Suppression increases with the increase of the exchange energy $`H_e`$ and of electrical quality of the interface; far from the interface, $`x>>\xi _S`$, the bulk superconductivity is restored.
Using $`\mathrm{\Phi }_S(\omega ,x)`$ (9) and the self-consistency condition (2) one can find the spatial variation of the order parameter in the $`S`$ layer $`\mathrm{\Delta }_S(x)`$ for different values of $`\gamma _B`$ and $`\gamma _M<<1`$. The exchange interaction influence on the spatial variation of the order parameter in the S layer is shown in Fig. 1. Namely, the dependence of difference of the order parameters for the case when magnetic interaction is turned off (i.e., an SN bilayer) and with ferromagnetic correlation (a SF bilayer) as function of distance from the interface is shown; the boundary parameters, $`\gamma _M`$ and $`\gamma _B`$, are fixed. It is seen, that influence of magnetism decreases with increasing the distance from the S/F boundary. The scale at which superconductivity reaches the value for a SN bilayer is $`\xi _S`$ from the interface. The curves in Fig. 1 illustrate the spatial dependencies of the induced exchange correlation in the superconductor for the case of vanishing interface resistance $`\gamma _B=0`$. With an increase of the SF boundary resistance the electrical coupling of the S and F metals decreases and in the limit $`\gamma _B\mathrm{}`$ the metals become decoupled.
### III.3 Exchange field spin-splitting of DOS and intra-gap states
Spin splitting of DOS and intergap states in the S layer are other manifestations of magnetic correlation leakage into a superconductor. Note that the magnetic layer does not influence the DOS of the normal metal. In this case the decay length is extremely small $`p_F^11\stackrel{ฬ}{A}`$ and the effect can be neglected.
The Green functions for the S layer $`G_S(\omega ,x)`$ and $`G_S(\omega ,x)`$ for both spin subbands can be obtained using solutions for the functions $`\mathrm{\Phi }_S(\omega ,x)`$ with $`\stackrel{~}{\omega }=\omega +iH_e`$ and $`\stackrel{~}{\omega }=\omega iH_e`$, respectively. Performing the analytical continuation to the complex plane by the substitution $`\omega i\epsilon `$ we calculate the spatial dependence of quasiparticle DOS for spin โupโ and โdownโ subband: $`N_S(\epsilon ,x)=ReG_S(\omega ,x)`$ and $`N_S(\epsilon ,x)=ReG_S(\omega ,x)`$, respectively The total density of states for quasiparticles, by definition, is given by $`N_S(\epsilon ,x)=N_S(\epsilon ,x)+N_S(\epsilon ,x)`$. Using Eqs. (9) and (10), one can obtain the explicit expressions for the total DOS, as well as for the specified spin subband. The resulting expressions, which are cumbersome to be presented here, imply that for $`H_e0`$, $`\gamma _M0`$, and $`\gamma _B0`$ the density of quasiparticle states is spin-splitted: $`N_S(\epsilon ,x)N_S(\epsilon ,x)`$. This is because of the initial exchange field splitting of the Fermi surface in the F metal, which is manifested in the characteristics of the united system โ the SF bilayer. The symmetry of the density of states with respect to the energy variable is also lost: $`N_S(\epsilon >0,x)N_S(\epsilon <0,x)`$. However, as one can expect from the fermionic symmetry, the spin-up particles and spin-down holes have the same DOS, and likewise for spin-down particles and spin-up holes; as a result the total density $`N_S(\epsilon ,x)`$ is symmetric: $`N_S(\epsilon >0,x)=N_S(\epsilon <0,x)`$.
In Fig. 2 representative $`N_S(\epsilon ,x)`$ dependences at different distances from the S/F interface are presented for $`H_e=5\pi T_C`$ and $`\gamma _M=0.1`$, and vanishing boundary resistance ($`\gamma _B=0`$). In Fig. 3 the same dependence is presented for $`x/\xi _S=1`$ and different values of the exchange energy. We find that all features mentioned above are saved on a distance of a scale $`\xi _S`$ from the SF boundary. The spin-splitting decreases with an increase of the distance from the boundary and vanishes in the bulk of the S layer (see curve 4 in Fig. 2).
Other important features, shown in Figs. 2 and 3, are the local states that appear inside the energy gap at the distances up of a few $`\xi _S`$ from the S/F boundary. These intergap states are absent far from the S/F interface, and also if $`H_e=0`$. For small values of $`\gamma _M`$ and $`\gamma _B=0`$, as follows from the expression (9), $`N_S(\epsilon ,x)`$ has singularity for
$$\epsilon =\pm \mathrm{\Delta }_0\{1\frac{\gamma _M\beta _\epsilon \stackrel{~}{\epsilon }}{\epsilon +\gamma _M\beta _\epsilon \stackrel{~}{\epsilon }}\mathrm{exp}(\beta _\epsilon x/\xi _S)\}$$
(12)
where $`\beta _\epsilon ^2=(\mathrm{\Delta }_0^2\epsilon ^2)^{1/2}/\pi T_C`$ and$`\stackrel{~}{\epsilon }=\epsilon H_e`$. We found the singularity inside the superconducting gap, $`\mathrm{\Delta }_0<\epsilon <\mathrm{\Delta }_0`$, by numerical calculations . The local states are definitely not due to the spatial variation of the pair potential, but due to Cooper pairs breaking in the superconductor by the exchange-induced magnetic correlation. The region of their existence increases with the increasing of $`H_e`$, or increasing pair breaking effects. In the absence of spin-flip (e.g., spin-orbit) scattering, the subgap bands accommodate quasiparticles with a definite (โupโ or โdownโ) spin direction. These bands bear superficial resemblance to both the bands observed at interface of superconductor and perfectly insulating ferromagnet and bulk superconductor containing finite concentrations of magnetic impurities .
### III.4 Induced magnetization of the S layer
As we saw above, the influence of the ferromagnet on the superconductor is reflected in a nonzero value of the difference in the DOS for spin-up and spin-down unpaired electrons, $`N_S(\epsilon ,x)`$ and $`N_S(\epsilon ,x)`$. This DOS splitting causes an effective magnetization $`M_S(x)`$ of the S layer, that can be found using the relation:
$$M_S(x)/M_O=_0^{\mathrm{}}๐\epsilon \{N_S(\epsilon ,x)N_S(\epsilon ,x)\}f(\epsilon )$$
(13)
where $`M_O=gS_e\mu _B(=\mu _B)`$ is a quasiparticle magnetic moment, $`S_e=1/2,g=2`$ and $`f(\epsilon )=1/[\mathrm{exp}(\epsilon /T)+1]`$ is the Fermi distribution function. Figure 4 illustrates the mechanism of proximity induced magnetization of the S layer. For $`T<T_C`$ we took $`f(\epsilon )=1`$, i.e., all states below Fermi level are filled (dashed regions in Fig. 4), while all states above Fermi energy are empty. One can directly see from the figure that the S layer acquires a nonzero magnetic moment. This suggestion is confirmed by numerical calculations of $`M_S(x)`$ Eq.(13) shown in Figs. 5, 6. Figure 5 shows the magnetization of the superconductor versus distance from the S/F interface for fixed boundary parameters. The same magnetic characteristics but for a SF sandwich with fixed exchange energy and boundary transparency, and different proximity effect strength are presented in Fig. 6.
### III.5 Experiment
There are only a few experimental reports devoted to the questions discussed here. Interplay between magnetism and superconductivity in Nb/Co multilayers was recently investigated by Ogrin et al. . The upper critical fields of the samples were measured for the field applied parallel to the plane, H$`_{C_2||}`$ and perpendicular to the plane H$`_{C_2}`$ of the films. Effective thickness of the Co layer, d<sub>eff</sub> , they define through the well known relation:
$$d_{eff}=\left(\frac{\mathrm{\Phi }_0}{2\pi H_{C_2}}\right)^{1/2}\frac{H_{C_2}}{H_{C_2||}}$$
where $`\mathrm{\Phi }_0`$ stands for flux quantum. Experiments revealed that the effective thickness of the magnetic layer in $`Nb/Co`$ structures is usually much larger than its physical thickness $`d_{Co}`$. For example, taking the data on sample with d<sub>Co</sub> = 1.8.nm, the authors obtained a value $`d_{eff}=`$ 12 nm, so that $`d_{eff}>>d_{Co}`$. The โincreaseโ of the thicknesses was so great that in all samples, except for those with extremely thin magnetic layers, the crossover to a 3D state superconductivity is never in fact observed experimentally. This is to be contrasted with the case of nonmagnetic spacer layers, where these two length scales are comparable. Taking into account our results, we explain the rise of the effective magnetic layer thickness in the $`Nb/Co`$ multilayer as an impact of proximity effect. Namely, the induced magnetic correlation into the $`S`$ layer depletes Cooper pairs density at the $`SF`$ boundary, which results in an excess thickness of the magnetic layer.
The modification of the DOS in mesoscopic superconducting strips of $`Al`$ under the influence of magnetic proximity effect of a classical ferromagnet $`Ni`$ has also been studied both theoretically and experimentally in . However, since the tunnel spectroscopy experiments were carried out with a nonmagnetic probe, the authors could not measure spin-denendent local DOS in the superconducting side.
The interest in the magnetic proximity effect has been increased with the development of experimental techniques like neutron reflectometry and muon spin rotation, which allow to determine accurately the spatial distribution of magnetic moments. For example, very recently multilayered system YBa<sub>2</sub>Cu<sub>3</sub>O<sub>7</sub>/La<sub>2/3</sub>Ca<sub>1/3</sub>MnO<sub>3</sub> have been studied by neutron reflectometry in . Evidence for a characteristic difference between the structural and magnetic depth profiles is obtained from the occurrence of a structurally forbidden Bragg peak is a ferromagnetic state. The authors discussed findings in two possible scenarios: a sizable magnetic moment within the Slayer antiparallel to one in the F layer (inverse proximity effect), or a โdeadโ region in the F layer with zero net magnetic moment.
## IV Conclusion
In recent years, advances in materials growth and fabrication techniques have made it possible to create heterostructures with high quality interfaces. Taking into account that ferromagnet-superconductor hybrid systems have great scientific importance, and are promising for application in spin-electronics, it is not surprising that interest to these hybrid materials has been renewed. As far as the thickness of superconducting and magnetic metals in such structures may be a few atomic periods, understanding of how the proximity effects modify electronic properties of S/F interfaces is growing in importance.
We have studied in the magnetic correlations acquired by a superconductor at S/F interface due to a proximity effect. We have found that an equilibrium exchange of electrons between the F and S metals results not only in proximity induced superconductivity of the F metal, as was found earlier, but in proximity induced magnetism of the S metal, too. The magnetic correlations spread over a large distance which is of the order of the superconducting coherence length $`\xi _S`$ and can exceed both the ferromagnetic and the superconducting films thicknesses. That is why the existence of these magnetic properties of the S metal is quite important for SF nanoscale structures and should be taken into account while comparing theoretical results with experimental data. Summarizing the results, we should stress that for SF nanoscopic hybrid structures both phenomena, โ the superconducting and the magnetic proximity effects, โ take place simultaneously, and both should be paid attention to.
*Acknowledgements*. We wish to dedicate this paper to V.G. Barโyakhtar, our Master who played a significant, exceptional role in our post-student life, on the occasion of his 75th birthday, and to wish him continuing health and vigour in pursuing his scientific interest. The authors would like to thank V. V. Ryazanov, A.I. Buzdin, L. Tagirov, and M. A. Belogolovskii for valuable discussions of some questions of proximity effect phenomena. We also acknowledge E. A. Koshina for performing numerical calculations.
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Figure Captures
Fig. 1. The difference of the superconducting order parameter in the S layer versus distance from the interface for SN and SF structures with the same boundary parameters ($`\gamma _M=0.1`$, $`\gamma _B=0`$), and different ferromagnetic field energy $`H_e/\pi T_C=`$ 8, 9, 10, 12 and 15.
Fig. 2. Normalized density of state for spin โupโ quasiparticles in the S layer of the SF sandwich for $`\gamma _M=0.1`$ , $`\gamma _B=0`$ and $`H_e=5\pi T_C`$, and various distances from the S/F interface: $`x/\xi _S=`$ 0, 1, 5, and 30 (curves 1, 2, 3, and 4, respectively).
Fig. 3. Same as in Fig. 4 for $`\gamma _M=0.1`$, $`\gamma _B=0.1`$ and $`x=\xi _S`$, and various ferromagnetic field energies: $`H_e/\pi T_C=`$ 1, 2, and 5 (curves 1, 2, and 3, respectively).
Fig. 4. Quasiparticle density of states in the S layer near the S/F interface; $`\gamma _M=0.1`$, $`\gamma _B=0.0`$, $`x=\xi _S`$, and $`H_e=5\pi T_C`$. All states above Fermi energy are empty; all states below Firmi level are filled (dashed regions in figure).
Fig. 5. Leakage of magnetization into the S material versus distance from the interface for SF sandwich for $`\gamma _M=0.1`$ , $`\gamma _B=0`$ , and different exchange energies $`H_e/\pi T_C=`$7, 5, and 3 (curves 1, 2, and 3, respectively).
Fig. 6. Same as in Fig. 6 for $`\gamma _B=0`$ , $`H_e=3.5\pi T_C`$ and different proximity effect strength $`\gamma _M`$ = 0.1, 0.15, 0.2 (curves 1, 2, and 3, respectively).
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# Dynamic regimes of fluids simulated by multiparticle-collision dynamics
## I Introduction
The dynamics of complex fluids such as colloidal suspensions, dilute or semi-dilute polymer solutions, biological macromolecules, membranes, and aqueous surfactant solutions, is often governed by the hydrodynamic behavior of the solvent. Due to a large separation of length and time scales between the atomic scale of the solvent molecules and the mesoscopic scale of the solute, direct simulation approaches with explicit atomistic solvent are prohibitively costly in computer time. Therefore, several mesoscale simulation techniques have been developed in recent years in order to bridge the length- and time-scale gap. In particular, lattice-gas automata (LGA) Frisch et al. (1986); Koelman (1990), lattice Boltzmann methods (LBM) McNamara and Zanetti (1988); Lallemand and Luo (2000); Succi (2001), smoothed-particle hydrodynamics (SPH) Gingold and Monaghan (1977); Kum et al. (1995), dissipative particle dynamics (DPD) Hoogerbrugge and Koelman (1992); Espaรฑol and Warren (1995); Groot and Warren (1997), direct simulation Monte Carlo (DSMC) Metropolis et al. (1953); Bird (1976), fluid particle dynamics Tanaka and Araki (2000), and others, have been investigated. The basic idea of all these approaches is very similar: To obtain hydrodynamic behavior on length scales much larger than the atomic scale, the detailed interactions and dynamics of the solvent molecules are not important; instead mass and momentum conservation are the essential ingredients to obtain the correct hydrodynamic behavior. Therefore, the dynamics on the microscopic scale can be strongly simplified, as long as the conservation laws are strictly satisfied. The different methods listed above differ in the way the solvent dynamics is implemented.
Two main classes of mesoscopic simulation techniques can be distinguished, which are lattice and off-lattice methods. Lattice gas and lattice Boltzmann methods fall into the first class, while direct simulation Monte Carlo, dissipative particle dynamics, and fluid particle dynamics fall into the second class. Off-lattice approaches have the advantage that Galilean invariance is typically satisfied. Moreover, the interaction of the off-lattice solvent with solutes such as colloids, polymers and membranes can be taken into account more naturally.
The mesoscale simulation technique, which we are investigating in this paper, was introduced by Malevanets and Kapral Malevanets and Kapral (1999) a few years ago. It is a variant of the DSMC method, in which binary collisions are replaced by multi-particle collisions in a prescribed collision volume. This method has been called multi-particle-collision dynamics (MPCD) or stochastic rotation dynamics (SRD). It employs a discrete-time dynamics with continuous velocities and local multi-particle collisions. Mass and momentum are conserved quantities and it has been demonstrated that the hydrodynamic equations are satisfied Malevanets and Kapral (1999, 2000).
Certain transport coefficients, in particular the viscosity, of this solvent model have been studied intensively. Analytical expressions have been derived from kinetic theory by generalizing point-like collisions to finite collision volumes Ihle and Kroll (2003a, b); Tรผzel et al. (2003); Kikuchi et al. (2003). The theoretical expressions describe numerical results very well.
In this article, we study the transport coefficients as a function of the parameters of the MPCD fluid, in particular the mean free path in units of the size of the collision volume. We find two distinct regimes, in which the dynamics is either gas-like or fluid-like. This behavior can be characterized by the Schmidt number, which measures the ratio between viscous and diffusive transport. We find that MPCD allows us to tune the fluid behavior such that large Schmidt numbers are obtained and momentum transport dominates over mass transport. Analytical expressions Ihle and Kroll (2003b); Tรผzel et al. (2003); Kikuchi et al. (2003) for the tracer diffusion coefficient, which have been derived on the basis of a molecular-chaos assumption, are found to describe the simulation data very well for large mean free paths, but fail in the fluid regime. The reason is a build-up of correlations among the fluid particles by hydrodynamic interactions, which leads to enhanced diffusion coefficients. We will show that the latter leads to non-exponentially decaying velocity-autocorrelation functions at small mean free paths. Independent of the mean free path, we find that the algorithm reproduces the algebraic long-time decay typical in fluids.
In a further step, we investigate the diffusion of a heavy tracer particle in a MPCD solvent. It is very important to understand the contribution of the solvent dynamics on the solute diffusion. Two limiting situations are found: either Brownian or hydrodynamic behavior, depending on the collision time and the rotation angle. We explore the range of parameters where these different dynamical behaviors appear, and show how they emerge from the mesoscopic dynamics.
Finally, we study self-diffusion in colloidal dispersions with excluded-volume interactions as a function of the volume fraction. To this end, the MPCD method is combined with molecular dynamic simulations. We find that such a hybrid model displays the proper dynamics for the same parameter regime where the hydrodynamic behavior is found for the fluid. Our results in the collective regime are in good agreement with previous theoretical predictions based on Stokes hydrodynamics and the Smoluchowski equation Dhont (1996).
## II The Model
The fluid is modeled by $`N`$ point particles, which are determined by their positions $`๐ซ_i`$ and velocities $`๐ฏ_i`$, with $`i=1,\mathrm{},N`$. Positions and velocities are continuous variables, which evolve in discrete increments of time. The mass $`m`$ associated with the particles is taken to be the same, but more generally, different masses can be assigned. The algorithm consists of two steps, streaming and collision. In the streaming step the particles move ballistically according to their velocities during a time increment $`h`$, to which we will refer as collision time. Thereby, the evolution rule is
$$๐ซ_i(t+h)=๐ซ_i(t)+h๐ฏ_i(t).$$
(1)
In the collision step, the particles are sorted into collision boxes, and interact with all other particles in the same collision box. The collision boxes are typically the unit cells of a $`d`$-dimensional cubic lattice with lattice constant $`a`$, although other geometries would be possible. The collision is then defined as a rotation of the velocities of all particles in a box in a co-moving frame with its center of mass. Thus, the velocity of the $`i`$-th particle after the collision is
$$๐ฏ_i(t+h)=๐ฏ_{cm,i}(t)+(\alpha )\left[๐ฏ_i(t)๐ฏ_{cm,i}(t)\right],$$
(2)
where $`(\alpha )`$ is a stochastic rotation matrix, and $`๐ฏ_{cm,i}(t)=_j^{(i,t)}(m๐ฏ_j)/_jm`$ is the velocity of the center of mass of all particles $`j`$, which are located in the collision box of particle $`i`$ at time $`t`$. The conservation of local momentum and kinetic energy is guaranteed by construction. In two dimensions, the rotation of the relative velocity is simply given by an angle $`\pm \alpha `$. Here $`\alpha `$ is a parameter of the model; the sign is chosen randomly for each cell. In three dimensions, various schemes for the random collisions are possible Malevanets and Kapral (1999); Allahyarov and Gompper (2002); Tรผzel et al. (2003). The one employed in this paper consist in choosing a random direction in space for each box around which the relative velocities are rotated by an angle $`\alpha `$. A detailed explanation of the implementation is given in Ref. Allahyarov and Gompper (2002).
In order to ensure Galilean invariance for the full range of parameters, a random shift of the collision grid has to be performed in the execution of the collision step Ihle and Kroll (2001, 2003a). As a consequence of such a shift, the collision environment of each particle is independent of the average local velocity, and no special reference frame exists. Random shifts also facilitate the transfer of momentum between neighboring particles.
In the simulations, $`N`$ particles are initially placed at random in a cubic system of linear extension $`L`$. The average number of particles in a collision box is $`\rho =N(a/L)^d`$, the scaled number density. Starting from an arbitrary distribution of velocities, only a few steps are required to reach the Maxwell Boltzmann velocity distribution. The equilibrium temperature $`T`$ is then given by the average kinetic energy $`m๐ฏ_i^2=3k_BT`$, where $`k_B`$ is the Boltzmann constant. In the simulations, we scale length and time according to $`\widehat{x}=x/a`$ and $`\widehat{t}=t\sqrt{k_BT/ma^2}`$, which corresponds to the choice $`m=1`$, $`a=1`$, and $`k_BT=1`$ of reference units. The scaled mean free path is then given by $`\lambda =\widehat{h}`$. Basic parameters and the definitions of dimensionless quantities are collected in Table 1.
## III Dynamical Properties
The kinematic viscosity $`\nu =\eta /\varrho `$ has been calculated theoretically Malevanets and Kapral (1999, 2000); Ihle and Kroll (2001, 2003a, 2003b); Tรผzel et al. (2003); Kikuchi et al. (2003); Ihle et al. (2004) by means of kinetic theory and its validity has been checked with simulations. The total kinematic viscosity, $`\nu =\nu _{kin}+\nu _{coll}`$, is the sum of two contributions, the kinetic viscosity $`\nu _{kin}`$ and the collisional viscosity $`\nu _{coll}`$, which have been calculated in two and three dimensions. In three dimensions, the expressions Tรผzel et al. (2003); Kikuchi et al. (2003)
$`{\displaystyle \frac{\nu _{coll}}{\sqrt{k_BTa^2/m}}}`$ $`=`$ $`{\displaystyle \frac{1}{\lambda }}{\displaystyle \frac{(1\mathrm{cos}\alpha )}{18}}\left(1{\displaystyle \frac{1}{\rho }}\right)`$ (3)
$`{\displaystyle \frac{\nu _{kin}}{\sqrt{k_BTa^2/m}}}`$ $`=`$ $`\lambda \left[{\displaystyle \frac{1}{(42\mathrm{cos}\alpha 2\mathrm{cos}2\alpha )}}{\displaystyle \frac{5\rho }{\rho 1}}{\displaystyle \frac{1}{2}}\right]`$
have been derived.
The total kinematic viscosity has been determined numerically by the procedure explained in Ref. Lamura et al. (2001). Briefly, a three-dimensional system is considered with periodic boundary conditions in two dimensions and planar walls in the third dimension. Stick boundary conditions at the walls are implemented by considering bounce-back collisions with the walls. A gravitational field is applied in one direction parallel to the walls. After a relaxation time, the system reaches a stationary state with a parabolic velocity profile between the walls and in the direction of the force. This is Poiseuille flow. It is known Tritton (1988) that the measured maximum velocity of the parabola is inversely proportional to the viscosity of the fluid. The viscosity data obtained in this way are presented in Fig. 1 together with the theoretical predictions of Eq. (3). The obtained agreement is quite remarkable, in contrast to the case of other mesoscopic simulation techniques such as dissipative particle dynamics Pagonabarraga et al. (1998). Density fluctuations can also be included in the theory Kikuchi et al. (2003), which noticeably improves the agreement with the simulations results for small number densities; for $`\rho =5`$ and $`\rho =10`$, these contributions are negligible.
Alternative methods to determine the viscosity from simulations have been employed in Refs. Kikuchi et al. (2003) and Ihle and Kroll (2003b), where a system under shear flow and vorticity correlations have been used, respectively.
The ratio between the kinetic and the collisional contributions to the kinematic viscosity varies considerably with the model parameters, as can be seen easily from the theoretical expressions (3). In Fig. 1 the total kinematic viscosity and its two contributions are plotted as a function of the rotation angle and the collision time step. The collisional contribution is dominant for large collision angles and small collision times, while the kinetic viscosity dominates in the opposite case of small collision angles and large collision times.
Kinetic transport is due to the movement of the particles themselves, i.e., when a particle moves it carries a certain amount of the relevant quantities as momentum and energy, while collisional transport is due to transfer of energy and momentum from one particle to another during collisions. In MPCD, kinetic transport is therefore dominant when the mean free path is larger than the size of the collision box and for small values of the rotation angle. If the rotation angle is small, there is little exchange of momentum between particles due to collisions. The situation where the kinetic transport dominates is characteristic for gases. In fluids the usual situation is the opposite, the transport of momentum is mainly due to collisions.
A convenient measure of the importance of hydrodynamics is the Schmidt number $`Sc=\nu /D`$, where $`\nu `$ is the kinematic viscosity and $`D`$ the diffusion coefficient. Thus, $`Sc`$ is the ratio between momentum transport and mass transport. It is known that this number for gases is smaller than but on the order of unity, while in fluids like water it is on the order of $`10^2`$ to $`10^3`$. A prediction for the Schmidt number of a MPCD fluid can be obtained from the theoretical expressions (3) for the kinematic viscosity, and the diffusion coefficient, see Eq. (17) below. In Fig. 2, we plot the theoretical prediction for $`Sc`$ as a function of the collision time for different values of the rotation angle. This shows that $`Sc`$ becomes considerably larger than unity for the same range of parameters where the collisional viscosity is considerably larger than the kinetic viscosity (Fig. 1). We will show that the dynamical behavior in the two limits is fundamentally different. We will call the parameter region of large rotation angles and small collision times the collective regime, and the opposite region the particle regime. This classification has similar consequences as the one introduced in dissipative particle dynamics (DPD) Ripoll et al. (2001), although we do not investigate wave-length dependent properties here.
## IV Simple Fluid Correlations
Correlations between particles are responsible for hydrodynamic interactions. Therefore, we are interested in characterizing the velocity correlations in a MPCD fluid.
### IV.1 Velocity Autocorrelation Functions
An analytical expression for the velocity autocorrelation function (VACF) has been derived in Refs. Ihle and Kroll (2003b); Tรผzel et al. (2003). The collision step in Eq. (2) can be rewritten as
$$\begin{array}{cc}\hfill ๐ฏ_i(nh)=& ๐ฏ_i((n1)h)+\hfill \\ & \left((\alpha )๐\right)\left[๐ฏ_i((n1)h)๐ฏ_{cm,i}((n1)h)\right],\hfill \end{array}$$
(4)
where $`๐`$ is the unit matrix and $`t=nh`$ is the discretized time, with $`n`$ the number of collision steps. By multiplying this expression with the velocity at time zero and taking thermal averages, we obtain
$$\begin{array}{c}\hfill ๐ฏ_i(nh)๐ฏ_i(0)=(1\gamma _\alpha )๐ฏ_i((n1)h)๐ฏ_i(0)\\ \hfill +\gamma _\alpha ๐ฏ_{cm,i}((n1)h)๐ฏ_i(0).\end{array}$$
(5)
Here, the rotational average over an arbitrary vector $`๐`$ in three dimensions is obtained from geometrical arguments to be
$$\left((\alpha )๐\right)๐=\frac{2}{3}(1\mathrm{cos}\alpha )๐\gamma _\alpha ๐.$$
(6)
This particular value of $`\gamma _\alpha `$ arises from the implementation of the rotation chosen in this paper. The remaining problem is to calculate the last term in Eq. (5). First, we neglect density fluctuations in the average of the center of mass velocity, which yields $`๐ฏ_{cm,i}(nh)_j^{(i,n)}๐ฏ_j/\rho `$. Furthermore, a molecular-chaos assumption implies that
$$๐ฏ_{cm,i}((n1)h)๐ฏ_i(0)\frac{1}{\rho }๐ฏ_i((n1)h)๐ฏ_i(0).$$
(7)
This approximation means that of all the particles in the collision box of particle $`i`$ after $`(n1)`$ collisions, only particle $`i`$ itself makes a non-zero contribution to the correlation function. This is the same as assuming that none of the other particles has any information about the state of particle $`i`$ at any time. The correlation at a certain time step can then be expressed in terms of the previous time step as
$$๐ฏ_i(nh)๐ฏ_i(0)\left(1\frac{\rho 1}{\rho }\gamma _\alpha \right)๐ฏ_i((n1)h)๐ฏ_i(0).$$
(8)
This implies that in this approximation, the VACF shows an exponential decay,
$$C_v(nh)\frac{๐ฏ_i(nh)๐ฏ_i(0)}{v_i^2(0)}(1\gamma )^n,$$
(9)
where the normalization factor follows from the equipartition theorem, $`๐ฏ_i^2(0)=3k_BT/m`$. The decorrelation factor $`\gamma `$ is defined as
$$\gamma =\frac{2}{3}(1\mathrm{cos}\alpha )\left(1\frac{1}{\rho }\right)\gamma _\alpha \gamma _\rho .$$
(10)
ยฟFrom Eq. (9), a characteristic time $`\tau _0=h/\mathrm{ln}(1\gamma )`$ can be extracted. Up to this time, the VACF follows the exponential decay for every set of parameters. However, the collective phenomena responsible for the hydrodynamic behavior appear at much later times.
In Fig. 3, simulation results of the VACF are presented for two different mean free paths $`\lambda `$. The theoretical prediction (9) is also displayed for both values of $`\lambda `$. For $`\lambda =1`$ the exponential decay is followed with very good accuracy until the crossover to the long-time tail behavior occurs. For $`\lambda =0.1`$ the purely exponential decay is followed only in the first collision; for long times, a long-time-tail behavior is observed similarly as for $`\lambda =1`$. What is different in this case is that after the first collision the system enters an intermediate regime where the VACF decay is significantly slower than the one described by the molecular-chaos approximation but is not yet the algebraic tail. Note that for the investigated rotation angle of $`\alpha =130`$, the mean free path $`\lambda =1`$ corresponds to the particle regime, while $`\lambda =0.1`$ corresponds to the collective regime.
It is interesting to note that for short times, the VACF decays monotonically only in the case that the correlation parameter $`\gamma `$ is smaller than unity. If $`\gamma 1`$, Eq. (9) predicts that the VACF exhibits damped oscillations. We have checked that this oscillatory behavior is indeed observed in the simulations. However, the viscosity curves show no particular features when this happens (compare Fig. 1, where the VACF for $`\rho =10`$ becomes oscillatory for $`\alpha 132`$).
### IV.2 Long-Time Tails
It is well known Alder and Wainwright (1970); Ernst et al. (1971); Ernst (2005) that the long-time behavior of the VACF in $`d`$-dimensional fluids in thermal equilibrium shows a universal behavior. This corresponds to a power-law tail, for which the explicit form can be calculated from a mode-coupling theory as Ernst et al. (1971),
$$C_v(t)\left(\frac{d1}{d\rho }\right)\frac{1}{[4\pi (D+\nu )t]^{d/2}},$$
(11)
where $`\nu `$ and $`D`$ are the transport coefficients of the fluid.
The results obtained for the long-time behavior of the VACF are consistent with the general prediction for fluids in thermal equilibrium in Eq. (11). The algebraic power $`t^{3/2}`$ is clearly reproduced in our simulations as can be seen in Fig. 4. The value of the amplitude in Eq. (11) is related to the kinematic viscosity $`\nu `$ and the diffusion coefficient $`D`$. Since both values are known for the MPCD fluid and discussed in this paper, quantitative comparison can also be performed. We find that the value for $`\lambda =1`$ is exactly reproduced by our simulations within the accuracy of the results, while the amplitude obtained for $`\lambda =0.1`$ is about $`10\%`$ smaller than the theoretic prediction. Ihle and Kroll Ihle and Kroll (2003a) obtain good agreement in a two-dimensional MPCD fluid with the expected $`t^1`$ behavior over a comparable time window.
The effect of finite system size can be seen in Fig. 4 for times $`\widehat{t}20`$, where the VACF crosses over from the algebraic to a faster, exponential decay. This effect is similar to that observed for the time dependence of the temperature autocorrelation function for a random-solid dissipative-particle-dynamics system Ripoll and Ernst (2004). There, it can be proved that the correlations decay faster after a time, where hydrodynamic modes become relevant which are truncated by the system size.
### IV.3 Importance of Many-Body Correlations
In the previous section, an exponential decay of the VACF has been theoretically predicted. This behavior is a consequence of the approximation in Eq. (7) which neglects any correlation among the particles in the same collision box at all times. In order to improve Eq. (7), we have to go beyond the molecular-chaos approximation. This is a formidable task. We start the procedure by calculating the center-of-mass correlation average for the first collision and, consecutively, the second and so on. For $`n=1`$ the approximation in Eq. (7) is exact, $`๐ฏ_{cm,i}(0)๐ฏ_i(0)=๐ฏ_i^2(0)/\rho `$. This is the reason why for the first time step, $`C_v(h)`$ agrees perfectly in all simulations. For $`n=2`$ it reads
$`๐ฏ_{cm,i}(h)๐ฏ_i(0)=`$ (12)
$`=`$ $`{\displaystyle \frac{1}{\rho }}{\displaystyle \underset{j}{\overset{(i,1)}{}}}\left\{๐ฏ_j(0)+\left((\alpha )๐\right)\left[๐ฏ_j(0)๐ฏ_{cm,j}(0)\right]\right\}๐ฏ_i(0)`$
$`=`$ $`{\displaystyle \frac{1}{\rho }}(1\gamma _\alpha )v_i^2(0)+{\displaystyle \frac{\gamma _\alpha }{\rho ^2}}{\displaystyle \underset{j}{\overset{(i,1)}{}}}{\displaystyle \underset{k}{\overset{(j,0)}{}}}๐ฏ_k(0)๐ฏ_i(0)`$
$``$ $`\left({\displaystyle \frac{1\gamma _\alpha }{\rho }}+{\displaystyle \frac{\gamma _\alpha }{\rho ^2}}\zeta _1\right)v_i^2(0)`$
where $`\zeta _1`$ denotes the number of particles that are neighbors of particle $`i`$ at both times $`t=h`$ and $`t=0`$. We use the term โneighborsโ for particles within the same collision box. The approximation in Eq. (7) is recovered for $`\zeta _1=1`$. This is the case when only the actual particle is considered to be in both collision boxes. As we have seen above, this is not a good approximation in the collective regime.
We denote the average number of remaining neighbors that one particle is revisiting after $`n`$ collisions as $`\zeta _n`$. This number could in principle be calculated analytically by probabilistic arguments, but in order to get a flavor of the improvement that such numbers produce in the theory, we determine $`\zeta _n`$ numerically in our simulations. As expected, these numbers strongly depend on the system parameters. A detailed study has not been performed, but we have observed that the number of remaining neighbors seems to be a universal function of the root-mean-square displacement of the tagged particle.
The measured numbers $`\zeta _n`$ are presented in Fig. 5 as a function of the root-mean-square displacement $`\left(๐ซ(t)๐ซ(0)\right)^2^{1/2}=\sqrt{6Dt}`$, where $`D`$ is the diffusion coefficient and $`t=nh`$. The diffusion coefficient is the one obtained from the analytical expression which will be deduced in the next section (see Eq. (17)). The data for different mean free paths seem to fall onto a single master curve with reasonable accuracy. When the numerical values of $`D`$ (discussed in the next section) are used instead of the theoretical result, the data collapse becomes even more accurate. For the large mean free path $`\lambda =1`$, the first collision takes place when $`\sqrt{6Dt}/a2`$, which implies that $`\zeta _11`$ is a good approximation. Note that in the representation chosen in Fig. 5, $`\zeta _n1`$ corresponds to the abscissa. The same displacement for a small mean free path $`\lambda =0.1`$ takes place when the particle has been involved in $`80`$ collisions on average. The first collision for $`\lambda =0.1`$ takes place when the average displacement is much smaller and many of the particles are still in the same collision box, which makes $`\zeta _11`$ a bad approximation. Indeed, we can infer from Fig. 5 that $`\zeta _12.1`$ for $`\lambda =0.1`$ and $`\rho =5`$, and $`\zeta _13.5`$ for $`\lambda =0.1`$ and $`\rho =10`$.
Following the same procedure as employed in Eq. (12), the velocity correlation function can be calculated for $`n=3`$,
$$\begin{array}{c}\hfill ๐ฏ_{cm,i}(2h)๐ฏ_i(0)=\frac{v_i^2(0)}{\rho }[(1\gamma )(1\gamma +\frac{\gamma }{\rho }\zeta _2)\\ \hfill +\frac{\gamma }{\rho }(1\gamma )\zeta _1+\frac{\gamma ^2}{\rho }(\zeta _2+\delta \zeta _2)]\end{array}$$
(13)
where $`\delta \zeta _2`$ is determined by
$$\zeta _2+\delta \zeta _2\frac{1}{v_i^2(0)}\underset{j}{\overset{(i,2)}{}}\underset{l}{\overset{(j,1)}{}}\underset{k}{\overset{(l,0)}{}}๐ฏ_k(0)๐ฏ_i(0).$$
(14)
This is the number of neighbors of particle $`i`$ at the two times $`t=2h`$ and $`t=0`$ together with the neighbors of the neighbors, or the result of ring collisions. Let us consider two particles $`i`$ and $`k`$, which are in the same collision box at $`t=2h`$ but not at $`t=0`$. If one, $`k`$, has been neighbor of a third particle $`j`$ at $`t=h`$ and this $`j`$ was neighbor of $`i`$ at $`t=0`$, then this combination also contributes to the correlation function. To obtain a reasonable prediction for this number is obviously not trivial. Furthermore, this relation will become more interconnected and difficult to predict for further time steps. It can be checked that with the approximations $`\zeta _n=1`$ and $`\delta \zeta _n=0`$, Eq. (13) reduces to Eq. (7), and consequently the exponential decay in Eq. (9) is recovered.
Now we come back to the correlation average in Eq. (5) which can be expanded with the help of Eq. (4) โ without any approximation โ
$$\begin{array}{c}\hfill ๐ฏ_i(nh)๐ฏ_i(0)=v_i^2(0)(1\gamma )^n\\ \hfill \gamma \underset{k=1}{\overset{n}{}}(1\gamma )^{nk}๐ฏ_{cm,i}((k1)h)๐ฏ_i(0).\end{array}$$
(15)
The predictions for short times can be improved compared to Eq. (9) by employing the results of Eqs. (12) and (13) on the right-hand side of Eq. (15), but setting $`\zeta _n1`$ and $`\delta \zeta _n0`$ for $`n3`$ as before. The result is shown in Fig. 6. We observe that the prediction for the second collision $`C_v(2h)`$ now agrees perfectly with the simulation data, which confirms our arguments. Nevertheless, the prediction for further steps is still only a small improvement compared to the exponential decay in Eq. (9).
The most relevant conclusion at this point is that in the collective regime the MPCD algorithm accounts for many-body collisions which are crucial for the build-up of correlations. This is known to be the origin of the hydrodynamic behavior in fluids.
## V Self-Diffusion
We study now the consequences of the different behavior in the two hydrodynamic regimes, which have been introduced in Sec. II, on the self-diffusion coefficient.
### V.1 Diffusion Coefficient
In the Green-Kubo formalism, the self-diffusion coefficient is given by $`D=\frac{1}{3}_0^{\mathrm{}}๐t๐ฏ(t)๐ฏ(0)`$. In the case that the time is discretized the integral has to be replaced by Ihle and Kroll (2001, 2003b)
$$D=\frac{1}{3}\left[\frac{1}{2}v^2(0)h+\underset{n=1}{\overset{\mathrm{}}{}}๐ฏ(nh)๐ฏ(0)\right]h.$$
(16)
In order to obtain an analytical prediction for the diffusion coefficient, an expression for $`๐ฏ(nh)๐ฏ(0)`$ is required. The Brownian approximation for the VACF given by Eq. (9) yields
$$D_0=\frac{k_BT}{m}h\left(\frac{1}{\gamma }\frac{1}{2}\right)$$
(17)
with $`\gamma `$ defined in Eq. (10). This expression coincides with that of Ref. Ihle and Kroll (2003b) with a different notation.
In the simulations, the diffusion coefficient is determined by a linear fit of the mean-square displacement for long times. We have checked that equivalent results for $`D`$ are also obtained directly from the VACF by employing Eq. (16).
Fig. 7 shows the relative deviation $`\mathrm{\Delta }D=(D_{sim}D_0)/D_0`$ of the diffusion coefficient from the expression (17). This expression should be a good approximation as long as the exponential decay (9) of the VACF applies. This is indeed the case for $`\lambda >0.6`$, what means that the long-time tail for these values has a negligible contribution for the diffusion coefficient. This is reasonable since the deviation from the exponential behavior appears when the VACF has decayed typically by three orders of magnitude (see $`\lambda =1.0`$ in Fig. 3). In contrast, Fig. 7 shows that the deviation from the Brownian behavior (17) increases with decreasing $`\lambda `$ for $`\lambda <0.5`$. This can be understood from the VACF since for small $`\lambda `$ the deviation from the exponential decay appears much earlier. Fig. 3 shows that for $`\lambda =0.1`$ the VACF has decayed only by about one order of magnitude when the deviation starts. This translates into a noticeable increment of the diffusion coefficient. This difference can be understood as a hydrodynamic enhancement of the diffusion coefficient for large values of the Schmidt number.
The diffusion coefficient for a simple MPCD fluid in two dimensions has been determined by Ihle and Kroll Ihle and Kroll (2003b). In their Fig. 15, results for $`\lambda =0.113`$ are presented as a function of the rotation angle; deviations from the theoretical prediction are found for large values of $`\alpha `$, which is in the range of parameters which we identify as the collective regime. They arrive at a similar conclusion that this is due to multiple encounters among particles. In three dimensions, some numerical results of the diffusion coefficient have been presented in Ref. Tucci and Kapral (2004), and good agreement with the molecular-chaos approximation has been found for a large range of number densities. However, the employed parameters (which correspond to $`\lambda >0.5`$) all belong to the particle regime, where we argue that a good agreement with the theory should be expected.
At this stage we come back to the discussion in Sec. III about the Schmidt number $`Sc=\nu /D`$. The analytic expression can be calculated from the viscosity $`\nu `$ in Eq. (3) and the diffusion coefficient in Eq. (17), as was already pointed out in Refs. Ihle and Kroll (2003b); Falck et al. (2004); Ripoll et al. (2004). Note that $`Sc`$ increases rapidly for small values $`\lambda 1`$ of the mean free path, where $`Sch^2`$. This allows arbitrary large values of the Schmidt number. Although very small values of the collision time significantly reduce the efficiency of the simulations, there is a range of $`\lambda `$-values which are not too small but still display fluid behavior corresponding to high $`Sc`$. On the other hand, the hydrodynamic enhancement of the diffusion coefficient in the collective regime leads to values of $`Sc`$ which are smaller than predicted by the analytical approximation. By substituting the numerically determined diffusion coefficient, it can be checked that $`Sc`$ is indeed smaller, but still large enough to display a fluid-like behavior.
### V.2 Continuum Time Limit
It is interesting to discuss the limit of small collision times $`h0`$, and small rotation angles $`\alpha 0`$. The leading contributions in the theoretical expressions (3) of the kinetic and collisional viscosity read in this limit
$$\nu _{coll}\frac{m\gamma _\rho }{36a}\left(\frac{\alpha ^2}{h}\right),\nu _{kin}\frac{k_BT}{a^3\gamma _\rho }\left(\frac{h}{\alpha ^2}\right),$$
(18)
with $`\gamma _\rho `$ defined in Eq. (10). This result shows that a finite viscosity is obtained in the continuum limit only if the ratio $`\alpha ^2/h`$ is kept constant. The additive term due to discrete times in Eq. (17) naturally vanishes in the continuum limit, because $`\gamma \alpha ^2`$.
The expressions (18) for the kinetic and collisional contributions to the viscosity show that the collective regime, where $`\nu _{coll}\nu _{kin}`$, corresponds to $`\alpha ^2/h1`$ in the continuum limit. In this regime, the leading contribution to the diffusion coefficient (17) is found to be
$$D\frac{3k_BT}{\gamma _\rho }\left(\frac{h}{\alpha ^2}\right).$$
(19)
The related Schmidt number
$$Sc=\frac{\nu }{D}\frac{1}{108}\frac{m^2\gamma _\rho ^2}{ak_BT}\left(\frac{\alpha ^2}{h}\right)^2$$
(20)
can be very large since $`\alpha ^2/h1`$. This shows that the model has a proper continuum limit. However, due to the requirement of very small collision times, this limit is not very convenient from a computational point of view.
It is very satisfactory to see that the Stokes-Einstein relation is satisfied in this case, since the diffusion coefficient is inversely proportional to the viscosity,
$$D\frac{k_BT}{6\pi \rho \nu _{coll}R}\mathrm{with}R=\frac{2a}{\pi \rho }$$
(21)
and defines an effective particle radius inversely proportional to the number density. We want to emphasize, however, that the Stokes-Einstein relation is not only satisfied in the continuum limit, but always when the additive term $`1/2`$ in Eq. (17) can be neglected and the collisional dominates the kinetic viscosity. In this case, Eq. (21) is also valid.
## VI Dynamics of Embedded Particles
After the behavior of a simple MPCD fluid has been characterized, the next important question is how complex fluids can be modeled. As first step, we investigate the behavior of a single heavy point-like particle, which could represent a solute particle or a colloidal sphere embedded in a simple fluid. Also, the monomers in a polymer chain can be represented as point particles Malevanets and Yeomans (2000); Ripoll et al. (2004); Winkler et al. (2004); Ali et al. (2004). This is a quite convenient strategy, since the solute-solvent interactions are modeled by just including the point-like solute particles in the collision step. Then we study different concentrations of these heavy particles.
### VI.1 Single Heavy Tracer Particle
For the simulation of heavy point-like particles embedded in a solvent, the algorithm is the same as described for the simple fluid in Sec. II. The only point where the higher mass plays a role is in the calculation of the velocity of the center of mass, where the different particle masses have to be taken into account via $`๐ฏ_{cm,i}(t)=_j^{(i,t)}(m_j๐ฏ_j)/_jm_j`$. In thermal equilibrium, the average kinetic energy of light and heavy particles is the same. Therefore, the average momentum of the heavy particle of mass $`M`$ is a factor $`(M/m)^{1/2}`$ larger than the average momentum of a light particle. This implies that a heavy particle has a larger contribution in the center-of-mass velocity than a light particle. Since the center-of-mass velocity and therefore also the velocities of all particles after the collision step depends on $`M`$ and the mass $`m\rho `$ of the solvent particles in a collision cell, the effective coupling between the solvent and the solute must depend in general on the ratio $`M/(m\rho )`$.
We denote the heavy particle position and velocity with capital letters $`๐`$ and $`๐`$. Of course, all types of particles are involved in the center-of-mass calculation or other sums over particles. The VACF can be calculated in the molecular-chaos approximation as explained in Sec. IV.1, except for the center-of-mass correlation in Eq. (7), which for the heavy particle yields
$$๐_{cm}((n1)h)๐(0)\frac{M}{m\rho +M}๐((n1)h)๐(0).$$
(22)
because in the collision box of the heavy particle the total mass is $`(M+m\rho )`$. The correlation at time zero depends now on the heavy particle mass,
$$๐^2(0)=3\frac{k_BT}{M}.$$
(23)
By inserting these results in the expression equivalent to Eq. (5), we obtain the molecular-chaos approximation for the normalized VACF of the heavy particle,
$$C_V(t)\frac{๐(nh)๐(0)}{V^2(0)}(1\gamma )^n,$$
(24)
where the decorrelation factor $`\gamma `$ is now given by
$$\gamma =\gamma _\alpha \frac{m\rho }{m\rho +M}\gamma _\alpha \gamma _1,$$
(25)
and $`\gamma _1`$ is defined for one heavy particle in the presence of $`\rho `$ fluid particles, in contrast to $`\gamma _\rho `$ in Eq. (10), where a fluid particle is surrounded by $`(\rho 1)`$ other fluid particles.
In Fig. 8 results for the normalized VACF of one heavy particle in the collective regime are presented for different values of its mass. The solvent mass density has been chosen to be equal to the solute mass, i.e., $`\rho =M/m`$. In this way, $`\gamma =\gamma _\alpha /2`$ and the analytical expression (24) is independent of the heavy particle mass. Fig. 8 shows that after the second collision all the simulation data exhibit a non-exponential decay. This is not very surprising, since a similar behavior was observed for the simple fluid in Fig. 3 for parameter values within the collective regime. A slightly slower decay is displayed at lower number density $`\rho `$, but an asymptotic curve is clearly approached for large values of $`\rho `$. The deviations for small $`\rho `$ are due to the presence of density fluctuations.
The dependence of the VACF of a single heavy tracer particle of mass $`M=5m`$ on the mean free path $`\lambda `$ of the solvent is shown in Fig. 9. Corresponding results for the simple fluid are shown in Fig. 3. For $`\lambda =0.1`$, the qualitative behavior of tracer particles with $`M=m`$ and $`M=5m`$ is very similar. The first collision perfectly follows the molecular-chaos approximation, followed by a slower-than-exponential decay for intermediate times and a crossover to a power-law decay for long times. However, note that since the exponential decay is slower for the heavy particle, the deviations from Brownian behavior appear when the VACF has decayed to approximately one third of its original value for the employed values of $`\rho `$ and $`\alpha `$, while for the simple fluid case the VACF has decayed to $`6\%`$ of its original value. This implies that the hydrodynamic enhancement is more pronounced for particles of larger mass. For $`\lambda =1`$, small deviations from the exponential decay are visible for short times; for long times, the crossover to the power-law behavior can be seen.
Analytical approximation for the diffusion coefficient can be calculated similar to Sec. V.1. It reads,
$$D_0=\frac{k_BT}{M}h\left(\frac{1}{\gamma }\frac{1}{2}\right)$$
(26)
where the decorrelation factor $`\gamma `$ is now given by Eq. (25).
Simulation results for the diffusion coefficient $`D_M`$ of a heavy tracer particle are plotted in Fig. 10 as a function of the mass $`M/m`$, for fixed solvent density $`\rho =5`$ and two different sets of parameters. The agreement of the simulations with the approximation (26) is again very good for parameter values within the particle regime, $`\lambda =1`$ and $`\alpha =45`$, but not within the collective regime, $`\lambda =0.1`$ and $`\alpha =130`$. This is the same behavior as observed in the simple fluid (see Fig. 7) and indicates again the presence of a hydrodynamic contribution to the diffusion coefficient in the collective regime.
In Fig. 11, the hydrodynamic contribution to the diffusion coefficient (in units of the Brownian contribution) is plotted as a function of the scaled mean free path $`\lambda `$ for a heavy tracer particle of mass $`M=5m`$ and for a simple fluid tracer particle (compare Fig. 7). It can be seen that $`D_H`$ increases considerably for small $`\lambda `$ in both cases. This increment is significantly more pronounced for the heavy particle, which corresponds to the slower decay of the VACF in Fig. 9 for the larger mass. A small deviation of the VACF from the exponential decay was observed in Fig. 9 at short times for $`\lambda =1`$. This deviation translates into the small hydrodynamic enhancement of the diffusion coefficient of the heavy particle that can be seen in Fig. 11, even at โlargeโ mean free paths $`\lambda 1`$.
Fig. 10 shows that for a fixed density $`\rho `$ in the collective regime, the hydrodynamic enhancement increases with increasing mass of the solute particle until $`M/m2\rho `$, and then levels off and becomes independent of the solute mass for $`M/m\rho `$. This is consistent with the diffusion behavior of colloidal spheres, where the diffusion coefficient is independent of the mass of the colloidal particles.
Kikuchi et al. Kikuchi et al. (2003) determine numerically the friction coefficient acting on a particle of mass $`M`$ and velocity $`๐ฏ`$ in a MPCD solvent. Their simulation results, for a fluid of $`\lambda 0.9`$, compare nicely with the analytical prediction, independently on the mass of the particle. However, we want to point out that this agreement is not very surprising, since their result is obtained from the velocity autocorrelation function after the first collision step, where the molecular-chaos approximation is always exact (see Sec. IV.3).
The increase of the hydrodynamic coupling of solute and solvent with increasing solute mass can be understood as follows. The relative mass of the solute and solvent particles appears in the collision step via the calculation of the center-of-mass velocity. If solute particles have the same mass as solvent particles and there is a large number of solvent particles per cell, the solvent particles transfer a large random momentum to the solute particle. Simultaneously, the effect of the solute particle momentum on the solvent is small. For this reason, the hydrodynamic contribution to the diffusion constant of a particle of equal mass, shown in Fig. 7, is only of the order of $`30\%`$ for the largest Schmidt number considered. In contrast, this hydrodynamic enhancement is $`65\%`$ when $`M/m\rho `$ and $`75\%`$ when $`M/m2\rho `$, as can be seen in Fig. 10. A very large mass of the solute particle is not very convenient either, because it implies a large ballistic regime and a long diffusion time. Therefore, we conclude that a mass $`M/m\rho `$ for the solute particle is a optimal choice to enhance the hydrodynamic coupling between solute and fluid particles.
### VI.2 Finite Concentration of Heavy Point-Like Particles
At a finite concentration of solute particles, an important question is to which extent solute particles build up hydrodynamic interactions among themselves through the fluid particles when simulated with MPCD. We study therefore systems with different concentrations of heavy particles for sets of parameters within the particle and the collective regimes, respectively. We address this question by investigating the tracer-diffusion coefficient.
Simulations with different heavy particle concentrations are performed by changing the total number $`N_M`$ of heavy particles but keeping fixed the volume $`V=L^3`$ and the number $`N`$ of solvent particles. The corresponding number density of heavy particles is defined as $`\varphi =N_M(a/L)^3`$. In Fig. 12, the diffusion coefficients for three different values of the mean free path are displayed. Very surprisingly, when the data are normalized by the corresponding diffusion coefficients in the limit of vanishing density $`\varphi `$, all three data sets, which are both in the particle and the collective regime, collapse onto a single curve. We recall that the hydrodynamic enhancement for the diffusion coefficient of a single heavy particle, here denoted as $`D_M(0)`$, is quite different among these three values of $`\lambda `$ (see Fig. 11). It can be inferred from the data collapse in Fig. 12 that there is no extra hydrodynamic contribution among these heavy particles, which is consistent with the idea that there is no hydrodynamic screening for point particles Ahlrichs et al. (2001); Ladd (1996).
The dependence of the diffusion coefficient on the heavy particle number density can be understood along the same lines as for the simple fluid or the single heavy particle. We assume that in each collision box there is a fixed number of fluid particles $`\rho `$, but that the number of heavy particles $`n`$, fluctuates from one collision box to another. The probability $`P(n)`$ of a given heavy particle to be found in a cell with a total of $`n1`$ other heavy particles is given by the Poisson distribution function, $`P(n)=e^\varphi \varphi ^{n1}/(n1)!`$. The corresponding decorrelation factor for a heavy particle in a collision box with $`(n1)`$ other heavy particles and $`\rho `$ fluid ones is
$$\gamma _n=1M/(\rho m+nM),$$
(27)
compare the definition of $`\gamma _1`$ in Eq. (25) for a single heavy particle in a collision box. The diffusion coefficient is then given by Eq. (26), where the decorrelation factor is now $`\gamma =\gamma _\alpha _{n=1}^{\mathrm{}}P(n)\gamma _n`$. In the regime of low number density, $`\varphi 1`$, this implies
$$\gamma =\gamma _\alpha \left[(1\varphi )\gamma _1+\varphi \gamma _2+๐ช\left(\varphi ^2\right)\right].$$
(28)
In the special case of $`\rho =M/m`$, the sum can be evaluated analytically and yields
$$\gamma =\gamma _\alpha \left(1\left(e^\varphi +\varphi 1\right)/\varphi ^2\right).$$
(29)
In Fig. 12 the simulation data for the normalized diffusion coefficient at different volume fractions are compared with the theoretical prediction obtained from Eq. (26) with the decorrelation function in Eq. (29). It can be seen that this prediction overestimates the values for the diffusion coefficients. Further studies are required to understand the origin of this deviation.
## VII Hybrid Dynamics
In order to go one step further in the development of an efficient simulation technique for suspensions of colloidal particles with MPCD, we next investigate the effect of excluded-volume interactions between the heavy particles. To this end, the MPCD algorithm has to be combined with standard molecular dynamics (MD) for the solute particles.
### VII.1 The Model
We consider a dispersion of spherical colloidal particles in three dimensions. The interactions of solvent particles among themselves and with colloids take place in the MPCD collisional step, exactly in the same way as described for the heavy point-like particles in Sec. VI.1. However, the streaming step (1) is used only for the solvent particles. The position update of the colloidal particles is performed in several MD steps between MPCD collisions. In these MD steps, colloids interact via an excluded-volume potential. We use the truncated repulsive Lennard Jones potential Andersen et al. (1976)
$$V^{RLJ}(r)=\{\begin{array}{ccc}\hfill 4\epsilon \left[\right(\frac{\sigma }{r})^{12}(\frac{\sigma }{r}\left)^6\right]+\epsilon \text{,}& & rr_{min}\hfill \\ \hfill 0\text{,}& & r>r_{min}\hfill \end{array}$$
(30)
where $`r`$ is the distance between the centers of the colloidal particles. The parameter $`\sigma `$ is related to the particle diameter; it is chosen to equal the collision box length, $`\sigma =a`$, so that there is typically no more than one colloid particle in each collision box. The potential strength is taken to be equal to the thermal energy, $`\epsilon =k_BT`$, the cut-off radius is $`r_{min}=2^{1/6}\sigma `$, and the mass of the particles is taken to be $`M=5m`$. The MD time steps are integrated with the velocity-Verlet algorithm Allen and Tildesley (1987) with a time step $`\mathrm{\Delta }t=0.002\sqrt{\epsilon /ma^2}`$.
In other words, we consider a system of colloidal particles interacting through repulsive Lennard Jones potentials whose positions and velocities evolve in discrete time intervals $`\mathrm{\Delta }t`$. This procedure is interrupted every $`h/\mathrm{\Delta }t`$ steps for the interaction with the fluid particles. This interaction is a MPCD event where solvent and solute particles interchange momentum. This implies that the solvent particles can enter the cores of the colloidal particles, but the colloids cannot interpenetrate each other.
The hybrid model described here is a variant of the model introduced previously by Malevanets and Kapral Malevanets and Kapral (2000, 2004). In their model, both the solute-solute and solute-solvent interactions were taken into account through excluded-volume potentials with MD, and only the solvent-solvent interactions were mesoscopically described through MPCD. The advantage of the model described here comes from the fact that in the MD steps just the solute particles are considered. This leads to a considerable speed up of the simulations.
### VII.2 Diffusion in Colloidal Dispersions
We measure the diffusion coefficient of the dispersion through the mean-square displacement of a tracer particle, as before. Simulations are performed for different colloidal concentrations. The volume fraction of colloidal particles $`\phi `$ is the fraction of the total volume $`V`$ occupied by the colloidal particles, $`\phi =(\pi /6)\sigma _{eff}^3\rho _M`$, where the effective diameter $`\sigma _{eff}`$ is determined by the Barker-Henderson expression Barker and Henderson (1976)
$$\sigma _{eff}=_0^{r_{min}}๐r\left[1\mathrm{exp}(V^{RLJ}(r)/k_BT)\right].$$
(31)
For our choice of Lennard-Jones parameters, this gives $`\sigma _{eff}=1.01\sigma `$. The number density of colloidal particles is $`\rho _M=(N_M1)/VN_M/V`$, where $`N_M`$is the number of heavy particles with excluded-volume interactions.
For later comparison and better understanding of our hybrid model results, we recall first the basic behavior of a system with excluded-volume interactions only. In Fig. 13 we show the results for the diffusion coefficient of a MD simulation of repulsive Lennard-Jones particles. Kinetic theory for hard spheres predicts in the low-density limit Boon and Yip (1980)
$$D_{\mathrm{MD}}(\phi )=\frac{3}{8\sigma _{eff}^2}\sqrt{\frac{k_BT}{\pi M}}\frac{1}{\rho _M},$$
(32)
This analytical prediction is depicted in Fig. 13 together with the simulation results. It can be seen that for small volume fraction, the $`\phi ^1`$ behavior is properly reproduced, while for large volume fractions a linear behavior can be inferred.
The density dependence of the self-diffusion coefficient of colloidal hard spheres in a hydrodynamic bath has been calculated in Ref. Dhont (1996),
$$D_S(\phi )=D_S(0)\left[12.1\phi +๐ช\left(\phi ^2\right)\right].$$
(33)
The diffusion coefficient now decreases linearly with the volume fraction, in contrast with the kinetic theory result (32) for a gas of hard spheres. In the calculation of Eq. (33), Brownian and hydrodynamic terms have to be considered, and it has been found that the hydrodynamic terms almost cancel. For a colloidal dispersion in a Brownian bath Dhont (1996) the first-order correction in Eq. (33) equals $`2.0\phi `$. Thus, no significant differences are expected between Brownian and hydrodynamic measurements of the diffusion coefficient.
Simulation results with the hybrid method are shown in Fig. 14. The simulations presented here are performed with rotation angle $`\alpha =130`$, fluid number density $`\rho =5`$, and mass $`M=5m`$ of the colloidal particle. We vary the mean free path between $`\lambda =0.02`$ and $`\lambda =2.0`$.
In the limit of very small volume fractions, the repulsive interactions between colloids are negligible, and the colloidal dispersion will behave as the dispersion of heavy point-like particles presented in Secs. VI.1 and VI.2. In this limit, we know from Eq. (26) that the diffusion coefficient $`D(0)`$ increases with the mean free path $`\lambda `$ (with $`D(0)\lambda `$ in the molecular-chaos approximation). The decrease of $`D(\phi )`$ with decreasing $`\lambda `$ displayed in Fig. 14 arises then as a natural consequence. Furthermore, the MPCD interactions of the colloid particles with the fluid imply that the self-diffusion coefficient at small densities does not diverge but goes to finite value dictated by Eq. (26).
For small but finite volume fractions, in the case of large values of $`\lambda `$, we observe a behavior reminiscent of the $`\phi ^1`$ decay of hard-sphere gases, instead of the linear decrease expected from Eq. (33). This can be understood since, in the limit of very large mean free paths, the colloids will essentially interact with each other rather than with the solvent. This behavior is not seen in experiments of colloidal dispersions, because the diffusive length scale is typically much smaller than the diameter of the particles.
Therefore, the appropriate parameters for the modeling of colloidal dispersions have again to be chosen in the collective regime. In Fig. 15, the normalized diffusion coefficient is shown, where $`D(0)`$ is extrapolated from the simulated data. The linear behavior in Eq. (33) is indeed observed for the smallest values of the mean free path, $`\lambda =0.02`$ and $`\lambda =0.1`$, within the accuracy of the simulations. Thus, we find that in order to obtain the theoretically predicted behavior (33) from simulation of the MD - MPCD hybrid model, small values of the mean free path and large values of the rotation angle $`\alpha `$ are required, i.e. parameters in the collective regime.
However, an almost identical dependence of the diffusion coefficient on the volume fraction is predicted theoretically in the absence of hydrodynamics interactions. In order to investigate this point in more detail, we have performed simulations of a hybrid model similar to the one presented here, but with a completely Brownian solvent. One way of transforming a MPCD fluid in a Brownian solvent has been introduced by Kikuchi et al. Kikuchi et al. (2002), where the velocities among all the fluid particles are randomly interchanged after each MPCD collision step. We propose an alternative method which does not consider any solvent particles. Instead, at every $`h/\mathrm{\Delta }t`$ steps the MD dynamics is interrupted for a rotation of the (full) velocity of each colloid around a random axis by and angle $`\alpha `$. In this case, the diffusion coefficient at zero volume fraction $`D(0)`$ is given by Eq. (26) but the decorrelation factor is $`\gamma \gamma _\alpha `$ with $`\gamma _\alpha `$ of Eq. (6). The simulation results of the dependence of the diffusion coefficient on the volume fraction are quite similar to those displayed in Fig. 15. The data for $`D(\phi )`$ follow a linear decay only for very small values of $`\lambda `$, where the friction is large and $`D(0)`$ is small enough to represent a fluid. For large values of $`\lambda `$, $`D(\phi )`$ has a concave shape, reminiscent of the $`\phi ^1`$ behavior of gases, similarly as observed for the hydrodynamic simulations.
Simulations with a similar hybrid method of a two-dimensional colloidal suspension have been reported by Falck et al. Falck et al. (2004). In the majority of the presented results, they consider excluded-volume interaction among colloids but not between colloids and solvent particles. They measure an apparent tracer diffusion coefficient (since the diffusion coefficient in two dimensions diverges with increasing system size) of the colloids for different concentrations. Three different Schmidt numbers are studied. A similar trend in the data is observed as in our simulations: the normalized diffusion coefficient increases with increasing Schmidt number, in particular for intermediate values of the volume fraction $`\phi `$. However, our interpretation is different. While Falck et al. attribute this effect to hydrodynamics, we believe that it is due to the crossover from gas-like to diffusive behavior of the colloidal dynamics.
In summary, our hybrid model describes the dynamics of a dispersion of hard-sphere colloids very well in the collective regime of the solvent. In the hydrodynamic interaction, only the leading contribution for large distances is included in our model. This implies that lubrication forces between neighboring particles at short distances, as well as the coupling between rotational degrees of freedom, are neglected. We conclude from the very weak dependence of our results for the normalized diffusion coefficients on the mean free path, which controls the strength of the hydrodynamic interaction, that our model works very well for not too concentrated colloidal dispersions.
## VIII Summary and Conclusions
In this paper we have performed a detailed analysis of the hydrodynamic properties of a fluid simulated with MPCD. We identify two hydrodynamic regimes in terms of the parameters of the MPCD algorithm. The particle regime is characterized by dynamical properties being closer to those of a gas than to those of a liquid. The Schmidt number is small and the dominant transport mechanism is kinetic transport. This is the regime obtained for large values of the collision time and/or small values of the rotation angle. The second and more relevant regime for fluid simulations is the collective regime. In this regime the Schmidt number is large and collisional transport dominates over kinetic transport โ this characterizes liquid-like behavior. These properties are obtained for large values of the rotation angle and small values of the collision time.
Different quantities have been measured in both regimes. The main conclusion is that the diffusion coefficient shows a hydrodynamic enhancement in the collective regime. In the study of the VACF we observe that the behavior can be understood in both regimes as an exponentially decay for short times and algebraic decay for long times. In the particle regime, a simple crossover between both behaviors is observed while an extra intermediate behavior is displayed in the collective regime. This intermediate behavior of the VACF is typically a slower than the initial exponential decay. We have shown that the origin of this intermediate decay region is due to the build-up of correlations by many-body collisions, which is in conceptual agreement with the hydrodynamic behavior. The theoretical predictions for the diffusion coefficient are based on a molecular-chaos assumption, which gives an exponential decay of the VACF. Consequently, a deviation from the theoretical prediction is found in the collective regime. This deviation can be understood as a hydrodynamic contribution to the Brownian value.
In a further step, we have investigated the differences between the particle and the collective regime for complex fluids. We have studied the behavior of heavy particles embedded in the MPCD fluid which can represent solute or colloidal particles dissolved in a simple fluid. This study demonstrates that optimal hydrodynamic coupling occurs when the mass of the tagged particle is on the order of the solvent mass in a collision cell.
In order to describe colloidal dispersions at finite volume fractions, it is necessary to account for excluded volume interactions among colloidal particles. To this end, a hybrid model was studied, which combines MPCD for the solvent with MD simulations for the colloidal particles. We show that only for parameters within the collective regime does the hybrid model reproduce the proper hydrodynamic behavior. In this case, the results agree well with the theoretical calculations with and without hydrodynamic interactions, as well as with experimental results.
A more precise modeling of colloidal particles would require new interactions among fluid and colloids such that fluid particles would not freely travel through colloidal particles and eventually angular momentum could be interchanged among them.
In the future, it will be interesting to explore in which applications a more detailed description of colloidal interactions is necessary, as compared to our simplified model which allows more particles and larger system sizes, and is therefore well suited to study cooperative phenomena.
## Acknowledgments
We thank G. Vliegenthart, H. Noguchi, D.M. Kroll, T. Ihle, N. Kikuchi and A. Lamura for helpful discussions. Financial support of this work by the German Research Foundation (DFG) within the SFB TR6, โPhysics of Colloidal Dispersions in External Fieldsโ, is gratefully acknowledged. M.R. also acknowledges partial support from the Projects No. BFM2001-0290 and FIS2004-01934.
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# The influence of strength of hyperon-hyperon interactions on neutron star properties.
## 1 Introduction
The analysis of the role of strangeness in nuclear structure in the aspect of multi-strange system is of great importance for relativistic heavy-ion collisions and for astrophysics for the description of hyperon star matter. At the core of a neutron star the matter density ranges from a few times the density of normal nuclear matter to about an order of a magnitude higher. Thus exotic forms of matter such as hyperons are expected to emerge in the interior of a neutron star. The appearance of these additional degrees of freedom and their impact on a neutron star structure have been the subject of extensive studies , , , . Properties of matter at such extreme densities are of particular importance in determining forms of equations of state relevant to neutron stars and successively in examining their global parameters.
The existence of bound strange hadronic matter which in addition to nucleons contains also hyperons has profound consequences for astrophysics. The starting point in studying the role of strangeness in nuclear structure is the knowledge of the properties of a single hypernucleus. Quantum chromodynamics should be applied to the theoretical description of hadronic systems owing to the fact that it constitutes the fundamental theory of strong interactions. However, at the hadronic energy scale where the experimentally observed degrees of freedom are not quarks but hadrons the direct description of nuclei in terms of QCD become inadequate and thus alternative approaches had to be formulated. One of them is quantum hadrodynamics (QHD) , which gives quantitative description of the nuclear many body problem. QHD is a relativistic quantum field theory in which nuclear matter description in terms of baryons and mesons is provided. The original model (QHD-I) contains nucleons interacting through the exchange of simulating medium range attraction $`\sigma `$ meson and $`\omega `$ meson responsible for short range repulsion. Extension (QHD-II) of this theory , includes also the isovector meson $`\rho `$. Theoretical description of strange hadronic matter, which satisfactorily reproduces nucleon-nucleon and hyperon-nucleon data, has been given within the non-relativistic and relativistic mean field models. This approach is based on the notion of the meson-exchange model in which baryons interact through the exchange of mesons. In addition to the $`\sigma `$, $`\omega `$ and $`\rho `$ mesons these models contain $`\sigma ^{}`$ and $`\varphi `$ mesons, introduced in order to reproduce the strong attractive hyperon-hyperon interactions .
The vector coupling constants are chosen according to SU(6) symmetry whereas the scalar coupling constants are fixed to hypernuclear data. Recent reports on observations of a $`{}_{\mathrm{\Lambda }\mathrm{\Lambda }}{}^{6}He`$ hypernucleus made by Takahashi et al. provided information on the $`\mathrm{\Lambda }\mathrm{\Lambda }`$ interaction energy $`\delta B_{\mathrm{\Lambda }\mathrm{\Lambda }}=1.01\pm 0.20_{0.11}^{+0.18}`$ MeV. This allows one to determine the value of the $`\mathrm{\Lambda }`$ well depth in $`\mathrm{\Lambda }`$ matter at density $`0.5\rho _0`$ ($`\rho _0`$ denotes the saturation density) at the level of $`U_\mathrm{\Lambda }^\mathrm{\Lambda }5`$ MeV . The results of earlier experiments , , , , give the value of $`U_\mathrm{\Lambda }^\mathrm{\Lambda }`$ estimated with the use of Nijmegen model D at the level of $`20`$ MeV . This paper examines the implications of the strength of hyperon-hyperon coupling constants on the $`\beta `$-equilibrated hyperon star matter, hence doing all calculations within the framework of relativistic mean field model with two parameterizations, namely the standard TM1 and TMA . The results have been obtained for two cases: the weak and strong $`YY`$ interactions. First, the properties of the isospin symmetric strange hadronic matter have been investigated. The obtained saturation curves resemble those obtained by Song et al.. The results for the two models (TM1 and TMA) are very similar. The extension of the considered model to the $`\beta `$ equilibrated asymmetric hyperon star model has been done in the subsequent section. As a result of this the composition, the equation of state and the hyperon star structure for the two parameterizations have been obtained.
## 2 Hypernuclei
Information on single hypernuclei can be summarized in the following points:
* $`\mathrm{\Lambda }`$ hypernuclei: there is a large amount of data on the binding energies $`B_\mathrm{\Lambda }`$ of $`\mathrm{\Lambda }`$โs bound in various single particle orbitals in hypernuclei. This enables us to study deeply bound states inside the nucleus over an extensive range of mass number. An analysis of these data with the use of Skyrme-Hartree-Fock model gives the potential depth of a single $`\mathrm{\Lambda }`$ in nuclear matter at the value of
$$U_\mathrm{\Lambda }^{(N)}2730MeV$$
(1)
which corresponds to 1/3-1/2 of the nucleon well depth $`U_N^{(N)}`$ (in this text all potentials are considered as attractive but the convention of positive sign has been used).
* $`\mathrm{\Sigma }`$ hypernuclei : the experimental status of $`\mathrm{\Sigma }`$ nucleus potential still remains controversial. The calculations of $`\mathrm{\Sigma }`$ hypernuclei have been based on analysis of $`\mathrm{\Sigma }^{}`$ atomic data. Phenomenological analysis of level shifts and widths in $`\mathrm{\Sigma }^{}`$ atoms made by Batty et al. indicates that the $`\mathrm{\Sigma }`$ potential is attractive only at the nuclear surface changing into repulsive one in increasing density. The small attractive component of this potential is not sufficient to form bound $`\mathrm{\Sigma }`$-hypernuclei. Balberg et al. in their paper show that the system which includes $`\mathrm{\Sigma }`$, $`\mathrm{\Lambda }`$ hyperons and nucleons will be unstable with respect to strong reactions $`\mathrm{\Sigma }+N\mathrm{\Lambda }+N`$ (78 MeV), $`\mathrm{\Sigma }+\mathrm{\Sigma }\mathrm{\Lambda }+\mathrm{\Lambda }`$ (156 MeV), $`\mathrm{\Sigma }+\mathrm{\Lambda }\mathrm{\Xi }+N`$ (50 MeV), $`\mathrm{\Sigma }+\mathrm{\Xi }\mathrm{\Lambda }+\mathrm{\Xi }`$ (80 MeV) , in parenthesis the $`Q`$ values for each reaction are given.
* $`\mathrm{\Xi }`$ hypernuclei - for the $`\mathrm{\Xi }`$ hypernuclei there exist a few emulsion events reported in literature indicating the existence of a bound system. The interpreted data give the potential of a $`\mathrm{\Xi }`$ in nuclear matter with a depth of
$$U_\mathrm{\Xi }^{(N)}2025MeV.$$
(2)
* $`\mathrm{\Lambda }\mathrm{\Lambda }`$ hypernuclei: the properties of single hypernuclei are one aspect of studying strangeness in nuclear systems; the other is connected with multi-strange systems. The extrapolation to a multi-strange system is based on the data concerning double $`\mathrm{\Lambda }`$-hypernuclei. Data on $`\mathrm{\Lambda }\mathrm{\Lambda }`$ hypernuclei are extremely scarce. Observation of double-strange hypernuclei $`\mathrm{\Lambda }\mathrm{\Lambda }`$ provide information about the $`\mathrm{\Lambda }\mathrm{\Lambda }`$ interaction. Several events have been identified which indicate an attractive $`\mathrm{\Lambda }\mathrm{\Lambda }`$ interaction. The analysis of the data allows one to estimate the binding energies of $`{}_{\mathrm{\Lambda }\mathrm{\Lambda }}{}^{6}He`$, $`{}_{\mathrm{\Lambda }\mathrm{\Lambda }}{}^{10}Be`$ and $`{}_{\mathrm{\Lambda }\mathrm{\Lambda }}{}^{13}B`$. The measurement of the masses of double-$`\mathrm{\Lambda }`$ hypernuclei gives information on the sum of the binding energy of the two $`\mathrm{\Lambda }`$ hyperons $`B_{\mathrm{\Lambda }\mathrm{\Lambda }}`$ and the $`\mathrm{\Lambda }\mathrm{\Lambda }`$ interaction energy $`\mathrm{\Delta }B_{\mathrm{\Lambda }\mathrm{\Lambda }}`$. These two quantities can be defined as
$`B_{\mathrm{\Lambda }\mathrm{\Lambda }}({}_{\mathrm{\Lambda }\mathrm{\Lambda }}{}^{A}Z)`$ $`=`$ $`B_\mathrm{\Lambda }({}_{\mathrm{\Lambda }\mathrm{\Lambda }}{}^{A}Z)+B_\mathrm{\Lambda }({}_{\mathrm{\Lambda }}{}^{A1}Z)`$ (3)
$`\mathrm{\Delta }B_{\mathrm{\Lambda }\mathrm{\Lambda }}({}_{\mathrm{\Lambda }\mathrm{\Lambda }}{}^{A}Z)`$ $`=`$ $`B_\mathrm{\Lambda }({}_{\mathrm{\Lambda }\mathrm{\Lambda }}{}^{A}Z)B_\mathrm{\Lambda }({}_{\mathrm{\Lambda }}{}^{A1}Z).`$
Table LABEL:tab:lambda compiles experimental values of the two observables mentioned above , , , .
The obtained values of $`\mathrm{\Delta }B_{\mathrm{\Lambda }\mathrm{\Lambda }}`$ indicate that the $`\mathrm{\Lambda }\mathrm{\Lambda }`$ interaction is attractive and rather strong. The value of $`\mathrm{\Delta }B_{NN}`$ equals 6-7 MeV for comparison. Following the estimation made by Schaffner et al. it is possible to approximately determine the ratio of the $`\mathrm{\Lambda }`$ well depth in $`\mathrm{\Lambda }`$ matter and nucleon well depth in nuclear matter
$$\frac{U_\mathrm{\Lambda }^{(\mathrm{\Lambda })}}{U_N^{(N)}}=\frac{V_{\mathrm{\Lambda }\mathrm{\Lambda }}}{V_{NN}}\frac{1/4}{3/8}\frac{1}{2}$$
(4)
where $`1/2`$ stands for a nucleon and $`\mathrm{\Lambda }`$ density ratio and $`(1/4)/(3/8)`$ denotes spin-isospin weights of spatially symmetric two-body configurations. Taking into account the value of $`V_{\mathrm{\Lambda }\mathrm{\Lambda }}\mathrm{\Delta }B_{\mathrm{\Lambda }\mathrm{\Lambda }}45`$ MeV and the value $`V_{NN}67`$ MeV one can estimate the ratio $`V_{\mathrm{\Lambda }\mathrm{\Lambda }}/V_{NN}3/4`$ and thus the equation (4) gives
$$\frac{U_\mathrm{\Lambda }^{(\mathrm{\Lambda })}}{U_N^{(N)}}\frac{1}{4}.$$
(5)
For $`U_N^N80`$ MeV the estimated value of $`U_\mathrm{\Lambda }^\mathrm{\Lambda }`$ potential is $`20`$ MeV.
However, the recent data analysis of double hypernucleus $`{}_{\mathrm{\Lambda }\mathrm{\Lambda }}{}^{6}He`$ done by Takahashi et al. gives the following value of the $`\mathrm{\Lambda }\mathrm{\Lambda }`$ interaction energy $`\mathrm{\Delta }B_{\mathrm{\Lambda }\mathrm{\Lambda }}=1.01\pm 0.20_{0.11}^{+0.18}`$ MeV what indicate that the interaction is much weaker. The potential well depth evaluated for this data has the value $`U_\mathrm{\Lambda }^{(\mathrm{\Lambda })}5`$ MeV.
## 3 The model
The theoretical description of the properties of strange hadronic matter is given within the relativistic mean field approach. The considered model involves baryons interacting through the exchange of simulating medium range attraction $`\sigma `$ meson and $`\omega `$ meson responsible for short range repulsion. The model also includes the isovector meson $`\rho `$. In order to reproduce attractive hyperon-hyperon interaction two additional hidden-strangeness mesons, which do not couple to nucleons, have been introduced, namely the scalar meson $`f_0(975)`$ (denoted as $`\sigma ^{}`$) and the vector meson $`\varphi (1020)`$.
The effective Lagrangian function for the system can be written as a sum of a baryonic part including the full octet of baryons together with baryon-meson interaction terms and a mesonic part
$$=_B+_M.$$
(6)
The interacting baryons are described by the Lagrangian function $`_B`$ which is given by
$$_B=\underset{B}{}\overline{\psi }_B(i\gamma ^\mu D_\mu m_B+g_{\sigma B}\sigma +g_{\sigma ^{}B}\sigma ^{})\psi _B.$$
(7)
where $`B`$ stands for $`N,\mathrm{\Lambda },\mathrm{\Sigma },\mathrm{\Xi }`$ and $`\mathrm{\Psi }_B^T=(\psi _N,\psi _\mathrm{\Lambda },\psi _\mathrm{\Sigma },\psi _\mathrm{\Xi })`$. The covariant derivative $`D_\mu `$ is defined as
$$D_\mu =_\mu +ig_{\omega B}\omega _\mu +ig_{\varphi B}\varphi _\mu +ig_{\rho B}I_{3B}\tau ^a\rho _\mu ^a.$$
(8)
The Lagrangian density for meson fields takes the form
$`_M`$ $`=`$ $`{\displaystyle \frac{1}{2}}_\mu \sigma ^\mu \sigma U(\sigma )+{\displaystyle \frac{1}{2}}_\mu \sigma ^{}^\mu \sigma ^{}{\displaystyle \frac{1}{2}}m_\sigma ^{}^2\sigma ^2+`$ (9)
$`+`$ $`{\displaystyle \frac{1}{2}}m_\varphi ^2\varphi _\mu \varphi ^\mu {\displaystyle \frac{1}{4}}\varphi _{\mu \nu }\varphi ^{\mu \nu }{\displaystyle \frac{1}{4}}\mathrm{\Omega }_{\mu \nu }\mathrm{\Omega }^{\mu \nu }+{\displaystyle \frac{1}{2}}m_\omega ^2\omega _\mu \omega ^\mu +`$
$``$ $`{\displaystyle \frac{1}{4}}R_{\mu \nu }^aR^{a\mu \nu }+{\displaystyle \frac{1}{2}}m_\rho ^2\rho _\mu ^a\rho ^{a\mu }+{\displaystyle \frac{1}{4}}c_3(\omega _\mu \omega ^\mu )^2.`$
The field tensors $`R_{\mu \nu }^a`$, $`\mathrm{\Omega }_{\mu \nu }`$, $`\varphi _{\mu \nu }`$ are defined as
$$R_{\mu \nu }^a=_\mu \rho _\nu ^a_\nu \rho _\mu ^a+g_\rho \epsilon _{abc}\rho _\mu ^b\rho _\nu ^c,$$
(10)
$$\mathrm{\Omega }_{\mu \nu }=_\mu \omega _\nu _\nu \omega _\mu ,\varphi _{\mu \nu }=_\mu \varphi _\nu _\nu \varphi _\mu $$
(11)
The potential function $`U(\sigma )`$ possesses a polynomial form
$$U(\sigma )=\frac{1}{2}m_\sigma ^2\sigma ^2+\frac{1}{3}g_3\sigma ^3+\frac{1}{4}g_4\sigma ^4.$$
(12)
The baryon mass is denoted by $`m_B`$ whereas $`m_M`$ $`(M=\sigma ,\omega ,\rho ,\sigma ^{},\varphi )`$ are masses assigned to the meson fields. The derived equations of motion constitute a set of coupled equations which have been solved in the mean field approximation. In this approximation meson fields are separated into classical mean field values and quantum fluctuations which are not included in the ground state
| $`\sigma =\overline{\sigma }`$ +$`s_0`$ | $`\sigma ^{}=\overline{\sigma }^{}`$ \+ $`s_0^{}`$ | |
| --- | --- | --- |
| $`\varphi _\mu =\overline{\varphi }_\mu +f_0\delta _{\mu 0}`$ | $`\omega _\mu =\overline{\omega }_\mu +w_0\delta _{\mu 0}`$ | $`\rho _\mu ^a=\overline{\rho }_\mu ^a+r_0\delta _{\mu 0}\delta ^{3a}`$ |
In the field equations the derivative terms are neglected and only time-like components of the vector mesons will survive if one assumes homogenous and isotropic infinite matter. The field equations derived from the Lagrange function at the mean field level are
$$m_\sigma ^2s_0+g_3s_0^2+g_4s_0^2=\underset{B}{}g_{\sigma B}m_{eff,B}^2S(m_{eff,B},k_{F,B})$$
(13)
$$m_\omega ^2w_0+c_3w_0^3=\underset{B}{}g_{\omega B}n_B$$
(14)
$$m_\rho ^2r_0=\underset{B}{}g_{\rho B}I_{3B}n_B$$
(15)
$$m_\sigma ^2s_0^{}=\underset{B}{}g_{\sigma ^{}B}m_{eff,B}^2S(m_{eff,B},k_{F,B})$$
(16)
$$m_\varphi ^2f_0=\underset{B}{}g_{\varphi B}n_B.$$
(17)
The function $`S(m_{eff,B},k_{F,B})`$ is expressed with the use of the integral
$$S(m_{eff,B},k_{F,B})=\frac{2J_B+1}{2\pi ^2}_0^{k_{F,B}}\frac{m_{eff,B}}{\sqrt{k^2+m_{eff,B}}}k^2๐k$$
(18)
where $`J_B`$ and $`I_{3B}`$ are the spin and isospin projections of baryon $`B`$, $`k_{F,B}`$ is the Fermi momentum of species $`B`$, $`n_B=\gamma _Bk_{F_B}^3/6\pi ^2`$ ($`\gamma _B`$ stands for the spin-isospin degeneracy factor which equals 4 for nucleons and $`\mathrm{\Xi }`$ hyperons and 2 for $`\mathrm{\Lambda }`$).
The obtained Dirac equation for baryons has the following form
$$(i\gamma ^\mu _\mu m_{eff,B}g_{\omega B}\gamma ^0\omega _0g_{\varphi B}\gamma ^0f_0)\psi _B=0$$
(19)
with $`m_{eff,B}`$ being the effective baryon mass generated by the baryon and scalar fields interaction and defined as
$$m_{B,eff}=m_B(g_{\sigma B}s_0+g_{\sigma ^{}B}s_0^{}).$$
(20)
The total energy of the system is given by
$`\epsilon `$ $`=`$ $`{\displaystyle \frac{1}{2}}m_\rho ^2r_0^2+{\displaystyle \frac{1}{2}}m_\varphi f_0^2+{\displaystyle \frac{1}{2}}m_\sigma ^{}(s_0^{})^2+{\displaystyle \frac{1}{2}}m_\omega ^2w_0^2`$
$`+`$ $`{\displaystyle \frac{3}{4}}c_3w_0^4+U(s_0)+ฯต_B`$
where $`\epsilon _B`$ equals
$$\epsilon _B=\underset{B}{}\frac{1}{3\pi ^2}_0^{k_{F,B}}k^2๐k\sqrt{k^2+m_{eff,B}^2}.$$
(22)
## 4 Parameters.
The considered model does not include $`\mathrm{\Sigma }`$ hyperons due to the remaining uncertainty of the form of their potential in nuclear matter at saturation density , , . The parameters that enter the Lagrangian function are collected in Tables 2 and 3. They are the standard TM1 parameter set supplemented by hyperon-meson coupling constants. In the scalar sector the scalar coupling of the $`\mathrm{\Lambda }`$ and $`\mathrm{\Xi }`$ hyperons requires constraining in order to reproduce the estimated values of the potential felt by a single $`\mathrm{\Lambda }`$ and a single $`\mathrm{\Xi }`$ in saturated nuclear matter
$`U_\mathrm{\Lambda }^{(N)}(\rho _0)`$ $`=`$ $`g_{\sigma \mathrm{\Lambda }}s_0(\rho _0)g_{\omega \mathrm{\Lambda }}w_0(\rho _0)2728MeV`$ (23)
$`U_\mathrm{\Xi }^{(N)}(\rho _0)`$ $`=`$ $`g_{\sigma \mathrm{\Xi }}s_0(\rho _0)g_{\omega \mathrm{\Xi }}w_0(\rho _0)1820MeV.`$
Assuming the SU(6) symmetry for the vector coupling constants and determining the scalar coupling constants from the potential depths, the hyperon-meson couplings can be fixed.
The strength of hyperon couplings to strange meson $`\sigma ^{}`$ is restricted through the following relation
$$U_\mathrm{\Xi }^{(\mathrm{\Xi })}U_\mathrm{\Lambda }^{(\mathrm{\Xi })}2U_\mathrm{\Xi }^{(\mathrm{\Lambda })}2U_\mathrm{\Lambda }^{(\mathrm{\Lambda })}.$$
(24)
which together with the estimated value of hyperon potential depths in hyperon matter provides effective constraints on scalar coupling constants to the $`\sigma ^{}`$ meson. The currently obtained value of the $`U_\mathrm{\Lambda }^{(\mathrm{\Lambda })}`$ potential at the level of 5 MeV permits the existence of additional parameter set which reproduces this weaker $`\mathrm{\Lambda }\mathrm{\Lambda }`$ interaction. In the text this parameter set is marked as weak, whereas strong denotes the stronger $`\mathrm{\Lambda }\mathrm{\Lambda }`$ interaction for $`U_\mathrm{\Lambda }^{(\mathrm{\Lambda })}20`$ MeV. The relation 24 indicates that for the weak $`\mathrm{\Lambda }\mathrm{\Lambda }`$ interaction the potential $`U_\mathrm{\Xi }^{(\mathrm{\Xi })}`$ takes the value $`10`$ MeV whereas for the strong interaction $`U_\mathrm{\Xi }^{(\mathrm{\Xi })}40`$ MeV.
The vector coupling constants for hyperons are determined from $`SU(6)`$ symmetry as
$`{\displaystyle \frac{1}{2}}g_{\omega \mathrm{\Lambda }}={\displaystyle \frac{1}{2}}g_{\omega \mathrm{\Sigma }}=g_{\omega \mathrm{\Xi }}={\displaystyle \frac{1}{3}}g_{\omega N}`$ (25)
$`{\displaystyle \frac{1}{2}}g_{\rho \mathrm{\Sigma }}=g_{\rho \mathrm{\Xi }}=g_{\rho N};g_{\rho \mathrm{\Lambda }}=0`$
$`2g_{\varphi \mathrm{\Lambda }}=2g_{\varphi \mathrm{\Sigma }}=g_{\varphi \mathrm{\Xi }}={\displaystyle \frac{2\sqrt{2}}{3}}g_{\omega N}.`$
## 5 Infinite symmetric strange hadronic matter.
For the symmetric strange hadronic matter there no isospin dependance, and there is no contribution coming from the $`\rho `$ meson field. There are only two conserved charges: the baryon number $`n_b=n_\mathrm{\Lambda }+n_N+n_\mathrm{\Xi }`$ and the strangeness number $`n_s=n_\mathrm{\Lambda }+2n_\mathrm{\Xi }`$. These two conserved charges allow one to define the parameter which specifies the strangeness contents in the system and is strictly connected to the appearance of particular hyperon species in the model
$$fs=\frac{n_s}{n_b}=\frac{n_\mathrm{\Lambda }+2n_\mathrm{\Xi }}{n_b}.$$
(26)
In a multi-strange system, for sufficient number density of $`\mathrm{\Lambda }`$ hyperons, the process $`\mathrm{\Lambda }+\mathrm{\Lambda }\mathrm{\Xi }+N`$, where $`N`$ stands for nucleon, becomes energetically allowed. Thus, beside $`N`$ and $`\mathrm{\Lambda }`$ also $`\mathrm{\Xi }^{}`$ and $`\mathrm{\Xi }^0`$ hyperons have to contribute to the composition of strange hadronic matter. In general the chemical equilibrium conditions for the processes $`\mathrm{\Lambda }+\mathrm{\Lambda }n+\mathrm{\Xi }^0`$ and $`\mathrm{\Lambda }+\mathrm{\Lambda }p+\mathrm{\Xi }^{}`$ are established by the following relations between chemical potentials
$$2\mu _\mathrm{\Lambda }=\mu _n+\mu _{\mathrm{\Xi }^0}2\mu _\mathrm{\Lambda }=\mu _p+\mu _\mathrm{\Xi }^{}.$$
(27)
This relation for symmetric matter can be rewritten as $`2\mu _\mathrm{\Lambda }=\mu _N+\mu _\mathrm{\Xi }`$. The chemical potential $`\mu _B`$ ($`B=N,\mathrm{\Lambda },\mathrm{\Xi }`$) through the Hugenholtz-van-Hove theorem is related to the Fermi energy of each baryon in the following way
$$\mu _B=\sqrt{m_{B,eff}^2+k_{F,B}^2}+g_{\omega B}w_0+g_{\varphi B}f_0.$$
(28)
The binding energy of the system can be obtained from the following relation
$$Eb=\frac{1}{n_b}(\epsilon Y_Nm_NY_\mathrm{\Lambda }m_\mathrm{\Lambda }Y_\mathrm{\Xi }m_\mathrm{\Xi })$$
(29)
where $`Y_B(B=N,\mathrm{\Lambda },\mathrm{\Xi }`$) denotes concentrations of particular baryons.
## 6 Results
The density dependance of binding energies of multi-strange system involving nucleons, $`\mathrm{\Lambda }`$ and $`\mathrm{\Xi }`$ hyperons for different parameter sets are presented in Fig.1 and Fig.2. For each parameterization two separated cases are considered, namely the strong and weak $`YY`$ interactions. Individual lines represent binding energies obtained for different values of strangeness fraction $`fs`$. In all cases the value $`fs=0`$ corresponds to the state when only nucleons are present in the system and the equation of state characteristic to nuclear matter is reproduced.
The equilibrium density $`n_{b_0}`$ can be obtained by minimizing the binding energy with respect to the baryon number density $`n_B`$. Increasing the value of the parameter $`f_s`$, what is equivalent with the increasing value of the strangeness contents in the matter, the binding energy for each fixed value of the parameter $`fs`$ has been calculated. The minimum values of binding energies for the fixed values of $`fs`$ have been determined and the results are plotted in Fig.3.
For both parameterizations (TM1 and TMA), in the case of strong $`YY`$ interaction the increasing value of $`fs`$ leads to more bound system with the minimum shifted towards higher densities. Contrary to this situation the increasing value of strangeness contents in hyperon matter characterized by weak $`YY`$ interaction gives shallower minima in the result. In Fig 4 the equilibrium density $`n_{b_0}`$ as a function of strangeness contents for the two cases mentioned above is presented. The results are depicted for the TM1 parametrization.
Relative concentrations of $`\mathrm{\Lambda }`$ and $`\mathrm{\Xi }`$ hyperons in the case of symmetric strange hadronic matter are presented in Fig.5. In both cases the onset of $`\mathrm{\Lambda }`$ hyperon is followed by the onset of $`\mathrm{\Xi }`$ hyperon. The populations of $`\mathrm{\Lambda }`$ hyperons obtained in the TM1 and TMA-weak models are reduced in comparison with those calculated for the strong $`YY`$ interaction models. Concentrations of $`\mathrm{\Xi }`$ hyperons are higher for the weak TM1 and TMA models. In the case of weak $`YY`$ interaction model the $`\mathrm{\Xi }`$ hyperon thresholds are shifted towards lower strangeness fractions. This has an influence on the properties of neutron star matter.
## 7 Hyperon star matter.
Neutron star interiors and relativistic heavy ion collisions offer suitable environment for the existence of multi-strange hyperon system. In the case of relativistic heavy ion collisions hot and dense nuclear matter is probed whereas neutron star interiors represent the density dominated scenario ($`T0`$). The analysis of the role of strangeness in nuclear structure in the aspect of multi-strange system is of great importance for neutron star matter. An imperfect knowledge of the neutron star matter equation of state, especially in the presence of hyperons, causes many uncertainties in determining neutron star structure. It seems to be very important to estimate the influence of nucleon-hyperon and hyperon-hyperon interaction on the equation of state. The comparison between weak interaction time scales ($`10^{10}`$ s) and a time scale connected with the lifetime of a relevant star indicates that there is a difference between the neutron star matter constrained by charge neutrality and generalized $`\beta `$ equilibrium and the matter in high energy collisions; the latter matter is constrained by isospin symmetry and strangeness conservation. The condition of $`\beta `$ equilibrium in the case of neutron star matter implies the presence of leptons and is realized by adding electrons and muons to the baryonic matter. The Lagrangian of free leptons has the following form (30)
$$L_l=\underset{L=e,\mu }{}\overline{\psi }_L(i\gamma ^\mu _\mu m_L)\psi _L.$$
(30)
In general, in neutron star matter muons start to appear after $`\mu _e`$ has reached the value equal to the muon mass. The appearance of muons not only reduces the number of electrons but also affects the proton fraction. The requirements of charge neutrality
$$n_p=n_e+n_\mu +n_\mathrm{\Xi }^{}$$
(31)
and equilibrium under the weak processes
$$B_1B_2+LB_2+LB_1$$
(32)
($`B_1`$ and $`B_2`$ denote baryons) is decisive in determining the composition of the hyperon star matter. The equilibrium conditions between baryonic and leptonic species, which are present in hyperon star matter, lead to the following relations between their chemical potentials
$`\mu _p=\mu _n\mu _e\mu _\mathrm{\Lambda }=\mu _{\mathrm{\Xi }^0}=\mu _n`$ (33)
$`\mu _\mathrm{\Xi }^{}=\mu _n+\mu _e\mu _\mu =\mu _e.`$
The relations which determine the chemical potentials (28) and effective baryon masses (20) indicate that the composition of hyperon star matter is altered when the strength of the hyperon-hyperon interaction is changed. Fig.6 presents fractions of particular strange baryon species $`Y_B`$ as a function of baryon number density $`n_B`$ for TM1 and TMA parameterizations. Starting the analysis of these graphs from moderate densities it is evident that $`\mathrm{\Lambda }`$ hyperons are the most abundant strange baryons. $`\mathrm{\Lambda }`$ is also the first strange baryon that emerges in hyperon star matter, it is followed by $`\mathrm{\Xi }^{}`$ and $`\mathrm{\Xi }^0`$ hyperons. For the weak and strong $`YY`$ interaction models the sequence of appearance of hyperons is the same however, in the case of the weak $`YY`$ interaction the threshold for the $`\mathrm{\Xi }^{}`$ and $`\mathrm{\Xi }^0`$ hyperons are shifted towards lower densities. Due to the repulsive potential of $`\mathrm{\Sigma }`$ hyperons their onset points are possible at very high densities which are not relevant for neutron stars.
In the case of weak $`YY`$ interaction the population of $`\mathrm{\Lambda }`$ hyperons is reduced whereas relative concentrations of $`\mathrm{\Xi }^{}`$ and $`\mathrm{\Xi }^0`$ hyperons are enhanced.
The populations of protons and leptons are also altered by the change of the strength of hyperon coupling constants. Through the requirement of charge neutrality and $`\beta `$-equilibrium condition the onset points and concentrations of hyperons affect the negatively charged lepton and proton abundance. Populations of protons for TM1 and TMA parameterizations increase for the weak $`YY`$ interaction models. Results are presented in Fig.7.
The hyperonization of matter in the presented models can be also analyzed through the density dependence of the strangeness fraction. The density dependance of the $`fs`$ parameter in the hyperon star matter is depicted in Fig.8.
For the TMA parameter set more hyperon rich matter is obtained for the strong $`YY`$ interaction model, the TM1 parameter set gives the same result for very high densities, for moderate densities the strangeness contents is higher for the weak model.
The equation of state of the system can be calculated from the stress-energy density tensor $`T_{\mu \nu }`$ which is defined as
$$T_{\mu \nu }=\underset{a}{}_\nu \mathrm{\Phi }^a(x)\frac{L(x)}{(^\mu \mathrm{\Phi }^a(x))}g_{\mu \nu }L(x)$$
(34)
where $`\mathrm{\Phi }^a=(\varphi _M,\psi _B,\psi _L)`$ and $`\varphi _M=(\sigma ,\omega ,\rho ,\sigma ^{},\varphi )`$. The equation of state also has been constructed for the strong and weak $`YY`$ interactions, for the chosen TM1 and TMA parameterizations. The results are presented in Fig.9 where besides the equations of state obtained for the parameterizations discussed in this paper two other selected equations of state are shown. The first represents the neutron star model (with zero strangeness) and the second the hyperon star model with hyperon couplings derived from the quark model .
Generally the equation of state obtained for the weak $`YY`$ interaction model is less stiff that for the strong model. However, there is an energy density range for which the weak model gives the stiffer equation of state than the strong $`YY`$ interaction model. The obtained form of the equation of state is influenced by considerably altered effective baryon masses. The results are presented in Fig. 10.
On specifying the equation of state the properties and structure of neutron stars can be obtained from hydrostatic equilibrium equations of Tolman, Oppenheimer and Volkoff. The application of the composite equation of state constructed by adding Bonn and Negele-Vauterin equations of state allows us to calculate the neutron star structure for the entire neutron star density span. At higher densities the equation of state depends on the nature of strong interactions hence the strength of hyperon-hyperon interactions should exert influence on neutron star parameters. The calculated mass-radius relations for the determined forms of the equations of state are presented in Fig.11.
In Table 4 the parameters of the maximum mass configurations are collected.
From Fig.6 it is evident that $`\mathrm{\Xi }^{}`$ and $`\mathrm{\Xi }^0`$ hyperons emerge at very high densities. Although their onset points in the case of weak $`YY`$ interaction models are shifted towards lower densities, still their presence in stable neutron star configuration is uncertain. In Table 5 the star parameters for the appearance of $`\mathrm{\Xi }`$ hyperons are collected.
The value of densities, masses and radii collected in the Table 5 indicate that $`\mathrm{\Xi }`$ hyperons do not appear in stable hyperon star configurations for both parameter sets. Thus only $`\mathrm{\Lambda }`$ hyperon will be present in the composition of hyperon star matter for the considered models. Fig.11 shows the position of the configurations, which parameters are presented in Table 5, on the mass-radius relations.
Neutron stars are purely gravitationally bound compact stars. The gravitational binding energy of a relativistic star is defined as a difference between its gravitational and baryon mass.
$$Eb_g=(M_pm(R))c^2$$
(35)
where
$$M_p=4\pi _0^R๐rr^2(1\frac{2Gm(r)}{c^2r})^{\frac{1}{2}}\rho (r)$$
(36)
Of considerable relevance is the numerical solution of the above equation for the selected equations of state. Fig. 12 depicts the gravitational mass which includes interactions versus the baryonic mass for all the considered models.
## 8 Conclusions.
The main goal of this paper was to study the influence of the strength of hyperon-hyperon interactions on the properties of the hyperon star matter and on a hyperon star structure. It has been shown that replacing the strong $`YY`$ interaction model by the weak one introduces large differences in the composition of a hyperon star matter both in the strange and non-strange sectors. There is a considerable reduction of $`\mathrm{\Lambda }`$ hyperon concentration whereas the concentrations of $`\mathrm{\Xi }^{}`$ and $`\mathrm{\Xi }^0`$ hyperons are enhanced. In the non-strange sector the populations of protons and electrons are changed. The weak model permits larger fractions of protons and electrons. The presence of hyperons in general leads to the softening of the equation of state. For the employed weak $`YY`$ interaction model there is a density range for which the obtained equation of state is stiffer than the one calculated with the use of the strong model. This is clearly visible in the case of TM1 parameter set. For higher densities the weak model gives less stiff equation of state. The behavior of the equation of state is directly connected with the value of the maximum star mass. Equilibrium conditions namely charge neutrality and $`\beta `$-equilibrium determine the composition of the star. For both parameterizations the onset points for $`\mathrm{\Xi }^{}`$ and $`\mathrm{\Xi }^0`$ hyperons are localized at high densities which are relevant to unstable branches of the mass-radius relations. Thus the obtained stable hyperon star model is composed of neutrons, protons, $`\mathrm{\Lambda }`$ hyperons and leptons. The appearance of $`\mathrm{\Xi }^{}`$ and $`\mathrm{\Xi }^0`$ hyperons in neutron star interior will be possible in a very special configuration. A protoneutron star model with very high central density can lead to a hyperon star with $`\mathrm{\Xi }`$ and $`\mathrm{\Lambda }`$ hyperons in its interior. The possibility of the existence of such protoneutron star model will be the subject of future investigations. One can compare the obtained results with those presented in the paper by Schaffner et al. , where the analysis of the existence of the third family of stable compact stars has been performed for highly attractive hyperon-hyperon interaction. However, recent experimental data indicate for much weaker strength of $`YY`$ interaction. Employing these data and the estimated value of the $`\mathrm{\Lambda }`$ well depth $`U_{\mathrm{\Lambda }\mathrm{\Lambda }}5`$ MeV the result of the paper can not be confirmed.
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# Static quark anti-quark interactions at zero and finite temperature QCD. II. Quark anti-quark internal energy and entropy
## I Introduction
This is the second part of our discussion of thermal modifications of the strong forces in finite temperature QCD Kaczmarek:2005ui (for a detailed introduction to this subject and further references see Kaczmarek:2005ui ; Karsch:2005ex ). At finite temperature, $`T0`$, the free energy of a static quark anti-quark pair McLerran:1981pb ; McLerran:1980pk , separated by distance $`r`$, is an important tool to analyze the in-medium modification of the QCD forces. Similar to the free energies also the internal energies have recently been introduced Kaczmarek:2002mc ; Zantow:2003ui and are expected to play an important role in the discussion of quarkonia binding properties Matsui:1986dk ; Brown:2004qi ; Wong:2004kn ; Shuryak:2004tx ; Park:2005nv ; Brambilla:2004wf ; Digal:2001iu ; Digal:2001ue ; Wong:2001kn ; Wong:2001uu ; Shuryak:2003ty . Moreover, the structure of these energies at large distances and high temperatures is important for our understanding of the bulk properties of the QCD plasma phase, e.g. the screening property of the quark gluon plasma Kaczmarek:1999mm ; Kaczmarek:2004gv , the equation of state Beinlich:1997ia ; Karsch:2000ps . They also provide important input to the construction of effective models based on properties of the Polyakov loop Pisarski:2002ji ; Kaczmarek:2002mc ; Kaczmarek:2005ui ; Kaczmarek:2003ph ; Dumitru:2003hp ; Dumitru:2004gd . Up to quite recently Karsch:2000ps ; Kaczmarek:2005ui ; Kaczmarek:2003ph ; Petreczky:2004pz ; Kaczmarek:2005uw ; Kaczmarek:2005uv most of these discussions concerned the quark anti-quark free energies in quenched QCD. Several qualitative differences, however, are to be expected when changing from free to internal energies and/or when taking into account the influence of dynamical fermions. The difference between free and internal energy arises from non-trivial entropy contributions Kaczmarek:2002mc ; Zantow:2003ui . Moreover, in QCD with light dynamic quarks the large distance behavior of the strong interaction will show a qualitative different behavior due to the possibility of string breaking.
Properties of the finite temperature quark anti-quark free energies and the heavy quark potential at $`T=0`$, $`V(r)`$, have been discussed for $`2`$-flavor QCD in Refs. Karsch:2000ps ; Kaczmarek:2005ui ; Kaczmarek:2005uw ; Kaczmarek:2005uv ; Kaczmarek:2003ph . Recently some results have also been reported for $`3`$-flavor QCD Petreczky:2004pz . Here we will continue our analysis of the fundamental forces of QCD at finite temperature. We analyze the partition function of $`2`$-flavor QCD in the presence of heavy quarks and extract the quark anti-quark internal energies and entropies. This paper is organized as follows: In section II we introduce the change in internal energy and entropy due to the presence of static quarks and anti-quarks in a thermal heat bath. We discuss their temperature dependence at large distances, in particular, their behavior in the vicinity of the transition, in section III. We finally discuss the qualitative and quantitative differences between free and internal energies in section IV and discuss their binding properties with respect to quarkonium binding in the vicinity of the transition. Section V contains our conclusions. Details on our simulation parameters, the lattice actions used in our calculations as well as details on the analysis of the quark anti-quark free energies are given in Ref. Kaczmarek:2005ui .
## II Partition function in the presence of heavy quarks
### II.1 Free energy, internal energy and entropy
As we are interested in the lattice formulation of QCD at finite temperature in thermal equilibrium, we consider the (Euclidean) path integral, i.e. we consider the partition function of the QCD heat bath,
$`๐ต(T,V)`$ $``$ $`{\displaystyle ๐A๐\mathrm{\Psi }๐\overline{\mathrm{\Psi }}e^{S[A,\mathrm{\Psi },\overline{\mathrm{\Psi }}]}}=e^{F(T,V)/T},`$ (1)
where $`T`$ ($`V`$) denotes the temperature (volume) and $`S[A,\mathrm{\Psi },\overline{\mathrm{\Psi }}]`$ the QCD action. We also investigate the corresponding system containing additional heavy quarks McLerran:1981pb , i.e.
$`๐ต_๐ช(r,T,V)`$ $``$ $`{\displaystyle ๐A๐\mathrm{\Psi }๐\overline{\mathrm{\Psi }}๐ช_r^{(c)}[W,W^{}]e^{S[A,\mathrm{\Psi },\overline{\mathrm{\Psi }}]}}`$ (2)
$`=`$ $`e^{\stackrel{~}{F}_๐ช(r,T,V)/T},`$
where $`๐ช_r^{(c)}[W,W^{}]`$ denotes an operator which introduces static color sources representing quarks and anti-quarks separated by distances $`r\{r_i\}`$ in a specific color representation $`c`$. A static color source appearing in $`๐ช_r^{(c)}[W,W^{}]`$ located at $`๐ฑ`$ is described by the thermal Wilson line,
$`W(๐ฑ)`$ $`=`$ $`๐ฏ\mathrm{exp}\left(i{\displaystyle _0^{1/T}}๐x_0\lambda ๐_0(x_0,๐ฑ)\right).`$ (3)
The expectation value of $`๐ช_r^{(c)}[W,W^{}]`$, i.e. the $`n`$-point Polyakov loop correlation function $`๐ช_r^{(c)}[W,W^{}]`$, is related to the change in free energy, $`F_๐ช(r,T)\stackrel{~}{F}(r,T,V)F(T,V)`$, due to the presence of static quark anti-quark sources in the heat bath,
$`F_๐ช(r,T)`$ $`=`$ $`T\mathrm{ln}๐ช_r^{(c)}[W,W^{}]+TC_๐ช`$
$`=`$ $`T\left(\mathrm{ln}๐ต_๐ช(r,T,V)\mathrm{ln}๐ต(T,V)\right)+TC_๐ช,`$
where $`C_๐ช`$ can be fixed through renormalization Note1 . For instance, in the case of $`2`$-point Polyakov loop correlation functions the singlet ($`1`$), averaged ($`\overline{q}q`$) and octet ($`8`$) color representations of the operator $`๐ช_r^{(c)}[W,W^{}]`$ can be specified as McLerran:1981pb ; Nadkarni:1986as ; Philipsen:2002az
$`๐ช_r^{(1)}[W,W^{}]`$ $`=`$ $`{\displaystyle \frac{1}{3}}\mathrm{Tr}W(0)W^{}(|r|),`$ (5)
$`๐ช_r^{(\overline{q}q)}[W,W^{}]`$ $`=`$ $`{\displaystyle \frac{1}{9}}\mathrm{Tr}W(0)\mathrm{Tr}W^{}(|r|),`$ (6)
$`๐ช_r^{(8)}[W,W^{}]`$ $`=`$ $`{\displaystyle \frac{1}{8}}\mathrm{Tr}W(0)\mathrm{Tr}W^{}(|r|)`$ (7)
$`{\displaystyle \frac{1}{24}}\mathrm{Tr}W(0)W^{}(|r|).`$
The color singlet, averaged and octet quark anti-quark free energies, i.e. $`F_1(r,T)`$, $`F_{\overline{q}q}(r,T)`$ and $`F_8(r,T)`$, respectively, have already been discussed extensively in quenched and full QCD Attig:1988ey ; Kaczmarek:2002mc ; Kaczmarek:2004gv ; Kaczmarek:2005ui ; Kaczmarek:2003ph ; Kaczmarek:2003dp ; Petreczky:2004pz ; Kaczmarek:2005uw ; Kaczmarek:2005uv ; Nakamura:2004wr .
For the purpose of discussing internal energies ($`U_๐ช(r,T)`$) and entropies ($`S_๐ช(r,T)`$), we follow the conceptual approach suggested in Refs. Kaczmarek:2002mc ; Zantow:2003ui and consider thermal derivatives of the QCD partition functions introduced above, i.e.
$`U_๐ช(r,T)`$ $`=`$ $`T^2{\displaystyle \frac{F_๐ช(r,T)/T}{T}},`$ (8)
which leads to
$`U_๐ช(r,T)`$ $`=`$ $`T^2\left({\displaystyle \frac{1}{๐ช_r^{(c)}[W,W^{}]}}๐ช_r^{(c)}[W,W^{}]{\displaystyle \frac{S[A,\mathrm{\Psi },\overline{\mathrm{\Psi }}]}{T}}{\displaystyle \frac{S[A,\mathrm{\Psi },\overline{\mathrm{\Psi }}]}{T}}\right)`$ (9)
$``$ $`\stackrel{~}{U}_๐ช(r,T,V)U(T,V).`$
Note here that the derivative of the operator $`๐ช_r^{(c)}`$ with respect to temperature, $`๐ช_r^{(c)}/T`$, vanishes due to (3). A similar relation can be derived for the entropies starting from
$`S_๐ช(r,T)`$ $`=`$ $`{\displaystyle \frac{F_๐ช(r,T)}{T}}`$ (10)
$``$ $`\stackrel{~}{S}_๐ช(r,T,V)S(T,V){\displaystyle \frac{TC_๐ช}{T}},`$
i.e. the observable $`TS_๐ช(r,T)`$ could be calculated from the difference, $`TS_๐ช(r,T)=U_๐ช(r,T)F_๐ช(r,T)`$, with $`U_๐ช(r,T)`$ and $`F_๐ช(r,T)`$ given in (9) and (II.1). We have also specified constant contributions, $`C_๐ช`$, which result from divergent contributions to the free energies and will control the internal energies and entropies at large quark anti-quark separations. Once the free energies are fixed through renormalization also the constant contributions to the internal energies and entropies are properly determined.
We note that Eq. (9), and similar (10), open the possibility for a direct calculations of $`U_๐ช(r,T)`$ and $`S_๐ช(r,T)`$ and define properly the quantities we aim to discuss here, i.e. the change in internal energies and entropies due to the presence of static quarks and anti-quarks in the QCD heat bath. Quite similar to the change in free energies, $`F_๐ช(r,T)=\stackrel{~}{F}_๐ช(r,T)F(T)`$, also the changes in internal energies and entropies, $`U_๐ช(r,T)=\stackrel{~}{U}_๐ช(r,T,V)U(T,V)`$ and $`S_๐ช(r,T)=\stackrel{~}{S}_๐ช(r,T,V)S(T,V)`$, are expected to behave like intensive observables and, in particular, will show no volume dependence in the thermodynamic limit. It, however, can no longer be assumed that the $`r`$-dependence of the quark anti-quark free energies (II.1) are given by the $`r`$-dependences of the internal energies (9) alone, i.e. we expect
$`F_๐ช(r,T)`$ $`=`$ $`U_๐ช(r,T)TS_๐ช(r,T).`$ (11)
This indicates a quite complicated relation between free energies, internal energies and entropies and stresses the important role of the entropy contribution which is still present in free energies. In the limit of vanishing temperature, $`T0`$, however, $`TS_๐ช(r,T0)`$ will vanish and
$`F_๐ช(r,T0)=U_๐ช(r,T0)V_๐ช(r),`$ (12)
i.e. $`F_๐ช(r,T0)`$ and $`U_๐ช(r,T0)`$ could be used<sup>1</sup><sup>1</sup>1We assume here the existence of an energy, $`V_๐ช(r)`$, at $`T=0`$ which corresponds to the expectation value of $`lim_{T0}๐ช_r^{(c)}`$. to define the corresponding energies at $`T=0`$.
Unfortunately, the non-perturbative formulation of energies given in (9) is still complicated and, in particular, it is complicated to be given in a form suitable for lattice simulations Karsch:2000ps ; Boyd:1996bx ; Boyd:1995zg ; Montvay:1994cy . In the thermodynamic limit, however, the internal energy and entropy described above can be calculated equally well from Eqs. (8, 10). We indeed have used these relations and performed the derivatives with respect to temperature based on finite difference approximations using the non-perturbative free energies at neighboring temperature as input (see Tab. 1 of Ref. Kaczmarek:2005ui ). With this method no perturbative uncertainties get introduced in our calculations and a test of confinement and deconfinement could be given. In fact, as we use in our calculations the renormalized free energy as input, both quantities, the internal energies and entropies, survive the continuum limit. We stress here, however, that we also provide non-perturbative renormalization prescriptions for the quark anti-quark internal energies and entropies which are independent from a renormalization of the quark anti-quark free energies.
### II.2 Theoretic expectations and renormalization
Following Kaczmarek:2005ui we consider mainly the quark anti-quark internal energies and entropies in the color singlet channel. At large distances and high temperatures, i.e. $`rT1`$, $`T`$ sufficiently above $`T_c`$, as well as at small distances and zero temperature, i.e. $`r\mathrm{\Lambda }_{\text{QCD}}1`$, the singlet free energies, $`F_1(r,T)`$, are indeed dominated by one gluon exchange Kaczmarek:2005ui . Using the perturbative small distance relation for the singlet free energy Brown:1979ya ; Nadkarni:1986as ; Nadkarni:1986cz ; McLerran:1981pb ,
$`F_1(r,T)`$ $``$ $`{\displaystyle \frac{4}{3}}{\displaystyle \frac{\alpha (r)}{r}},r\mathrm{\Lambda }_{\text{QCD}}1,`$ (13)
standard thermodynamic relations (Eqs. (8) and (10)) lead to
$`U_1(r1/\mathrm{\Lambda }_{\text{QCD}},T)`$ $``$ $`{\displaystyle \frac{4}{3}}{\displaystyle \frac{\alpha (r)}{r}},`$ (14)
and
$`S_1(r1/\mathrm{\Lambda }_{\text{QCD}},T)`$ $``$ $`0.`$ (15)
We thus expect that at small distances the singlet free and internal energies are controlled to a large extent by energy, i.e. in the limit of small distances both will smoothly approach the zero temperature heavy quark potential, $`V(r)`$. At larger distances, however, the quark anti-quark free energies are strongly temperature dependent Kaczmarek:2005ui ; Kaczmarek:2005uv ; Kaczmarek:2005uw and thus non-vanishing entropy contributions arise. In this case differences between free and internal energies are expected to become important and the quark anti-quark free energy will be to a large extent controlled by $`TS_1(r,T)`$.
To clarify the role of the entropy we compare<sup>2</sup><sup>2</sup>2Details on our lattice simulations, in particular, on the calculation of $`F_1(r,T)`$, are given in Ref. Kaczmarek:2005ui . Details on the computation of $`U_1(r,T)`$ and $`S_1(r,T)`$ will be given in Sec. III and IV. in Fig. 1(a) the short and large distance parts of the singlet free and internal energies at $`T1.3T_c`$ and show in Fig. 1(b) the corresponding entropy contribution, $`TS_1(r,T1.3T_c)`$. We also indicate in both figures the small distance behavior expected from Eqs. (14) and (15) as solid lines, i.e. the line in Fig. 1(a) indicates the heavy quark potential from Refs. Kaczmarek:2005ui ; Kaczmarek:2005uw , and in Fig. 1(b) it indicates the zero level. As both, the internal energy and entropy, have been calculated using renormalized free energies proper renormalization of both observables is already incorporated by construction. It can clearly be deduced from Fig. 1(a) that the singlet free and internal energies smoothly approach $`V(r)`$ at small distances and the free energy thus indeed is dominated by the energy contribution. In fact, the entropy contribution shown in Fig. 1(b) is quite small at small distances and indicates a vanishing entropy contribution in the limit $`r0`$. Unfortunately we could not go to smaller distances to clearly demonstrate this behavior. Moreover, at small distances $`U_1(r,T)`$ and $`TS_1(r,T)`$ suffer from lattice artifacts which result from small distance lattice artifacts in the free energies. At intermediate and large distances, however, the free and internal energies shown in Fig. 1(a) deviate from each other and $`TS_1(r,T)`$ indeed plays an important role for the behavior of $`F_1(r,T)`$. At asymptotic large distances the internal energies and entropies approach temperature dependent constant values, i.e. $`U_{\mathrm{}}(T)lim_r\mathrm{}U_1(r,T)`$ and $`S_{\mathrm{}}(T)lim_r\mathrm{}S_1(r,T)`$ are finite for finite temperatures. These values are indicated by the arrows in Fig. 1 and, in particular, at finite temperature
$`U_{\mathrm{}}(T)`$ $`>`$$``$ $`F_{\mathrm{}}(T)`$ (16)
is evident.
A similar behavior of $`U_{\mathrm{}}(T)`$ and $`S_{\mathrm{}}(T)`$ can be deduced from high temperature perturbation theory. To be more precise, high temperature perturbation theory suggests a cubic leading order dependence<sup>3</sup><sup>3</sup>3Here and in what follows we already have anticipated the running of the coupling with the expected dominant scale $`T`$. Of course, the running of the coupling appears only beyond leading order. of the free energy on the coupling Gava:1981qd , i.e.
$`F_{\mathrm{}}(T){\displaystyle \frac{4}{3}}m_D(T)\alpha (T)๐ช(g^3T).`$ (17)
Here $`\alpha (T)=g^2(T)/4\pi `$ and $`m_D(T)`$ denotes the Debye mass which in re-summed leading order is given by
$`m_D(T)`$ $`=`$ $`\left(1+{\displaystyle \frac{N_f}{6}}\right)^{1/2}g(T)T.`$ (18)
We note that this leading order result is gauge invariant. In the following we use the renormalization group $`\beta `$-function to evaluate the derivatives of the coupling, i.e. for an arbitrary function $`(g,T)`$ we use
$`T{\displaystyle \frac{d(g,T)}{dT}}`$ $`=`$ $`T{\displaystyle \frac{(g,T)}{T}}+\beta (g){\displaystyle \frac{d(g,T)}{dg}},`$ (19)
where $`\beta (g)=\beta _0g^3+๐ช(g^5)`$ in perturbation theory. Assuming this behavior the internal energy, $`U_{\mathrm{}}(T)`$, and entropy, $`S_{\mathrm{}}(T)`$, are expected to behave like
$`U_{\mathrm{}}(T)`$ $``$ $`4m_D(T)\alpha (T){\displaystyle \frac{\beta (g)}{g(T)}}๐ช(Tg^5),`$ (20)
and
$`S_{\mathrm{}}(T)`$ $``$ $`+{\displaystyle \frac{4}{3}}{\displaystyle \frac{m_D(T)}{T}}\alpha (T)+4{\displaystyle \frac{m_D(T)}{T}}\alpha (T){\displaystyle \frac{\beta (g)}{g(T)}}`$ (21)
$``$ $`+๐ช(g^3).`$
The leading contribution to $`TS_{\mathrm{}}(T)`$ is similar to the free energy in Eq. (17) and thus at high temperatures and large distances the free energy is indeed expected to be to large extent dominated by the entropy contribution, i.e. at leading order $`S_{\mathrm{}}(T)F_{\mathrm{}}(T)/T`$. Although the entropy itself will vanish logarithmically in the high temperature limit, i.e. $`S_{\mathrm{}}(T\mathrm{})=0`$, the contribution $`TS_{\mathrm{}}(T)`$ will clearly dominate the differences between free and internal energy at high temperatures,
$`U_{\mathrm{}}(T)F_{\mathrm{}}(T)`$ $`=`$ $`TS_{\mathrm{}}(T)+๐ช(g^3T).`$ (22)
Thus, the difference between free and internal energy is expected to increase continuously with increasing temperature when approaching the perturbative high temperature regime. Only in the limit of zero temperature, $`T0`$, the observable $`TS_{\mathrm{}}(T)`$ will vanish as $`S_{\mathrm{}}(T)`$ is a dimension less quantity which due to string breaking stays finite in QCD. Any qualitative change in the observable $`TS_{\mathrm{}}(T)`$ as function of temperature between both limits, i.e. $`T0`$ and $`T\mathrm{}`$, is not quite obvious and, if present, may signal the phase change from the chiral symmetry broken phase at low to the deconfinement phase at high temperatures.
We may finally note that although we discuss in the following the internal energies and entropies calculated from renormalized free energies, it is conceptually quite satisfying that both observables could equally well be renormalized by matching their short distance parts to the heavy quark potential (14) and zero (15), respectively. This is indeed evident from Figs. 1(a, b). As no additional divergences get introduced at finite temperature also their large distance properties are properly fixed in the continuum limit, in particular, also the manifestly gauge invariant observables $`U_{\mathrm{}}(T)`$ and $`S_{\mathrm{}}(T)`$.
## III The confinement deconfinement transition
We begin our discussion of the finite temperature energies and entropies at asymptotic large distances, $`r\mathrm{}`$. Actually, to avoid any fit we again Kaczmarek:2005ui approximated the value of the free energy at infinite distance, $`F_1(r\mathrm{},T)`$, by the value of the quark anti-quark free energy calculated at the largest possible separation on the lattice, i.e. $`F_{\mathrm{}}(T)F_{\overline{q}q}(N_\sigma /2,T)`$, and calculated separately the internal energy, $`U_{\mathrm{}}(T)`$, and entropy, $`S_{\mathrm{}}(T)`$. Both quantities are obtained from derivatives of the color averaged free energy, $`F_{\overline{q}q}(r,T)`$, which is a manifestly gauge invariant observable. Our results for $`U_{\mathrm{}}(T)`$ and $`S_{\mathrm{}}(T)`$ are summarized in Tab. 1. It is quite satisfying that the values obtained for $`U_{\mathrm{}}(T)`$ and $`S_{\mathrm{}}(T)`$ reproduce the free energy given in Tab. 2 of Ref. Kaczmarek:2005ui , i.e. the quantity $`U_{\mathrm{}}TS_{\mathrm{}}`$ matches to the value for $`F_{\mathrm{}}=2T\mathrm{ln}|L|`$ for all temperatures.
In parts of our analysis we are also interested in the flavor and quark mass dependence of the finite temperature energies and entropies. We thus again Kaczmarek:2005ui compare our results ($`N_f=2`$) to results from quenched QCD ($`N_f=0`$) Kaczmarek:2002mc ; Phd and also in parts to a recent study of $`3`$-flavor QCD Petreczky:2004pz . To convert the observables to physical units we use $`T_c=270`$ MeV in quenched, $`T_c=200`$ MeV in $`2`$-flavor ($`m_\pi /m_\rho 0.7`$) and $`T_c=193`$ MeV in $`3`$-flavor ($`m_\pi /m_\rho 0.4`$) QCD. It should be obvious, however, that a comparison of free and internal energies crucially depends on the relative normalization of the corresponding zero temperature heavy quark potentials used for renormalization and thus a comparison could be affected by flavor and/or quark mass dependent (over-all) constant contributions. Here, and in what follows, the relative normalization of the heavy quark potentials in quenched and full QCD is such that there is no constant contribution in the Cornell Ansatz for $`V(r)`$ at large distances. We also note that any undetermined constant contribution to the heavy quark potential at zero temperature will add a non-perturbative over-all constant to the free and internal energies which would also affect the comparison of these observables with perturbation theory Zantow:2003uh . We stress again, however, that the quark anti-quark entropy is unaffected by any undetermined finite renormalization of $`V(r)`$ at zero temperature, i.e. $`S_{\mathrm{}}(T)`$ does not dependent on any flavor and/or quark mass dependent normalization terms that could contribute to $`V(r)`$ at $`T=0`$.
### III.1 The free energy
Our results for $`F_{\mathrm{}}(T)`$ are summarized in Fig. 2 as function of $`T/T_c`$ and compared to $`F_{\mathrm{}}(T)`$ obtained in quenched and $`3`$-flavor QCD. While $`F_{\mathrm{}}(T)`$ in quenched QCD exhibits a singularity at $`T_c`$ due to the first order phase transition and is infinite below $`T_c`$, it is well-defined and finite in full QCD at all temperatures due to string breaking below and color screening above the transition. In this case $`F_{\mathrm{}}(T)`$ is steadily decreasing with increasing temperatures in the whole temperature range analyzed by us. A discussion of $`F_{\mathrm{}}(T)`$, in particular for $`T<T_c`$, has already been given Kaczmarek:2005ui ; Kaczmarek:2005uv . We may add here that the slope of $`F_{\mathrm{}}(T)`$ as function of temperature indeed turns out to be maximal in the vicinity of the transition. Although flavor and/or quark mass dependences of the observable $`F_{\mathrm{}}(T)`$ can clearly be seen when comparing the data for $`2`$\- and $`3`$-flavor QCD, at temperatures about $`1.2T_c<T<\mathrm{\hspace{0.33em}2}T_c`$ no or only little differences between quenched and full QCD could be identified. Flavor dependences are, however, to be expected when reaching temperatures in the perturbative regime.
### III.2 The entropy
The entropy contribution, $`TS_{\mathrm{}}(T)`$, obtained in $`2`$-flavor QCD is shown in Fig. 3 as function of $`T/T_c`$ and is again compared to results from quenched and $`3`$-flavor QCD. We indeed find $`TS_{\mathrm{}}(T)\mathrm{\hspace{0.33em}\hspace{0.25em}\hspace{0.33em}0}`$ at all temperatures analyzed here. Moreover, the data for $`TS_{\mathrm{}}(T)`$ at low temperatures also suggest a vanishing contribution in the zero temperature limit, $`T0`$. Unfortunately we could not go to smaller temperatures to clearly demonstrate this behavior. On the other hand, in the high temperature phase, i.e. at $`T>\mathrm{\hspace{0.33em}2}T_c`$, we indeed find a tendency for an increase of $`TS_{\mathrm{}}(T)`$ with temperature as expected from (22). Actually, this small increase is consistent with the rise given in (21). To demonstrate this we also compared $`TS_{\mathrm{}}(T)`$ in $`2`$-flavor QCD to Eq. (21) using the perturbative $`2`$-loop coupling, i.e.
$`g_{2loop}^2(T)`$ $`=`$ $`2\beta _0\mathrm{ln}\left({\displaystyle \frac{\mu T}{\mathrm{\Lambda }_{\overline{\mathrm{MS}}}}}\right)+{\displaystyle \frac{\beta _1}{\beta _0}}\mathrm{ln}\left(2\mathrm{ln}\left({\displaystyle \frac{\mu T}{\mathrm{\Lambda }_{\overline{\mathrm{MS}}}}}\right)\right),`$
with
$`\beta _0`$ $`=`$ $`{\displaystyle \frac{1}{16\pi ^2}}\left(11{\displaystyle \frac{2N_f}{3}}\right),`$
$`\beta _1`$ $`=`$ $`{\displaystyle \frac{1}{(16\pi ^2)^2}}\left(102{\displaystyle \frac{38N_f}{3}}\right),`$
assuming vanishing quark masses. We used $`T_c/\mathrm{\Lambda }_{\overline{\mathrm{MS}}}=0.77(15)`$ Karsch:2000ps ; Gockeler:2005rv ; Laine:2005ai and the ambiguity in fixing the scale in perturbation theory, $`\mu =\pi ,\mathrm{},4\pi `$. This estimate is shown within the dashed lines and qualitatively agrees with the lattice data for $`T>\mathrm{\hspace{0.33em}2}T_c`$. However, to clearly establish the perturbative increase of $`TS_{\mathrm{}}(T)`$ with increasing temperature will require the analysis of significantly higher temperatures. We also note that within the statistical accuracy of the data we find for $`T>\mathrm{\hspace{0.33em}2}T_c`$ the tendency,
$`S_{\mathrm{}}^{N_f=0}(T)<S_{\mathrm{}}^{N_f=2}(T)<S_{\mathrm{}}^{N_f=3}(T).`$ (24)
It appears indeed quite reasonable that introducing additional flavor degrees of freedom may enhance the finite temperature quark anti-quark entropy in the deconfined phase. It is, however, quite difficult to separate clearly the different effects from flavor and quark mass dependence in full QCD. In particular, at temperatures below $`T_c`$ the tendency given in (24) may change as can be seen from the temperature dependence of $`F_{\mathrm{}}(T)`$ shown in Fig. 2.
In contrast to the small temperature dependence of $`TS_{\mathrm{}}(T)`$ at low and high temperatures, $`TS_{\mathrm{}}(T)`$ shows qualitatively and quantitatively significant differences at temperatures in the vicinity of the transition. In fact, $`TS_{\mathrm{}}(T)`$ obtained in $`2`$-flavor QCD exhibits a sharp peak at $`T_c`$. This behavior signals the high temperature phase transition/crossover in QCD. As the peak is so sharp we may introduce $`T_l`$ ($`T_u`$) defined as the lower (upper) temperature at which $`S_{\mathrm{}}(T)`$ approaches about half of the peak value, i.e. $`S_{\mathrm{}}(T_{l,u})S_{\mathrm{}}(T_c)/2`$. We find for $`2`$-flavor QCD $`T_l`$ about $`0.89T_c`$ and $`T_u`$ about $`1.07T_c`$ using $`S_{\mathrm{}}(T_c)16.5`$. This temperature range is shown at the bottom of Fig. 2 as thick line. A similar behavior is also apparent in $`3`$-flavor QCD. This indicates that the crossover from the low to the high temperature phase in QCD takes place in a small temperature range around $`T_c`$. We stress, however, that the temperature dependence of $`F_{\mathrm{}}(T)`$ also at temperatures above $`1.07T_c`$ is still to a large extent dominated by non-perturbative effects.
### III.3 The internal energy
The internal energy, $`U_{\mathrm{}}(T)`$, in $`2`$-flavor QCD is shown in Fig. 4 as function of temperature and is compared to the corresponding free energy, $`F_{\mathrm{}}(T)`$ (solid line), already shown in Fig. 2. We indeed find $`U_{\mathrm{}}(T)>F_{\mathrm{}}(T)`$ at all temperatures analyzed here. It can clearly be seen that the temperature dependence of $`F_{\mathrm{}}(T)`$ and $`U_{\mathrm{}}(T)`$ is qualitatively and quantitatively different. While the free energy steadily decreases with increasing temperatures the internal energy exhibits a pronounced peak. Again this peak is sharply localized at the (pseudo-) critical temperature. At the temperatures analyzed by us the internal energy below $`T_c`$ is rapidly increasing with increasing temperatures. Again we indicate by dotted lines the plateau value of the heavy quark potential at zero temperature, $`V(r_{\text{breaking}})1000\mathrm{\hspace{0.33em}1200}`$ MeV, using $`r_{\text{breaking}}1.2\mathrm{\hspace{0.33em}1.4}`$ fm Pennanen:2000yk . A comparison of $`U_{\mathrm{}}(T)`$ with this value shows again that most of the temperature dependence of $`U_{\mathrm{}}(T)`$ is sharply localized at temperatures in the vicinity of the transition. A qualitatively similar behavior is also apparent in $`3`$-flavor QCD Petreczky:2004pz . Comparing the available data some flavor or quark mass dependence, $`U_{\mathrm{}}^{N_f=2}(T_c)4000`$ MeV $`>`$$``$ $`U_{\mathrm{}}^{N_f=3}(T_c)3000`$ MeV, can be observed and also the value $`U_{\mathrm{}}(T_c^+)`$ in quenched QCD Phd is of similar magnitude. As noted in Sec. II.2, however, the values for $`U_{\mathrm{}}(T)`$ may depend on the relative normalization of $`V(r)`$ at $`T=0`$ used for renormalization.
At temperatures above $`T_c`$, $`U_{\mathrm{}}(T)`$ rapidly drops while at higher temperatures, i.e. $`T>\mathrm{\hspace{0.33em}1.3}T_c`$, the temperature dependence of $`U_{\mathrm{}}(T)`$ turns out to be much weaker than in the vicinity of phase transition. However, contact with the perturbative relation, Eq. (20), is not expected at those temperatures as $`U_{\mathrm{}}(T)`$ is still positive. In fact, when approaching the perturbative high temperature regime also $`U_{\mathrm{}}(T)`$ is expected to exhibit a change in sign and will slowly diverge with respect to (20).
## IV Finite temperature energies and potential models
The analysis of bound state problems has been quite successful in terms of potential theory at $`T=0`$ Eichten:1978tg ; Eichten:1979ms ; Jacobs:1986gv . For the discussion of quarkonium suppression patterns at finite temperature one also often resorts to potential models Digal:2001iu ; Digal:2001ue ; Wong:2004kn ; Shuryak:2003ty ; Park:2005nv . Of course, the strong interaction remains unaffected by temperature and the modeling of thermal modifications of heavy quark bound states requires the definition of an effective potential, $`V_{\text{eff}}(r,T)`$ Karsch:2005ex , which can be given only phenomenologically, for instance, by using the modifications of the free and internal energies. It is thus important to understand the binding properties of the different finite temperature energies. For this purpose we also calculated the quark anti-quark internal energies for the temperatures given in Tab. 1 at several finite distances. Parts of our results for $`U_1(r,T)`$ are summarized in Fig. 5 at temperatures below (a) and above (b) the transition. Similar results have been obtained for internal energies in the averaged and octet channels.
To gain some insight into the consequences these energies have for quarkonium dissociation we compare the asymptotic ($`r\mathrm{}`$) energies with the potential energies at a distance corresponding to the size of some quarkonium states,
$`\mathrm{\Delta }E_i(T=0)V(\mathrm{})V(r_i),`$ (25)
where the radii $`r_i`$ ($`i=J/\psi ,\chi _c,\mathrm{}`$) are listed in the first row of Tab. 2. At zero temperature $`V(\mathrm{})`$ is taken to be twice the energy needed to create the lowest heavy-light meson. For our purpose we consider $`V(\mathrm{})V(r_{\text{breaking}})1000`$ MeV where $`r_{\text{breaking}}`$ is the distance at which the string is expected to break at zero temperature, $`r_{\text{breaking}}>\mathrm{\hspace{0.33em}1.2}`$ fm Pennanen:2000yk . This energy is shown in Fig. 6 as horizontal line. The resulting energy for $`J/\psi `$ is also shown. The energies for some charmonium and bottomonium states are summarized in Tab. 2 and compared to the mass difference obtained from $`2M_{D,B}m_i`$ where $`M_{D,B}`$ denotes the $`D`$\- and $`B`$-meson masses and $`m_i`$ the masses of the different quarkonium states Eidelman:2004wy . Of course, the wave functions for the different quarkonium states will also reach out to larger distances Jacobs:1986gv and thus our estimate for the different energy levels $`E_i(T=0)`$ can only be taken as indicative for the relevant energies. Potential model analysis, using for instance the Schrรถdinger equation, will do better in this respect.
Similarly we can estimate the temperature dependence of the energy levels for the different quarkonium states from $`E_i(T)V_{\text{eff}}(r_i,T)`$. Again these energy levels will only characterize the relevant energies and the sizes of these states may also become temperature dependent. At finite temperature, however, the values for these levels are expected to depend crucially also on the specific modeling of the effective potential, $`V_{\text{eff}}(r,T)`$. This is obvious from the different energy levels for the $`J/\psi `$ shown in Fig. 6 which we obtained by using as $`V_{\text{eff}}(r,T)`$ the singlet free energy (lower dashed line) and singlet internal energy (upper dashed line). Due to the steeper rise of the internal energy compared to free energy $`E_{J/\psi }(T)`$ is enhanced compared to the energy level obtained from the internal energy. It is interesting to note here that $`E_{J/\psi }(T=1.3T_c)`$ deduced from the internal energy is even larger than at zero temperature while in terms of the free energy it is smaller than $`E_{J/\psi }(T=0)`$. For the characterization of the relevant energies needed to dissociate the bound state we again consider $`\mathrm{\Delta }E_i(T)V_{\text{eff}}(\mathrm{},T)V_{\text{eff}}(r_i,T)`$, which depends on the definition of $`V_{\text{eff}}(r,T)`$. This is evident from Fig. 7 where we show the temperature dependence of $`\mathrm{\Delta }E_{J/\psi }(T)`$ estimated from $`V_{\text{eff}}(r,T)U_1(r,T)`$ (filled circles) and $`V_{\text{eff}}(r,T)F_1(r,T)`$ (open circles). While $`\mathrm{\Delta }E_{J/\psi }(T)`$ is continuously decreasing with increasing temperatures when using $`F_1(r,T)`$ as effective potential, its binding pattern appears quite different from what one obtains by using $`V_{\text{eff}}(r,T)U_1(r,T)`$. Actually, in the latter case the binding pattern exhibits a maximum in the vicinity of the transition while it rapidly drops above $`T_c`$. Similar results have been obtained also for other charmonium and bottomonium states and are summarized in Tab. 2.
Of course, the model dependences of the dissociation energies at finite temperature also affect the analysis of suppression patterns and corresponding dissociation temperatures Digal:2001iu ; Digal:2001ue ; Wong:2004kn . Actually, using $`V_{\text{eff}}(r,T)U_1(r,T)`$ suggests that suppression of $`J/\psi `$ may occur only at temperatures close but above the transition while from $`V_{\text{eff}}(r,T)F_1(r,T)`$ one finds that $`J/\psi `$ dissolves already at temperatures below the crossover. Similar model dependences enter also the analysis of excited quarkonium states and the corresponding estimates for the dissociation temperatures are summarized in Tab. 2 using four different definitions for $`V_{\text{eff}}(r,T)`$, i.e. we used the singlet free and internal energies ($`F_1(r,T)`$, $`U_1(r,T)`$) as well as the finite temperature energies in the color averaged channel ($`F_{\overline{q}q}(r,T)`$, $`U_{\overline{q}q}(r,T)`$).
## V Summary and conclusions
Following Kaczmarek:2002mc ; Zantow:2003ui we introduced and analyzed the change in internal energy and entropy due to the presence of a static quark anti-quark pair in a QCD heat bath. Both observables are introduced as intensive observables as appropriate derivatives of the renormalized free energy. Similar to the singlet quark anti-quark free energies Kaczmarek:2002mc ; Zantow:2003uh ; Zantow:2001yf ; Kaczmarek:2005ui also the singlet internal energies become temperature independent in the limit of small distances and are controlled by the zero temperature running coupling.
We analyzed qualitative and quantitative differences that appear when changing from free energies to internal energies as observable that defines an effective potential that can be used in model calculations. We discussed the important role of the entropy contribution at finite temperature. At short distances, at intermediate and at large distances the entropy contribution is non-zero and shows non-trivial $`r`$-dependences. We find positive entropy contributions, $`S_1(r,T)>\mathrm{\hspace{0.33em}0}`$, and thus $`U_1(r,T)>F_1(r,T)`$. Similar to the free energies, also the large distance properties of the internal energies and entropies are controlled by string breaking below and color screening above deconfinement and both approach temperature dependent constant values, which define $`U_{\mathrm{}}(T)`$ and $`S_{\mathrm{}}(T)`$, at asymptotic large distances. Actually, the difference between free and internal energies at high temperatures, $`TS_{\mathrm{}}(T)4m_D(T)\alpha (T)/3`$, is supposed to increase with increasing temperatures. In particular, $`U_{\mathrm{}}(T)`$ and $`S_{\mathrm{}}(T)`$, are, similar to $`F_{\mathrm{}}(T)`$ Kaczmarek:2002mc , again introduced as manifest gauge invariant observables and clearly signal the QCD plasma transition. In fact, while the plateau values which are approached by the free energies, $`F_{\mathrm{}}(T)`$, are rapidly decreasing in the vicinity of the transition Kaczmarek:2005ui ; Kaczmarek:2005uv , the values approached by the internal energies, $`U_{\mathrm{}}(T)`$, and entropies, $`S_{\mathrm{}}(T)`$, show both a sharp peak at the (pseudo-) critical temperature. Similar results are also obtained in quenched and $`3`$-flavor QCD Kaczmarek:2002mc ; Zantow:2003ui ; Phd ; Petreczky:2004pz . However, qualitative differences become quite transparent in the vicinity and below the transition when comparing these observables obtained in quenched and full QCD. In quenched QCD the first order phase transition is related to singularities in thermodynamic observables which can indeed be seen after renormalization in the temperature dependence of the finite temperature energies, entropies and the renormalized Polyakov loop (see also Kaczmarek:2002mc ; Kaczmarek:2005ui ; tobep ). In contrast to quenched QCD, in full QCD the phase change is a crossover and we consequently do not see any singularities in the finite temperature energies and entropies nor in the temperature dependence of the renormalized Polyakov loop Kaczmarek:2005ui .
We also investigated the temperature dependence of quarkonium binding in the vicinity of the transition. For this purpose we used the finite temperature energies to define appropriate effective potentials, $`V_{\text{eff}}(r,T)`$, as one would do in potential models. As effective potentials we used the singlet and averaged free and internal energies, i.e. we defined $`V_{\text{eff}}(r,T)`$ through $`F_1(r,T)`$, $`F_{\overline{q}q}(r,T)`$, $`U_1(r,T)`$ and $`U_{\overline{q}q}(r,T)`$, and estimated the binding energies at temperatures below and above the transition. In all cases the binding energies of the quarkonium states indeed become weaker with increasing temperatures above the transition and this may lead to dissociation of parts of these states at temperatures close but above $`T_c`$. Our analysis, however, shows strong dependencies of the binding energies and dissociation temperatures on the specific modeling of $`V_{\text{eff}}(r,T)`$. To some extent these model dependencies enter from quite general grounds when exchanging the definition of $`V_{\text{eff}}(r,T)`$ from free energies into internal energies. When using free energies the binding energies continuously and rapidly decrease when crossing the transition and most of the quarkonium bound states may thus indeed dissolve at temperatures in the vicinity and below the transition Digal:2001iu ; Digal:2001ue . The temperature dependence of the binding energies deduced from internal energies, however, turns out to be more complicated in the vicinity of the transition. An effective potential defined through quark anti-quark internal energies suggests increasing binding energies below the transition which exhibit a peak at $`T_c`$. This may imply that all quarkonium states analyzed here are still bound at the transition. The binding energies of the different states rapidly decrease above $`T_c`$ leading again to quarkonium dissociation, however, the temperatures which are relevant for dissociation are shifted to larger temperatures than those deduced from free energies. The recent potential model calculations Shuryak:2004tx ; Wong:2004kn using the properties of $`U_1(r,T)`$ at temperatures above the transition as well as the lattice analysis of quarkonium bound states in quenched QCD Asakawa:2003re ; Datta:2003ww support our findings.
###### Acknowledgements.
We thank the Bielefeld-Swansea collaboration for providing us their configurations with special thanks to S. Ejiri. We would like to thank E. Laermann, F. Karsch and H. Satz for many fruitful discussions. F.Z. thanks P. Petreczky for his continuous support. We thank K. Petrov and P. Petreczky for sending us the data of Ref. Petreczky:2004pz . This work has partly been supported by DFG under grant FOR 339/2-1 and by BMBF under grant No.06BI102 and partly by contract DE-AC02-98CH10886 with the U.S. Department of Energy. At an early stage of this work F.Z. has been supported through a stipend of the DFG funded graduate school GRK881.
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# Nonlinear evolution equations in QCD and effective Hamiltonian at high energy
## 0.1 Introduction: high energy limit and the Pomeron
One of the central problems in the theory of strong interactions is the understanding of the behaviour of the hadronic cross sections in the limit of high energies. The experimental data on the total cross section show slow but distinct increase with energy. It has been suggested that this rise, which can be parametrized by a power-like form,
$$\sigma (s)s^{0.08},$$
(1)
is mediated by the โsoft Pomeronโ DL . Perturbative calculation summing the leading logarithmic contributions of large logs $`\mathrm{ln}s/t`$ results in the BFKL Pomeron, which however leads to a stronger rise of the cross section with energy BFKL
$$\sigma (s)s^{4\mathrm{ln}2\alpha _sN_c/\pi },$$
(2)
with the leading exponent being much larger than in (1) when evaluated at physically interesting values of the strong coupling constant ($`\alpha _s0.2`$). In any of these cases, the power-like increase of the cross section is in contradiction with the Froissart bound Froissart which allows at most logarithmic increase with energy
$$\sigma (s)\sigma _0\mathrm{ln}^2s,$$
(3)
with a normalization coefficient $`\sigma _0`$ related to the inverse of the pion mass squared. Froissart bound stems from the very general principles: unitarity and the finite range of the strong interactions. Thus, the important question arises whether it is possible to identify and calculate other type of Feynman diagrams which will lead to the restoration of the unitarity.
## 0.2 Color Glass Condensate and the JIMWLK equation
The Color Glass Condensate is an effective theory of the strong interactions at very high energies. Apart from the standard BFKL evolution of the gluon density it also contains the recombination diagrams which are important when the gluon density becomes very high. It therefore describes the phenomenon known as perturbative parton saturation. When the density of gluons becomes very high due to their enhanced splitting described by the BFKL Pomeron, the recombination effects start to become important and reduce the growth. It is believed that this phenomenon is important for the restoration of the unitarity. The Color Glass Condensate theory describes the scattering of the projectile off a target which constitutes a dense system of soft gluons. The S-matrix for the scattering of the quark-antiquark dipole in the target field $`\alpha `$ is described by
$$S(๐,๐)=\frac{1}{N_c}\mathrm{Tr}(V_๐V_๐^{}),$$
(4)
where Wilson line
$$V_๐=P\mathrm{exp}\left(ig๐x^{}\alpha _a(x^{},๐)t^a\right),$$
(5)
is the path ordered exponential along the trajectory of the projectile. The physical scattering matrix between the $`q\overline{q}`$ dipole and the hadron target is obtained by taking the average of (4) over the color field of the target
$$S(๐,๐)_\tau =D[\rho ]๐ต[\rho ]_\tau S(๐,๐),$$
(6)
where $`\rho `$ is a color charge which generates field $`\alpha `$. One cannot compute the weight function $`๐ต[\rho ]`$ since it contains nonperturbative information about the hadron but one can compute the evolution of this weight function with increasing rapidity $`\tau \mathrm{ln}s`$. The basic equation of the Color Glass Condensate theory is the JIMWLK equation JIMWLK which governs the evolution of the weight function $`๐ต[\rho ]_\tau `$ with rapidity
$$\frac{๐ต[\rho ]_\tau }{\tau }=H_{JIMWLK}๐ต[\rho ]_\tau ,$$
(7)
where $`H`$ is the JIMWLK Hamiltonian. The JIWMLK equation is a very complicated functional evolution equation, which however reduces to one closed and relatively simple equation, the Balitsky-Kovchegov equation, in the large $`N_c`$ limit and in the dipole picture. Diagramatically JIMWLK equation contains BFKL Pomeron ladder diagrams as well as the triple Pomeron interaction diagrams which reduce the growth of the gluon density. More precisely it contains Pomeron merging diagrams, however it misses the Pomeron splittings. One can understand it by looking at the structure of the operator $`H_{JIMWLK}`$
$$H_{JIMWLK}=\frac{1}{2\pi }_{๐,๐,๐}K_{๐๐๐}\frac{\delta }{\delta \alpha ^a(๐)}\left[1+\stackrel{~}{V}_๐^{}\stackrel{~}{V}_๐\stackrel{~}{V}_๐^{}\stackrel{~}{V}_๐\stackrel{~}{V}_๐^{}\stackrel{~}{V}_๐\right]^{ab}\frac{\delta }{\delta \alpha ^b(๐)},$$
(8)
which is second order in derivatives $`\frac{\delta }{\delta \alpha }`$ and all orders (through $`\stackrel{~}{V}`$โs)<sup>1</sup><sup>1</sup>1$`\stackrel{~}{V}`$ denotes the Wilson line in the adjoint representation. in the field $`\alpha `$. That means, that the evolution of the correlation functions in the field $`\alpha `$ can be (by using (6,7,8)) schematically represented as
$$\frac{\stackrel{n}{\stackrel{}{\alpha \mathrm{}\alpha }}}{\tau }=\underset{mn}{}๐ฆ_m\stackrel{m}{\stackrel{}{\alpha \mathrm{}\alpha }}.$$
(9)
The evolution of $`n`$ point functions is coupled to the higher order correlation functions but not to the lower correlation functions.
## 0.3 Dual version of JIMWLK and effective Hamiltonian
In order to include the additional diagrams corresponding to the Pomeron splittings in the Color Glass Condensate formalism, terms which are higher order in derivatives have to be incorporated MSW . By including these terms in the evolution, the Pomeron loops are generated, which are known to be important contributions at very high energy. It has been suggested KL , that in general there is a duality relation between the evolution of the dense system, governed by the JIWMLK equation, and the dilute regime. In particular, the Pomeron mergings and splittings are believed to be the dominant diagrams in the dense and the dilute regimes, respectively, of the gluon field. In particular the evolution in the dilute regime is governed by the dual equation to the JIMWLK evolution. Formally, the duality can be expressed as the transformation KL ; HIMST
$$x^{}x^+,\frac{\delta }{\delta \alpha ^a(x^{},๐)}i\rho ^a(x^+,๐),\alpha ^a(x^{},๐)i\frac{\delta }{\delta \rho ^a(x^+,๐)},$$
(10)
where $`\rho `$ is the charge of the target and the $`\alpha `$ is the field generated by this charge measured by the projectile. The dual version of the JIMWLK equation is then
$$\frac{๐ต[\rho ]_\tau }{\tau }=\stackrel{~}{H}๐ต[\rho ]_\tau ,$$
(11)
where the dual Hamiltonian
$$\stackrel{~}{H}=\frac{1}{2\pi }_{๐,๐,๐}K_{๐๐๐}\rho ^a(๐)\left[1+\stackrel{~}{W}_๐^{}\stackrel{~}{W}_๐\stackrel{~}{W}_๐^{}\stackrel{~}{W}_๐\stackrel{~}{W}_๐^{}\stackrel{~}{W}_๐\right]^{ab}\rho ^b(๐).$$
(12)
The Wilson line $`\stackrel{~}{W}`$ is given by
$$\stackrel{~}{W}_๐=P\mathrm{exp}\left(g๐x^+\frac{\delta }{\delta \rho _a(x^+,๐)}T^a\right),$$
(13)
where now the integration and ordering is in the other light-cone variable $`x^+`$. The above evolution equation thus includes the process of Pomeron splitting by evolving the weight function of the target. The duality can be also rexpressed in term of the symmetry between the projectile and target, what is the splitting from the target side, can be interpreted as a Pomeron merging from the projectile side. The general Hamiltonian at high energy should therefore include both of these processes at the same time. It can be shown HIMST that it is entirely expressed in terms of Wilson lines along $`x^{}`$ and $`x^+`$ directions
$`H_{\mathrm{eff}}={\displaystyle \frac{1}{2\pi g^2N_c}}{\displaystyle _๐}\mathrm{Tr}`$ $`[\stackrel{~}{V}_{\mathrm{}}(^i\stackrel{~}{W}_{\mathrm{}})(^i\stackrel{~}{V}_{\mathrm{}}^{})\stackrel{~}{W}_{\mathrm{}}^{}+\stackrel{~}{V}_{\mathrm{}}\stackrel{~}{W}_{\mathrm{}}(^i\stackrel{~}{V}_{\mathrm{}}^{})(^i\stackrel{~}{W}_{\mathrm{}}^{})`$
$`+(^i\stackrel{~}{W}_{\mathrm{}}^{})(^i\stackrel{~}{V}_{\mathrm{}})\stackrel{~}{W}_{\mathrm{}}\stackrel{~}{V}_{\mathrm{}}^{}+\stackrel{~}{W}_{\mathrm{}}^{}(^i\stackrel{~}{V}_{\mathrm{}})(^i\stackrel{~}{W}_{\mathrm{}})\stackrel{~}{V}_{\mathrm{}}^{}],`$ (14)
where the infinity limits in the Wilson lines $`\stackrel{~}{V}`$ and $`\stackrel{~}{W}`$ correspond to the limits in the $`x^+`$ and $`x^{}`$ lightcone variables correspondingly. This Hamiltonian is reminiscent of the nonlinear sigma model for the effective theory at high energy, proposed in VV . The important point to note here however, is that $`\stackrel{~}{V}`$ and $`\stackrel{~}{W}`$ are interpreted as operators which are the functions of fields and derivatives of fields, respectively, with very nontrivial commutation relations. The effective Hamiltonian (0.3) has an interesting symmetry property, namely it is invariant under the following transformation
$`\stackrel{~}{W}_{\mathrm{}}\stackrel{~}{V}_{\mathrm{}},\stackrel{~}{V}_{\mathrm{}}\stackrel{~}{W}_{\mathrm{}}^{},\stackrel{~}{W}_{\mathrm{}}^{}\stackrel{~}{V}_{\mathrm{}}^{},\stackrel{~}{V}_{\mathrm{}}^{}\stackrel{~}{W}_{\mathrm{}},`$ (15)
which is refferred to as the selfduality property. The effective Hamiltonian reduces to JIMWLK (8) and its dual (12) when the Wilson lines $`\stackrel{~}{W}`$ and $`\stackrel{~}{V}`$ are expanded to the first nontrivial order in derivatives and fields, respectively.
The results presented in this talk have been obtained in the collaboration with Y.Hatta, E.Iancu, L.McLerran and D.Triantafyllopoulos HIMST . This research has been supported by the U. S. Department of Energy, Contract No. DE-AC02-98CH10886 and by the Polish Committee for Scientific Research, KBN Grant No. 1 P03B 028 28.
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# Equations of Motion with Multiple Proper Time: A New Interpretation of Spin
## I Introduction
There are three major properties of quantum physics which are different from classical physics:
1) Non-local property of single quantum particle: a single particle can stay in different places at the same time. Non-local property is the most important part which quantum physics is distinguish itself from classical physics. Without non-local property, the wave of single particle becomes an oscillator, which can be found in classical physics.
2) Statistical effect in the measurement of single particle.
3) Spin and related statistics: Bose-Einstein statistics and Fermi-Dirac statistics.
In previous paper chen1 we demonstrate that: if a classical particle moves under three proper time (three dimensional time), it will give the same non-local property of quantum particle and the same statistical effect of measurement as quantum particle. The properties 1) and 2) of quantum particle that we mentioned above can be derived from classical particle under three proper time model. For instance, Fig1. draw the trajectory of single free particle under two independent proper time $`\tau `$ and $`\sigma `$ on $`x_0x_i`$ plane in Riemann Space.
In Fig1. a single particleโs trajectory is determined by two proper time ($`\tau `$ and $`\sigma `$) where world line $`\tau `$ and $`\sigma `$ are orthogonal to each other. At each points on world line $`\tau `$, particle will also move along world line $`\sigma `$. As the result: At $`t=0`$, the single particle will stay at many positions: $`x_1,x_2,\mathrm{}x_n`$ with different values of $`\tau `$ and $`\sigma `$ $`(\tau _1,\sigma _1)..(\tau _n,\sigma _n)`$. Also the single particle will stay at $`x=0`$ at different time: $`t_1,t_2,..t_n`$ with different values of $`\tau ,sigma`$; where $`x_n=h/mu`$ and $`t_n=h/mc^2`$ which are de Broglie wavelength and period. By adding periodic conditions for $`\tau `$ and $`\sigma `$, Fig1. becomes de Broglie plane wave for quantum particle. Briefly speaking, in paper chen1 , we built 3 proper time model for free particle as below:
1)Along world line $`\tau `$, free particle moves the same as classical particle with classical energy and momentum.
2)Free particle can move along world line $`\sigma `$ and $`\varphi `$(where $`\sigma `$ and $`\varphi `$ are second and third proper time respectively) while first proper time $`\tau `$ unchanged. World lines of $`\sigma `$ and $`\varphi `$ are also straight lines. Relations between $`\sigma `$ and 2nd time dimension $`x_4`$, $`\varphi `$ and third time dimension $`x_5`$ are $`x_4=e^{i\sigma }=e^{\frac{im_0\tau }{\mathrm{}}}`$, $`x_5=e^{i\varphi }=e^{\frac{im_0\tau }{\mathrm{}}}`$, These equations come from the geometry of $`x_4`$ and $`x_5`$, which are loops in complex plane.
The statistical effect of measurement for single particle is derived from the nature character of time: the apparatus can meet a particle at a spatial location if and only if all three values of three dimensional time are equals between particle and apparatus: $`t_i^{apparatus}=t_i^{particle}`$. But because we only have knowledge of one dimensional time, we donโt know how to synchronize other two time dimensions, the results of measurement become statistical. (see paper chen1 for detail).
This paper will give more clear picture of 3-proper time model than previous paper chen1 . We will extend above 3-proper time model to include particle with spin. We will derive spin, Bose-Einstein statistics and Fermi-Dirac statistics from classical physics under three proper time. In section II, 3-proper time model is introduced for spinless particle. In section III, 3-proper time model for Boson with spin $`>0`$ is derived. In section IV, 3-proper time for Fermion is discussed. In section V, we derived Bose-Einstein statistics and Fermi-Dirac statistics for Boson and Fermion respectively. In last section, we will give the interpretation of what is spin. The purpose of this paper is to show that: when a single particle moving under 3-proper time, the trajectories of a classical particle are equivalent to a quantum field with spin.
## II 3-proper time models for free spinless single particle
We need clarify some terminologies which will be used in this paper. we will use both words: three dimensional time and three proper time. Proper time is the same meaning as it is in relativity: Proper time is a special affine parameter of space-time; Each proper time has one related world line, i.e. each proper time corresponding to one individual movement. Time dimension means the time value in a chosen time coordinates system. We will use word : โthe direction of speed of particleโ which means the spatial moving direction of the particle with proper time $`\tau `$. โWorld line $`\tau `$โ means the trajectory of particleโs motion by proper time $`\tau `$. Similarly we will use โWorld line $`\sigma `$โ and โWorld line $`\varphi `$ โ.
Through out this paper, we keep some reasonable assumptions:
1)Particleโs world lines by different proper time are orthogonal to each other.
2)All three proper time are independent.
3)We implement cylinder condition on 2nd and 3rd proper time $`\sigma `$, and $`\varphi `$: Each of them is angle of a loop with value from $`0`$ to $`2\pi `$.
For spinless particle, the 3-proper time model is:
1) World line $`\tau `$ is the same as the world line in relativity; The velocity of particle in world line $`\tau `$ is $`u`$, which is the same as velocity in relativity, it is also the group velocity of de Broglie wave.
2) The projection of world line $`\sigma `$ in 4-dimensional time-space $`(x_0,\stackrel{}{x})`$ is orthogonal to world line $`\tau `$; The velocity of particle on world line $`\sigma `$ is $`v`$, which is the phase velocity of de Broglie wave.
3) The projection of world line $`\varphi `$ in 4-dimensional time-space $`(x_0,\stackrel{}{x})`$ is coincident with world line $`\tau `$.
Considering particle under rest reference frame, velocity on world line $`\tau `$ is zero: $`u_i=0`$. In this case, phase velocity of de Broglie wave v becomes infinite since $`v=\omega \lambda =c^2/u`$. To understand this, considering in rest frame, proper time $`\tau =t`$, world line $`\tau `$ becomes $`X_0`$ axis in Fig1., world line $`\sigma `$ becomes $`X_i`$ axis in Fig1.; Particle does not move with $`\tau `$, but it still moves with second proper time $`\sigma `$. The definition of velocity with $`\sigma `$ is $`v_i=cdx_i/dx_0`$ where c is speed of light, since this motion does not change first dimensional time t: $`dx_0=0`$, so $`v_i`$ becomes infinite. To obtain the equation of motion under $`\sigma `$, we need know the relationship between $`\sigma `$ and 4-dimensional space-time coordinates. In previous paper chen1 , world line $`\sigma `$ can go infinite since plane wave can be everywhere in the space. In reality, the inital position of particle is always under certain boundary condition. Now we localize the particleโs wave-length up to one wave-length. Let second time dimension $`x_4=\sigma `$, the spatial direction of moving by $`\sigma `$ is $`x_1`$, and choose equation of motion in rest frame be:
$$\sigma =\pi \mathrm{cos}\frac{m_0c}{\mathrm{}}x_1\mathrm{\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}0}x<\frac{h}{m_0c}$$
(1)
and
$`x_1=x_5`$ (2)
$`x_j=0j=0,2,3`$ (3)
$`\sigma =x_4`$ (4)
(5)
where $`m_0`$ is rest mass of particle, c is speed of light, h is planck constant; $`x_4`$ is 2nd time dimension; $`x_5`$ is 3rd time dimension. Equation (1) tells us that in the rest reference frame, particle moves like an oscillator with proper time $`\sigma `$. The equation contains $`m_0`$, so we can understand that the motion of oscillating comes from the rest energy: $`m_0c^2`$. Back to general reference frame with velocity $`u=u_1>0`$, by using Lorentz transformation, we turn back to Fig1., where particle is still non-localized within one wave length. Using Lorenze transformation, and noticing that $`v_1=c^2/u_1`$, Equation (1) becomes:
$$\sigma =\pi \mathrm{cos}\frac{im_0c}{\mathrm{}}\frac{(v_1tx_1)}{\sqrt{1v_1^2/c^2}}=\pi \mathrm{cos}\frac{1}{\mathrm{}}(Etp^1x_1)$$
(6)
where $`x_1`$ and t are coordinates in new reference frame. The imaginary number i is to keep $`\sigma `$ being real number in Lorenze transformation. We can rewite equation (6) as
$`\sigma ={\displaystyle \frac{\pi }{2}}(ae^{\frac{i}{\mathrm{}}p_1x_1}+a^{}e^{\frac{i}{\mathrm{}}p_1x_1})`$ (7)
where $`a^{}`$ is complex conjugate of a, which is
$$a=\frac{\pi }{2}e^{\frac{i}{\mathrm{}}Et}$$
(8)
Equation (7) is derived under the condition that the direction of motion is point to $`x_1`$. In general case, it can be written as:
$`\sigma ={\displaystyle \frac{\pi }{2}}(ae^{\frac{i}{\mathrm{}}p^ix_i}+a^{}e^{\frac{i}{\mathrm{}}p^ix_i})`$ (9)
where i= 1,2,3. In quantum field theory, $`a`$ in equation(9) is annihilation operator of scalar field, $`a^{}`$ is creation operator. Here we see that $`a^{}`$ means particle moves toward positive t direction, $`a`$ means particle moves toward negative t direction.
Let the relation between second proper time $`\sigma `$ and second time dimension $`x_4`$ be:
$$d\sigma =\mathrm{cos}\frac{1}{\mathrm{}}(p^\alpha x_\alpha m_0x_5)dx_4$$
(10)
The difference between $`\sigma `$ and $`x_4`$ comes from the geometry of time-space. On world line $`\sigma `$, it becomes equation(5).
The geometry of 6-dimensional time-space is:
$$ds^2=dx_\alpha dx^\alpha +\psi ^2dx_4dx^4dx_5dx^5$$
(11)
where $`\psi `$ can be chosen as
$$\psi =e^{\frac{i}{\mathrm{}}(p^\alpha x_\alpha m_0x_5)}$$
(12)
or
$$\psi =\mathrm{cos}\frac{1}{\mathrm{}}(p^\alpha x_\alpha m_0x_5)$$
(13)
The metric of 6-dimensional space-time is
$$\left(\widehat{g}_{AB}\right)=\left(\begin{array}{cc}g_{\alpha \beta }& \\ \psi & \\ 1& \end{array}\right)$$
(14)
where metric elements $`g_{\alpha \beta }`$ is 4-dimensional metric. Substitute above metric into Einstein field equation chen2 chen3 :
$$\widehat{G}_{AB}=\kappa \widehat{T}_{AB},$$
(15)
We can derive Klein-Golden equation:
$$_\alpha ^\alpha \psi +m_0^2\psi =0$$
(16)
where $`\alpha `$ is 0..3. From above discussion, we see that non-local property and field equation of quantum scalar field can be derived by classical particle moving under 3 proper time. Look at โtimeโ part of equation (11):
$$ds_t^2=dx_0dx_5dx^5+\psi dx_4dx^4$$
(17)
Compare to geometry of sphere:
$$ds_t^2=dr^2d\theta ^2r^2\mathrm{sin}\theta ^2d\varphi ^2$$
(18)
The 2nd and 3rd time dimension is similar to sphere with unit radius.
## III 3-proper time models for free Boson with spin one and spine $`>1`$
In section II, for spinless particle, the projection of world line $`\varphi `$ on 4-dimensional space-time $`(t,\stackrel{}{x})`$ is coincident with world line $`\tau `$. In this section we will see that, when the projections of world lines $`\tau `$, $`\sigma `$, $`\varphi `$ on 4-dimensional space-time $`(t,\stackrel{}{x})`$ are separated, we get equations of quantum particle with integer spin $`>1`$.
### III.1 3-proper time model for photon
Since the rest mass of photon is zero, we can not use equation (1). Instead, the oscillation for photon comes from electric energy and magnetic energy. For free photon with single frequency, Electric field and magnetic field both are perpendicular to the direction of wave-vector $`\stackrel{}{k}`$. Now we choose $`x_1`$ axis as direction of electric field $`\stackrel{}{E}`$, $`x_2`$ axis as direction of magnetic field $`\stackrel{}{B}`$, $`x_3`$ axis as direction of wave-vector $`\stackrel{}{k}`$. First, we build 3-proper time model in rest frame where Photonโs speed with world line $`\tau `$ is zero. In world line $`\sigma `$, we choose the equation of motion as:
$`x_1=\lambda (E_0)\mathrm{arccos}{\displaystyle \frac{\sigma }{\pi }}`$ (19)
$`x_3=\lambda \mathrm{arccos}{\displaystyle \frac{\sigma }{\pi }}`$ (20)
where k is wave-vector of photon, $`\lambda (E_0)`$ is coefficient dependent on magnitude of electric field $`E_0`$. So $`\sigma `$ can be derived from either $`x_1`$ or $`x_3`$. From equation(20), we have:
$`\sigma =\pi \mathrm{cos}kx_3`$ (21)
Similarly for $`\varphi `$:
$`x_2=\lambda (B_0)\mathrm{arccos}{\displaystyle \frac{\varphi }{\pi }}`$ (22)
$`x_3=\lambda \mathrm{arccos}{\displaystyle \frac{\sigma }{\pi }}`$ (23)
$$\varphi =\pi \mathrm{cos}kx_3$$
(24)
Where $`\lambda (B_0)`$ is a variable the dependent on magnitude of magnetic field $`B_0`$.
In regular reference frame, photon moves at direction $`x_3`$ with speed c. The equation for $`\sigma `$ is:
$$\sigma =\pi \mathrm{cos}(\omega tkx_3)$$
(25)
and the equations of motions for $`\varphi `$ is:
$$\varphi =\pi \mathrm{cos}(\omega tkx_3)$$
(26)
We know that electric field E and magnetic field B of photon perpendicular to wave-vector $`\stackrel{}{k}`$. Let
$$\stackrel{}{E}=\alpha _e\sigma \stackrel{}{e_1}$$
(27)
$$\stackrel{}{B}=\alpha _b\varphi \stackrel{}{e_2}$$
(28)
where $`\alpha _{e(b)}`$ is constant; $`\stackrel{}{e_1}`$, $`\stackrel{}{e_2}`$ are unit vector point to $`x_1`$ and $`x_2`$ axis. Now we derived photon plane wave by using classical motion of single photon under 3-proper time.
### III.2 3-proper time model for free massive Boson
As we have seen in this paper: besides the classical motion on proper time $`\tau `$; For spinless particle, it has extra oscillation by proper time $`\sigma `$, the oscillation is caused by rest energy; for photon, it has two extra oscillations by $`\sigma `$ and $`\varphi `$ which are perpendicular to the direction of wave-vector, they are caused by electric field E and magnetic field B. Now for a massive Boson with spin one, it can have three separated oscillations. Also in previous sections, we assume that, in rest frame, the world line of first proper time $`\tau `$ is only moving toward $`x_0`$ dirction, $`t=\tau `$. That is not a required condition. In general, we can assume that: in rest frame, world line $`\tau `$ oscillates around t, while the average effect of motion is the same as classical proper time $`\tau `$.
For free particle with spin one, let:
1) Three independent proper time be: $`\tau `$, $`\sigma `$, $`\varphi `$.
2) Each world line perpendicular to each other in 3-dimensional space.
3) We localized each world lines so that, they are all cosine functions.
In rest frame, Let the equations of motion for world line $`\sigma `$ be
$`x_1=\lambda (V_{10})\mathrm{arccos}{\displaystyle \frac{\sigma }{\pi }}`$ (29)
$`x_3={\displaystyle \frac{\mathrm{}}{m_0c}}\mathrm{arccos}{\displaystyle \frac{\sigma }{\pi }}`$ (30)
and from equation(30), we can represent $`\sigma `$ by $`x_3`$:
$$\sigma =\pi \mathrm{cos}(\frac{m_0c}{\mathrm{}}x3)$$
(31)
Where $`\lambda (V_{10})`$ is a variable the dependent on magnitude of vector field $`V_{10}`$. The equations of motion for world line $`\varphi `$ be:
$`x_2=\lambda (V_{20})\mathrm{arccos}{\displaystyle \frac{\varphi }{\pi }}`$ (32)
$`x_3={\displaystyle \frac{\mathrm{}}{m_0c}}\mathrm{arccos}{\displaystyle \frac{\varphi }{\pi }}`$ (33)
Where $`\lambda (V_{20})`$ is a variable the dependent on magnitude of vector field $`V_{20}`$. And from equation(33), we can represent $`\sigma `$ by $`x_3`$:
$$\varphi =\pi \mathrm{cos}(\frac{m_0c}{\mathrm{}}x3)$$
(34)
The equations of motion for world line $`\theta `$ be:
$$x_3=\frac{\mathrm{}}{m_0c}\mathrm{arccos}\frac{\tau }{\pi }$$
(35)
And:
$$\tau =\pi \mathrm{cos}(\frac{m_0c}{\mathrm{}}x3)$$
(36)
In the reference frame with speed of particle $`u_3>0`$ toward direction of $`x_3`$, by using Lorentz transformation: we have:
$$\sigma =\pi \mathrm{cos}\frac{1}{\mathrm{}}(Etp^ix_i)$$
(37)
$`\varphi =\pi \mathrm{cos}{\displaystyle \frac{1}{\mathrm{}}}(Etp^ix_i)`$ (38)
$`\theta =\pi \mathrm{cos}{\displaystyle \frac{1}{\mathrm{}}}(Etp^ix_i)`$ (39)
When $`m_0`$ is zero, above equations becomes the equations of photon. Equations(37)(39)(39) can also be written as:
$$\sigma =a_1e^{\frac{i}{\mathrm{}}p^ix_i}+a_1^{}e^{\frac{i}{\mathrm{}}p^ix_i}$$
(40)
$$\varphi =a_2e^{\frac{i}{\mathrm{}}(p^ix_i)}+a_2^{}e^{\frac{i}{\mathrm{}}(p^ix_i)}$$
(41)
$$\tau =a_3e^{\frac{i}{\mathrm{}}(p^ix_i)}+a_3^{}e^{\frac{i}{\mathrm{}}(p^ix_i)}$$
(42)
where $`a_i^{}`$ is complex conjugate of a which is
$$a_i=\frac{\pi }{2}e^{\frac{i}{\mathrm{}}Et}$$
(43)
$`a_i`$ is equivalent to annihilation operator of scalar field, $`a^{}`$ is creation operator. The model assume that vector $`V_1`$, $`V_2`$ perpendicular to the direction of speed of particle
$$\stackrel{}{V_1}=\alpha _e\sigma \stackrel{}{e_1}$$
(44)
$$\stackrel{}{V_2}=\alpha _b\varphi \stackrel{}{e_2}$$
(45)
where $`\alpha _{e(b)}`$ is constant.
We can also see that the components of world lines on $`x_3`$ axis (the direction of speed) gives wave functions for vector field; the components of world lines on $`x_1`$ and $`x_2`$ axis gives the mangnitude of vector field.
### III.3 Lorentz Covariant Vector Fields
The above equations(40)(41)(42) contains three independent motions, they are not invariant under Lorentz transformation. To obtain Lorentz invariant vector fields: using the new fields $`A_\alpha `$ ($`\alpha =0\mathrm{}3`$), where
$`\widehat{A_0}+_{x_0}\stackrel{}{\widehat{A}}=\sigma \stackrel{}{e_1}`$ (46)
$`\times \stackrel{}{\widehat{A}}=\varphi \stackrel{}{e_2}`$ (47)
$`_{x_5}\stackrel{}{\widehat{A}}=\theta \stackrel{}{e_3}`$ (48)
where $`\widehat{A}(\widehat{A_0},\stackrel{}{\widehat{A}})`$ is 4-vector which is Lorentz invariant, $`\stackrel{}{e_i}`$ is unit vector point to the same direction of $`x_i`$ axis. When $`m_0`$ is zero, then all $`_{x_4}`$ becomes zero, above equations becomes the equations of photon.
In paper chen2 chen3 , I approved that by chosen time-space metric:
$$\left(\widehat{g}_{AB}\right)=\left(\begin{array}{cc}g_{\alpha \beta }+\widehat{A}_\alpha \widehat{A}_\beta & \widehat{A}_\alpha \\ 1& \widehat{A}_\beta \\ & 1\end{array}\right)$$
(49)
where A, B equals 0โฆ5. $`\alpha `$, $`\beta `$ equals 0โฆ3. And put above metric into Einstein field equation(15), one can get the same quantum field equations of vector field:
$$\frac{1}{4}\widehat{F}_{\alpha \beta }\widehat{F}^{\alpha \beta }\frac{1}{2}m_0^2\widehat{A}_\alpha \widehat{A}^\alpha =0$$
(50)
and Proca equation:
$$_\alpha \widehat{F}_{\alpha \beta }+m_0^2\widehat{A}_\beta =0$$
(51)
Here we obtain field equations for quantum vector field by using pure classical method (general relativity equations). When $`m_0=0`$, above equations become Maxwell-equation for free photon in vacuum.
Under time-space metric, we can write world line as:
$$ds^2=dx_\alpha dx^\alpha dx_5dx^5+(\widehat{A}_\alpha dx^\alpha +dx_4)^2$$
(52)
where $`\alpha `$ is 0โฆ3. From equation(52) we can see that, if we choose 2nd time coordinate as
$$dx_{new}^4=\widehat{A}_\alpha dx^\alpha +dx^4$$
(53)
We got local flat time-space.
For general boson particle with spin $`>0`$, we can get the field equations by using similar methods: i.e. first, create three independent motion vectors by three independent proper time. Then convert them to 4-vectors.
### III.4 Lorentz Covariant Vector Fields for general electro-magnetic fields
When $`m_0=0`$, equation(50) is only satisfied for free plane wave photon. Generally, it is not satisfied by regular electro-magnetci fields: for instance, static electric field. But static electric field is not particle. There are two matters in relativity: particle and curved space-time. Static electric field can be treated as curved space-time. The 6-dimensional time-space metric for regular electro-magnetic fields is:
$$\left(\widehat{g}_{AB}\right)=\left(\begin{array}{cc}g_{\alpha \beta }+A_\alpha A_\beta & A_\alpha \\ 1& A_\beta \\ A_5A_\beta & A_51+A_5A_5\end{array}\right)$$
(54)
where A, B equals 0โฆ5. $`\alpha `$, $`\beta `$ equals 0โฆ3. put above metric into Einstein field equation(15), we get:
$$\frac{1}{2}F_{\alpha \beta }F^{\alpha \beta }_\alpha A_5^\alpha A^5=0$$
(55)
$$_\alpha F_{\alpha \beta }=0$$
(56)
and equation
$$_\alpha ^\alpha A_5=0$$
(57)
For static electric field, let
$$A_5=\frac{e}{r}$$
(58)
$`A_5`$ satisfied equations (55) and (57).
## IV 3-proper time model for free single particle with half integer spin
In section II and III, the trajectories of particle with 2nd and 3rd proper time are straight lines. In this section, we will see that: when the trajectories of particle with 2nd and 3rd proper time become circles, we will obtain equations of motion for particle with the half-integer spin.
For free particle with half spin, in rest frame (i.e. in the fram $`t=\tau `$), we build 3-proper time model for $`\tau `$, $`\sigma `$, $`\varphi `$ as below:
1) Choosing a reference direction for the rotations by $`\sigma `$ and $`\varphi `$. Let trajectory of $`\sigma `$ parallel to $`x_1x_2`$ plane, $`x_3`$ perpendicular to world lines $`\sigma `$ (this is equivalent to choose $`x_3`$ representation for spin).
2) Let trajectory of $`\varphi `$ be a rotation from $`x_3`$ axis to $`\stackrel{}{s}`$ where $`\stackrel{}{s}`$ is unit vector in $`x_1x_2`$ plane.
The circles of rotation must be small. It is resonable to assume that the diameter of circle equals $`h/m_0c`$ , so the radius of circle is
$$r_0=\frac{h}{2m_0c}$$
(59)
Let the equations of motion of $`\sigma `$ be:
$`x_1=r_0\mathrm{cos}{\displaystyle \frac{m_0c}{\mathrm{}}}\sigma `$ (60)
$`x_2=r_0\mathrm{sin}{\displaystyle \frac{m_0c}{\mathrm{}}}\sigma `$ (61)
We can see that $`x_1`$ and $`x_2`$ has a phase difference $`\pi /2`$ . Let the equations of motion for $`\varphi `$ be:
$`x_3=r_0\mathrm{cos}{\displaystyle \frac{m_0c}{4\mathrm{}}}\varphi `$ (62)
$`x_s=r_0\mathrm{sin}{\displaystyle \frac{m_0c}{4\mathrm{}}}\varphi `$ (63)
let the equations of motion for $`\tau `$ be:
$`x_0=\tau (1+\mathrm{cos}({\displaystyle \frac{m_0c}{\mathrm{}}}\tau ))`$ (64)
$`x_5=\tau \mathrm{sin}({\displaystyle \frac{m_0c}{\mathrm{}}}\tau )`$ (65)
The motions of world lines $`\sigma `$ and $`\varphi `$ are oriented, suppose that on world line $`\varphi `$ at point $`\alpha _0`$, particleโs moving direction of world line $`\sigma `$ is perpendicular to $`\varphi `$ towards its right, and assume the direction of rotation towards positive $`x_3`$ direction. When it moves to the opposite point of $`\alpha _0`$ at sphere, its motion of world line $`\sigma `$ is still perpendicular to $`\varphi `$ towards its right, but the direction of rotation towards negative $`x_3`$ direction which is opposite to previous motion. So, to build a stable two-proper time motion on spherical surface, we have to limit all the motions on half sphere. Thatโs why we have factor $`\frac{1}{4}`$ inside equations (63) (63)(65): $`\varphi `$ is from 0 to $`2\pi `$, 1/4 factor to make the rotation from $`x_3`$ to $`x_s`$ with angel from 0 to $`\pi /2`$.
For the reference frame with velocity of particle $`u>0`$, define one vector for each motion. For world line $`\sigma `$:
$$\stackrel{}{V_1}=(\stackrel{}{e_1}+i\stackrel{}{e_2})e^{\frac{i}{\mathrm{}}(Etp^ix_i)}$$
(66)
where i comes from the $`\pi /2`$ phase difference of $`x_1`$ and $`x_2`$ since world line $`\sigma `$ is a rotation. For world line $`\varphi `$:
$$\stackrel{}{V_2}=\stackrel{}{e_3}e^{\frac{i}{\mathrm{}}(Etp^ix_i)}$$
(67)
We ignored the contribution for $`\stackrel{}{e_s}`$ since the direction of $`\stackrel{}{e_s}`$ is from 0 to $`2\pi `$, the sum of contribution will be zero. For world line $`\tau `$, define:
$$\stackrel{}{V_3}=\stackrel{}{e_\tau }+(\stackrel{}{e_0}+i\stackrel{}{e_5})e^{\frac{i}{\mathrm{}}(Etp^ix_i)}$$
(68)
where $`\stackrel{}{e_\tau }`$ is the direction of speed of particle. In rest frame, it is the same as $`x_0`$ axis. In previous section, the components of each world lines on direction of speed of particle give wave-function. Here we let:
$`\psi _1=\stackrel{}{V_1}\stackrel{}{e_n}=\mathrm{sinh}{\displaystyle \frac{\alpha }{2}}({\displaystyle \frac{u_1}{u}}+i{\displaystyle \frac{u_2}{u}})e^{\frac{i}{\mathrm{}}(Etp^ix_i)}`$ (69)
$`\psi _2=\stackrel{}{V_2}\stackrel{}{e_n}=\mathrm{sinh}{\displaystyle \frac{\alpha }{2}}{\displaystyle \frac{u_3}{u}}e^{\frac{i}{\mathrm{}}(Etp^ix_i)}`$ (70)
$`\psi _3=\stackrel{}{V_3}\stackrel{}{e_n}=\mathrm{cosh}{\displaystyle \frac{\alpha }{2}}e^{\frac{i}{\mathrm{}}(Etp^ix_i)}+1`$ (71)
where $`u=\sqrt{u_1^2+u_2^2+u_3^2}`$ is the speed of particle; $`\stackrel{}{e_n}`$ is the direction of $`\widehat{n}`$ shown in Fig2.
$$\stackrel{}{e_n}=\stackrel{}{e_0}+\stackrel{}{e_\tau }$$
(72)
and
$`\mathrm{cosh}\alpha ={\displaystyle \frac{1}{\sqrt{1\frac{u^2}{c^2}}}}`$ (73)
$`\mathrm{sinh}\alpha ={\displaystyle \frac{\frac{u}{c}}{\sqrt{1\frac{u^2}{c^2}}}}`$ (74)
It is easy to see that:
$`\psi _1={\displaystyle \frac{p_1+ip_2}{m_0}}e^{\frac{i}{\mathrm{}}(Etp^ix_i)}=C_0\psi _4^D`$ (75)
$`\psi _2={\displaystyle \frac{p_3}{m_0}}e^{\frac{i}{\mathrm{}}(Etp^ix_i)}=C_0\psi _3^D`$ (76)
$`\psi _3={\displaystyle \frac{p_0}{m_0}}e^{\frac{i}{\mathrm{}}(Etp^ix_i)}+1=C_0\psi _0^D+1`$ (77)
where $`\psi _0^D`$,$`\psi _2^D`$,$`\psi _3^D`$ are three non-zero components of the solution of Dirac equation in $`x_3`$ representation, with positive energy and spin. $`C_0`$ is normalization constant. The constant item 1 can be ignored. Therefore, $`\psi _i`$ corresponding to the three non-zero components of Dirac wave-function.
Now we define 5-vector $`\widehat{K}`$, the first 4 components are:
$`\widehat{K_0}=C\psi _3e^{im_0x_5}\widehat{K_1}=C\psi _1e^{im_0x_5}`$
$`\widehat{K_2}=iC\psi _1\varphi _3e^{im_0x_5}\widehat{K_3}=C\psi _2e^{im_0x_5}`$ (78)
where C is constant. If we modify equation(72) as:
$$\stackrel{}{e_n}=\stackrel{}{e_0}+\stackrel{}{e_\tau }+\stackrel{}{e_5}$$
(79)
Then
$$\psi _3=\stackrel{}{V_3}\stackrel{}{e_n}=\mathrm{cosh}\frac{\alpha }{2}(1+i)e^{\frac{i}{\mathrm{}}(Etp^ix_i)}+1$$
(80)
The extra item $`ie^{\frac{i}{\mathrm{}}(Etp^ix_i)}`$ gives us 5th components of $`\widehat{K}`$:
$$\widehat{K_5}=Ce^{\frac{i}{\mathrm{}}(Etp^ix_im_0x_5)}$$
(81)
The geometry of 6-dimensional time-space for Fermion is chen3 :
$$\left(\widehat{g}_{AB}\right)=\left(\begin{array}{cc}g_{\alpha \beta }+\widehat{K}_\alpha \widehat{K}_\beta & \widehat{K}_\alpha \widehat{K}_\alpha \widehat{K}_5\\ \widehat{K}_\beta & \mathrm{\hspace{0.33em}\hspace{0.33em}1}\widehat{K}_5\\ \widehat{K}_5\widehat{K}_\beta & \widehat{K}_51+\widehat{K}_5\widehat{K}_5\end{array}\right)$$
(82)
Paper chen2 chen3 show that: put above $`\widehat{g}_{AB}`$ into Einstein field equation, We can get Dirac equation. Under metric equation (82), the interval $`ds`$ can be written as:
$$ds^2=dx_\alpha dx^\alpha dx_5dx^5+(\widehat{K}_\alpha dx^\alpha +\widehat{K}_5dx^5+dx^4)^2$$
(83)
we can see that, if we choose 2nd local time coordinate $`x_4`$ as
$$dx_{new}^4=\widehat{K}_\alpha dx^\alpha +\widehat{K}_5dx^5+dx^4$$
(84)
We obtain local flat time-space. I.e. Local inertia reference frame exists in this 6-dimesional time-space.
## V Interpretation of Bose-Einstein Statistics and Fermi-Dirac Statistics
Consider an event A: particle 1 and particle 2 interact each other at spatial point $`O(x_1,x_2,x_3)`$. In classical 4-dimensional time-space, event A can happen if and only if both particle meet at $`O(x_1,x_2,x_3)`$ at the same time t. In another words: their world lines across each other at point $`O(t,\stackrel{}{x})`$. In 3+3 dimensional time-space, event A can happen if and only if the world lines of particle 1 and particle 2 across each other at $`O(t,\stackrel{}{x},x_4,x_5)`$. If world lines of particle 1 and world lines of particle 2 parallel each other, then there is no interaction between them (Here we neglect the case of 3rd particle involving interaction).
For Boson, Fig3. shows a distribution of world lines of particle 1 and particle 2, where their world lines $`\tau `$, $`\sigma `$ and $`\varphi `$ parallel each other respectively, particle 2 always slightly โlaterโ than particle 1. Because particle is โpoint particleโ, the size of particle is infinitesimal, then inside a spatial volume with the size of one wave-length, we can put infinite particles inside where their world lines are parallel to each other. No interactions will be occurs inside this volume. We get Bose-Einstein condensation. Also because we can not synchronize the 2nd and 3rd time dimension, then we can not identify each particle individually โ they are identical particles.
For Fermion, section IV shows that, the world lines $`\sigma `$ and $`\varphi `$ filled half spherical surface with diameter equals $`h/m_0c`$. In the spatial volume with the size of $`h/m_0c`$, we can not find two parallel spherical without crossing each other. But we can have another particle on the other half spherical surface with $`x_1x_2`$ rotation point to negative $`x_3`$ direction. That is, within the volume with the size of $`h/m_0c`$, there can only have two non-interactive Fermions with rotations towards opposite directions. This is Fermi-Dirac statistics.
Above interpretation is based on the equations for free Boson and Fermion. I believe that the basic concepts can be extended to cases of non-free particles.
## VI Summary and Discussions
Now we are back to the initial topic of this paper: what is spin? As we have seen in previous sections
1) For spinless particle, the projections of world lines $`\tau `$, $`\sigma `$ and $`\varphi `$ on 3-dimensional space are coincident on the same spatial straight lines.
2) For particle with integer spin $`>0`$, the projections of world lines $`\tau `$, $`\sigma `$ and $`\varphi `$ on 3-dimensional space are three separate lines. Both $`\sigma `$ and $`\varphi `$ have spatial components which perpendicular to spatial direction of world line $`\tau `$.
3) For particle with half-integer spin, world lines $`\sigma `$ and $`\varphi `$ become rotations on a spherical surface (like a real spin).
As the summary, the spin is inner 6-dimensional geometry properties of particle, it determines the motion of three-proper time of the particle.
This paper only study the case of free particle, it will be interesting to see the behavior of 3-proper time motions for interactive particles. I believe that 3-proper time model will not only reproduce the physics of quantum field theory, but also give more detail and more physics for quantum particle. To find the difference between 3-proper time model and quantum physics will be another interesting topic.
Finally, Iโd like to say that, there will be many different 3-proper time models. In future, there will be 3-proper time models which may better fit into quantum physics than the model proposed by this paper. Again, the purpose of this paper is to show that, 3-proper time model provides a possible way to interpret the physical effect of spin particles and basic quantum physics.
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# Photometric observations from theoretical flip-flop models
## 1 Introduction
In many active stars the spots concentrate on two permanent active longitudes which are $`180^{}`$ apart. In some of these stars the dominant part of the spot activity changes the longitude every few years. This so-called flip-flop phenomenon was first reported by Jetsu et al. (jetsu1 (1991), jetsu2 (1993)) in the single, late type giant FK Com. Berdyugina & Tuominen (ber\_tuo (1998)) reported periodic flip-flops between permanent active longitudes in four RS CVn binaries. Their results were confirmed in the case of II Peg by Rodonรฒ et al. (rod (2000)). The persistent active longitude structures and flipping between two active longitudes have also been reported over the years based on photometric observation (e.g. Berdyugina et al. ber\_lqhya (2002); Korhonen et al. kor4 (2002); Jรคrvinen et al jarv (2005)) and on Doppler images (Berdyugina et al. ber2 (1998); Korhonen et al. kor\_ff (2001)). After its discovery in cool stars, the flip-flop phenomenon has also been reported in the Sun (Berdyugina & Usoskin ber\_sun (2003)). A review on the flip-flop phenomenon in cool stars and the Sun is given by Berdyugina (ber\_ff (2004)).
In order to explain this phenomenon, a non-axisymmetric dynamo component, giving rise to two permanent active longitudes 180ยฐ apart, is needed together with an oscillating axisymmetric magnetic field. Fluri & Berdyugina (flu\_ber (2004)) suggest also another possibility with a combination of stationary axisymmetric and varying non-axisymmetric components. Unfortunately, no simple dynamo mechanism is yet known for such a configuration. There are models with anisotropic $`\alpha `$-effect or a weak differential rotation, which could produce non-axisymmetric components, but only recently Moss (moss1 (2004), moss2 (2005)) reported coexisting mixed components with a differential rotation that depends on the distance to the rotation axis. This differential rotation configuration is a very plausible state for fast rotators. Also, a weak non-axisymmetric field coexisting with a dominant axisymmetric field, assuming a solar-like rotation law, was found by Moss (moss99 (1999)). This solution was not analysed for flip-flops, as the possibility of them being present on the Sun was not being discussed at that time. Flip-flop solutions for a rotation law similar to the solar one and anisotropic $`\alpha `$ have also been found by Elstner & Korhonen (elstner (2005)).
According to the calculations by Elstner & Korhonen (elstner (2005)) the shift of the spots in a flip-flop event is 180ยฐ only in some cases, mainly for the stars with thin convective zones. In stars with thick convective zones they found a shift that is close to 90ยฐ. Similar results were reported by Moss (moss2 (2005)). Recently, Olรกh et al. (olah (2005)) re-analysed some of the old photometric observations of FK Com and found a flip-flop event in which spots on both active longitudes vanished briefly, and one of the new spots appeared on an old active longitude and the other one 90ยฐ away from that. This is the first evidence suggesting that flip-flops where the spots shift only by 90ยฐcan also occur.
In this paper we use the model calculations presented by Elstner & Korhonen (elstner (2005)) and convert them into synthetic photometric observations. This is used to investigate the expected long-term photometric behaviour of active stars showing the flip-flop phenomenon. Hopefully, this will help us in identifying new stars exhibiting this intriguing phenomenon. At the moment only few stars showing it are known and no statistically significant correlation between the stellar parameters and the flip-flop phenomenon can therefore be carried out.
## 2 Model
The model consists of a turbulent fluid in a spherical shell of inner radius $`r_{\mathrm{in}}`$ and outer radius $`r_{\mathrm{out}}`$.
We solve the induction equation
$$\frac{B}{t}=\mathrm{curl}(\alpha _qB\eta _\mathrm{T}\mathrm{curl}B),$$
(1)
in spherical coordinates ($`\mathrm{r},\theta ,\phi `$) for an $`\alpha ^2\mathrm{\Omega }`$-dynamo. Here we used a solar-like rotation law in the corotating frame of the core
$`\mathrm{\Omega }(r,\theta )={\displaystyle \frac{1}{2}}\mathrm{\Omega }_0\left[1+\mathrm{erf}\left({\displaystyle \frac{rr_c}{d_1}}\right)\right](\mathrm{\Omega }_s\mathrm{\Omega }_c)`$ (2)
where $`\mathrm{\Omega }_s=\mathrm{\Omega }_{eq}a\mathrm{cos}^2\theta `$ with $`\mathrm{\Omega }_{eq}/2\pi =460.7\mathrm{nHz}`$, $`\mathrm{\Omega }_c/2\pi =432.8\mathrm{nHz}`$ and $`a/2\pi =125.82\mathrm{nHz}`$ are used. Normalising length with stellar radius $`R_{}`$ and time with diffusion time $`t_d=R_{}^2/\eta `$ we can define the dimensionless dynamo numbers
$$C_\alpha =\alpha _0R_{}/\eta $$
(3)
and
$$C_\mathrm{\Omega }=(\mathrm{\Omega }_{eq}\mathrm{\Omega }_c)R_{}^2/\eta .$$
(4)
For $`R_{}=R_{}`$, the diffusivity $`\eta =510^{12}\mathrm{cm}^2\mathrm{s}^1`$ and $`\mathrm{\Omega }_0=1`$ we get $`C_\mathrm{\Omega }`$=172. Notice that a definition with $`\mathrm{\Omega }_{eq}\mathrm{\Omega }_{pol}`$ gives $`C_\mathrm{\Omega }=780`$. Because of these different definitions of $`C_\mathrm{\Omega }`$ in the literature we describe the strength of the differential rotation in our models with the value of $`\mathrm{\Omega }_0`$.
In order to identify the lifetime and the maximal possible pole to equator difference of the angular velocity for a flip-flop solution also for models with isotropic $`\alpha `$, we performed several calculations with models similar to those presented in Moss (moss2 (2005)). Here we used the differential rotation law (Eq. 2) with $`\mathrm{\Omega }_s=ar^2\mathrm{sin}^2\theta `$ and $`\mathrm{\Omega }_c=0`$. The same normalisation was used.
Only the symmetric part
$`\alpha _{rr}=\alpha _0\mathrm{cos}\theta (1.2\mathrm{c}\mathrm{o}\mathrm{s}^2\theta )`$
$`\alpha _{\theta \theta }=\alpha _0\mathrm{cos}\theta (1.2\mathrm{s}\mathrm{i}\mathrm{n}^2\theta )`$
$`\alpha _{\phi \phi }=\alpha _0\mathrm{cos}\theta `$
$`\alpha _{r\theta }=\alpha _{\theta r}=2\alpha _0\mathrm{cos}^2\theta \mathrm{sin}\theta `$ (5)
of the $`\alpha `$-tensor is included (cf. Rรผdiger et al rued (2002)). For the isotropic model $`\alpha _{r\theta }=\alpha _{\theta r}=0`$ and the diagonal terms $`\alpha _{ii}=\alpha _0\mathrm{cos}\theta `$. In order to saturate the dynamo we choose a local quenching of
$$\alpha _q=\frac{\alpha }{1+B^2/B_{\mathrm{eq}}^2}.$$
(6)
The inner boundary is a perfect conductor at $`r_{\mathrm{in}}=0.3`$ and the outer boundary resembles a vacuum condition, by including an outer region up to 1.2 stellar radii into the computational grid with 10 times higher diffusivity. At the very outer part the pseudo vacuum condition (tangential component of the magnetic field and vertical component of the electric field vanish on the surface) is used. In order to see the influence of the thickness of the convection zone we have chosen $`r_\mathrm{c}=0.4`$ for a thick (results shown in Fig. 3a) and $`r_\mathrm{c}=0.7`$ for a thin (Fig. 3b) convection zone. Below the convection zone from $`r_{\mathrm{in}}`$ to $`r_\mathrm{c}`$ the diffusivity was reduced by a factor 1000 and $`\alpha `$ was set to zero. The critical $`C_\alpha `$ and periods for the models are given in Table 1.
In all the models we used a normalised $`\eta =0.5`$, $`C_\alpha `$ and $`\mathrm{\Omega }_0`$ for the thin and thick convection zone models are given in Table 2 and the parameters for the model used in Fig. 4 are $`\mathrm{\Omega }_0=0.3`$, $`r_{\mathrm{in}}=0.2`$ and $`r_\mathrm{c}=0.3`$ $`C_\alpha =8`$.
## 3 From the model to observations
For converting the possible spot pattern from the model calculation into synthetic photometric observations, we first have to decide which value of the magnetic pressure on the stellar surface results in a spot. For doing this a three temperature model was chosen in which the values of magnetic pressure that are $`70`$ % of the maximum value are considered to form the โumbraโ and the values $`<70`$ % and $`30`$ % of the maximum form the โpenumbraโ. The values $`<30`$ % of the maximum denote the unspotted surface. For investigating the long-term changes in the spot strength the maximum value of the magnetic pressure was taken from the whole run, not from the individual maps.
A typical star showing flip-flops is a cool giant or a zero age main sequence object. For describing the realistic spot temperatures on such stars, 5000 K was chosen as the unspotted surface temperature and 3500 K & 4250 K as the umbral and penumbral temperatures, respectively. After the assignment of the spot temperatures to the magnetic pressure maps, synthetic light-curves were calculated from the maps. The limb-darkening coefficient from Al-Naimyi (aln (1978)) for 5000 K at the central wavelength of the Johnson V band was used for all the three temperatures. Fig. 1 shows the magnetic pressure map obtained from the dynamo calculations, the corresponding spot configuration and the synthetic light-curve calculated from the temperature map.
## 4 Results
The results from the thick and thin convection zone flip-flop models, that were first discussed in Elstner & Korhonen (elstner (2005)), have been converted into light-curves as described in the previous section. The model parameters are given in Table 2.
In Fig. 2 an example of a sequence of temperature maps exhibiting a migrating spot pattern and a flip-flop event is shown. The features are symmetric with respect to the equator. It is clearly seen that in this thick convection zone model the spot shift in the flip-flop is $`90^{}`$ and not 180ยฐ, as is more commonly seen in the observations.
For investigating the long-term photometric behaviour obtained from the models, calculations were done starting around 80 diffusion times, running 50000 timesteps (one timestep is $`1.910^5`$ diffusion times or approximately 1/5 days for our models) and calculating a map of magnetic pressure at the surface every 100 steps. These maps were then converted into temperature maps and synthetic light-curves, and the light-curves were plotted against time to see the long-term behaviour. Figs. 3a & b show the calculated light-curve behaviour for three different inclination angles for the thick and thin models, respectively.
In the case of the thick model (Fig. 3a) we see that the minimum magnitude is very strongly modulated during the flip-flop cycle, whereas the maximum magnitude is much less affected, except when viewed from small inclinations. The minimum and maximum in the maximum magnitude occur around the time of the minimum and maximum in the minimum magnitude. This behaviour is the same as seen by Fluri & Berdyugina (flu\_ber (2004)) in their first case (sign-change of the axisymmetric component). The inclination of the star affects mainly the amplitude of the variation, as it determines how close the spots are to the centre of the visible disk (location of the maximum effect on the light-curve). In the case of high-latitude spots, as seen in these models, this effect totally dominates over the other inclination effect, i.e. how much of the spots on the โsouthernโ hemisphere are visible. The changes in the inclination affect the behaviour of the maximum magnitude more strongly than the minimum magnitude.
The behaviour seen in the thin convection zone model (Fig. 3b) is different from the one seen in the thick model. The amplitude of the variation is much smaller because the active longitudes are further apart than in the thick case. It is also worth noting, that in this case the variation in the maximum and minimum magnitudes is very different from the thick case; here the maximum of the maximum magnitude occurs near the minimum of the minimum magnitude. There is a small shift towards the later timesteps for the maximum of the maximum magnitude in comparison to the minimum of the minimum magnitude.
## 5 Discussion
### 5.1 Dependence of the solution on the model parameters
The differential rotation needed for a flip-flop solution with a period of about 5 years is rather small. For a solar sized star the latitudinal difference in $`\mathrm{\Omega }`$ ($`\mathrm{\Delta }\mathrm{\Omega }`$) should be about 10% of the solar value, independent of the global rotation. For $`R_{}=2R_{}`$ the $`\mathrm{\Delta }\mathrm{\Omega }`$ can be about 40%. The situation seems similar in all cases, i.e. for isotropic $`\alpha `$ tensor used together with a differential rotation depending on the distance to the rotation axis and with an anisotropic $`\alpha `$ tensor used with both a solar-like rotation law and axis distance dependent rotation law. An example of the time evolution of the magnetic field energy in the m=0,1 components with isotropic $`\alpha `$, $`\mathrm{\Delta }\mathrm{\Omega }`$ 30% of the solar value, and a rotation law depending on the axis distance (cf. Moss 2005) is shown in Fig. 4. As can be seen, with relatively strong differential rotation the non-axisymmetric component, m=1, is initially excited, but the mixed component solution survives only approximately 5 diffusion times.
Increasing the diffusivity gives a strong flip-flop phenomenon also for higher pole to equator differences of $`\mathrm{\Omega }`$, but with a smaller flip-flop period. The flip-flop solutions appear preferential for a positive $`\alpha `$ in the northern hemisphere in the case of radially increasing $`\mathrm{\Omega }`$ at the equator. This leads to a poleward migration of the spots. Observations indicate solar-like equatorward migration pattern in solar-like stars (see e.g. Katsova et al. kats (2003)), but there is also a detection of poleward migration of the spots in the RS CVn binary HR 1099 (Vogt et al. vogt (1999); Strassmeier & Bartus str\_bar (2000)).
### 5.2 Comparing the synthetic light-curves with the observations
Many active stars show long-term light-curves where time periods with small and large amplitude in the photometry alternate, as seen in our flip-flop models. For the behaviour seen in the thick model (Fig. 3a) a good stellar counterpart is DX Leo (see e.g. Messina & Guinan mes (2002)). It is easier to find stellar counterparts for the thin case (Fig. 3b). Some examples, like LQ Hya and EI Eri, can be seen for instance in Olรกh & Strassmeier (olah\_str (2002)). The fact that the thin case seems to be dominant implies that the flip-flops where the spots shift 180ยฐ are more common.
As seen in the case of FK Com (Olรกh et al olah (2005)), some stars can show both 90ยฐ and 180ยฐ shifts in the spots during a flip-flop event. For investigating what alternating 90ยฐ and 180ยฐ flip-flops would look like in the long-term photometry, we combine the calculated light-curves from the models showing 90ยฐ (from the thick model) and 180ยฐ (from the thin model) flip-flops, taking alternatively one 90ยฐ flip-flop and one 180ยฐ flip-flop. The result of combining the two types of flip-flops is shown in Fig. 5a. The long-term light-curve behaviour produced by this is similar, but not identical, to the second case of Fluri & Berdyugina (flu\_ber (2004)), which shows alternatively small and large amplitude changes that are symmetric with respect to the mean magnitude, i.e. the maximum of the maximum magnitude occurs at the time of the minimum of the minimum magnitude. Our combination of the models with different spot shifts in the flip-flop event is symmetrical only during the smaller amplitude phase. During the larger amplitude phase the minimum of the maximum magnitude occurs at the time of the minimum of the minimum magnitude. For example, the light-curve of $`\sigma `$ Gem (see Fig. 5b) shows indications of this kind of behaviour.
It is often difficult to see in the real observations the activity pattern caused by the flip-flop behaviour. This is partly due to the years long time series of observations needed to see the pattern and partly due to the solar-like cyclic changes in the over-all activity level of many active stars. Quite drastic changes in the brightness of some stars can be seen on top of the possible flip-flop signature (see e.g. HK Lac and IL Hya in Olรกh & Strassmeier olah\_str (2002)) and these large changes are likely to mask the patterns caused by the flip-flops.
## 6 Summary
We have synthesised photometry from Dynamo calculations exhibiting flip-flop behaviour. This was done for investigating the long-term changes in the photometric behaviour seen over several flip-flop cycles. On the whole, the activity patterns discussed in this paper imply flip-flop phenomenon and stars showing these patterns should be further investigated for checking if they really show flip-flops. A statistically significant sample of stars is needed for deeper understanding of this phenomenon. Further, more effort should be put to measuring the $`\mathrm{\Delta }\mathrm{\Omega }`$, meridional flow, and latitude migration of the spots on different types of stars. All these parameters have important implications for the dynamo calculations.
###### Acknowledgements.
We would like to thank the referee Dr. Moss for his very useful comments on this paper. HK acknowledges the support from the German *Deutsche Forschungsgemeinschaft, DFG* project KO 2320/1.
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# The Extended Chandra Deep Field-South Survey: X-ray Point-Source Catalog
## 1 Introduction
Wide-area X-ray surveys have played a fundamental role in understanding the nature of the sources that populate the X-ray universe. Early surveys like the Einstein Medium Sensitivity Survey (Gioia et al., 1990), the ROSAT International X-ray/Optical Survey (Ciliegi et al., 1997), and the ASCA Large Sky Survey (Akiyama et al., 2000) showed that the vast majority of the bright X-ray sources are active galactic nuclei (AGNs). More specifically, shallow wide-area surveys in the soft (0.5โ2.0 keV) X-ray band yield mostly unobscured, broad-line AGNs, which are characterized by a soft X-ray spectrum with a photon index of $`\mathrm{\Gamma }`$ 1.9 (Nandra & Pounds, 1994). In contrast, deep X-ray surveys โ particularly surveys that make use of the unprecedented, sub-arsecond spatial resolution of the Chandra X-ray Observatory โ find AGN with harder X-ray spectra ($`\mathrm{\Gamma }`$ 1.4) at fainter fluxes, more like the hard spectrum of the X-ray background.
Deep Chandra surveys have thus opened a new vista on resolving the X-ray background and identifying the role and evolution of accretion power in all galaxies. The cosmic X-ray background is now almost completely resolved ($``$ 70โ90%) into discrete sources in the deep, pencil beam surveys like the Chandra Deep Fields (CDF-N, Brandt et al., 2001; CDF-S, Giacconi et al., 2002). To understand the composition of the sources that make up the X-ray background, population synthesis models have been constructed (Madau et al., 1994; Comastri et al., 1995; Gilli et al., 1999; Gilli et al., 2001; Treister & Urry, 2005) which typically require approximately 3 times as many obscured AGN as traditional Type 1 (unobscured) AGN.
While the deep fields provide the deepest view of the X-ray universe and have generated plentiful AGN samples at lower luminosities, the small area covered by pencil-beam surveys means luminous sources are poorly sampled. In an attempt to determine the luminosity function of X-ray emitting AGN up to $`z5`$, as well as to leverage existing deep multiwavelength data in the extended $`30^{}\times 30^{}`$ field centered on the CDF-S, the region surrounding the CDF-S was recently observed by $`Chandra`$. Covering $``$ 1100 square arcminutes ($``$ 0.3 deg<sup>2</sup>), the Extended Chandra Deep Field-South (ECDFS) survey is the largest Chandra survey field at this depth ($``$ 230 ks), and is the second deepest and widest survey ever conducted in the X-rays (the XMM-Newton survey of the Lockman Hole is deeper in the hard band and has $``$ 30% more area; Hasinger 2004).
In this paper, we present the X-ray catalog for the ECDFS and the number counts in two energy bands. In subsequent papers, we will present the optical and near-IR properties of these X-ray sources, including first results from our deep optical spectroscopy campaign obtained as part of the one-square-degree MUltiwavelength Survey by Yale/Chile (MUSYC) (Gawiser et al., 2005).
In Section 2, we describe our data reduction procedure. In Section 3, we describe the point source detection and astrometry. The X-ray source catalog and basic survey results are presented in Section 4 and the conclusions are given in Section 5. The average Galactic column density along this line of sight for the four pointings is $`9.0\times 10^{19}`$ cm<sup>-2</sup> (Stark et al., 1992). $`H_0=70`$ km s<sup>-1</sup> Mpc<sup>-1</sup>, $`\mathrm{\Omega }_\mathrm{M}=0.3`$, and $`\mathrm{\Omega }_\mathrm{\Lambda }=0.7`$ are adopted throughout this paper which is consistent with the cosmological parameters reported by Spergel et al. (2003). All coordinates throughout this paper are J2000.
## 2 Observations and Data Reduction
### 2.1 Instrumentation and Diary of Observations
All nine observations of the ECDFS survey field were conducted with the Advanced CCD Imaging Spectrometer (ACIS) on-board the Chandra X-ray Observatory<sup>1</sup><sup>1</sup>1For additional information on the ACIS and Chandra see the Chandra Proposersโ Guide at http://cxc.harvard.edu/proposer/POG/html. as part of the approved guest observer program in Cycle 5 (proposal number 05900218 โ PI Niel Brandt; Lehmer et al. 2005). ACIS consists of 10 CCDs, distributed in a 2$`\times `$2 array (ACIS-I) and a 1$`\times `$6 array (ACIS-S). All 4 of the ACIS-I CCDs are front-illuminated (FI) CCDs; 2 of the 6 ACIS-S CCDs are back-illuminated CCDs (S1 and S3). Of these 10 CCDs, at most 6 can be operated at any one time. Table 1 presents a journal of the Chandra observations of the ECDFS. All nine observations were conducted in VERY FAINT mode (See the Chandra Proposersโ Guide, pg. 95) so that the pixel values of the $`5\times 5`$ event island are telemetered rather than just the $`3\times 3`$ event island as in FAINT mode. This telemetry format offers the advantage of further reducing the instrument background after ground processing (see Section 2.2). Observation ids (ObsIds) 5019โ5022 and 6164 also had the ACIS-S2 CCD powered on (see Table 1). However, due to the large off-axis angle of the S2 CCD during these observations, it has a much broader point spread function (PSF) and hence lower sensitivity so we exclude data from this CCD and focus only on data collected from the ACIS-I CCDs. The on-axis CCD for each ACIS-I observation is I3 and the ACIS-I field of view is $`16\stackrel{}{\mathrm{.}}9\times 16\stackrel{}{\mathrm{.}}9`$.
### 2.2 Data Reduction
All data were re-processed using the latest version of the Chandra Interactive Analysis of Observations<sup>2</sup><sup>2</sup>2See http://cxc.harvard.edu/ciao/. (CIAO; Version 3.2.1; released 10 February 2005) software as well as version 3.0.0 of the calibration database (CALDB; released 12 December 2004). We chose to re-reduce all nine datasets rather than simply use the standard data processing (SDP) level 2 event files because we wanted the datasets to be reduced in a consistent manner, and more importantly, we wanted to take advantage of better Chandra X-ray Center (CXC) algorithms to reduce the ACIS particle background by using the information located in the outer 16 pixels of the $`5\times 5`$ event island<sup>3</sup><sup>3</sup>3See โReducing ACIS Quiescent Background Using Very Faint Modeโ, http://cxc.harvard.edu/cal/Acis/Cal\_prods/vfbkgrnd/index.html., as well as executing a new script that is much more efficient at identifying ACIS โhot pixelsโ and cosmic ray afterglow events<sup>4</sup><sup>4</sup>4See http://cxc.harvard.edu/ciao/threads/acishotpixels/.. This new hot pixel and cosmic ray afterglow tool is now implemented in the standard data processing pipeline at the CXC but was not applied in the SDP pipeline of the present observations.
The following procedure was used to arrive at a new level 2 event file. Before applying the new CIAO tool to identify ACIS hot pixels and cosmic ray afterglow events, the pixels identified in the CXC-provided level 1 event file as being due to a cosmic ray afterglow event were reset. An afterglow is the residual charge from the interaction of a cosmic ray in a front-side illuminated CCD frame. Some of the excess charge is captured by charge traps created by the radiation damage suffered early in the mission (see Townsley et al., 2000 and references therein) and released over the next few to a few dozen frames. If these afterglow events are not removed from the data, they can result in the spurious detection of faint sources. To better account for such events, a new, more precise method for identifying afterglow events was developed by the CXC and has now been introduced into the standard data processing pipeline. This CIAO tool, โacis\_run\_hotpix,โ was then run on the reset level 1 event file to identify and flag hot pixels and afterglow events in all 9 ACIS observations of the ECDFS. The last step in producing the new level 2 event file was to run the CIAO tool โacis\_process\_events.โ In addition to applying the newest gain map supplied in the latest release of the CALDB, this tool also applies the pixel randomization and the ACIS charge transfer inefficiency (CTI) correction. The latter corrects the data for radiation damage sustained by the CCDs early in the mission. All of these corrections are part of the standard data processing and are on by default in acis\_process\_events. The time-dependent gain correction is also applied to the event list to correct the pulse-invariant (PI) energy channel for the secular drift in the average pulse-height amplitude (PHA) values. This drift is caused primarily by gradual degradation of the CTI of the ACIS CCDs (e.g., Schwartz & Virani, 2004). Finally, the observation-specific bad pixel map created by โacis\_run\_hotpixโ was supplied, and the option to clean the ACIS particle background by making use of the additional pixels telemetered in VERY FAINT mode was turned on.
Once a new level 2 event file was produced, we applied the standard grade filtering to each observation, choosing only event grades 0, 2, 3, 4, and 6 (the standard ASCA grade set), and the standard Good Time Intervals supplied by the SDP pipeline. We also restrict the energy range to 0.5โ8.0 keV, as the background rises steeply below and above those limits<sup>5</sup><sup>5</sup>5See http://cxc.harvard.edu/contrib/maxim/stowed/.. Lastly, we examined the background light curves for all 9 observations as the ACIS background is known to vary significantly. For example, Plucinsky & Virani (2000) found that the front-illuminated CCDs can show typical increases of 1 - 5 cts s<sup>-1</sup> above the quiescent level, while the back-illuminated CCDs can show large excursions โ as high as 100 cts s<sup>-1</sup> above the quiescent level โ during background flares. The durations of these intervals of enhanced background are highly variable, ranging from 500 s to 10<sup>4</sup> s. The cause of these background flares is currently not known (see Grant, Bautz, & Virani, 2002); however, they may be caused by low-energy protons ($`<`$100 keV; e.g., Plucinsky & Virani, 2000; Struder et al., 2001). The time periods corresponding to these background flares are generally excised from the data before proceeding with further analysis although not always (see Brandt et al., 2001; Nandra et al., 2005). In fact, Kim et al. (2004) find that the source detection probability depends strongly on the background rate. To examine our observations for such periods, we used the CIAO script ANALYSE\_LTCRV.SL which identifies periods where the background is $`\pm `$3$`\sigma `$ above the mean. All 9 observations were filtered according to this prescription (see Table 1 for a comparison of raw exposure time vs. filtered exposure time), resulting in only $``$40.6 ks ($``$4%) being lost due to background flares (954.2 ks vs. 913.6 ks). Of this, 20.3 ks were excluded from the end of ObsID 5017 due to a flare in which the count rate increased by a factor of $``$2. Table 2 lists the net exposure time for each of the 4 pointings used to image the ECDFS region. The net exposure time for each of the 4 pointings varies from a low of 205 ks to a high of 239 ks, with the mean net exposure time for the entire survey field of 228 ks. These extra steps in processing help remove spurious sources and result in fewer catalog sources than if the standard processing or pipeline products were used.
## 3 Data Analysis
In this paper, we report on the sources detected in three standard X-ray bands (see Table 3): 0.5โ8.0 keV (full band), 0.5โ2.0 keV (soft band), and 2.0โ8.0 keV (hard band). The raw ACIS resolution is 0.492 arcsec pixel<sup>-1</sup>, however, source detection and flux determinations were performed on the block 4 images, i.e., 1.964 arcsec pixel<sup>-1</sup>, as the source detect tool and exposure map generation require significant computer resources for full size images; for greater accuracy, source positions were determined from the block 1 images.
### 3.1 Image and Exposure Map Creation
Observations at each of the four pointings were combined via the CIAO script โmerge\_all.โ This script was executed using CIAO version 2.3 because of a known bug in the โasphistโ tool under CIAO version 3.2.1; this bug results in incorrect exposure maps for the merged image<sup>6</sup><sup>6</sup>6See the usage warning at http://cxc.harvard.edu/ciao/threads/merge\_all/.. At each pointing, the observation with the longest integration time was used for co-ordinate registration. For example, when merging ObsIds 5015 and 5016, the merged event list was registered to ObsId 5015 as it has approximately twice the integration time as 5016. Table 2 lists the ObsIds for each pointing, as well as the raw and the net integration time. For each pointing, we constructed images in the three standard bands: 0.5โ8.0 keV (full band), 0.5โ2.0 keV (soft band), and 2.0โ8.0 keV (hard band); see Table 3. The full band exposure-corrected image for the entire survey field<sup>7</sup><sup>7</sup>7Raw and smoothed ASCA -grade images for all three standard bands (See Table 3) are available from http://www.astro.yale.edu/svirani/ecdfs/. is presented in Figure 1.
We constructed exposure maps in these three energy bands for each pointing and for the entire survey field<sup>8</sup><sup>8</sup>8Exposure maps for all three standard bands (see Table 3) are available from http://www.astro.yale.edu/svirani/ecdfs/.. These exposure maps were created in the standard way and are normalized to the effective exposure of a source location at the aim point. The procedure used to create these exposure maps accounts for the effects of vignetting, gaps between the CCDs, bad column filtering, and bad pixel filtering. However, it should be noted that charge blooms caused by cosmic rays can reduce the detector efficiency by as much as few percent<sup>9</sup><sup>9</sup>9See http://cxc.harvard.edu/ciao/caveats/acis\_caveats\_050620.html.. There is currently no way to account for such charge cascades; however, when a tool becomes available, we will correct for this effect as necessary and make the new exposure maps publicly available at the World Wide site listed in Footnote 7. The exposure maps were binned by 4 so that they were congruent to the final reduced images. A photon index of $`\mathrm{\Gamma }`$ = 1.4, the slope of the X-ray background in the 0.5โ8.0 keV band (e.g. Marshall et al., 1980; Gendreau et al., 1995; Kushino et al., 2002) was used in creating these exposure maps.
In order to calculate the survey area as a function of the X-ray flux in the soft and hard bands, we used the exposure maps generated for each band and assumed a fixed detection threshold of 5 counts in the soft band and 2.5 in the hard band ($``$2$`\sigma `$). Dividing these counts by the exposure map, we obtain the flux limit at each pixel for each band. The pixel area is then converted into a solid angle and the cumulative histogram of the flux limit is constructed (Figure 2). The total survey area is $``$ 1100 arcmin<sup>2</sup> ($``$ 0.3 deg<sup>2</sup>). A more precise method of determining the survey area as a function of the X-ray flux is described by Kenter & Murray (2003); however, this would affect only the faint tail of the sample and would not significantly alter the present results. Therefore, a more sophisticated treatment is deferred to a later paper.
### 3.2 Point Source Detection
To perform X-ray source detection, we applied the CIAO wavelet detection algorithm wavdetect (Freeman et al., 2002). Although several other methods have been used in other survey fields to find sources in Chandra observations (e.g., Giacconi et al., 2002; Nandra et al., 2005), we chose wavdetect to allow a straightforward comparison between sources found in our catalog with those found in the CDF-S (Giacconi et al., 2002; Alexander et al., 2003). Moreover, wavdetect is more robust in detecting individual sources in crowded fields and in identifying extended sources than the other CIAO detection algorithm, celldetect. Point-source detection was performed in each standard band (see Table 3) using a โ$`\sqrt{2}`$ sequenceโ of wavelet scales; scales of 1, $`\sqrt{2}`$, 2, $`2\sqrt{2}`$, 4, $`4\sqrt{2}`$, and 8 pixels were used. Brandt et al. (2001), for example, showed that using larger scales can detect a few additional sources at large off-axis angles but found that this โ$`\sqrt{2}`$ sequenceโ gave the best overall performance across the CDF-N field. Moreover, as Alexander et al. (2003) point out, sources found with larger scales tend to have source properties and positions too poorly defined to give useful results.
Our criterion for source detection is that a source must be found with a false-positive probability threshold ($`p_{thresh}`$) of $`1\times 10^7`$ in at least one of the three standard bands. This false-positive probability threshold is typical for point-source catalogs (e.g., Alexander et al., 2003; Wang et al., 2004), although Kim et al. (2004) found that a significance threshold parameter of $`1\times 10^6`$ gave the most efficient results in the Chandra Multiwavelength Project (ChaMP) survey. We ran wavdetect using both probability thresholds and found that using the lower significance threshold (i.e., $`1\times 10^6`$) results in only an additional 64 unique sources. Visual inspection of each of these sources suggest they are bona fide X-ray sources. However, because these are sources found with the lower significance threshold, we present them in a separate table (the secondary catalog; Table 5). The primary catalog (Table 4) is a compilation of 587 unique sources found using the higher significance threshold in at least one of the three energy bands. For the remaining source detection parameters, we used the default values specified in CIAO which included requiring that a minimum of 10% of the on-axis exposure was needed in a pixel before proceeding to analyze it. We also applied the exposure maps generated for each pointing (see Section 3.1) to mitigate finding spurious sources which are most often located at the edge of the field of view.
The number of spurious sources per pointing is approximately $`p_{thresh}\times N_{pix}`$, where $`N_{pix}`$ is the total number of pixels in the image, according to the wavdetect documentation. Since there are approximately 2 $`\times `$ 10<sup>6</sup> pixels in each image for each pointing, we expect $``$0.2 spurious sources per pointing per band for a probability threshold of $`1\times 10^7`$. Therefore, treating the 12 images searched as independent, we expect $``$ 2-3 false sources in our primary catalog (Table 4) for the case of a uniform background. Of course the background is neither perfectly uniform nor static as the level decreases in the gaps between the CCDs and increases slightly near bright point sources. As mentioned by Brandt et al. (2001) and Alexander et al. (2003), one might expect the number of false sources to be increased by a factor of $``$ 2โ3 due to the large variation in effective exposure time across the field and the increase in background near bright sources due to the point-spread function (PSF) wings. But our false-source estimate is likely to be conservative by a similar factor since wavdetect suppresses fluctuations on scales smaller than the PSF. That is, a single pixel is unlikely to be considered a source detection cell โ particularly at large off-axis angles (Alexander et al., 2003).
The source lists generated by the procedure above for each of the standard bands in each of the pointings of the ECDFS were merged to create the point-source catalogs presented in Tables 4 and 5. The source positions listed in each catalog are the full band wavdetect-determined positions except when the source was detected only in the soft or hard bands. To identify the same source in the different energy bands, a matching radius of 2.โณ5 or twice the PSF size of each detect cell, whichever was the largest, was used. For comparison, Alexander et al. (2003) and Nandra et al. (2005) used a matching radius of 2.โณ5 for sources within 6โฒ of the aimpoint, and 4.โณ0 for sources with larger off-axis angles. With our method, 9 and 3 soft- and hard-band sources, respectively, have more than 1 counterpart, so we took the closest one. Note that both Tables 4 and 5 excludes sources found by wavdetect in which one or both of the axes of the โsource ellipseโ collapsed to zero. Over the survey field, 70 such sources are found; in general, these are unusual sources and although the formal probability of being spurious is low, there may be problems with these detections. Hornschemeier et al. (2001) found that using the wavdetect-determined counts for such objects as we do results in a gross underestimate of the number of counts even though the source was detected with a probability threshold of $`1\times 10^7`$. Since these sources would appear in catalogs that do circular aperture photometry instead, we present this list in a separate catalog (Table 6) for completeness.
Below we define the columns in Tables 4 and 5, our primary and secondary source catalogs for the ECDFS survey.
* Column 1 gives the ID number of the source in our catalog.
* Column 2 indicates the International Astronomical Union approved names for the sources in this catalog. All sources begin with the acronym โCXOYECDFโ (for โYale E-CDFโ)<sup>10</sup><sup>10</sup>10Name registration submitted to http://cdsweb.u-strasbg.fr/viz-bin/DicForm..
* Columns 3 and 4 give the right ascension and declination, respectively. These are wavdetect-determined positions for the unbinned images. If a source is detected in multiple bands, then we quote the position determined in the full band; when a source is not detected in the full band, we quote the soft-band position or the hard-band position.
* Column 5 gives the PSF cell size, in units of arcseconds, as determined by wavdetect. The farther off-axis a source lies, the larger the PSF size.
* Columns 6, 7, and 8 give the count rates (in units of cts s<sup>-1</sup>) in the full band and the corresponding upper and lower errors estimated according to the prescription of Gehrels (1986). If a source is undetected in this band, no count rate is tabulated.
* Columns 9, 10, and 11 give the count rates (in units of cts s<sup>-1</sup>) in the soft band and the corresponding upper and lower errors estimated according to the prescription of Gehrels (1986). If a source is undetected in this band, no count rate is tabulated.
* Columns 12, 13, and 14 give the count rates (in units of cts s<sup>-1</sup>) in the hard band and the corresponding upper and lower errors estimated according to the prescription of Gehrels (1986). If a source is undetected in this band, no count rate is tabulated.
* Column 15 lists the full band flux (in units of erg cm<sup>-2</sup> s<sup>-1</sup>) calculated using a photon slope of $`\mathrm{\Gamma }=`$1.4 and corrected for Galactic absorption. If a source was undetected in the full band but was detected in the hard or soft band, the hard- or soft- band flux (in that order of priority) was used to extrapolate to the full band assuming a photon slope of 1.4.
* Column 16 lists the soft band flux (in units of erg cm<sup>-2</sup> s<sup>-1</sup>) calculated using a photon slope of $`\mathrm{\Gamma }=`$1.4 and corrected for Galactic absorption. If a source was undetected in the soft band but was detected in the full or hard band, the full- or hard- band flux (in that order of priority) was used to extrapolate to the soft band assuming a photon slope of 1.4.
* Column 17 lists the hard band flux (in units of erg cm<sup>-2</sup> s<sup>-1</sup>) calculated using a photon slope of $`\mathrm{\Gamma }=`$1.4 and corrected for Galactic absorption. If a source was undetected in the hard band but was detected in the full or soft band, the full- or soft- band flux (in that order of priority) was used to extrapolate to the hard band assuming a photon slope of 1.4.
* Column 18 provides individual notes for each source. Examples include the catalog ID (c#) if detected in the CDF-S by Alexander et al. (2003), or if the source was selected from a band other the full band (โhโ or โsโ) or only detected in the full band (โfโ).
To determine source counts for each of our sources, we extracted counts in the standard bands from each of the images using the geometry of the wavdetect source cell and the wavdetect-determined source position. For example, to determine the counts in the soft band, we used the position and geometry determined by wavdetect in the soft band image to extract soft band counts. Some studies use circular aperature photometry to extract sources counts. However, as both Hornschemeier et al. (2001) and Yang et al. (2004; see their Figure 5) demonstrate, both techniques generally return the same result. Net count rates were then calculated using the effective exposure (which includes vignetting) for each pointing (exposure maps generated as described in Section 3.1). Errors were derived following Gehrels (1986), assuming an 84% confidence level. Note that the exposure maps do account for the degradation of the soft X-ray response of ACIS due to the build-up of a contamination layer on the ACIS optical blocking filter (Marshall et al., 2004; see Section 3.4). Therefore, the count rates reported in Table 4 are exposure- and contamination-corrected.
In Table 7 we summarize the source detections in the three standard bands, and in Table 8 we summarize the number of sources detected in one band but not in another. To convert the count rates to flux, we determined the conversion factor for each band assuming a photon slope of $`\mathrm{\Gamma }`$ = 1.4 and the mean Galactic $`N_H`$ absorption along the line-of-sight for each of the 4 pointings (N<sub>H</sub> = 9 $`\times `$ 10<sup>19</sup> cm<sup>-2</sup>; Stark et al., 1992).
Our faintest soft-band sources have $`4`$ counts (about one every 1.5 days), and our faintest hard-band sources have $`6`$ counts; these sources are detected near the aim point. The corresponding 0.5โ2.0 keV and 2โ8 keV flux limits, corrected for the Galactic column density, are $`1.7\times 10^{16}`$ erg cm<sup>-2</sup> s<sup>-1</sup> and $`3.9\times 10^{16}`$ erg cm<sup>-2</sup> s<sup>-1</sup>, respectively. Of course, these flux limits vary and generally increase across the field of view.
Undoubtedly, there are some sources in Table 4 that are extended sources (i.e., resolved by Chandra). Giacconi et al. (2002) find 18 extended sources in their 1 Ms catalog of the CDF-S out of 346 unique sources. The ECDFS survey has approximately 25% the integration time of the CDF-S but is approximately 3 times larger in area. Therefore, we expect roughly the same fraction of our sources reported in Table 4 are likely to be extended. The identification, X-ray, and optical properties of these sources will be presented in a later paper.
### 3.3 Astrometry
Given the superb Chandra spatial resolution, the on-axis positional accuracy is often quoted as being accurate to within 1โณ (e.g., Kim et al., 2004); in fact, the overall 90% uncertainty circle of a Chandra X-ray absolute position has a radius of 0.6 arcsec, and the 99% limit on positional accuracy is 0.8 arcsec<sup>11</sup><sup>11</sup>11See http://cxc.harvard.edu/cal/ASPECT/celmon/.. Nevertheless, as the off-axis angle increases, the PSF broadens and becomes circularly asymmetric (see Chandra Proposerโs Guide; URL listed in Footnote 1). Therefore, source positions for faint sources at large off-axis angles may not be accurate. In order to test the astrometry of the wavdetect-determined positions, we have matched our full-band X-ray positions provided in Table 4 against deep $`BVR`$-band imaging produced by the MUltiwavelength Survey by Yale/Chile (MUSYC<sup>12</sup><sup>12</sup>12For more information: http://www.astro.yale.edu/musyc/.; Gawiser et al., 2005). The 5$`\sigma `$ depth of the MUSYC optical imaging of this field is 27.1 AB mag with approximately 0.โณ85 seeing. Correlating the X-ray positions reported in Tables 4 and 5 with the optical positions found for sources in the ECDFS field, we find that approximately 72% of the sources reported in Table 4 and 41% of the sources reported in Table 5 have an optical counterpart within 1.โณ5 of the X-ray position. Furthermore, comparing the X-ray positions with the optical positions for these matched sources, we find a mean offset of -0.โณ08 in RA and +0.โณ28 in Dec. (We do not correct the X-ray positions for these offsets.) The optical properties of these X-ray sources will be presented in a forthcoming paper (Virani et al. 2005b, in prep.).
### 3.4 Accuracy of Source Detections and Fluxes
Approximately one third of the ECDFS field overlaps with the 1 Ms Chandra Deep Field South (see Figure 1 for the field layout). This is very useful as it allows us to compare our results with the properties of the overlapping sources already published. In particular, we used the catalog of Alexander et al. (2003), who re-analyzed the original CDF-S data. In Figure 3 we show the ratio of our fluxes to those reported by Alexander et al. (2003) for the overlapping sources. For this comparison, neither the CDF-S nor the ECDFS sources were corrected for intrinsic Galactic absorption. (This correction is $``$4% in the soft band and is negligible in the hard band.) Error bars are calculated by adding in quadrature the statistical (Poisson) uncertainties in the counts plus a 10% error arising from the likely range in spectral slopes (see Section 3.5).
Sources were matched using the closest CDF-S counterpart to each ECDFS source, using a maximum search radius of $`2^{\prime \prime }`$. To compare the fluxes of matched sources in the two data sets, we excluded the most discrepant top and bottom 15% of the flux ratios, and found our fluxes are $``$14% higher in the soft band and $``$11% higher in the hard band. In the first case, the difference can be explained by the different treatment of the contamination layer, which is particularly important in the soft band. The Alexander et al. (2003) catalog used ACISABS<sup>13</sup><sup>13</sup>13Available at http://www.astro.psu.edu/users/chartas/xcontdir/xcont.html to correct their fluxes for the presence of a contamination layer in the ACIS instrument. This tool assumes a spatially-uniform contamination layer composed of hydrogen, carbon, nitrogen, and oxygen. However, recent analysis of grating data (Marshall et al., 2004) shows that the amount of contamination correction depends on the spatial position on the instrument, and that the actual composition of the contamination is hydrogen, carbon, oxygen, and fluorine (P. Plucinsky, priv. comm.). These two new discoveries may have caused Alexander et al. (2003) to underestimate the contamination correction, thus making their fluxes lower in the soft band. In the hard band, the discrepancy can be explained by our assumed value of $`\mathrm{\Gamma }=`$ 1.4 for the spectral slope to calculate fluxes, while Alexander et al. (2003) used individual spectral fits for most of these overlapping sources. We conclude that the fluxes are broadly consistent and that systematic uncertainties in their average values are $``$ 15%, although individual fluxes have larger uncertainties (and some AGN may have actually varied).
### 3.5 Simulations
We performed extensive XSPEC and MARX simulations to investigate the statistical properties of the catalog, its completeness, and its flux limits. First, in order to investigate the effect of a fixed photon slope on the true flux of sources found in the ECDFS, we simulated 2000 sources with extreme photon spectral slopes, $`\mathrm{\Gamma }`$=1 and $`\mathrm{\Gamma }`$=2, and with fluxes distributed smoothly from the minimum to the maximum in our sample. We then computed their count rates in a typical ECDFS pointing ($``$ 230 ks). Using a fixed photon slope of $`\mathrm{\Gamma }`$=1.4 to compute fluxes then results in systematic flux errors of $``$ 10% in both the hard and soft bands.
To investigate the completeness of our catalog, we used MARX to simulate X-ray images of sources with known properties, including the range of count rates from just below our threshold to just above our highest count rate, and a generous range of spectral slopes (1 $`\mathrm{\Gamma }`$ 2) drawn from the observed $`\mathrm{\Gamma }`$ distribution observed in the 1 Ms CDFS survey (Alexander et al., 2003). We positioned 1000 sources of known fluxes (consistent with an exposure time of $``$ 230 ks) randomly within the ECDFS survey field, so the background and noise properties of the data are real. We then analyzed these simulated data with the same procedures used on the real ECDFS data; that is, we performed source detection on the resulting event list via wavdetect. This resulted in $``$90% of the sources being recovered overall, with incompleteness becoming important below $``$2$`\times 10^{16}`$ erg cm<sup>-2</sup> s<sup>-1</sup> and $``$2$`\times 10^{15}`$ erg cm<sup>-2</sup> s<sup>-1</sup>, in the soft and hard bands, respectively.
## 4 Results and Discussion
We found 651 unique sources in the Extended Chandra Deep Field-South survey field, which spans $``$ 0.3 deg<sup>2</sup> on the sky. Of these, 561 were detected in the 0.5โ8.0 keV full band, 529 in the 0.5โ2.0 keV soft band, and 335 in the 2.0โ8.0 keV hard band. There are 9 hard-band sources that are not detected in either the soft or full bands, 81 soft-band sources are not detected in either the hard or full bands, and 56 full-band sources are not detected in either the soft or hard bands (see Table 8). Of the 335 hard-band sources, 83 were not detected in the soft band ($``$20%); these are candidates for highly absorbed sources. Of the 529 and 335 sources detected in the soft and hard bands, respectively, 118 and 73 are detected in the CDF-S itself. Over this 0.11 deg<sup>2</sup> area, with an exposure time of $``$ 1 Ms, Giacconi et al. (2002) found 346 unique sources, of which 307 were detected in the 0.5โ2.0 keV band and 251 in the 2โ10 keV band. In the CDF-N, with an area similar to the CDF-S but with twice the exposure, Alexander et al. (2003) found 503 X-ray sources in the 2 Ms exposure. The number of sources found in the ECDFS is consistent with these two pencil beam surveys, given an approximate slope of unity for the X-ray counts in this flux range.
The cumulative distribution of sources for the soft and hard bands is shown in Figure 4. Error bars for a given bin were calculated by adding in quadrature the error bars from the previous bin to the 84% confidence error bars appropriate to the additional number of sources in the present bin, following the procedure described in Gehrels (1986). The observed distribution is compared to the compilation of Moretti et al. (2003) and to the $`\mathrm{log}`$ Nโ$`\mathrm{log}`$ S for the Chandra deep fields reported by Bauer et al. (2004). In the soft band there is very good agreement with the comparison sample in the flux range from $`4\times 10^{14}`$ to $`2\times 10^{16}`$ erg cm<sup>-2</sup> s<sup>-1</sup>. At the bright end, the discrepancy is not statistically significant, $``$1$`\sigma `$, because there are few bright X-ray sources in our field. At fluxes below $``$2$`\times 10^{16}`$ erg cm<sup>-2</sup> s<sup>-1</sup>, the observed $`\mathrm{log}`$ Nโ$`\mathrm{log}`$ S in the ECDFS flattens relative to the comparison samples because of incompleteness near the flux limit. Sources with soft fluxes of $``$2$`\times 10^{16}`$ erg cm<sup>-2</sup> s<sup>-1</sup> are only detected at the $``$ 2$`\sigma `$ level, and thus not all sources will be recovered.
The $`\mathrm{log}`$ Nโ$`\mathrm{log}`$ S relation for the hard band is shown in the right panel of Figure 4 and is compared again with the distributions of Moretti et al. (2003) and Bauer et al. (2004). Moretti et al. (2003) used 2-10 keV instead of 2-8 keV for the hard band. To convert 2-10 keV fluxes to the 2-8 keV band, we used a factor of 0.8, corresponding to the flux conversion assuming a $`\mathrm{\Gamma }`$=1.4 spectral slope. Bauer et al. (2004) quote 2-8 keV but appear to have used 2-10 keV, so we also converted their fluxes by the same factor (which reproduces their curve in Figure 4 of their paper). As in the soft band, very good agreement with previously reported $`\mathrm{log}`$ Nโ$`\mathrm{log}`$ S relations is seen for the 4$`\times 10^{14}`$ to 2$`\times 10^{15}`$ erg cm<sup>-2</sup> s<sup>-1</sup> range, and again, incompleteness at the faint end explains the observed discrepancy.
The differential $`\mathrm{log}`$ Nโ$`\mathrm{log}`$ S for both the soft and hard bands is shown in Figure 5. These observed distributions are compared to the predictions of the AGN population synthesis model of Treister & Urry (2005) which explains the X-ray background as a superposition of mostly obscured AGN. This model also explains the multiwavelength number counts of AGN in the Chandra Deep Fields (Treister et al., 2004). Given that these models match very well to the observed cumulative flux distributions from existing surveys, it is not surprising that this model also successfully explains the $`\mathrm{log}`$ Nโ$`\mathrm{log}`$ S distributions in the ECDFS field. Discrepancies can be found only at the fainter end, where incompleteness causes the number of observed sources to fall below the model prediction.
One of the early Chandra results was the finding that fainter X-ray sources have in general harder spectra (Giacconi et al., 2001), represented by higher values of the hardness ratio. Figure 6 shows that this effect is also observed in the ECDFS field, for a much larger number of sources. This trend is explained by obscuration since the soft band count rate is relatively more affected than the hard band, creating a harder observed X-ray spectrum while at the same time reducing the observed soft flux. This is in accordance with the general picture of AGN unification, although the precise geometry is not constrained, and it is as expected from population synthesis models (e.g., Treister & Urry, 2005 and references therein) which require a large number of obscured AGN at moderate redshift to explain the spectral shape of the X-ray background.
## 5 Conclusions
We present here the X-ray properties of sources detected in deep Chandra observations of the ECDFS field, the largest Chandra survey ever performed in terms of both area and depth. This survey covers a total of 0.3 square degrees, roughly 3 times the area of each very deep Chandra Deep Field. A total of 651 unique sources were detected in the four ACIS-I pointings in this field; 81 sources were detected in the soft but not in the full band, while 9 were detected only in the hard band. Roughly 15% of these 651 unique sources โ 118 sources in the soft band and 73 in the hard band โ were previously detected in the CDF-S. The fluxes derived for these sources agree well with the fluxes obtained from the CDF-S observations.
The X-ray $`\mathrm{log}`$ Nโ$`\mathrm{log}`$ S in the soft and hard bands agree well with those derived from other X-ray surveys and with predictions of the most recent AGN population synthesis models for the X-ray background.
As first discovered in early deep Chandra observations, we find in this sample that faint X-ray sources have in general harder spectra, indicating that these sources are likely obscured AGN at moderate redshifts. This is predicted by AGN unification models that explain the properties of the X-ray background. A future paper will discuss the optical and near-IR properties of these objects. This field was observed with the Spitzer Space Telescope by the MIPS GTO team and will also be observed by Spitzer as part of an approved program related to the MUSYC survey (PI: P. van Dokkum).
The source catalogs and images presented in this paper are available in electronic format on the World Wide Web (http://www.astro.yale.edu/svirani/ecdfs). We will continue to improve the source catalog as better calibration information, analysis methods, and software become available. For example, we plan to optimize the searching for variable sources and to study the multiwavelength properties of these X-ray sources.
Note: After this paper was submitted, another catalog paper by Lehmer et al. (2005) appeared on astro-ph. Our catalogs are similar but the analysis assumptions are different and therefore the source catalogs differ, as do the papers. We expect the comparison to be useful.
We thank the referee for helpful comments that improved the manuscript and are grateful to Samantha Stevenson of the CXC Help Desk for her help and patience in answering our many questions regarding CIAO-related tools. We also acknowledge the help of Jeffrey Van Duyne in cross-correlating the X-ray and optical positions. This work was supported in part by NASA grant HST-GO-09425.13-A. ET would like to thank the support of Fundaciรณn Andes, Centro de Astrofรญsica FONDAP and the Sigma-Xi foundation through a Grant in-aid of Research. EG acknowledges support by the National Science Foundation under Grant No. AST-0201667, an NSF Astronomy and Astrophysics Postdoctoral Fellowship (AAPF).
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# Power-law Tails in Non-stationary Stochastic Processes with Asymmetrically Multiplicative Interactions
## I Introduction
A mechanism of power-law emergence has been taken a strong interest for many researchers because of its ubiquity in nature. A number of attempts have been developed to explain the advent of various power-law distributions BTW1987 ; Jensen ; Sornette ; AB2002 , but underlying physics has not been clarified yet.
Recently, inelastic Maxwell models(IMM) BCG2000 ; EB2002 ; BK2002 have been studied extensively and provides a new mechanism of power laws. In IMM, randomly chosen particles undergo binary inelastic collisions and non-integer power-law tails appear in probability distribution function(PDF) of particle velocities. This processes are described by the Boltzmann equation with a velocity independent collision rate and are analytically tractable. The moment equations of IMM form a closed set which is able to be solved sequentially as an initial value problem. A nontrivial power-law exponent which is the function of a dissipation parameter is determined from a transcendental equation. This power law differs from usual critical phenomena in one major point that the system has a power-law region rather than a ciritical point. Therefore fine-tuning of parameters is not necessary and the power law is easy to be observed inside the region everywhere. ben-Avraham et al. BBLR2003 extended IMM to a symmetrically linear collision rule.
On the other hand, power laws are widely observed in social and economic phenomenaPareto1897 ; Zipf1949 ; MS2000 ; Takayasu2002 . It is well-known that wealth distributions such as capitals and incomes obey a power law in high-wealth range, which is called Paretoโs lawPareto1897 ; Champernowne1953 ; DGP2003condmat ; FujiGAGScondmat ; Yakocondmat . It is also recognized that the size $`S`$ of cities satisfies power law $`1/S`$ in the cumulative distribution function, that is, Zipfโs lawZipf1949 ; ZM1997 ; Gabaix1999 . A multiplicatively interacting stochastic process is one candidate to explain the lawsZM1997 ; Slacondmat . In economic phenomena, the Matthew effect (the rich gets richer and the poor gets poorer) plays an important role, which results in asymmetry in interactions. Ispolatov et al. IKR1998 have discovered that a power-law distribution with the exponent of unity in PDF arises in multiplicative processes of greedy exchange, where the rich always gets richer. The purpose of this paper is to study a model of asymmetrically multiplicative interactions and to clarify the influence of asymmetry. Note that the model includes IMM BCG2000 ; EB2002 ; BK2002 , symmetirically multiplicative interaction(SMI) model BBLR2003 , and greedy multiplicative exchange(GME) model IKR1998 as special cases.
The paper is organized as follows. In Sec II, we begin with introducing the model. The Fourier transform of the master equation is performed. PDF is assumed to be sum of the regular term and the singular term. Then, a transcendental equation is derived from the singular terms. In Sec. III, the growth rate $`\gamma `$ of the processes and the power-law exponent $`s`$ of tails are discussed in three cases I, I, and II. In case I and I, $`\gamma `$ and $`s`$ are computed explicitly by solving the transcendental equation. In case II, moment equations do not form a closed set, and $`\gamma `$ and $`s`$ cannot be calculated analytically. Thus numerical simulations are performed. The good agreement between the theory and simulations is achieved. In the last section, we discuss our results.
## II Transcendental equation
Let us consider the stochastic processes that distinguishes two particles in the manner of the magnitude of quantity asymmetrically: When a particle of positive quantity $`x(>0)`$ interact with a particle of quantity $`y(>0)`$, post-interaction quantities $`x^{}`$ and $`y^{}`$ are given by
$`\left(\begin{array}{c}x^{}\\ y^{}\end{array}\right)=\left(\begin{array}{cc}c(1a)& cb\\ da& d(1b)\end{array}\right)\left(\begin{array}{c}x\\ y\end{array}\right)(xy)`$ (7)
where $`0a,b1`$ and $`c,d>0`$ are interaction parameters representing amplification rates and exchange rates of larger and smaller quantities, respectively. Two particles are selected randomly. There are two trivial cases (i) $`a=b=0`$ or $`a=b=1`$, (ii) $`a=0,b=1`$ or $`a=1,b=0`$. In the former case, two particles experience no interactions and PDF becomes log-normal. In the latter, one particle ends up taking all the quantity and the others having null quantities. Hereafter we omit these cases.
The normalized distribution function $`f(z,t)`$ obeys the master equation.
$`{\displaystyle \frac{f(z,t)}{t}}+f(z,t)`$ $`=`$ $`{\displaystyle _0^{\mathrm{}}}๐y{\displaystyle _y^{\mathrm{}}}๐xf(x,t)f(y,t)`$ (8)
$`\times [\delta (z(c(1a)x+cby))`$
$`+\delta (z(dax+d(1b)y))].`$
The integration is performed along the solid line illustrated in Fig. 1. The kink at $`y=x`$ stems from asymmetry of the model.
The Fourier transform $`g(k,t)=_0^{\mathrm{}}๐ze^{ikz}f(z,t)`$ of Eq. (8) is performed by changing variables $`p=xy`$, $`q=y`$ and using the formula
$`{\displaystyle _0^{\mathrm{}}}๐pe^{ikp}=\pi \delta (k)+iv.p.{\displaystyle \frac{1}{k}}`$
where $`v.p.`$ means the Cauchy principal value. It follows that
$`{\displaystyle \frac{}{t}}g(k,t)+g(k,t)`$ $`=`$ $`{\displaystyle \frac{1}{2}}[g(c(1a)k,t)g(cbk,t)+g(dak,t)g(d(1b)k,t)]`$ (9)
$`{\displaystyle \frac{i}{2\pi }}v.p.{\displaystyle _{\mathrm{}}^{\mathrm{}}}๐k^{}{\displaystyle \frac{1}{k^{}c(1a)k}}g(k^{},t)g(c(1a+b)kk^{},t)`$
$`{\displaystyle \frac{i}{2\pi }}v.p.{\displaystyle _{\mathrm{}}^{\mathrm{}}}๐k^{}{\displaystyle \frac{1}{k^{}dak}}g(k^{},t)g(d(a+1b)kk^{},t).`$
The $`v.p.`$ terms in r.h.s. of Eq (9) come from the integration of kinked lines, and represent asymmetry of the model. Next, $`g(k,t)`$ is expanded as
$`g(k,t)=1+{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(ik)^n}{n!}}m_n(t)`$ (10)
where $`m_n(t)_0^{\mathrm{}}๐zz^nf(z,t)`$ is the n-th order moment. Even though the moment equations do not form a closed set generally, they can be derived formally as
$`{\displaystyle \frac{d}{dt}}m_n(t)\lambda _n(t)m_n(t)`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{l=1}{\overset{n1}{}}}\left(\begin{array}{c}n\\ l\end{array}\right)\{c^n(1a)^lb^{nl}+d^na^l(1b)^{nl}\}m_l(t)m_{nl}(t)(n1),`$ (13)
$`\lambda _n(t)`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left[c^n\left\{(1a)^n+b^n\right\}+d^n\left\{a^n+(1b)^n\right\}\right]1`$ (16)
$`{\displaystyle \frac{i}{2\pi i^n}}{\displaystyle \frac{1}{m_n(t)}}{\displaystyle \underset{l=0}{\overset{n}{}}}\left(\begin{array}{c}n\\ l\end{array}\right)\{c^n(1a)^lb^{nl}+d^na^l(1b)^{nl}\}v.p.{\displaystyle _{\mathrm{}}^{\mathrm{}}}๐k^{}{\displaystyle \frac{1}{k^{}}}g^{(l)}(k^{},t)g^{(nl)}(k^{},t)`$
where $`g^{(l)}(k,t)`$ denotes the $`l`$th order deirivative
$`g^{(l)}(k,t)=i^l{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(ik)^n}{n!}}m_{l+n}(t).`$ (17)
In real $`x`$-$`y`$ space, the pseudo eigenvalue $`\lambda _n(t)`$ is expressed as
$`\lambda _n(t)`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left[c^n\left\{(1a)^n+b^n\right\}+d^n\left\{a^n+(1b)^n\right\}\right]1`$ (20)
$`{\displaystyle \frac{1}{2}}{\displaystyle \frac{1}{m_n(t)}}{\displaystyle \underset{l=0}{\overset{n}{}}}\left(\begin{array}{c}n\\ l\end{array}\right)\{c^n(1a)^lb^{nl}+d^na^l(1b)^{nl}\}{\displaystyle _0^{\mathrm{}}}๐y{\displaystyle _y^{\mathrm{}}}๐x(x^ly^{nl}y^lx^{nl})f(x,t)f(y,t).`$
Since $`f(x,t)`$ includes information of all-order of moments, the moment equations do not form a closed set in the presence of the integral terms of Eq. (20). In IMM and SMI, however, the integral terms vanish. Therefore, $`\lambda _n`$ is constant and the moment equations (13) form a closed set.
We are interested in similarity solutions of the form
$`\{\begin{array}{ccc}f(z,t)& =& e^{\gamma t}\mathrm{\Psi }(\xi ),\\ g(k,t)& =& \mathrm{\Phi }(\eta )\end{array}`$ (23)
where $`\gamma `$ is the growth rate (scaling parameter) of the system, and $`\xi =ze^{\gamma t}`$ and $`\eta =ke^{\gamma t}`$ are scaled variables. The scaled PDF $`\mathrm{\Phi }(\eta )`$ satisfies
$`\gamma \eta {\displaystyle \frac{d\mathrm{\Phi }(\eta )}{d\eta }}+\mathrm{\Phi }(\eta )`$ $`=`$ $`{\displaystyle \frac{1}{2}}[\mathrm{\Phi }(c(1a)\eta )\mathrm{\Phi }(cb\eta )+\mathrm{\Phi }(da\eta )\mathrm{\Phi }(d(1b)\eta )]`$ (24)
$`{\displaystyle \frac{i}{2\pi }}v.p.{\displaystyle _{\mathrm{}}^{\mathrm{}}}๐k^{}{\displaystyle \frac{1}{\eta ^{}c(1a)\eta }}\mathrm{\Phi }(\eta ^{})\mathrm{\Phi }(c(1a+b)\eta \eta ^{})`$
$`{\displaystyle \frac{i}{2\pi }}v.p.{\displaystyle _{\mathrm{}}^{\mathrm{}}}๐k^{}{\displaystyle \frac{1}{\eta ^{}da\eta }}\mathrm{\Phi }(\eta ^{})\mathrm{\Phi }(d(a+1b)\eta \eta ^{}).`$
Here we assume that the function $`\mathrm{\Phi }(\eta )`$ is described by the sum of the regular and singular components
$`\mathrm{\Phi }(\eta )=\mathrm{\Phi }_{regular}(\eta )+\mathrm{\Phi }_{singular}(\eta ),`$ (25)
$`\mathrm{\Phi }_{regular}(\eta )=1+{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(i\eta )^n}{n!}}\mu _n,`$ (26)
$`\mathrm{\Phi }_{singular}(\eta )=C\mathrm{exp}\left(i{\displaystyle \frac{\pi s}{2}}sgn(\eta )\right)\mathrm{\Gamma }(s)|\eta |^s`$ (27)
where $`\mu _n`$ is the $`n`$-th order moment of $`\mathrm{\Phi }(\eta )`$, $`s`$ a non-integer exponent, $`C`$ the normalization constant, $`\mathrm{\Gamma }(s)(s0)`$ the gamma function, and $`sgn(\eta )`$ the signature of $`\eta `$. The form of the singular component Eq. (27) is given by the Fourier transform of $`\mathrm{\Psi }(\xi )=C/\xi ^{1+s}`$. The leading small-$`\eta `$ behavior of the singular component $`\mathrm{\Phi }_{singular}(\eta )|\eta |^s`$ reflects the tail of the scaled PDF $`\mathrm{\Psi }(\xi )1/\xi ^{1+s}`$ as $`\xi \mathrm{}`$. Substituting Eq. (25) into Eq. (24), we obtain relations:
$`O(\eta )`$ $`\gamma =\lambda _1=`$ $`{\displaystyle \frac{1}{2}}[c(1a+b)+d(a+1b)]+{\displaystyle \frac{1}{2}}(cd)(1ab)A1,`$ (28)
$`O(\eta ^n)`$ $`(n\gamma \lambda _n)\mu _n=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{l=1}{\overset{n1}{}}}\left(\begin{array}{c}n\\ l\end{array}\right)\{c^n(1a)^lb^{nl}+d^na^l(1b)^{nl}\}\mu _l\mu _{nl}(n2),`$ (31)
$`O(\eta ^s)`$ $`\gamma s=\lambda _s=`$ $`\{c(1a)\}^s+\{da\}^s1`$ (32)
where
$`A={\displaystyle \frac{1}{\mu _1}}{\displaystyle _0^{\mathrm{}}}๐\xi _2{\displaystyle _{\xi _2}^{\mathrm{}}}๐\xi _1(\xi _1\xi _2)\mathrm{\Psi }(\xi _1)\mathrm{\Psi }(\xi _2)(0<A1),`$ (33)
and the eigenvalue $`\lambda _n`$ is given by
$`\lambda _n`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left[c^n\left\{(1a)^n+b^n\right\}+d^n\left\{a^n+(1b)^n\right\}\right]1`$ (36)
$`{\displaystyle \frac{1}{2}}{\displaystyle \frac{1}{\mu _n}}{\displaystyle \underset{l=0}{\overset{n}{}}}\left(\begin{array}{c}n\\ l\end{array}\right)\{c^n(1a)^lb^{nl}+d^na^l(1b)^{nl}\}{\displaystyle _0^{\mathrm{}}}๐\xi _2{\displaystyle _{\xi _2}^{\mathrm{}}}๐\xi _1(\xi _1^l\xi _2^{nl}\xi _2^l\xi _1^{nl})\mathrm{\Psi }(\xi _1)\mathrm{\Psi }(\xi _2).`$
In the process of deriving Eq. (32), we use a convergence condition; $`\mathrm{\Phi }(\eta )0(|\eta |\mathrm{})`$ and the formulae to calculate regular-singular integrals
$`v.p.{\displaystyle _{\mathrm{}}^{\mathrm{}}}๐\eta ^{}{\displaystyle \frac{1}{\eta ^{}\alpha \eta }}|\eta ^{}|^s\mathrm{\Phi }_{regular}((\alpha +\beta )\eta \eta ^{})`$ $`=`$ $`\mathrm{tan}({\displaystyle \frac{\pi s}{2}})\pi sgn(\alpha \eta )|\alpha \eta |^s\mathrm{\Phi }_{regular}(\beta \eta ),`$
$`v.p.{\displaystyle _{\mathrm{}}^{\mathrm{}}}๐\eta ^{}{\displaystyle \frac{1}{\eta ^{}\alpha \eta }}sgn(\eta ^{})|\eta ^{}|^s\mathrm{\Phi }_{regular}((\alpha +\beta )\eta \eta ^{})`$ $`=`$ $`\mathrm{cot}({\displaystyle \frac{\pi s}{2}})\pi |\alpha \eta |^s\mathrm{\Phi }_{regular}(\beta \eta ).`$
In the caluculation of deriving moment equations, singular-singular integrals are able to be neglected because the higer-order singular terms in Fourier space are less singular in real space. That is, it is possible to linearize the moment equations for the singular term, which is the key ingredient of derivation. Notice that moment equations (28) and (31) for regular parts can be directly derived from the master equation (8) without recourse to Fourier transform.
The transcendental equation of the model with asymmetrically multiplicative interactions is formally obtained from Eqs. (28) and (32) by eliminating $`\gamma `$
$`\lambda _1s`$ $`=`$ $`\{c(1a)\}^s+\{da\}^s1(s>1).`$ (38)
When Eq. (38) has a non-trivial solution in $`0s1`$, the singular component dominates over the regular one because of the divergence of all the moments. Therefore the parameter $`\gamma `$ is not determined by Eq. (28). As originally reported by ben-Avraham et al.BBLR2003 instead, $`\gamma `$ takes the minimum value, which is realized when the line $`\gamma s`$ comes into contact with the curve $`\{c(1a)\}^s+\{da\}^s1`$. Consequently, the transcendental equation in the case $`0s1`$ becomes
$`\{c(1a)\}^s\mathrm{ln}\left({\displaystyle \frac{e}{\{c(1a)\}^s}}\right)+\{da\}^s\mathrm{ln}\left({\displaystyle \frac{e}{\{da\}^s}}\right)`$
$`=1(0s1).`$ (39)
## III Growth rate and Exponent
In this section, the growth rate $`\gamma `$ of the processes and the exponent $`s`$ of power-law tails are investigated. Three cases I, I, and II are discussed separately in accordance with the type of the transcendental equation. Equation (28) tells us that the first-order moment equation admits a closed form solution and the first-order eigenvalue $`\lambda _1`$ becomes constant in the case $`c=d`$ (I) and $`a+b=1(a0,1)`$ (I), while it does not otherwise (II). When $`\lambda _1`$ is constant, $`\gamma `$ and $`s`$ can be calculated analytically, which are done in subsections A and B. Case II is treated in subsection C. In this case, explicit calculation of $`\gamma `$ and $`s`$ becomes impossible at $`s>1`$. When $`0s1`$, on the contrary, the selection of the minimum growth rate leads to the irrelevance of $`\lambda _1`$ in determining $`\gamma `$ and $`s`$. In this situation, therefore, the transcendental equation is given by Eq. (39) and $`\gamma `$ and $`s`$ are computed analytically. First, the continuity of $`s`$ at $`s=1`$ is discussed. Then, numerical simulations are carried out and compared with the theory.
### III.1 Case I : $`c=d`$
When $`c=d`$, first-order eigenvalue $`\lambda _1`$ is given by $`c1`$. Then, the transcendental equation reads
$`(c1)s=\{c(1a)\}^s+\{ca\}^s1(s>1),`$ (40)
$`\{c(1a)\}^s\mathrm{ln}\left({\displaystyle \frac{e}{\{c(1a)\}^s}}\right)+\{ca\}^s\mathrm{ln}\left({\displaystyle \frac{e}{\{ca\}^s}}\right)`$
$`=1(0s1).`$ (41)
It should be emphasized that the transcendental equation, the growth rate $`\gamma `$, and the exponent $`s`$ are independent of $`b`$, exchange rate of smaller quantities. When amplification rates of larger and smaller quantities are same ($`c=d`$), only exchange of larger quantities determine similarity properties of processes. Although Eqs. (40) and (41) are formally same as those for symmetric interactions ($`c=d,a=b`$)BBLR2003 , the physical meaning of the parameters is completely different. A phase diagram for $`s`$ is illustrated in Fig. 2.
Power-law tails exist outside the hatched region. Inside the region, $`s`$ diverges and power-law tails disappear. The exponent $`s`$ varies continuously as a function of parameters $`a,c`$. The diagram is symmetric with respect to the line $`a=0.5`$. The two points $`c=1,a=0`$ and $`c=a=1`$ are singular points where the transcendental equations become identities. Actually, these points correspond to greedy multiplicative exchange(GME) processes and the exponent $`s`$ is given by zero IKR1998 . In the limit $`a0`$ or $`a1`$, $`s`$ also goes to zero at $`c1`$, which suggests that the cases $`a=0`$ and $`a=1`$ for arbitrary values of $`c`$ belong to the same universality class as GME. In Appendix, we show this analytically.
### III.2 Case I : $`a+b=1(a0,1)`$
At $`a+b=1`$, $`\lambda _1=c(1a)+da1`$ and the transcendental equation becomes
$`\{c(1a)+da1\}s=\{c(1a)\}^s+\{da\}^s1`$ (42)
for $`s>1`$. The equation for $`0s1`$ is given by Eq. (39). Cross-sectional phase diagrams at fixed $`d`$ for $`s`$ are plotted in Fig. 3. Boundaries between the regions with and without power-law tails are described by three curves
$`c={\displaystyle \frac{1}{1a}},a={\displaystyle \frac{1}{d}},c={\displaystyle \frac{1da}{1a}}.`$ (43)
Therefore the area without power-law tails is unbounded when $`d1`$, while it is bounded when $`d>1`$. Asymmetry between $`c`$ and $`d`$ gives rise to asymmetry of phase diagrams.
### III.3 Case II : $`cd`$ and $`a+b1`$
In this case, the first-order moment equation (28) does not allow a closed form solution. However a transcendental equation can be derived formally and possibility for the presence of power-law tails can be shown. The results are given by Eqs. (28) and (32). The value of $`A`$ in Eq. (33) cannot be obtained analytically. Therfore $`\gamma `$ and $`s`$ are unable to be calculated at $`s>1`$. When $`0s1`$ and the selection of the minimum growth rate takes place, all of $`\lambda _i(i1)`$ are irrelevant and $`\gamma `$ and $`s`$ are determined from Eq. (39).
First we discuss the continuity of $`s(a,b,c,d)`$ at $`s=1`$. Suppose that the power-law tail dominates in determining $`A`$ and we set $`\stackrel{~}{\mathrm{\Psi }}(\xi )=(1+ฯต)L^{1+ฯต}/\xi ^{2+ฯต}(s=1+ฯต)`$, where $`\stackrel{~}{\mathrm{\Psi }}(\xi )`$ is a normalized power-law PDF and $`L`$ is a lower cutoff in order to avoid the divergence of the integral. We obtain
$`A`$ $``$ $`{\displaystyle \frac{_L^{\mathrm{}}๐\xi _2_{\xi _2}^{\mathrm{}}๐\xi _1(\xi _1\xi _2)\stackrel{~}{\mathrm{\Psi }}(\xi _1)\stackrel{~}{\mathrm{\Psi }}(\xi _2)}{_L^{\mathrm{}}๐\xi \xi \stackrel{~}{\mathrm{\Psi }}(\xi )}}`$
$`=`$ $`{\displaystyle \frac{1}{1+2ฯต}}1(ฯต0).`$
Thus
$`\lambda _{s=1}=\lambda _{s1}=c(1a)+da1.`$ (45)
This fact indicates $`s`$ is continuous and satisfies Eq. (32) identically at $`s=1`$.
In the general case, $`0<A1`$ and the transcendental equation (32) has two nontrivial solutions, $`s_1<1`$ and $`s_2>1`$ generally. Ernst and Brito EB2002 considered the case and argued that each solution has a different role. That is, $`s_1`$ is a expansion variable and $`s_2`$ controls the singularity of the tail,
$`\mathrm{\Phi }(\eta )=1+{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(i\eta ^{s_1})^n}{n!}}\mu _n+\mu _{s_2}\eta ^{s_2}.`$ (46)
We find one more reason to explain why $`s_2`$ determines the power-law tail. In the limit $`cd`$ or $`a+b1`$, $`s_11`$ and $`s_2s`$. Assumption that $`s`$ is a continuous function of parameters $`a,b,c,d`$ leads to the conclusion that $`s_2`$ corresponds to a nontrivial solution in the cases I and I.
Now we compare the theory and numerical simulations. The value of $`A`$ is computed numerically. Then the growth rate $`\gamma `$ and the exponent $`s_2`$ (and $`s_1`$) are calculated via Eqs. (28) and (32). Obtained results are compared with those of simulations in Figs. 4-5. We find that Eqs. (28) and (32) hold well.
## IV Summary and Discussions
In this work, we have investigated asymmetrically multiplicative interactions (AMI) processes analytically and numerically. The model includes greedy multiplicative exchange (GME) model IKR1998 , inelastic Maxwell model (IMM) BCG2000 ; EB2002 ; BK2002 , and symmetrically multiplicative interactions (SMI) processes BBLR2003 as special cases. The relations are summarized in Table 1.
In the processes of deriving the transcendental equation, we encounter the fact that the first-order moment equation does not admit a closed form solution generally. However, it is possible to linearize singular terms. Thus we can formally construct the transcendental equation and show the possibility for the existence of non-integer power-law tails in PDF. It becomes evident that when $`c=d`$, the exchange parameter $`b`$ of small quantities is irrelevant to power-law tails. In the general case $`cd`$ and $`a+b1`$, the growth rate $`\gamma `$, the exponent $`s`$, and PDF are determined consistently via Eqs. (28) and (32). The exponent $`s`$ turns out to be a continuous function of parameters $`a,b,c,d`$.
This model is applicable to many kinds of asymmetrically interacting processes such as particle systems and biological phenomena. The most important would be wealth distributions in economic systems, where the Matthew effect gives rise to asymmetry of interactions. The model is quite general if terms are replaced as โparticleโ $``$ โagentโ, โquantityโ $``$ โcapitalโ or โincomeโ, and โinteractโ $``$ โtradeโ or โdealโ. Power laws are widely observed in economic phenomenaMS2000 ; Takayasu2002 . Some of them might be explained by this model.
*
## Appendix A General Greedy exchange Processes
Ispolatov et al.IKR1998 have reported that greedy multiplicative exchange processes where $`a=0`$ and the rich always gets richer exhibit a power-law behavior $`1/z`$ in PDF. They refered only to the case $`c=d=1`$, but as discussed in the case I, the tail $`1/z`$ with $`s=0`$ is also realized at $`c1`$ and $`d1`$. Here we extend their approach to general greedy exchange processes $`a=0,c1,d1`$ and examine the condition for $`f(z,t)1/z`$. The master equation is
$`{\displaystyle \frac{f(z,t)}{t}}`$ $`=`$ $`f(z,t)`$
$`+{\displaystyle \frac{1}{cb}}{\displaystyle _{\frac{z}{c(1a+b)}}^{\frac{z}{c(1a)}}}๐xf(x,t)f({\displaystyle \frac{zc(1a)x}{cb}},t)`$
$`+{\displaystyle \frac{1}{d(1b)}}{\displaystyle _{\frac{z}{d(a+1b)}}^{\frac{z}{da}}}๐xf(x,t)f({\displaystyle \frac{zdax}{d(1b)}},t)`$
In the limit of $`a0`$, Eq. (A) reduces to
$`{\displaystyle \frac{f(z,t)}{t}}`$ $`=`$ $`f(z,t)`$
$`+{\displaystyle \frac{1}{cb}}{\displaystyle _{\frac{z}{c(1+b)}}^{\frac{z}{c}}}๐xf(x,t)f({\displaystyle \frac{zcx}{cb}},t)`$
$`+{\displaystyle \frac{1}{d(1b)}}f({\displaystyle \frac{z}{d(1b)}},t){\displaystyle _{\frac{z}{d(1b)}}^{\mathrm{}}}๐xf(x,t).`$
Performing the transformation $`(zcx)/cbx`$ and using the normalization condition $`_0^{\mathrm{}}๐xf(x,t)=1`$, we have
$`{\displaystyle \frac{f(z,t)}{t}}`$ $`=`$ $`{\displaystyle _0^{\frac{z}{c(1+b)}}}๐xf(x,t)\left[{\displaystyle \frac{1}{c}}f({\displaystyle \frac{zcbx}{c}},t)f(z,t)\right]`$
$`f(z,t){\displaystyle _{\frac{z}{c(1+b)}}^{\mathrm{}}}๐xf(x,t)`$
$`+{\displaystyle \frac{1}{d(1b)}}f({\displaystyle \frac{z}{d(1b)}},t){\displaystyle _{\frac{z}{d(1b)}}^{\mathrm{}}}๐xf(x,t).`$
In the limit $`a0`$, divergent terms vanishes and the solution of Eq. (A) is obtained by
$`f(z,t)={\displaystyle \frac{D}{zt}},D={\displaystyle \frac{1}{\mathrm{ln}\left(\frac{d(1b)}{c}\right)}}.`$ (50)
Equation (50) shows that the condition $`d(1b)/c<1`$ is necessary to ensure $`f(z,t)0`$ . It becomes evident that the condition that general greedy exchange processes has the distribution $`1/z`$ is given by
$`d(1b)<c.`$ (51)
We find that the condition (51) is always satisfied at $`c=d`$.
Next, we estimate the lower and upper cutoffs $`z_1`$, $`z_2`$ of the distribution Eq. (50) in order to maintain the normalization condition of PDF. The moments are evaluated as
$`m_0(t)`$ $`{\displaystyle _{z_1}^{z_2}}๐zf(z,t)`$ $`{\displaystyle \frac{D}{t}}\mathrm{ln}\left({\displaystyle \frac{z_2}{z_1}}\right)=1,`$
$`m_1(t)`$ $`{\displaystyle _{z_1}^{z_2}}๐zzf(z,t)`$ $`{\displaystyle \frac{D}{t}}(z_2z_1)=m_1(0)e^{\gamma t}.`$
Using $`z_2>>z_1`$, we get
$`z_1{\displaystyle \frac{m_1(0)}{D}}te^{(\gamma 1/D)t},z_2{\displaystyle \frac{m_1(0)}{D}}te^{\gamma t}.`$ (52)
We have carried out numerical simulations of general greedy exchange processes. The results are shown in Fig. 6.
The distribuitons $`\mathrm{\Psi }(\xi )1/\xi `$ is observed for several tens order of magnitude. It should be emphasized that in the case $`c>1`$, the tails with $`s>0`$ appear asymptotically ($`\xi \mathrm{}`$). With increasing $`c`$, $`s`$ decreases and fluctuations grow. This indicates that the tails come from fluctuation due to amplification of large quantities. However, detailed analysis of the tails remains for a future work.
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# Planck Distribution in Noncommutative Space
## I Introduction
In quantum mechanics, the phase space is defined by replacing the canonical position and momentum variables with the Hermitian operators which obey the well-known Heisenberg commutation relations $`[\widehat{x}_i,\widehat{p}_j]=i\mathrm{}\delta _{ij}`$. Later, inspired by the quantum mechanics, it was suggested that one could use the idea of space-time noncommutativity at very small length scales to introduce an effective ultraviolet cutoff
$`[x^\mu ,x^\nu ]=i\mathrm{\Theta }^{\mu \nu },`$ (1)
where $`\mathrm{\Theta }^{\mu \nu }`$ is an antisymmetric tensor describing the strength of the noncommutative effects and plays an analogous role to $`\mathrm{}`$ in usual quantum mechanics. Note that the matrix $`\mathrm{\Theta }^{\mu \nu }`$ is not a tensor since its elements are identical in all reference frames.
Noncommutative field theories are constructed from the standard field theories just by replacing the usual multiplication with the Moyal $``$-product in the Lagrangian density. The $``$-product of the two fields up to first order is defined by
$$\left(\mathrm{\Psi }\mathrm{\Phi }\right)(x)=\mathrm{\Psi }(x)\mathrm{\Phi }(x)+\frac{i}{2}\mathrm{\Theta }^{\mu \nu }_\mu \mathrm{\Psi }(x)_\nu \mathrm{\Phi }(x).$$
(2)
Now, let us write Hamiltonian of the particle interacting with the radiation field in noncommutative space. The interaction Hamiltonian in noncommutative space is assumed to be obtainable from the standard prescription if we replace the ordinary product with the $``$-product. The idea is to formulate the theories in noncommutative space as theories in commutative space and to express the noncommutativity by an appropriate $``$-product. Then, the interaction Hamiltonian in noncommutative space is given by
$$\widehat{H}_{int.}=\frac{e}{mc}\widehat{\text{A}}\text{p}.$$
(3)
Note also that $`\widehat{\text{A}}`$ in the above Hamiltonian is the noncommutative gauge field operator. Here, caret indicates the noncommutative quantities.
Gauge theory was formulated in noncommutative space . The Seiberg-Witten map allows us to write the noncommutative fields in terms of the ordinary fields. The Seiberg-Witten map for the gauge field is known to lowest order in $`\mathrm{\Theta }^{\mu \nu }`$
$$\widehat{A_\nu }=A_\nu \frac{1}{2}\mathrm{\Theta }^{\alpha \beta }A_\alpha (_\beta A_\nu +F_{\beta \nu }).$$
(4)
Substitution of the equation (4) into the Hamiltonian (3) and applying the definition (2) yields
$$\widehat{H}_{int.}=\frac{e}{mc}\text{A}.\text{p}\frac{ie}{2mc}\mathrm{\Theta }^{\mu \nu }_\mu \text{A}._\nu \text{p}+\frac{e}{2mc}\mathrm{\Theta }^{\mu \nu }\underset{j=1}{\overset{3}{}}A_\mu (_\nu A_j+F_{\nu j})p_j,$$
(5)
where $`p_j`$ is the $`j`$-th component of the momentum operator and $`A_j=(0,\text{A})`$. The first term in the above Hamiltonian is the standard one and the other two terms are due to the noncommutative nature of the space.
A few more words about the contribution to the interaction Hamiltonian are in order. Since A contains only the operators $`a^{}`$ and $`a`$ for photons linearly, the second term in Eq. (5) induces transitions for which one photon is either produced or annihilated. On the other side, the last term in the equation (5) induces transitions where two photons are involved. Here, we are interested only in one photon processes, so we restrict our consideration to the interaction Hamiltonian which includes A linearly. So, we omit the last term in Eq. (5).
If we use the plane wave expansion of A with the box normalization condition, the second term in the above Hamiltonian can be rewritten as
$$\frac{ie}{2mc}\mathrm{\Theta }^{\mu \nu }_\mu \text{A}._\nu \text{p}=\frac{ie}{2mc\mathrm{}}\mathrm{\Theta }^{\mu \nu }\underset{k\rho }{}N_kk_\mu (a_{k\rho }e^{i(k.xwt)}a_{k\rho }^{}e^{i(k.xwt)})\widehat{ฯต}_{k\rho }.\text{p}p_\nu ,$$
(6)
where $`\widehat{ฯต}_{k\rho }`$ are the polarization vectors ($`\rho =1,2`$) and $`N_k=\sqrt{{\displaystyle \frac{2\pi \mathrm{}c^2}{L^3\omega _k}}}`$. Here, we used the following identity: $`_\nu \text{p}={\displaystyle \frac{i}{\mathrm{}}}p_\nu \text{p}`$.
Having obtained the contribution to the interaction Hamiltonian, now we can calculate the contribution of this noncommutative term to the first order transition matrix element.
$$<f,n_{k\rho }\pm 1\frac{ie}{2mc\mathrm{}}\mathrm{\Theta }^{\mu \nu }\underset{k\rho }{}N_kk_\mu (a_{k\rho }e^{i(k.x\omega t)}a_{k\rho }^{}e^{i(k.x\omega t)})\widehat{ฯต}_{k\rho }.\text{p}p_\nu |i,n_{k\rho }>$$
(7)
An atom is in the initial state $`|i>`$ and decays into the final state $`|f>`$ by emitting a photon with the wave vector k and polarization $`\rho `$. It can also change itโs state by absorbing a photon with energy $`\mathrm{}\omega `$. The initial and the final of the total atom+radiation field are denoted by $`|i,n_{k\rho }>`$ and $`<f,n_{k\rho }\pm 1|`$, respectively. The number of photons is increased or decreased by one unit. Note that there is a minus sign in front of the operator $`a^{}`$. This is because of the existence of the derivative of the gauge field operator in Hamiltonian (5). This leads to some differences between the two theories. For example, the transition matrix elements in electric dipole approximation are the same both for absorbtion and emission processes in usual quantum mechanics. However, they are not equal in noncommutative theories.
Let us calculate the first order transition matrix element for emission of a photon.
$$\frac{e}{mc}\sqrt{n_{k\rho }+1}N_k<fe^{ik.x}\widehat{ฯต}_{k^{}\rho ^{}}.\text{p}(1+\frac{i}{2\mathrm{}}\mathrm{\Theta }^{\mu \nu }k_\mu p_\nu )i>.$$
(8)
Similarly, the transition matrix element for the absorbtion of a photon is given by
$$\frac{e}{mc}\sqrt{n_{k\rho }}N_k<fe^{ik.x}\widehat{ฯต}_{k^{}\rho ^{}}.\text{p}(1\frac{i}{2\mathrm{}}\mathrm{\Theta }^{\mu \nu }k_\mu p_\nu )i>.$$
(9)
## II Noncommutative Planck Distribution
In this section, we derive Planckโs radiation formula. It was first derived by Einstein in 1917. We consider a sample of atoms in thermal equilibrium. The number of atoms in state $`|f>`$ is denoted by $`N_f`$ and for those in state $`|i>`$ the number $`N_i`$. Transitions occur between the two states; photons absorbed or emitted from the radiation field.
Equilibrium requires that time derivation of the number of atoms is zero.
$$\dot{N_f}=\dot{N_i}=0,$$
(10)
The number of particles with energy $`E`$, in thermal equilibrium at temperature $`T`$, is proportional to the Boltzmann factor, so
$$\frac{N_f}{N_i}=\mathrm{exp}(\frac{E_fE_i}{k_BT}),$$
(11)
where $`T`$ represents temperature and $`k_B`$ is Boltzmannโs constant.
We know from quantum statistical mechanics that this ratio is also equal to
$$\frac{N_f}{N_i}=\frac{(Trans.prob./time)_{emis.}}{(Trans.prob./time)_{abs.}}.$$
(12)
Now, we are in a position to calculate this ratio by taking the absolute square of the transition matrix elements Eqs.(8,9).
$`{\displaystyle \frac{(Trans.prob./time)_{emis.}}{(Trans.prob./time)_{abs.}}}=`$ (13)
$`{\displaystyle \frac{n_{k\rho }+1}{n_{k\rho }}}\times {\displaystyle \frac{|<fe^{ik.x}\widehat{ฯต}_{k\rho }.\text{p}i>+(i/2\mathrm{})\mathrm{\Theta }^{\mu \nu }k_\mu <fe^{ik.x}\widehat{ฯต}_{k\rho }.\text{p}p_\nu i>|^2}{|<fe^{ik.x}\widehat{ฯต}_{k\rho }.\text{p}i>(i/2\mathrm{})\mathrm{\Theta }^{\mu \nu }k_\mu <fe^{ik.x}\widehat{ฯต}_{k\rho }.\text{p}p_\nu i>|^2}}.`$ (14)
Let us compute the term with $`\mathrm{\Theta }^{\mu \nu }`$ further by using the electric dipole approximation method $`(e^{\pm ikx}1)`$.
$`\mathrm{\Theta }^{\mu \nu }k_\mu <f\widehat{ฯต}_{k\rho }.\text{p}p_\nu i>=\mathrm{\Theta }^{\mu \nu }k_\mu {\displaystyle \underset{l}{}}<f\widehat{ฯต}_{k\rho }.\text{p}l><lp_\nu i>,`$ (15)
where $`p_\nu =(p_0,\text{p})`$.
If $`\nu =0`$, then $`<lp_0i>={\displaystyle \frac{E_i}{c}}\delta _{li}`$. If $`\nu 0`$, then $`<lp_ji>=im\omega _{li}<l|x_j|i>`$, where $`j=1,2,3`$. Then, the equation (15) can be rewritten
$`{\displaystyle \frac{im\omega _{fi}E_i}{c}}\mathrm{\Theta }^{\mu 0}k_\mu <f\widehat{ฯต}_{k\rho }.\text{x}i>m^2{\displaystyle \underset{l}{}}\mathrm{\Theta }^{\mu j}k_\mu \omega _{fl}\omega _{li}<f\widehat{ฯต}_{k\rho }.\text{x}l><lx_ji>.`$ (16)
The elements of $`\mathrm{\Theta }^{\mu \nu }`$ are very small ($`\mathrm{\Theta }^{\mu \nu }<<1`$). So, we can neglect the square terms of $`\mathrm{\Theta }^{\mu \nu }`$. Let us define $`ฯต`$ which includes $`\mathrm{\Theta }^{\mu \nu }`$ linearly as
$$ฯต=\frac{\frac{m}{2\mathrm{}}\mathrm{\Theta }^{\mu j}k_\mu _l\omega _{fl}\omega _{li}<f\widehat{ฯต}_{k\rho }.\text{x}l><lx_ji>}{\omega _{fi}<f\widehat{ฯต}_{k\rho }.\text{x}i>}.$$
(17)
Then the ratio Eq. (LABEL:mamik) is approximated as
$$\frac{(Trans.prob./time)_{emis.}}{(Trans.prob./time)_{abs.}}\frac{n_{k\rho }+1}{n_{k\rho }}\times \frac{12ฯต}{1+2ฯต}\frac{n_{k\rho }+1}{n_{k\rho }}\times (14ฯต),$$
(18)
At the last step, we used the following relation $`({\displaystyle \frac{1}{1+ฯต}}1ฯต)`$. Substitution of the equations (11,18) into the definition (12) leads to the Planck distribution function in noncommutative space.
$$n_{k\rho }=\frac{1}{\kappa \mathrm{exp}(\frac{\mathrm{}\omega _k}{k_BT})1},$$
(19)
where $`\kappa `$ is given by
$$\kappa =\frac{1}{14ฯต}1+4ฯต.$$
(20)
In commutative limit ($`\mathrm{\Theta }^{\mu \nu }0`$, $`\kappa 1`$), the Eq. (19) is reduced to the usual Planck distribution function.
Cosmic black body radiation may be one of the experimental test of noncommutative theories. The remnants of the intense radiation field produced in the beginning can be presented as a black body radiation. The value of the physical parameters like energy and temperature should be shifted because of the parameter $`\kappa `$ in Eq. (19), if noncommutativity of the space plays a role in the first times of big bang. One can study many new interesting physical phenomena with this noncommutative Planck distribution function Eq. (18).
Up to now, we have not said anything about the elements of the noncommutative parameter $`\mathrm{\Theta }^{\mu \nu }`$. Actually, there is no agreement in the literature on how these elements are constructed explicitly. While many researchers set $`\mathrm{\Theta }^{0i}=0`$ to avoid problems with unitarity and casuality , the others assume that the components of $`\mathrm{\Theta }^{\mu \nu }`$ are constant over cosmological scales, in any given frame of reference there is a special noncommutative direction . Hewett-Petriello-Rizzo parametrized them with the three different angles. In , $`\mathrm{\Theta }^{\mu \nu }`$ is interpreted as a background B-field.
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# Hi-Res Spectroscopy of a Volume-Limited Hipparcos Sample within 100 parsec
## 1. Introduction
Accurate parallax measurements of solar neighborhood stars from the Hipparcos catalog have enabled detailed studies of stellar properties including studies of the star formation history (Hernandez, Valls-Gabaud, & Gilmore 2000; Bertelli & Nasi 2001). By looking at stellar abundances at different ages, we can constrain the star formation history of the solar neighborhood and study the history of chemical evolution. Subgiants are of particular interest because they lie in an area where isochrones are well separated, enabling more accurate age determination.
We have initiated a project to obtain high-resolution echelle spectra of evolved Hipparcos stars in a volume-limited sample within 100 parsec. We are performing detailed abundance analysis to investigate the chemical evolution and the star formation history of the solar neighborhood.
Currently, we have obtained spectra for about 50% of our sample. We are attempting automated analysis to obtain fundamental attributes such as \[Fe/H\], effective temperature, surface gravity, and age. Here, we present initial results for 10 subgiants over a range of effective temperatures. We compare our results using several different methods to literature values, and discuss future development for this project.
## 2. Sample
We select evolved Hipparcos stars with acceptable parallax uncertainties in the Northern Hemisphere within 100 parsec that have M$`{}_{V}{}^{}`$ 4 and B-V $``$ 0.38, as shown in Fig. 1(a). To probe histories of star formation and chemical evolution, abundance analysis of subgiants is ideal because we can use the well-separated isochrones of the subgiant branch to obtain accurate ages. Therefore we start our analysis on subgiants first, of which the 10 that we are presenting here are also shown in Fig. 1(a). They are chosen to cover a wide range of effective temperatures to test the accuracy of our automated method.
We are randomly selecting about 500 stars to observe as a representative sample. We exclude known binaries and multiple-star systems to avoid complications such as abundance ambiguity in spectroscopic binaries. We will also exclude stars with high rotational velocities and significant on-going mass loss.
By collecting our own data, we ensure a uniform sample, consistency in data reduction, and sufficient SNR to obtain abundances of various individual elements. We obtain high-resolution echelle spectra with the ARC 3.5-m telescope at the Apache Point Observatory, and use standard techniques for spectrum extraction and calibration. To date, we have obtained over 300 spectra, tallied in Fig. 1(b), and are performing initial abundance analysis.
## 3. Analysis
Due to the number of stars in our sample, we are attempting to automate the analysis wherever possible. We measure equivalent widths ($`W`$) of selected โcleanโ Fe I and Fe II lines using nonlinear least squares fitting with a gaussian profile. An example is shown in Fig. 2. For the 10 stars discussed here, we also manually measure $`W`$ for comparison with the automatic results.
We use these measurements in MOOG (Sneden 2002) with a grid of Kurucz (1993) model atmospheres based on initial estimates of stellar properties. Our grid is interpolated from existing model atmospheres and has effective temperature ($`T_{\text{eff}}`$), logarithmic surface gravity ($`\mathrm{log}g`$), and mircoturbulence velocity ($`\xi `$) steps of 50 K, 0.1 cm s<sup>-2</sup>, and 0.1 km s<sup>-1</sup> respectively. For each star, we obtain \[Fe/H\], $`T_{\text{eff}}`$, $`\mathrm{log}g`$, and $`\xi `$ by finding zero slopes of Fe I abundance versus excitation potential (EP) and reduced equivalent width (RW), and ionization balance between Fe I and Fe II. We can compare the $`\mathrm{log}g`$ obtained from parallax and $`T_{\text{eff}}`$ from color with those from spectroscopy.
With a derived metallicity, we fit Padova isochrones (Girardi et al. 2002) through a star to obtain its mass ($``$), age ($`\tau `$), and โtrigonometricโ $`\mathrm{log}g`$. These isochrones are interpolated for sufficient fitting via closest-point match in the $`L`$,$`T_{\text{eff}}`$-plane and have \[Fe/H\] intervals of 0.05 dex. For the stellar luminosities, we apply corrections to the absolute magnitudes as suggested by Lutz & Kelker (1973). In future work, we will improve isochrone fitting by including probability weighting to allow for uncertainties in the photometry and distances, and overlap between isochrones.
## 4. Initial Results - 10 Subgiants
We present preliminary results for 10 subgiants. They are selected from the color-magnitude diagram to span a range of temperatures in order to evaluate the accuracy of our method over varying stellar properties (in this case, $`T_{\text{eff}}`$). We perform equivalent width measurement, abundance analysis, and isochrone fitting as described in the Analysis section. In Table 1, we tabulate the results from the automated $`W`$ fitting.
HIP13679 lies at the hook of main sequence turn-off, HIP39681 has considerable parallax uncertainty ($`\frac{\sigma }{\pi _o}=`$ 0.16), and HIP61995 turns out to be on the red giant branch. These cause substantial uncertainties in $``$ and $`\tau `$ values. In the future, we will refine our selection criteria to minimize these issues.
We compare our iron abundance from automated and manual $`W`$ measurements with literature values (Cayrel de Strobel et al. 1997; Cayrel de Strobel, Soubiran, & Ralite 2001; Nordstrรถm et al. 2004) in Fig. 3. On average, the automated measurements agree better with published values, except for HIP13679. More data on the lower and higher $`T_{\text{eff}}`$ ends are required to determine if there is any \[Fe/H\] trend with effective temperature.
We compare the derived stellar parameters from the automatic and manual $`W`$ measurements in Fig. 4. We use $`T_{\text{eff}}`$ from $`VI`$ (Bessel, Castelli, & Plez 1998) for visual convenience. The โautomaticโ method underestimates \[Fe/H\] and overestimates $`\xi `$ compared to the โmanual.โ However the results for $`\mathrm{log}g`$ and $`T_{\text{eff}}`$ are comparable except at $`>`$6000ร
.
Two possible reasons for these deviations are systematic differences in the two $`W`$ measurement methods and slopes affected by outliers at the upper or lower limits of EP and RW during abundance analysis. This calls for more data and careful visual inspection of the results. For the former, we could fine-tune the automation, such as implementing dynamic continuum determination and line templates as a function of temperature. For the latter, we could add more lines near the EP and RW edges, or formulate an analysis method less dependent of the slopes.
From our initial results of 10 stars from automated $`W`$ fitting, we show the age-metallicity relationship (AMR) in Fig. 5. Edvardsson et al. (1993) reported a real scatter of metallicity with respect to age, and much published literature states no tight correlation in AMR for \[Fe/H\]. Our present results here show no surprising contradictions.
## 5. Future Work
We will continue to obtain data until our targeted sample size of $``$500 stars is achieved. We will improve our current automation algorithm for more robust $`W`$ measurements and better accuracy in derived stellar properties. In addition to Fe, we will extend abundance analysis to C, O, Mg, and many other elements made possible by high-resolution spectroscopy.
Once we have sufficient abundance information, we will be able to incorporate this into a study of the star formation history of the solar neighborhood, using observed values to constrain the models. We will also investigate the abundance distribution of heavy elements and the age-metallicity relation.
By selecting a volume-limited sample, we are biased towards thin-disk stars that have smaller velocity dispersion from the plane. We are also biased towards more metal-rich stars due to a lower-limit cutoff in the color. In the sample of subgiants, we are biased towards less massive stars as they have a longer timescale in that evolutionary stage than their more massive counterparts. We will attempt to correct for these biases in our analysis.
A useful by-product of this project will be a database of high resolution spectra of volume-limited sample of solar neighborhood stars, yielding a catalog of fundamental stellar properties, kinematics, and detailed abundances.
### Acknowledgments.
Support for this work is provided by the National Science Foundation. We also thank the organizing committee of Cozumel Resolved Stellar Populations Meeting 2005 for the wonderful time in Cancรบn, Mรฉxico.
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# References
USTC-ICTS-05-6
hep-th/0506115
Tachyon condensation on the intersecting brane-antibrane system
Hua Bai<sup>a</sup><sup>1</sup><sup>1</sup>1E-mail: huabai@mail.ustc.edu.cn, J. X. Lu<sup>a,b,c</sup><sup>2</sup><sup>2</sup>2E-mail: jxlu@ustc.edu.cn and Shibaji Roy<sup>d</sup><sup>3</sup><sup>3</sup>3E-mail: roy@theory.saha.ernet.in
<sup>a</sup> Interdisciplinary Center for Theoretical Study
University of Science and Technology of China, Hefei, Anhui 230026, P. R. China
<sup>b</sup> Center for Mathematics and Theoretical Physics
Shanghai Institute for Advanced Study
University of Science and Technology of China, Shanghai 201315, P. R. China
<sup>c</sup> Interdisciplinary Center of Theoretical Studies
Chinese Academy of Sciences, Beijing 100080, P. R. China
<sup>d</sup> Saha Institute of Nuclear Physics, 1/AF Bidhannagar, Calcutta-700 064, India
Abstract
We generalize our study of the tachyon condensation on the brane-antibrane system \[hep-th/0403147\] to the intersecting brane-antibrane system. The supergravity solutions of the intersecting brane-antibrane system are characterized by five parameters. We relate these parameters to the microscopic physical parameters, namely, the number of D$`p`$-branes ($`N_1`$), the number of $`\overline{\mathrm{D}}p`$-branes ($`\overline{N}_1`$), the number of D$`(p4)`$-branes ($`N_2)`$, the number of $`\overline{\mathrm{D}}(p4)`$-branes ($`\overline{N}_2`$) and the tachyon vev $`T`$. We show that the solution and the ADM mass capture all the required properties and give a correct description of the tachyon condensation for the intersecting brane-antibrane system.
A coincident D-braneโantiD-brane pair (or a non-BPS D-brane) in Type II string theories is known to be unstable which is characterized, for example, by the presence of tachyonic mode on the D-brane world-volume . As a result, these systems decay and the decay occurs by a process known as tachyon condensation . Tachyon condensation is well understood in the open string description using either the string field theory approach or the tachyon effective action approach on the brane. In we have obtained a closed string (or supergravity) understanding of this process. We have interpreted the previously known non-supersymmetric, three parameter supergravity solutions with a symmetry ISO($`p,1`$) $`\times `$ SO($`9p`$) in ten space-time dimensions as the coincident D$`p`$$`\overline{\mathrm{D}}p`$ system. Then the three parameters in this solution were related to the microscopic physical parameters, namely, the number of D$`p`$-branes ($`N`$), the number of $`\overline{\mathrm{D}}p`$-branes ($`\overline{N}`$) and the tachyon vev of the D$`p`$$`\overline{\mathrm{D}}p`$ system. Using these relations we have calculated the ADM mass and have shown that the solution and the ADM mass capture all the required properties and give a correct description of the tachyon condensation advocated by Sen on the D-$`\overline{\mathrm{D}}`$ system.
In this paper we generalize our previous study of the tachyon condensation on the brane-antibrane system to the intersecting brane-antibrane system. The supergravity solution in this case is again non-supersymmetric and has the isometry ISO($`p4,1`$) $`\times `$ SO(4) $`\times `$ SO($`9p`$) which can obviously be identified as the intersecting brane-antibrane system or more precisely the intersecting D$`p`$-$`\overline{\mathrm{D}}p`$ and D$`(p4)`$-$`\overline{\mathrm{D}}(p4)`$ system for $`3p6`$. We find that the solution is characterized by five parameters and once the solution is realized to represent intersecting D$`p`$-$`\overline{\mathrm{D}}p`$ and D$`(p4)`$-$`\overline{\mathrm{D}}(p4)`$ system, we can immediately infer that the parameters must be related to the microscopic physical parameters $`N_1`$, the number of D$`p`$-branes, $`\overline{N}_1`$, the number of $`\overline{\mathrm{D}}p`$-branes, $`N_2`$, the number of D$`(p4)`$-branes, $`\overline{N}_2`$, the number of $`\overline{\mathrm{D}}(p4)`$-branes and the tachyon vev $`T`$<sup>4</sup><sup>4</sup>4So long as the bulk configuration is concerned, the worldvolume fields (in particular the tachyon) donโt need to satisfy their respective worldvolume equations of motion (for example, we can put worldvolume scalars and tachyon to constants and other fields to zero) in the spirit, for example, of . In other words, they can be off-shell. In this way, the tachyon vev will appear as a parameter labelling the solution.. The more exact relationships between the supergravity parameters and the microscopic physical parameters can be obtained by examining how the solution reduces to the supersymmetric configuration which corresponds to the four cases as follows. (a) The intersecting $`N_1`$ D$`p`$-branes and $`N_2`$ D$`(p4)`$-branes (when both $`\overline{N}_1`$ and $`\overline{N}_2`$ are zero). If in addition to $`\overline{N}_1=\overline{N}_2=0`$ we also have $`N_2=0`$ (or $`N_1=0`$), we still get a BPS configuration of half susy instead of quarter susy (as is the case for the intersecting D$`p`$/D$`(p4)`$ brane configuration) which is $`N_1`$ D$`p`$-branes (or delocalized $`N_2`$ D$`(p4)`$-branes). When all the $`N`$โs are zero, which is a trivial case, we get the maximally supersymmetric flat space-time. (b) The intersecting $`\overline{N}_1`$ $`\overline{\mathrm{D}}p`$-branes and $`\overline{N}_2`$ $`\overline{\mathrm{D}}(p4)`$-branes (when both $`N_1`$ and $`N_2`$ are zero). Similar to the previous case here also if in addition to $`N_1=N_2=0`$, we have $`\overline{N}_2=0`$ (or $`\overline{N}_1=0`$) we get half susy configuration which is $`\overline{N}_1`$ $`\overline{\mathrm{D}}p`$-branes (or delocalized $`\overline{N}_2`$ $`\overline{\mathrm{D}}(p4)`$-branes). (c) The intersecting $`N_1`$ D$`p`$-branes and $`\overline{N}_2`$ $`\overline{\mathrm{D}}(p4)`$-branes (when both $`\overline{N}_1`$ and $`N_2`$ are zero). This case actually preserves also one quarter of spacetime supersymmetry which seems impossible from a first look<sup>5</sup><sup>5</sup>5The ($`p,p^{}`$) system of D-branes with $`pp^{}=4`$ is very special and when one calculates the one-loop interaction amplitude between them (for example, see ), the R-R amplitude has zero contribution while the NS-NS amplitude has two terms cancelled for this case. Since the NS-NS amplitude is insensitive to the sign of the charges carried by the individual D-brane, therefore the force acting between the two branes vanishes, independent of the sign of their charges. Such a no-force condition indicates that there is a certain fraction of spacetime supersymmetry preserved. Further when two kinds of D-branes with different dimensionality intersect, the spacetime supersymmetry preserved is determined by the following two equations: $`ฯต_1=\eta \mathrm{\Gamma }^0\mathrm{\Gamma }^1\mathrm{}\mathrm{\Gamma }^pฯต_2`$ and $`ฯต_1=\eta ^{}\mathrm{\Gamma }^0\mathrm{\Gamma }^1\mathrm{}\mathrm{\Gamma }^p^{}ฯต_2`$ with $`\mathrm{\Gamma }^\mu `$ the ten dimensional $`\gamma `$-matrix, and $`ฯต_1`$ and $`ฯต_2`$ the two Majorana-Weyl supersymmetry parameters in IIA/IIB string theory (In IIA, $`p`$ and $`p^{}`$ are both even while in IIB they are odd). Here $`\eta =\pm `$ and $`\eta ^{}=\pm `$ label the sign of the charge carried by the corresponding branes, respectively. Assuming $`p>p^{}`$ and defining $`\mathrm{\Gamma }=\mathrm{\Gamma }^{p^{}+1}\mathrm{}\mathrm{\Gamma }^p`$, one can show that only for $`\mathrm{\Gamma }^2=I`$ with $`I`$ the unit matrix, i.e., for, $`(1)^{(pp^{})(pp^{}+1)/2}=1`$, a quarter of susy is preserved, independent of the values of $`\eta `$ and $`\eta ^{}`$. In a given theory, $`pp^{}`$ is even and the above condition implies $`pp^{}=4,8`$. As $`p=p^{}`$, $`\mathrm{\Gamma }`$ is a unit matrix, then only for $`\eta =\eta ^{}`$, one half of susy can be preserved. Here we discuss the so-called threshold bound states. For nonthreshold bound state, one half of susy can be preserved for $`pp^{}=2`$ based on U-duality.. Again similar to the previous case here also if in addition to $`\overline{N}_1=N_2=0`$, we have $`\overline{N}_2=0`$ (or $`N_1=0`$) we get half susy configuration which is $`N_1`$ D$`p`$-branes (or delocalized $`\overline{N}_2`$ $`\overline{\mathrm{D}}(p4)`$-branes). (d) This is similar in sprit to the case (c) but now with $`\overline{N}_1`$ and $`N_2`$ non-zero. Other special cases can be discussed accordingly.
When none of the $`N`$โs are zero, in general, the solution is not supersymmetric and there is a tachyon on the world-volume of the intersecting branes<sup>6</sup><sup>6</sup>6The tachyon can be expressed as $`T=\left(\begin{array}{cc}T_1& T_2\\ T_3& T_4\end{array}\right)`$. Here $`T_1`$ is in the bi-fundamental $`(N_1,\overline{N}_1)`$ of gauge group $`U(N_1)\times U(\overline{N}_1)`$ , $`T_2`$ in $`(N_1,\overline{N}_2)`$ of gauge group $`U(N_1)\times U(\overline{N}_2)`$, $`T_3`$ in $`(N_2,\overline{N}_1)`$ of gauge group $`U(N_2)\times U(\overline{N}_1)`$ and $`T_4`$ in $`(N_2,\overline{N}_2)`$ of gauge group $`U(N_2)\times U(\overline{N}_2)`$. The tachyon used in the discussion is actually $`\left(Tr(T\overline{T})\right)^{1/2}`$ with a proper normalization.. However we can always get BPS configuration at the end of tachyon condensation given the above discussion and the one in footnote 5. When $`N_1=\overline{N}_1`$ or $`N_2=\overline{N}_2`$, we will restore one half of susy at the end of tachyon condensation. When the two hold simultaneously, we end up with a flat spacetime at the end of tachyon condensation with maximal supersymmetry. For all the other cases, one quarter of susy is preserved at the end of tachyon condensation. For example, when $`N_1>\overline{N}_1`$ and $`N_2>\overline{N}_2`$ we expect to have intersecting $`(N_1\overline{N}_1)`$ D$`p`$-branes with $`(N_2\overline{N}_2)`$ D$`(p4)`$-branes at the end of tachyon condensation which is supersymmetric (quarter BPS), on the other hand, when $`N_1=\overline{N}_1`$ we expect to have delocalized $`(N_2\overline{N}_2)`$ D$`(p4)`$-branes at the end of tachyon condensation which is half BPS and so on. The recognition for having a supersymmetric background at the end of tachyon condensation is crucial and we use all these information to relate the supergravity parameters to the microscopic physical parameters.
Now in order to understand the tachyon condensation, we look at the expression of the total ADM mass of the intersecting brane-antibrane supergravity solution representing the total energy of the system. We then express this total energy in terms of the five microscopic physical parameters namely, $`N_1`$, $`\overline{N}_1`$, $`N_2`$, $`\overline{N}_2`$ and $`T`$ of the intersecting D$`p`$-$`\overline{\mathrm{D}}p`$ and D$`(p4)`$-$`\overline{\mathrm{D}}(p4)`$ system using the aforementioned relations. The total energy can be seen to be equal to or less than the sum of the masses of $`N_1`$ D$`p`$-branes, $`\overline{N}_1`$ $`\overline{\mathrm{D}}p`$-branes, $`N_2`$ D$`(p4)`$-branes and $`\overline{N}_2`$ $`\overline{\mathrm{D}}(p4)`$-branes indicating the presence of tachyon contributing the negative potential energy to the system. We will see that the energy expression gives the right picture of tachyon condensation as is expected of an intersecting brane-antibrane system. We will reproduce all the expected results from this general mass formula under various special limits at the top and at the bottom of the tachyon potential. We will also show how the various known BPS supergravity configurations can be reproduced in these special limits.
The non-supersymmetric intersecting D$`p`$/D$`(p4)`$ supergravity solution having isometry ISO($`p4,1`$) $`\times `$ SO(4) $`\times `$ SO($`9p`$) which can also be identified as the intersecting D$`p`$-$`\overline{\mathrm{D}}p`$ and D$`(p4)`$-$`\overline{\mathrm{D}}(p4)`$ has the form in space-time dimension $`d=10`$ as<sup>7</sup><sup>7</sup>7This configuration has also been considered previously in but in a rather different notations.,
$`ds^2`$ $`=`$ $`F_1^{\frac{7p}{8}}F_2^{\frac{11p}{8}}\left(dt^2+{\displaystyle \underset{i=1}{\overset{p4}{}}}dx_i^2\right)+F_1^{\frac{7p}{8}}F_2^{\frac{p3}{8}}{\displaystyle \underset{j=p3}{\overset{p}{}}}dx_j^2`$
$`+\left(H\stackrel{~}{H}\right)^{\frac{2}{7p}}F_1^{\frac{p+1}{8}}F_2^{\frac{p3}{8}}\left(dr^2+r^2d\mathrm{\Omega }_{8p}^2\right)`$
$`e^{2\varphi }`$ $`=`$ $`F_1^{\frac{3p}{2}}F_2^{\frac{7p}{2}}\left({\displaystyle \frac{H}{\stackrel{~}{H}}}\right)^{2\delta }`$
$`F_{[8p]}`$ $`=`$ $`b\mathrm{Vol}(\mathrm{\Omega }_{8p})`$
$`F_{[12p]}`$ $`=`$ $`c\mathrm{Vol}(\mathrm{\Omega }_{8p})dx_{p3}\mathrm{}dx_p`$ (1)
Note that in the above we have written the metric in the Einstein frame. The Vol($`\mathrm{\Omega }_{8p}`$) represents the volume form of the unit $`(8p)`$-dimensional sphere. Also the various functions appeared in the solution are defined below,
$`F_1`$ $`=`$ $`\mathrm{cosh}^2\theta _1\left({\displaystyle \frac{H}{\stackrel{~}{H}}}\right)^{\alpha _1}\mathrm{sinh}^2\theta _1\left({\displaystyle \frac{\stackrel{~}{H}}{H}}\right)^{\beta _1}`$
$`F_2`$ $`=`$ $`\mathrm{cosh}^2\theta _2\left({\displaystyle \frac{H}{\stackrel{~}{H}}}\right)^{\alpha _2}\mathrm{sinh}^2\theta _2\left({\displaystyle \frac{\stackrel{~}{H}}{H}}\right)^{\beta _2}`$
$`H`$ $`=`$ $`1+{\displaystyle \frac{\omega ^{7p}}{r^{7p}}},\stackrel{~}{H}=\mathrm{\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}1}{\displaystyle \frac{\omega ^{7p}}{r^{7p}}}`$ (2)
Here $`\alpha _1`$, $`\beta _1`$, $`\alpha _2`$, $`\beta _2`$, $`\delta `$ and $`\omega `$ are integration constants. $`b`$ and $`c`$ are the charge parameters related to the net charges of the D$`p`$-$`\overline{\mathrm{D}}p`$ and D$`(p4)`$-$`\overline{\mathrm{D}}(p4)`$-brane systems respectively. However, not all the parameters are independent and they are related as,
$`\alpha _1\beta _1`$ $`=`$ $`{\displaystyle \frac{p3}{2}}\delta ,b=(7p)(\alpha _1+\beta _1)\omega ^{7p}\mathrm{sinh}2\theta _1`$
$`\alpha _2\beta _2`$ $`=`$ $`{\displaystyle \frac{p7}{2}}\delta ,c=(7p)(\alpha _2+\beta _2)\omega ^{7p}\mathrm{sinh}2\theta _2`$ (3)
and from the consistency of equations of motion we also get
$$(\alpha _1+\beta _1)^2+(\alpha _2+\beta _2)^2+\left(4\frac{(p3)^2}{4}\frac{(p7)^2}{4}\right)\delta ^2=8\frac{8p}{7p}$$
(4)
So, treating $`(\alpha _1+\beta _1)\mathrm{\Delta }_1`$ and $`(\alpha _2+\beta _2)\mathrm{\Delta }_2`$ as independent, there are actually five independent parameters in the solution, namely, $`\theta _1`$, $`\theta _2`$, $`\omega `$, $`\mathrm{\Delta }_1`$ and $`\mathrm{\Delta }_2`$. We will relate each of these parameters to the physical microscopic parameters of the intersecting D$`p`$-$`\overline{\mathrm{D}}p`$ and D$`(p4)`$-$`\overline{\mathrm{D}}(p4)`$ system. But before we do that we would like to make some comments. First of all, we mention that as $`r\mathrm{}`$, both $`H`$ and $`\stackrel{~}{H}`$ $``$ 1, and so, $`F_1,F_21`$. The solution is therefore asymptotically flat. Also, note that the solution (1) represents magnetically charged intersecting brane-antibrane system. To obtain the electrically charged ones we just make a transformation $`F_{[8p]}e^{\frac{3p}{2}\varphi }F_{[8p]}`$ and $`F_{[12p]}e^{\frac{7p}{2}\varphi }F_{[12p]}`$, where $``$ represents the Hodge dual. Using these the gauge fields for the electrically charged solution can be written as,
$`A_{[p+1]}`$ $`=`$ $`\mathrm{sinh}\theta _1\mathrm{cosh}\theta _1\left({\displaystyle \frac{C_1}{F_1}}\right)dtdx_1\mathrm{}dx_p`$
$`A_{[p3]}`$ $`=`$ $`\mathrm{sinh}\theta _2\mathrm{cosh}\theta _2\left({\displaystyle \frac{C_2}{F_2}}\right)dtdx_1\mathrm{}dx_{p4}`$ (5)
where $`C_1`$ and $`C_2`$ are defined as,
$`C_1`$ $`=`$ $`\left({\displaystyle \frac{H}{\stackrel{~}{H}}}\right)^{\alpha _1}\left({\displaystyle \frac{\stackrel{~}{H}}{H}}\right)^{\beta _1}`$
$`C_2`$ $`=`$ $`\left({\displaystyle \frac{H}{\stackrel{~}{H}}}\right)^{\alpha _2}\left({\displaystyle \frac{\stackrel{~}{H}}{H}}\right)^{\beta _2}`$ (6)
We note from (2) that the solution has a potential singularity at $`r=\omega `$ and we will work only in the physically relevant region $`r>\omega `$. Also, without any loss of generality we will choose all of the parameters $`\mathrm{\Delta }_1`$, $`\mathrm{\Delta }_2`$ to be $`0`$. We remark that the solution is non-supersymmetric can be seen from the $`(H\stackrel{~}{H})^{2/(7p)}`$ factor in the last term of the metric in eq.(1). This is consistent with our interpretation of the solution to be intersecting brane-antibrane system. We can now express the parameter $`b`$ and $`c`$ in terms of the number of branes-antibranes using eq.(1) as follows,
$`Q_0^p(N_1\overline{N}_1)`$ $`=`$ $`{\displaystyle \frac{b\mathrm{\Omega }_{8p}}{\sqrt{2}\kappa _0}}b={\displaystyle \frac{\sqrt{2}\kappa _0Q_0^p(N_1\overline{N}_1)}{\mathrm{\Omega }_{8p}}}`$
$`Q_0^{p4}(N_2\overline{N}_2)`$ $`=`$ $`{\displaystyle \frac{c\mathrm{\Omega }_{8p}V_4}{\sqrt{2}\kappa _0}}c={\displaystyle \frac{\sqrt{2}\kappa _0Q_0^{p4}(N_2\overline{N}_2)}{\mathrm{\Omega }_{8p}V_4}}`$ (7)
where $`Q_0^p=(2\pi )^{(72p)/2}\alpha ^{(3p)/2}`$ is the unit charge on the D$`p`$-brane and similarly $`Q_0^{p4}`$ is the unit charge on the D$`(p4)`$-brane. $`\sqrt{2}\kappa _0=(2\pi )^{7/2}\alpha ^2`$ is related to 10 dimensional Newtonโs constant. $`V_4`$ is the compact volume of the four directions $`x_{p3}`$ to $`x_p`$. $`\mathrm{\Omega }_n=2\pi ^{(n+1)/2}/\mathrm{\Gamma }((n+1)/2)`$. Note that $`b0`$ as $`N_1\overline{N}_1`$ and $`c0`$ as $`N_2\overline{N}_2`$ as expected.
For the solution (1), the supersymmetry will be restored if and only if $`H\stackrel{~}{H}1`$ which always requires $`\omega ^{7p}0`$. As we have already stated in the beginning we have the following cases for which supersymmetry will be restored: (1) $`N_1=N_2=0`$, or, $`\overline{N}_1=\overline{N}_2=0`$, or, $`N_1=\overline{N}_2=0`$, or, $`\overline{N}_1=N_2=0`$ (or, $`N_1=N_2=\overline{N}_1=\overline{N}_2=0`$ which is the trivial case). (2) When none of the groups are zero, we can still have supersymmetric configuration at the end of tachyon condensation for which the number of susy restored depends on the initial values of $`N_1,N_2,\overline{N}_1`$ and $`\overline{N}_2`$ as discussed earlier. For example, when $`N_1=\overline{N}_1`$ and $`N_2=\overline{N}_2`$, we expect to get an empty space-time at the end of tachyon condensation with maximal supersymmetry. What is important is the recognition that for all these cases, $`\omega ^{7p}0`$. This observation will be crucial to relate $`\omega `$ as well as other parameters in terms of the physical microscopic parameters of the brane-antibrane system. We also point out that at the top of the tachyon potential when $`\omega ^{7p}`$ does not go to zero, it must give the correct number of branes so that the mass formula can be correctly reproduced. All these information can be captured in a single formula for $`\omega ^{7p}`$ as follows,
$$(7p)\omega ^{7p}=\sqrt{\frac{7p}{2(8p)}}\frac{2\kappa _0^2}{\mathrm{\Omega }_{8p}}T_p\left[\sqrt{N_1\overline{N}_1}+a\sqrt{N_2\overline{N}_2}\right]\mathrm{cos}T$$
(8)
In the above $`T_p=(2\pi )^p\alpha ^{(p+1)/2}`$ is the tension of a D$`p`$-brane. The factor $`2\kappa _0^2T_p/\mathrm{\Omega }_{8p}`$ in front is kept so that it will give the correct mass of the system. Also we have defined $`a=T_{p4}/(V_4T_p)=(2\pi \sqrt{\alpha ^{}})^4/V_4`$ relating the tensions of a D$`p`$-brane and a D$`(p4)`$-brane. We also point out that $`T`$ in the above represents the tachyon vev and we have taken $`T=0`$ as the top of the tachyon potential and $`T=\pi /2`$ as the bottom of the potential. Thus all the cases we have discussed above is contained in $`\omega ^{7p}`$ in eq.(8).
Now having obtained the form of $`\omega ^{7p}`$ in terms of the microscopic physical parameters, we would like to obtain similar relations for the other parameters as well. But first we will try to relate $`\delta `$, $`\mathrm{\Delta }_1`$ and $`\mathrm{\Delta }_2`$ in terms of the microscopic parameters. But before that we observe that the total ADM mass of the intersecting D$`p`$-$`\overline{\mathrm{D}}p`$ and D$`(p4)`$-$`\overline{\mathrm{D}}(p4)`$ can be obtained from the supergravity configuration (1) using the formula given in as,
$$M=\frac{\mathrm{\Omega }_{8p}}{2\kappa _0^2}(7p)\omega ^{7p}\left[\left(\mathrm{\Delta }_1\mathrm{cosh}2\theta _1+\frac{p3}{2}\delta \right)+\left(\mathrm{\Delta }_2\mathrm{cosh}2\theta _2+\frac{p7}{2}\delta \right)\right]$$
(9)
This clearly shows that the ADM mass or the total energy of the system breaks up into two parts corresponding to the D$`p`$-$`\overline{\mathrm{D}}p`$ and the D$`(p4)`$-$`\overline{\mathrm{D}}(p4)`$-brane systems. From our earlier experience on the tachyon condensation on the brane-antibrane system we find that in this case the parameters $`\mathrm{\Delta }_1`$ and $`\mathrm{\Delta }_2`$ satisfy the following relations in terms of $`\delta `$
$`\mathrm{\Delta }_1^2`$ $`=`$ $`{\displaystyle \frac{(p3)^2}{4}}\delta ^2{\displaystyle \frac{(p3)\delta }{\gamma A}}\sqrt{{\displaystyle \frac{(N_1\overline{N}_1)^2}{\mathrm{cos}^2T}}+4N_1\overline{N}_1\mathrm{cos}^2T}+{\displaystyle \frac{4N_1\overline{N}_1\mathrm{cos}^2T}{\gamma ^2A^2}}`$
$`\mathrm{\Delta }_2^2`$ $`=`$ $`{\displaystyle \frac{(p7)^2}{4}}\delta ^2{\displaystyle \frac{(p7)a\delta }{\gamma A}}\sqrt{{\displaystyle \frac{(N_2\overline{N}_2)^2}{\mathrm{cos}^2T}}+4N_2\overline{N}_2\mathrm{cos}^2T}+{\displaystyle \frac{4a^2N_2\overline{N}_2\mathrm{cos}^2T}{\gamma ^2A^2}}`$ (10)
where $`\gamma =\sqrt{(7p)/2(8p)}`$, $`\mathrm{\Delta }_1=\alpha _1+\beta _1`$, $`\mathrm{\Delta }_2=\alpha _2+\beta _2`$ and $`A=\sqrt{N_1\overline{N}_1}+a\sqrt{N_2\overline{N}_2}`$. Then using the parameter relation (4) we find that the parameter $`\delta `$ satisfy the following quadratic relation where only the $`\delta <0`$ root will be relevant for our discussion
$`\delta ^2{\displaystyle \frac{\delta }{4\gamma A}}\left[(p3)\sqrt{{\displaystyle \frac{(N_1\overline{N}_1)^2}{\mathrm{cos}^2T}}+4N_1\overline{N}_1\mathrm{cos}^2T}+a(p7)\sqrt{{\displaystyle \frac{(N_2\overline{N}_2)^2}{\mathrm{cos}^2T}}+4N_2\overline{N}_2\mathrm{cos}^2T}\right]`$
$`{\displaystyle \frac{(N_1\overline{N}_1+a^2N_2\overline{N}_2)\mathrm{sin}^2T+2a\sqrt{N_1N_2\overline{N}_1\overline{N}_2}}{\gamma ^2A^2}}=\mathrm{\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}0}`$ (11)
Thus $`\delta `$ is given entirely in terms of the microscopic physical parameters. As $`T0`$, the $`\delta `$ approaches a finite negative value while as $`T\pi /2`$, $`\delta =4(A/\gamma )\mathrm{cos}T/[(p3)|N_1\overline{N}_1|+a(p7)|N_2\overline{N}_2|]`$ if not all the four $`N_1,N_2,\overline{N}_1,\overline{N}_2`$ are equal and if so, then $`\delta =1/\gamma `$. Once $`\delta `$ is known we can determine $`\mathrm{\Delta }_1`$ and $`\mathrm{\Delta }_2`$ in terms of the microscopic parameters as well. After we determine the forms of $`\mathrm{\Delta }_1`$ and $`\mathrm{\Delta }_2`$ then using the form of $`\omega ^{7p}`$ given in (8) we can easily relate the parameters $`\theta _1`$ and $`\theta _2`$ to the microscopic physical parameters using (3) and (7) as,
$`\mathrm{sinh}2\theta _1`$ $`=`$ $`{\displaystyle \frac{|N_1\overline{N}_1|}{\gamma \mathrm{\Delta }_1A\mathrm{cos}T}},`$
$`\mathrm{sinh}2\theta _2`$ $`=`$ $`{\displaystyle \frac{a|N_2\overline{N}_2|}{\gamma \mathrm{\Delta }_2A\mathrm{cos}T}},`$ (12)
where we have assumed both $`\theta _1`$ and $`\theta _2`$ to be $`0`$. We note from (12) that as $`N_1\overline{N}_1`$, $`\theta _10`$. We have also seen it before from eq.(7) that $`N_1\overline{N}_1`$ implies $`b0`$ and so, $`\theta _10`$ implies the charge $`b0`$. Similarly, $`\theta _20`$ implies the charge $`c0`$. Note from eq.(4) that the values of $`|\mathrm{\Delta }_1|,|\mathrm{\Delta }_2|`$ and $`|\delta |`$ are all bounded from the above, therefore from (12) we see that both $`\theta _1`$ and $`\theta _2`$ blow up at the end of the tachyon condensation.
Now all the quantities in the ADM mass expression given in eq.(9) are known and so, substituting $`\omega ^{7p}`$, $`\mathrm{\Delta }_1`$, $`\mathrm{\Delta }_2`$ and $`\delta `$ we find that the mass expression has a very simplified form given by
$`M`$ $`=`$ $`T_p\left[\sqrt{(N_1\overline{N}_1)^2+4N_1\overline{N}_1\mathrm{cos}^4T}+a\sqrt{(N_2\overline{N}_2)^2+4N_2\overline{N}_2\mathrm{cos}^4T}\right]`$ (13)
$``$ $`T_p\left[(N_1+\overline{N}_1)+a(N_2+\overline{N}_2)\right]`$
We thus note that the total mass per unit $`p`$-brane volume of the system is less or equal to the sum of the those of $`N_1`$ D$`p`$-branes, $`\overline{N}_1`$ $`\overline{\mathrm{D}}p`$-branes, $`N_2`$ D$`(p4)`$-branes, $`\overline{N}_2`$ $`\overline{\mathrm{D}}(p4)`$-branes and the difference is the tachyon potential energy per unit $`p`$-brane volume $`V(T)`$ which is negative. We can easily check that $`T=0`$ gives the maximum of the energy, therefore the maximum of the tachyon potential (which is zero) while $`T=\pi /2`$ gives the corresponding minima. We want to point out that eq.(13) is consistent with our previous experience for a simple brane-antibrane system discussed in and the special feature for the present system under consideration is that no interaction exists between two D-branes with their dimensionality differing by four as discussed in footnote 5.
We will now check one by one whether the above mass formula produces all the required properties of the solution and the tachyon condensation. At $`T=0`$, i.e. at the top of the tachyon potential, $`\mathrm{cos}T=1`$ and we have from above $`M=T_p\left[(N_1+\overline{N}_1)+a(N_2+\overline{N}_2)\right]`$, producing the expected result. It is also easy to see from (8) that at that point $`\omega ^{7p}`$ does not go to zero indicating that the corresponding solution breaks all the supersymmetry as it should be. Note that in this case it does not matter whether we have (i) $`N_1\overline{N}_1`$ and $`N_2\overline{N}_2`$, (ii) $`N_1\overline{N}_1`$ and $`N_2\overline{N}_2`$, (iii) $`N_1\overline{N}_1`$ and $`N_2\overline{N}_2`$ or (iv) $`N_1\overline{N}_1`$ and $`N_2\overline{N}_2`$, in all four cases $`\omega ^{7p}`$ has the same value. Also at $`T=\pi /2`$ i.e. at the bottom of the tachyon potential we get from above $`M=T_p\left(|N_1\overline{N}_1|+a|N_2\overline{N}_2|\right)`$ again producing the expected result. Here also it does not matter whether we have cases (i) to (iv) above, we always have $`\omega ^{7p}`$ going to zero, indicating that we have a supersymmetric configurations. In fact, for (i) we have intersecting $`(N_1\overline{N}_1)`$ D$`p`$ branes and $`(N_2\overline{N}_2)`$ D$`(p4)`$ branes (when none of the inequalities are saturated), $`(N_1\overline{N}_1)`$ D$`p`$-branes (when the second inequality is saturated) and $`(N_2\overline{N}_2)`$ delocalized D$`(p4)`$-branes (when the first inequality is saturated). Similarly for case (ii). For case (iii) we have intersecting $`(N_1\overline{N}_1)`$ D$`p`$ branes and $`(\overline{N}_2N_2)`$ $`\overline{\mathrm{D}}(p4)`$ branes (when none of the inequalities are saturated), $`(N_1\overline{N}_1)`$ D$`p`$-branes (when the second inequality is saturated) and $`(\overline{N}_2N_2)`$ delocalized $`\overline{\mathrm{D}}(p4)`$-branes (when the first inequality is saturated). Similarly for case (iv).
Now we will discuss in a bit detail how one can recover the supersymmetric intersecting brane configuration at the end of tachyon condensation. Let us first discuss the case $`N_1>\overline{N}_1`$ and $`N_2>\overline{N}_2`$. Since for $`T\pi /2`$, $`\omega ^{7p}0`$, both $`H`$ and $`\stackrel{~}{H}`$ goes to unity as can be seen from eq.(2). Also we notice from (12) that both $`\theta _1`$ and $`\theta _2`$ goes to $`\mathrm{}`$. So, if we take a limit $`(\alpha _1+\beta _1)\mathrm{sinh}\theta _1ฯต_1^1`$, $`(\alpha _2+\beta _2)\mathrm{sinh}\theta _2ฯต_2^1`$ and $`\omega ^{7p}ฯต_1\overline{\omega }_1^{7p}`$ and $`\omega ^{7p}ฯต_2\overline{\omega }_2^{7p}`$ for some dimensionless parameters $`ฯต_1,ฯต_20`$, with $`\overline{\omega }_1^{7p}`$ and $`\overline{\omega }_2^{7p}`$ = finite, then we get from (3) $`b=(7p)\overline{\omega }_1^{7p}`$ and $`c=(7p)\overline{\omega }_2^{7p}`$. Both $`F_1`$ and $`F_2`$ in eq.(2) would then reduce to some harmonic functions
$`F_1`$ $``$ $`\overline{H}_1=\mathrm{\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}1}+{\displaystyle \frac{\overline{\omega }_1^{7p}}{r^{7p}}}`$
$`F_2`$ $``$ $`\overline{H}_2=\mathrm{\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}1}+{\displaystyle \frac{\overline{\omega }_2^{7p}}{r^{7p}}}`$ (14)
with $`\overline{\omega }_1^{7p}=b/(7p)`$ and $`\overline{\omega }_2^{7p}=c/(7p)`$. The corresponding configuration as can be seen from (1) with eq.(14) is an intersecting BPS configuration of $`(N_1\overline{N}_1)`$ D$`p`$ branes with $`(N_2\overline{N}_2)`$ D$`(p4)`$ branes. One quarter of susy is restored. The other three cases $`N_1<\overline{N}_1`$ and $`N_2<\overline{N}_2`$, or, $`N_1>\overline{N}_1`$ and $`N_2<\overline{N}_2`$, or, $`N_1<\overline{N}_1`$ and $`N_2>\overline{N}_2`$ will be very similar and we will not repeat the discussion. However when one of them is equal i.e. for the case $`N_1>\overline{N}_1`$ (or $`N_1<\overline{N}_1`$) and $`N_2=\overline{N}_2`$, we find from eq.(12) that $`\theta _2=0`$ which implies from eq.(2) that $`F_21`$ (because $`H`$ and $`\stackrel{~}{H}`$ go to unity). So, in this case we will have just one harmonic function $`\overline{H}_1`$ rather than two. The corresponding solution can be easily checked from (1) to get reduced to $`(N_1\overline{N}_1)`$ D$`p`$-branes (or $`(\overline{N}_1N_1)`$ $`\overline{\mathrm{D}}p`$-branes). Now one half of susy is restored. Similarly for the other case i.e. when $`N_1=\overline{N}_1`$ and $`N_2>\overline{N}_2`$ (or $`N_2<\overline{N}_2`$). The tachyon condensation can also be seen for the special case of $`N_1=\overline{N}_1`$ and $`N_2=\overline{N}_2`$. For this case, at the top of the potential ($`T=0`$) we get from (13) $`M=2T_p(N_1+aN_2)`$ as expected and the corresponding configuration breaks all the supersymmetry. However, at the end of tachyon condensation i.e. at $`T=\pi /2`$, $`M=0`$ corresponding to an empty space-time preserving all supersymmetry. In fact in this case we see from (8) that $`\omega ^{7p}`$ vanishes and so $`\overline{H}_1`$, $`\overline{H}_2`$ = 1. For this case, all susy is restored.
It is also easy to check that the mass formula produces the correct result when $`N_1=\overline{N}_1=0`$ or $`N_2=\overline{N}_2=0`$ or both. For the first case, we have from (13) $`M=a(N_2+\overline{N}_2)T_{p4}/V_4`$ at the top of the potential and $`M=a|N_2\overline{N}_2|T_{p4}/V_4`$ at the bottom. For the second case we have $`M=(N_1+\overline{N}_1)T_p`$ at the top and $`M=|N_1\overline{N}_1|T_p`$ at the bottom as expected. When both $`N_1=\overline{N}_1=0`$ and $`N_2=\overline{N}_2=0`$ we get $`M=0`$ as expected. It is not difficult to check that the space-time configuration at the end of tachyon condensation in these three cases indeed reduce to BPS $`(N_2\overline{N}_2)`$ delocalized D$`(p4)`$-branes (or $`(\overline{N}_2N_2)`$ delocalized $`\overline{\mathrm{D}}(p4)`$-branes), $`(N_1\overline{N}_1)`$ $`Dp`$-branes (or $`(\overline{N}_1N_1)`$ $`\overline{\mathrm{D}}p`$-branes) or an empty space-time respectively. The tachyon vev decouples in all these three cases at the bottom of the tachyon potential as expected.
To summarize, we have interpreted the supergravity solution given in (1) with the metric having the isometry ISO($`p4,1`$) $`\times `$ SO(4) $`\times `$ SO($`9p`$) as the intersecting D$`p`$-$`\overline{\mathrm{D}}p`$ and D$`(p4)`$-$`\overline{\mathrm{D}}(p4)`$-brane system. The five parameters appeared in the supergravity solution were then naturally interpreted as related to the five microscopic physical parameters of the system namely, the number ($`N_1`$) of D$`p`$-branes, the number ($`\overline{N}_1`$) of $`\overline{\mathrm{D}}p`$-branes, the number ($`N_2`$) of D$`(p4)`$-branes, the number ($`\overline{N}_2`$) of $`\overline{\mathrm{D}}(p4)`$-branes and the tachyon vev $`T`$. Based on the physical properties of the solution and the characteristic behavior of tachyon condensation, we have related the supergravity parameters with the microscopic physical parameters of the system. We have obtained the ADM mass representing the total energy per unit p-brane volume of the system from the supergravity configuration and related it with the microscopic physical parameters using the previously mentioned relations. We have shown that the ADM mass as well as the solution produce all the required properties of the tachyon condensation showing that the proposed relations capture the right picture of the tachyon condensation of the intersecting brane-antibrane system in the supergravity or closed string description.
Acknowledgements
One of us (SR) would like to thank the members of the Interdisciplinary Center for Theoretical Study at the University of Science and Technology of China at Hefei, where part of this work was done, for warm hospitality. We also acknowledge support by grants from the Chinese Academy of Sciences and the grants from the NSF of China with 90303002.
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# Low energy hadronic contribution to the QED vacuum polarization
(June 30, 2005)
## Abstract
Recent improvements in the low energy e<sup>+</sup>e<sup>-</sup> annihilation data and their influence on the determination of the hadronic contribution to the running of the QED fine structure constant at m<sub>Z</sub> are discussed. Using CMD-2 and KLOE measurements in the $`\rho `$ region we obtain $`\mathrm{\Delta }\alpha _{\mathrm{had}}^{(5)}(s)`$ = 0.02758 $`\pm `$ 0.00035 at s = $`m_Z^2`$.
In the year 2001, we published an updated evaluation of the hadronic contribution to the running of the QED fine structure constant BurPie2001 , based on a dispersion integral using a parametrization of the measured cross section of e<sup>+</sup>e$`{}_{}{}^{}`$ hadrons. We obtained a hadronic contribution of $`\mathrm{\Delta }\alpha _{\mathrm{had}}^{(5)}(s)`$ = 0.02761 $`\pm `$ 0.00036 at s = $`m_Z^2`$.
Our parametrization in the c.m.s. energy region of the $`\rho `$, the contribution of the $`\pi ^+\pi ^{}`$ final state from threshold to 1.8 GeV, was based on a pion form factor parametrization obtained by the CMD-2 Collaboration which used results of their measurements in the c.m.s. energy region between 0.61 and 0.96 GeV at the VEPP-2M collider CMD2-1999 . The overall uncertainty of the $`\rho `$ region integral, including the statistical uncertainty, was 2.3% (that of $`\mathrm{\Gamma }_{ee}`$ in CMD2-1999 ) in our analysis.
Since then, the CMD-2 collaboration improved the treatment of radiative corrections twice. An intermediate improvement has appeared in the published document CMD2-2001 and an additional improvement has become available in 2004 CMD2-2003 . We have concluded that the most recent CMD-2 results imply only a small change in the estimate of the hadronic contribution BurPie2003 .
Recently, the KLOE collaboration KLOE2004 has measured the cross section of e<sup>+</sup>e$`{}_{}{}^{}\pi ^+\pi ^{}`$ with high statistical accuracy in small energy bins using the โradiative returnโ from the $`\varphi `$ resonance to the $`\rho `$ in the $`\pi ^+\pi ^{}`$ mass range between 0.59 and 0.97 GeV.
We have been repeatedly asked to update our previous analysis and to comment on and quantify the influence of recent low energy measurements by KLOE and CMD-2 on our results. We find that the actual change turns out to be very small. Since the change is very small we have decided to submit this work as a brief report. To the extent that this report is an update of a previously published article, the choice is not to unnecessarily repeat the discussion for energy regions which did not change.
In our 2001 analysis, we used the parametrization of the pion form factor obtained by the CMD-2 collaboration. The contribution of the new results on $`\mathrm{\Delta }\alpha _{\mathrm{had}}^{(5)}(m_Z^2)`$ is now obtained by direct integration between measured KLOE and CMD-2 data points separately. For CMD-2 we use the โbareโ cross section and for KLOE the pion form factor data. These are quantities, in which the vacuum polarization corrections have been removed. The small $`\rho `$ contribution from lower and higher energies, not covered by new data, is evaluated as previously using the CMD-2 parametrization of the pion form factor. We treat the systematic uncertainties as fully correlated between different c.m.s. energies within the CMD-2 experiment. For the integration of the KLOE data, we constructed a covariance matrix based on the statistical covariance matrix with the addition of fully correlated systematic uncertainties as provided by the KLOE collaboration.
The results obtained from the dispersion integration of the KLOE and CMD-2 data at $`m_Z^2`$ are in good agreement with each other. The systematic uncertainty of the CMD-2 integration (0.6%) is smaller than the corresponding uncertainty of the KLOE integration (1.4%). On the other hand, the statistical uncertainty of the CMD-2 integration is slightly larger than the systematic one, while the statistical uncertainty of the KLOE integration is negligible. The integration results are combined as independent measurements in the evaluation of the $`\rho `$ contribution to $`\mathrm{\Delta }\alpha _{\mathrm{had}}^{(5)}(m_Z^2)`$.
We obtain a value of the hadronic contribution to the running of the QED fine structure constant of $`\mathrm{\Delta }\alpha _{\mathrm{had}}^{(5)}(s)`$ = 0.02758 $`\pm `$ 0.00035 at s = $`m_Z^2`$ corresponding to $`1/\alpha ^{(5)}(m_Z^2)=128.940\pm 0.048`$. The value of the $`\rho `$ contribution has changed from 0.00350 in BurPie2001 to 0.00347 and the relative uncertainty has decreased from 2.3% to 0.9%. The change of the uncertainty corresponds to the change of precision from the preliminary CMD-2 CMD2-1999 data to the combination of published CMD-2 CMD2-2003 and KLOE KLOE2004 data. The change of the value and the uncertainty of the hadronic contribution to the running of the QED fine structure constant at $`m_Z^2`$ is very small. In fact the $`\rho `$ region contributes to less than 13% to the dispersion integral and is known to much better precision than many of the other energy domains as can be concluded from the Table 1, which is the updated version from the Ref. BurPie2001 .
We note however, that the shape of the hadronic cross-sections measured by the KLOE and CMD-2 collaborations differ for some individual points by more than the systematic uncertainty would indicate Beijing . There also appears to be a small, but systematic energy shift in the observed cross sections between the KLOE and the CMD-2 data, which at present is not understood. The effect on the integrated cross sections which contribute to $`\mathrm{\Delta }\alpha _{\mathrm{had}}^{(5)}(m_Z^2)`$ is however negligible.
The situation is different for the hadronic contribution to the anomalous magnetic moment of the muon (g-2)<sub>ฮผ</sub> DEHZ ; Jegerlehner ; Teubner ; Davier . There the $`\rho `$ region provides the dominant contribution. The dispersion integral involves a different kernel which gives more weight to lower energies and larger sensitivity on systematic energy shifts.
Fig. 1, which is the updated version from the Ref. BurPie2001 , gives the summary of $`R_{\mathrm{had}}`$ measurements by different experiments and the current precision in different e<sup>+</sup>e<sup>-</sup> center-of-mass (cms) energy regions. $`R_{\mathrm{had}}`$ is the measured QED cross-section of the process e<sup>+</sup>e$`{}_{}{}^{}`$ hadrons, normalized to the QED cross-section for lepton-pair production. The uncertainty in the 1-2 GeV energy region is 15%. This region contributes to about 40% to the uncertainty on dispersion integral at $`m_Z^2`$, as can be seen from the Table 1 and Fig. 2 in the Ref. BurPie2001 . We would like to strongly encourage efforts to measure precisely $`R_{\mathrm{had}}`$ in this cms energy region.
We would like to thank S.I. Eidelman from the CMD-2 Collaboration and A. Denig and G. Venanzoni from the KLOE Collaboration for useful discussions.
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# Hydrodynamic Modes for a Granular Gas from Kinetic Theory
## I Introduction
The use of hydrodynamic equations to describe granular fluids in rapid flow has been in practice for many years Ha83 . The justification for this fluid-like description and prediction of the transport coefficients appearing in these equations has been the focus of attention for some time as well JyR85 ; LSJyCh84 ; GyS95 . In recent years, an accurate derivation of Navier-Stokes order hydrodynamics has been given from the Boltzmann equation for granular gases using an adaptation of the Chapman-Enskog method for normal gases BDKyS98 ; GyS95 . The expressions for the transport coefficients as a function of the degree of inelasticity have been confirmed by both Monte Carlo and molecular dynamics simulation ByC01 . Successful application of these Navier-Stokes equations to a number of different states also supports their validity BRyC99 ; Go03 ; ByR04 . However, the context of the hydrodynamic equations remains uncertain. What are the relevant space and time scales? How much inelasticity can be described in this way?
Such questions can be addressed for gases using the Boltzmann kinetic equation to describe the complete dynamics for properties of interest. Then it can be asked under what conditions do the hydrodynamic excitations emerge as the dominant dynamics. The analysis of this problem for normal gases is quite complete and precise Mc89 ; Scharf . The objective here is to initiate a similar formulation of the problem for granular gases DyB03 . First, solutions to the Boltzmann equation are considered for states that deviate from spatial homogeneity only by small perturbations. The dynamics in this case is governed by the *linear* Boltzmann operator for spatially inhomogeneous states. The spectrum of this operator determines all possible excitations on all space and time scales, and for all degrees of inelasticity. The first problem is to identify points in this spectrum corresponding to hydrodynamics. This is one of the main results reported here. Both the hydrodynamic eigenvalues and eigenfunctions are calculated for long wavelength excitations corresponding to Navier-Stokes order hydrodynamics. Their agreement with the corresponding results from the Chapman-Enskog method is established.
Next, the issue of conditions for the dominance of the hydrodynamic excitations, or modes, is considered. This entails showing that the hydrodynamic eigenvalues are smaller than all other parts of the spectrum, such that there is a long enough time scale for the latter to decay relative to the hydrodynamic modes. There are two new difficulties for granular gases. First, the hydrodynamic eigenvalues cannot be made arbitrarily small since they do not all scale with the inverse wavelength of the perturbation, as for normal gases. Second, there is a new characteristic frequency, the cooling rate, in addition to the collision frequency to set the scale of the spectrum. The cooling rate can be relatively large or small depending on the inelasticity of the gas particle collisions.
Analysis of the non-hydrodynamic spectrum of the Boltzmann equation for a granular gas remains a difficult unsolved problem. Instead, we consider here a simpler kinetic model BDyS99 that retains the exact hydrodynamic spectrum and allows identification of the entire spectrum as well. It is found that the time scale for the hydrodynamic excitations is longer than that for all other excitations, for any degree of inelasticity. A brief report of these results has been given in reference DyB03 , with the details and elaboration given below.
The plan of the paper is as follows. In the next section, a short summary of the Boltzmann equation is given, and the Navier-Stokes equations derived from it by the Chapman-Enskog procedure are recalled. More details are given in Appendix A. Also, the hydrodynamic equations are linearized and the hydrodynamic modes to second order in the wavevector are identified.
In Sec. III, the structure of the linear Boltzmann equation is discussed, emphasizing the relevance of the spectrum of the operator generating the linear dynamics to determine the existence and validity of hydrodynamics. The first of the above issues is addressed in Sec. IV, where the kinetic theory hydrodynamic modes are identified in the long wavelength limit. A technical point associated with the non Hermitian character of the operator is discussed in Appendix B. The results are extended to Navier-stokes order in Sec. V, and the obtained expressions are shown to agree with those derived in Sec. II from the hydrodynamic equations. Details of the calculations are given in Appendix C. The possibility of a description in terms of the hydrodynamic modes is studied in Sec. VI, by means of a model kinetic equation. It is established that there are length and time scales on which only the excitations associated to the hydrodynamic modes persist. The mathematical details are given in Appendix D. Finally, the last section contains a short summary of the main points and conclusions in the paper.
## II Boltzmann Equation and Hydrodynamic Modes
In all of the following, the simplest model of a granular fluid is considered: a low density gas of smooth and inelastic hard spheres ($`d=3`$) or disks ($`d=2`$) of mass $`m`$ and diameter $`\sigma `$. The binary collision rule given below is characterized by a constant coefficient of restitution $`\alpha `$, defined in the interval $`0<\alpha 1`$, and measuring the loss of energy in each collision. At sufficiently low density, the distribution function $`f(๐,๐,t)`$ is determined from the Boltzmann equation GyS95 ; vNyE01 ; ByP04
$$\left(_t+๐\mathbf{}\right)f=J[f,f],$$
(1)
where $`J`$ is the inelastic Boltzmann collision operator defined by
$$J[X,Y]\sigma ^{d1}d๐_1d\widehat{๐}\mathrm{\Theta }(\widehat{๐}๐)\widehat{๐}๐[\alpha ^2X(๐,๐^{},t,)Y(๐,๐_1^{},t)X(๐,๐,t)Y(๐,๐_1,t)],$$
(2)
for arbitrary functions $`X(๐)`$ and $`Y(๐)`$. Here, $`\widehat{๐}`$ is a unit vector along the line joining the centers of the colliding pair, $`\mathrm{\Theta }`$ is the Heaviside step function, and $`๐=๐๐_1`$. The primes on the velocities denote the initial values $`\{๐_1^{},๐_2^{}\}`$ that lead to $`\{๐_1,๐_2\}`$ following a โrestitutingโ binary collision,
$$๐^{}=๐\frac{1}{2}(1+\alpha ^1)(\widehat{๐}๐)\widehat{๐},๐_1^{}=๐_1+\frac{1}{2}(1+\alpha ^1)(\widehat{๐}๐)\widehat{๐}.$$
(3)
The usual Boltmann collision operator is recovered from Eq. (2) in the elastic limit $`\alpha =1`$.
The macroscopic variables of interest are the hydrodynamic fields, i.e. the density $`n(๐,t)`$, the flow velocity $`๐(๐,t)`$, and the (granular) temperature $`T(๐,t)`$. They are defined as moments of the solution to the Boltzmann equation,
$$\left(\begin{array}{c}n(๐,t)\\ n(๐,t)๐(๐,t)\\ \frac{d}{2}n(๐,t)T(๐,t)\end{array}\right)=๐๐\left(\begin{array}{c}1\\ ๐\\ \frac{1}{2}m\left(๐๐\right)^2\end{array}\right)f(๐,๐,t).$$
(4)
An exact set of equations for these variables is obtained from the following properties of the collision operator
$$๐๐\left(\begin{array}{c}1\\ ๐\\ \frac{1}{2}m\left(๐๐\right)^2\end{array}\right)J[f,f]=\left(\begin{array}{c}0\\ \mathrm{๐}\\ \frac{d}{2}nT\zeta \end{array}\right).$$
(5)
The first two equations follow from conservation of mass and momentum in the particle collisions. The last equation reflects the loss of energy in collisions due to the inelasticity. This appears through the โcooling rateโ $`\zeta (๐,t)0`$, defined by this equation.
The macroscopic balance equations resulting from the above properties are:
$$_tn+\mathbf{}\left(n๐\right)=0,$$
(6)
$$\left(_t+๐\mathbf{}\right)๐+(mn)^1\left[\mathbf{}(nT)+\mathbf{}\mathsf{\Pi }\right]=0,$$
(7)
$$(_t+๐\mathbf{}+\zeta )T+\frac{2}{nd}(nT\mathbf{}๐+\mathsf{\Pi }:\mathbf{}๐ฎ+\mathbf{}๐)=0.$$
(8)
The functionals giving the dissipative part of the pressure tensor, $`\mathsf{\Pi }`$, and the heat flux, $`๐ช`$, are also moments of the solution to the Boltzmann equation,
$$\mathsf{\Pi }(๐ซ,t)=๐๐m\left(๐ฝ๐ฝ\frac{V^2}{d}๐จ\right)f(๐,๐,t),$$
(9)
$$๐(๐,t)=๐๐\left(\frac{m}{2}V^2\frac{d+2}{2}T\right)๐ฝf(๐,๐,t),$$
(10)
where $`๐จ`$ is the unit tensor of dimension $`d`$ and $`๐ฝ๐๐`$ the so-called peculiar velocity. Equations (6)-(8) are the basis for a hydrodynamic description, once $`\mathsf{\Pi }(๐,t)`$, $`๐(๐,t)`$, and $`\zeta (๐,t)`$ are specified. These can be obtained from their definitions using a solution of the Boltzmann equation generated by the Chapman-Enskog method. This method assumes the existence of a solution whose space and time dependence is given entirely through the hydrodynamic fields and their gradients (i.e., a โnormalโ solution). As a result, $`\mathsf{\Pi }(๐,t)`$ and $`๐(๐,t)`$ are given in terms of these variables, and Eqs. (6)-(8) become a closed set of hydrodynamic equations. The primary results of this method are recalled in Appendix A. To leading order in the spatial gradients, the dissipative fluxes are found to be given by BDKyS98 ; ByC01
$$\mathsf{\Pi }_{ij}(๐,t)=\eta \left(_iu_j+_ju_i\frac{2}{d}\delta _{ij}\mathbf{}๐\right),$$
(11)
$$๐ช(๐ซ,t)=\kappa \mathbf{}T\mu \mathbf{}n.$$
(12)
For the cooling rate the result is:
$$\zeta (๐,t)=n\sigma ^{d1}\left(\frac{2T}{m}\right)^{1/2}\zeta _0^{}+\zeta _1^2T+\zeta _2^2n+\text{bilinear in }\mathbf{}n,\mathbf{}T,\mathbf{}๐\text{ terms},$$
(13)
with $`\zeta _0^{}`$ being a dimensionless positive constant proportional to $`\left(1\alpha ^2\right)`$. The explicit forms for the shear viscosity $`\eta `$, the thermal conductivity $`\kappa `$, and the new transport coefficients $`\mu `$, $`\zeta _1`$, and $`\zeta _2`$, peculiar to granular gases, are given in Appendix A. The nonlinear contributions to $`\zeta (t)`$ indicated in (13) play no role in the following linear analysis and will be not discussed further in this paper.
Equations (6)-(8) together with the โconstitutive relationsโ (11)-(13) are the Navier-Stokes hydrodynamic equations for a granular gas. A special solution of these equations is that for a homogeneous state, characterized by a constant density $`n_{hcs}`$, a vanishing velocity flow $`๐ฎ_{hcs}=0`$, and a uniform temperature $`T_{hcs}(t)`$ determined from
$$_tT_{hcs}\left(t\right)=n_{hcs}\sigma ^{d1}\left(\frac{2}{m}\right)^{1/2}\zeta _0^{}T_{hcs}^{3/2}\left(t\right).$$
(14)
This is referred to as the homogeneous cooling state (HCS), since the only macroscopic dynamics is the monotonic decrease in the temperature field. It is easily seen that the HCS as defined above is also a solution of the exact balance equations, taking into account that $`\zeta _{hcs}n_{hcs}\sigma ^{d1}(T_{hcs}/m)^{1/2}`$ on dimensional grounds. In the following, the solution to the Boltzmann equation will be considered for small perturbations of the HCS, and it will be useful to have the corresponding results from Navier-Stokes hydrodynamics. The linearized Navier-Stokes equations have time dependent coefficients, since the distribution function of the reference HCS depends on time. This complication can be overcome by the introduction of suitable dimensionless space and time scales, as well as scaled hydrodynamic fields. Thus we define $`๐^{}`$ and $`s`$ by
$$๐^{}=\frac{๐}{\mathrm{}},ds=\frac{v_{hcs}(t)}{\mathrm{}}dt.$$
(15)
Here, $`v_{hcs}\left(t\right)\left[2T_{hcs}(t)/m\right]^{1/2}`$ is the thermal velocity in the HCS, and $`\mathrm{}(n_{hcs}\sigma ^{d1})^1`$ is proportional to the mean free path of the particles. The dimensionless fields $`\delta y_j(๐^{},s)`$ are chosen as
$$\rho (๐^{},s)=\frac{n(๐,t)n_{hcs}}{n_{hcs}},\theta (๐^{},s)=\frac{T(๐,t)T_{hcs}(t)}{T_{hcs}(t)},๐(๐^{},s)=\frac{๐(๐,t)}{v_{hcs}(t)}.$$
(16)
Besides, since the equations are linear and the reference state is homogeneous, it is sufficient to consider a single Fourier mode,
$$\delta y_j(๐^{},s)=e^{i๐๐^{}}\delta \stackrel{~}{y}_j(๐,s).$$
(17)
The velocity components are chosen as a longitudinal component relative to $`๐`$, $`\stackrel{~}{\omega }_{}(๐,s)=\widehat{๐}\stackrel{~}{๐}(๐,s)`$, and $`d1`$ transverse components $`\stackrel{~}{\omega }_{,i}(๐,s)=\widehat{๐}^i\stackrel{~}{๐}(๐,s)`$, where $`\{\widehat{๐}๐/k,\widehat{๐}^{(i)},i=1,\mathrm{},d1\}`$ are a set of $`d`$ pairwise orthogonal unit vectors . The dimensionless linear Navier-Stokes equations then become a system of ordinary differential equations with constant coefficients that can be expressed in the compact form BDKyS98
$$_s\delta \stackrel{~}{๐}(๐,s)+๐ช\left(k\right)\delta \stackrel{~}{๐}(๐,s)=0,$$
(18)
where we use a $`d+2`$ dimensional space representation with $`\delta \stackrel{~}{๐}(๐,s)`$ being the vector
$$\delta \stackrel{~}{๐}(๐,s)=\left(\begin{array}{c}\stackrel{~}{\rho }(๐,s)\hfill \\ \stackrel{~}{\theta }(๐,s)\hfill \\ \stackrel{~}{\omega }_{}(๐,s)\hfill \\ \stackrel{~}{๐}_{}(๐,s)\hfill \end{array}\right),$$
(19)
and $`\stackrel{~}{๐}_{}(๐,s)`$ denoting the vector formed by the $`d1`$ components $`\stackrel{~}{\omega }_{,i}(๐,s)`$. The matrix $`๐ช\left(k\right)`$ is block diagonal, expressing the decoupling of transverse and longitudinal modes,
$$๐ช(k)=\left(\begin{array}{cc}๐ช_1& 0\\ 0& ๐ช_2\end{array}\right),$$
(20)
$$๐ช_1(k)=\left(\begin{array}{ccc}0& 0& ik\\ \zeta _0^{}+\left(\mu ^{}\zeta _2^{}\right)k^2& \frac{\zeta _0^{}}{2}+\left(\kappa ^{}\zeta _1^{}\right)k^2& \frac{2i}{d}k\\ \frac{i}{2}k& \frac{i}{2}k& \frac{\zeta _0^{}}{2}+\frac{2(d1)}{d}\eta ^{}k^2\end{array}\right),$$
(21)
$$๐ช_2(k)=\left(\frac{\zeta _0^{}}{2}\eta ^{}k^2\right)๐จ.$$
(22)
In the above expressions, $`๐ช_2`$ and $`๐จ`$ are matrices of dimension $`d1`$. The dimensionless transport coefficients used in the above equations are defined by
$$\eta ^{}=\frac{\eta }{\mathrm{}mn_{hcs}v_{hcs}},\kappa ^{}=\frac{2\kappa }{d\mathrm{}n_{hcs}v_{hcs}},\mu ^{}=\frac{2\mu }{d\mathrm{}T_{hcs}v_{hcs}},$$
(23)
$$\zeta _1^{}=\frac{T_{hcs}\zeta _1}{\mathrm{}v_{hcs}},\zeta _2^{}=\frac{n_{hcs}\zeta _2}{\mathrm{}v_{hcs}}.$$
(24)
The formal solution to the initial value problem (18) is
$$\delta \stackrel{~}{๐}(๐,s)=e^{๐ช\left(k\right)s}\delta \stackrel{~}{๐}(๐,0).$$
(25)
The eigenvalues and eigenvectors of the generator $`๐ช`$ for this dynamics define the $`d+2`$ Navier-Stokes order hydrodynamic modes. They are given by the solutions of the equation
$$๐ช(k)๐_j(๐)=\lambda _j(k)๐_j(๐),j=1,\mathrm{},d+2.$$
(26)
A simple calculation provides the expressions for the eigenvalues to order $`k^2`$. They are given by
$$\lambda _1(k)=\frac{k^2}{\zeta _0^{}},\lambda _2(k)=\frac{\zeta _0^{}}{2}\left(\frac{d+1}{d\zeta _0^{}}\kappa ^{}+\zeta _1^{}\right)k^2,$$
$$\lambda _{}(k)=\frac{\zeta _0^{}}{2}+\left[\frac{1}{d\zeta _0^{}}+\frac{2(d1)\eta ^{}}{d}\right]k^2,\lambda _{}(k)=\frac{\zeta _0^{}}{2}+\eta ^{}k^2,$$
(27)
the eigenvalue $`\lambda _{}(k)`$ being $`(d1)`$-fold degenerate. The corresponding eigenvectors to leading order in $`k`$ are
$$๐_1(๐)=\left(\begin{array}{c}1\\ 2\\ 0\\ \mathrm{๐}\end{array}\right),๐_2(๐)=\left(\begin{array}{c}0\\ 1\\ 0\\ \mathrm{๐}\end{array}\right),๐_{}(๐)=\left(\begin{array}{c}0\\ 0\\ 1\\ \mathrm{๐}\end{array}\right),$$
$$๐_{,i}(๐)=\left(\begin{array}{c}0\\ 0\\ 0\\ \widehat{๐}\end{array}\right).$$
(28)
Here,
* $`\mathrm{๐}=0`$, and $`\widehat{๐}=\widehat{\mathrm{๐}}=1`$, for $`d=2`$.
* $`\mathrm{๐}=\left(\begin{array}{c}0\\ 0\end{array}\right)`$, $`\widehat{\mathrm{๐}}=\left(\begin{array}{c}1\\ 0\end{array}\right)`$, and $`\widehat{\mathrm{๐}}=\left(\begin{array}{c}0\\ 1\end{array}\right)`$, for $`d=3`$.
The first of these modes is excited by the condition $`\stackrel{~}{\theta }(๐,0)=2\stackrel{~}{\rho }(๐,0)`$ at zero flow velocity. The second is produced by a temperature perturbation at constant density and also zero velocity, while the third one corresponds to a longitudinal velocity perturbation at constant temperature and density. There is a $`(d1)`$-fold degeneracy for the shear modes of eigenvalue $`\lambda _{}`$. These diffusive modes are excited by a perturbation of the velocity field in the transverse plane orthogonal to $`๐`$.
It should be noted that while the above analysis is restricted to the Navier-Stokes equations, derived from the Boltzmann equation by the Chapman-Enskog method, the eigenvalues and eigenvectors to order $`k`$ follow more generally from the exact macroscopic balance equations and do not depend on the approximate constitutive equations (11) and (12). The hydrodynamic modes sought by kinetic theory in the subsequent sections can therefore be defined as those excitations due to small perturbations which agree with the above in the long wavelength limit. Analyticity then allows extension of that identification to shorter wavelengths. A consistency check of the Chapman-Enskog method is agreement with the above results at order $`k^2`$. However, the concept of hydrodynamic modes in this context does not require the validity of the Chapman-Enskog method nor the limitation to the Navier-Stokes approximation.
## III Linear Boltzmann Equation
A more complete and accurate description of the response to small spatial perturbations of the density, temperature, and flow velocity is obtained directly from the Boltzmann equation. Consider first an isolated system. As already indicated in the previous section, the exact balance equations have a solution describing the HCS, with a monotonically decreasing temperature obeying Eq. (14). The solution to the Boltzmann equation corresponding to this macroscopic state is characterized by the scaling form GyS95
$$f_{hcs}(๐,t)=n_{hcs}v_{hcs}^d(t)\varphi \left(V^{}\right),$$
(29)
with
$$๐ฝ^{}=\frac{๐ฝ}{v_{hcs}(t)}=\frac{๐๐_{hcs}}{v_{hcs}(t)}.$$
(30)
Substitution into the Boltzmann equation gives
$$\frac{\zeta _{hcs}}{2}\frac{}{๐ฝ}\left(๐f_{hcs}\right)=J[f_{hcs},f_{hcs}].$$
(31)
For later convenience, a constant velocity $`๐ฎ_{hcs}`$ for the system as a whole has been included, although this can always be removed by means of a Gallilean transformation. Then, in the following it will be considered that $`๐_{hcs}=0`$, unless it be explicitly established otherwise. An exact and explicit solution of Eq. (31) is not known yet, but the behavior of $`\varphi `$ at large and small velocities has been determined vNyE01 and the results obtained by the direct simulation Monte Carlo method strongly supports the existence of such a scaling form BRyC96 .
The HCS distribution function is a โuniversalโ homogeneous solution in the same sense as the Maxwellian for elastic collisions. An arbitrary homogeneous state is expected to approach the HCS after a few collisions. Therefore, in discussing response of any homogeneous state to small spatial perturbations, it is sufficient to consider the HCS as the reference state. All the other cases will simply induce additional short time transients.
Consider then small spatial perturbations of the HCS
$$f(๐,๐,t)=f_{hcs}(๐,t)\left[1+\mathrm{\Delta }(๐,๐,t)\right],|\mathrm{\Delta }(๐,๐,t)|1.$$
(32)
To linear order in $`\mathrm{\Delta }(๐,๐,t)`$, the Boltzmann equation becomes
$$\left[_t+๐\mathbf{}+L(t)\right](f_{hcs}\mathrm{\Delta })=0,$$
(33)
where $`L(t)`$ is the linearized Boltzmann collision operator given by
$$L(t)X(๐)=J[f_{hcs},X]J[X,f_{hcs}],$$
(34)
for arbitrary $`X(๐)`$. Just as for the analysis of the Navier-Stokes equations in the previous section, the above linear kinetic equation takes a simpler form when expressed in terms of the dimensionless variables defined in Eq. (15), and considering a single Fourier mode
$$\mathrm{\Delta }(๐,๐,t)=e^{i๐๐^{}}\stackrel{~}{\mathrm{\Delta }}(๐,๐^{},s),$$
(35)
where $`๐^{}=๐/v_{hcs}(t)`$. Then, Eq. (33) becomes
$$\left(_s+i๐๐^{}+^{}\right)\stackrel{~}{\mathrm{\Delta }}(๐,๐^{},s)=0.$$
(36)
The dimensionless operator $`^{}`$ is now time independent and it is given by
$$^{}X\frac{\zeta _0^{}}{2}\varphi ^1\frac{}{๐^{}}\left(๐^{}\varphi X\right)+L^{}X,$$
(37)
where $`L^{}`$ is the dimensionless linear Boltzmann collision operator
$$L^{}X\varphi ^1\left(J^{}[\varphi ,\varphi X]+J^{}[\varphi X,\varphi ]\right),$$
(38)
$$J^{}[X,Y]๐๐_1^{}๐\widehat{๐}\mathrm{\Theta }(\widehat{๐}๐^{})\widehat{๐}๐^{}\left[\alpha ^2X(๐^{})Y(๐_1^{})X(๐^{})Y(๐_1^{})\right].$$
(39)
Here, $`๐^{}`$ and $`๐_1^{}`$ are related with $`๐^{}`$ and $`๐_1^{}`$ by Eqs. (3). The operator $`^{}`$ differs from the linearized Boltzmann collision operator $`L^{}`$ by terms representing the cooling effects of the inelastic collisions. The latter arise because the derivative with respect to $`s`$ is taken at constant $`๐^{}`$ rather than $`๐`$.
Solutions to the dimensionless, linear kinetic equation (36) for $`\stackrel{~}{\mathrm{\Delta }}`$ are sought in a Hilbert space whose scalar product is defined by
$$(X,Y)=๐๐ฏ^{}\varphi \left(v^{}\right)X^{}\left(๐^{}\right)Y(๐^{}),$$
(40)
with the dagger denoting complex conjugation. The formal solution to the kinetic equation is
$$\stackrel{~}{\mathrm{\Delta }}(๐,๐^{},s)=\frac{1}{2\pi i}๐ze^{zs}(z)\mathrm{\Delta }^{}(๐,๐^{},0),$$
(41)
$$(z)\left(zi๐๐^{}^{}\right)^1,$$
(42)
where the contour encloses the entire spectrum of $`i๐๐^{}+^{}`$ , both point and residual, counterclockwise. It is important to realize that *all the linear excitations of the granular gas are determined from this spectrum.* This formulation of the problem for small spatial perturbations provides a precise context for addressing many questions regarding hydrodynamics for a granular gas. The existence of the hydrodynamic modes and their role relative to other dynamical processes are determined by the characterization of the above spectrum. To see how this occurs, it is useful to recall briefly the status of the corresponding problem for gases with elastic collisions Mc89 ; Scharf . In that case, it has been proved that the hydrodynamic modes exist as $`d+2`$ poles located at the origin in the long wavelength limit and corresponding to the local conserved quantities. Furthermore, the remainder of the spectrum is bounded away from these poles, and the spectrum is analytic in $`๐`$ about $`๐=0`$, so this isolation of the hydrodynamic modes is preserved at finite wavelengths.
## IV Existence of Hydrodynamic Modes
The spectrum of $`i๐๐^{}+^{}`$ is expected to be quite complex, based on the special case of elastic collisions, with points spectra, continua, and limit points. The hydrodynamic excitations, whenever they exist, are part of the point spectrum so in order to investigate them it suffices to consider the eigenvalue problem
$$\left(i๐๐^{}+^{}\right)\mathrm{\Psi }_i=\lambda _i\left(k\right)\mathrm{\Psi }_i.$$
(43)
The search for hydrodynamic excitations can be carried out by assuming they are analytic in $`k`$ and looking first for the $`๐=0`$ solutions of Eq. (43). The practical issue of constructing these modes at finite $`k`$ is addressed in the next section.
The central idea for constructing the hydrodynamic eigenvalues and eigenvectors at $`k=0`$ is the note that the HCS is parameterized by the hydrodynamic fields $`n_{hcs},`$ $`T_{hcs}`$, and $`๐_{hcs}`$, which is now considered different from zero, as discussed at the beginning of the previous section. Therefore, differentiating the Boltzmann equation for the distribution function of the HCS, Eq. (31), with respect to these fields gives exact properties of the linearized Boltzmann collision operator. For example,
$$\frac{}{n_{hcs}}\left\{\frac{\zeta _{hcs}}{2}\frac{}{๐ฝ}\left(๐f_{hcs}\right)J[f_{hcs},f_{hcs}]\right\}=0$$
(44)
gives directly
$$\frac{1}{2}\frac{\zeta _{hcs}}{n_{hcs}}\frac{}{๐ฝ}\left(๐f_{hcs}\right)+\frac{\zeta _{hcs}}{2}\frac{}{๐ฝ}\left(๐ฝ\frac{f_{hcs}}{n_{hcs}}\right)+L\frac{f_{hcs}}{n_{hcs}}=0.$$
(45)
The derivatives of $`\zeta _{hcs}`$ and $`f_{hcs}`$ are easily calculated from the properties $`\zeta _{hcs}n_{hcs}`$ and $`f_{hcs}n_{hcs}`$. In terms of the dimensionless variables, and setting $`๐ฎ_{hcs}=0`$, Eq. (45) becomes
$$^{}1=\frac{\zeta _0^{}}{2}\varphi ^1\frac{}{๐^{}}\left(๐^{}\varphi \right).$$
(46)
Next, calculate the derivative of Eq. (31) with respect to $`T_{hcs}`$,
$$\frac{1}{2}\frac{\zeta _{hcs}}{T_{hcs}}\frac{}{๐ฝ}\left(๐ฝf_{hcs}\right)+\frac{\zeta _{hcs}}{2}\frac{}{๐ฝ}\left(๐\frac{f_{hcs}}{T_{hcs}}\right)+L\frac{f_{hcs}}{T_{hcs}}=0.$$
(47)
Since $`f_{hcs}`$ has the scaling form (29), it is
$$\frac{f_{hcs}}{T_{hcs}}=\frac{f_{hcs}}{2T_{hcs}}\left(d+๐ฝ\frac{\mathrm{ln}f_{hcs}}{๐ฝ}\right),$$
(48)
and, taking into account that $`\zeta _{hcs}T_{hcs}^{1/2}`$, Eq. (47) becomes
$`{\displaystyle \frac{\zeta _{hcs}}{2}}{\displaystyle \frac{}{๐ฝ}}\left(๐ฝf_{hcs}\right)`$ $``$ $`{\displaystyle \frac{\zeta _{hcs}}{2}}{\displaystyle \frac{}{๐ฝ}}\left[๐ฝf_{hcs}\left(d+๐ฝ{\displaystyle \frac{\mathrm{ln}f_{hcs}}{๐ฝ}}\right)\right]`$ (49)
$``$ $`L\left[f_{hcs}\left(d+๐ฝ{\displaystyle \frac{\mathrm{ln}f_{hcs}}{๐ฝ}}\right)\right]=0.`$
Setting $`๐_{hcs}=0`$ and transforming to dimensionless variables, the above equation yields
$$^{}\left(d+๐^{}\frac{\mathrm{ln}\varphi }{๐^{}}\right)=\frac{\zeta _0^{}}{2}\left(d+๐^{}\frac{\mathrm{ln}\varphi }{๐^{}}\right).$$
(50)
Finally, differentiating Eq. (31) with respect to $`๐ฎ_{hcs}`$, afterwards setting $`๐ฎ_{hcs}=0`$, and introducing dimensionless variables leads in a similar way to
$$^{}\left(\frac{\mathrm{ln}\varphi }{๐^{}}\right)=\frac{\zeta _0^{}}{2}\left(\frac{\mathrm{ln}\varphi }{๐^{}}\right).$$
(51)
The $`d+2`$ equations (46), (50), and (51) provide exact properties of $`^{}`$. In fact, Eqs. (50) and (51) are of the form of the eigenvalue problem to be solved at $`k=0`$, with eigenvalues given by $`\zeta _0^{}/2`$ and $`\zeta _0^{}/2`$, respectively. It is straightforward to construct linear combinations of (46) and (50) to obtain an additional eigenvalue and eigenvector. The results can then be expressed as DyB03
$$^{}\mathrm{\Psi }_i(0)=\lambda _i\left(0\right)\mathrm{\Psi }_i(0),i=1,\mathrm{},d+2,$$
(52)
with
$$\left\{\lambda _i(0)\right\}=\{0,\frac{\zeta _0^{}}{2},\frac{\zeta _0^{}}{2}\},$$
(53)
$$\left\{\mathrm{\Psi }_i(0)\right\}=\{d+1+๐^{}\frac{\mathrm{ln}\varphi }{๐^{}},d๐^{}\frac{\mathrm{ln}\varphi }{๐^{}},\widehat{๐}\frac{\mathrm{ln}\varphi }{๐^{}},\frac{\mathrm{ln}\varphi }{๐_{}^{}}\}.$$
(54)
The eigenvalue $`\zeta _0^{}/2`$ is $`d`$-fold degenerated. For convenience for the perturbation calculation to be carried out in the following section, the subspace associated to it has been rearranged. The velocity $`๐^{}`$ has been decomposed into its component in the direction of $`๐`$, $`v_{}^{}=\widehat{๐}๐^{}`$, and the remaining $`d1`$ ones, forming with it a set of $`d`$ pairwise orthogonal components, i.e. $`๐_{}^{}`$ is defined by the $`d1`$ components $`๐^{}\widehat{๐}^{(i)}`$. When appropriate, the set of the $`d`$ eigenfunctions associated with the eigenvalue $`\zeta _0^{}/2`$ will be denoted by $`๐ฟ_3(0)`$.
Clearly, the above are the long wavelength limit of the hydrodynamic modes defined by Eq. (27). This is a the first primary result of the analysis developed here, demonstration of the existence of hydrodynamic excitations in the spectrum of the linearized Boltzmann equation. The results are exact and apply for arbitrary degree of dissipation.
These long wavelength eigenfunctions are determined from the HCS distribution which depends only on the magnitude of $`v^{}`$, so the terms on the right side of Eq. (54) are all determined from $`\mathrm{ln}\varphi \left(v^{}\right)/v^{}`$. In the elastic limit, $`\varphi \left(v^{}\right)`$ becomes Gaussian, and the hydrodynamic eigenfunctions become linear combinations of $`1,`$ $`๐^{}`$, and $`v^2`$, i.e. of the summational invariants for the conservation laws of mass, momentum, and energy, as expected. For inelastic collisions, the eigenfunctions are quite different, particularly at large $`v^{}`$ due to an overpopulation in the HCS distribution relative to the Gaussian. This is illustrated in Figs. 1 and 2, where $`\mathrm{ln}\varphi \left(v\right)/v^{}`$ has been obtained from Direct Simulation Monte Carlo solution of the Boltzmann equation for the HCS Bi94 .
The set of eigenfunctions $`\mathrm{\Psi }_i(0)`$ span a $`d+2`$ dimensional subspace of the Hilbert space, but they are not orthogonal. Consistently, it is easily verified that the operator $`^{}`$ is not Hermitian. Therefore, it is useful to introduce a biorthogonal set of functions $`\mathrm{\Phi }_i,i=1,\mathrm{},d+2`$, with the requirement
$$(\mathrm{\Phi }_i,\mathrm{\Psi }_j(0))=\delta _{ij}.$$
(55)
To identify the appropriate set of functions, first note that
$$(1,^{}X)=0,(๐^{},^{}X)=\frac{\zeta _0^{}}{2}(๐^{},X),$$
(56)
as a consequence of the number of particles and momentum conservation in the moment equations (5). Therefore, $`1`$ and $`๐^{}`$ are eigenfunctions of the adjoint of $`^{}`$, $`^{}`$, with eigenvalues $`0`$ and -$`\zeta _0^{}/2`$, respectively. This gives the set
$$\left\{\mathrm{\Phi }_i\right\}=\{1,\mathrm{\Phi }_2,\widehat{๐}๐^{},๐_{}^{}\}.$$
(57)
The final choice of $`\mathrm{\Phi }_2`$ does not appear to be unique. This is discussed further in Appendix B. For the purposes of the next section, it suffices to make the choice
$$\left\{\mathrm{\Phi }_i\right\}=\{1,\frac{v^2}{d}+\frac{1}{2},\widehat{๐}๐^{},๐_{}^{}\}.$$
(58)
The function $`\mathrm{\Phi }_2`$ is not an eigenfunction of the adjoint operator $`^{}`$, but it is easily verified that the biorthogonality conditions (55) are satisfied.
## V Navier-Stokes Order Modes
In the previous section, the hydrodynamic modes were identified in the long wavelength limit. Assuming analyticity, their existence at finite wavevectors can be inferred. Furthermore, their explicit construction is possible by perturbation theory. This construction provides a critical test of the internal consistency of other quite different approaches (e.g., the Chapman-Enskog method discussed above and in Appendix A). In particular, the detailed form of the eigenvalues and the dependence of the associated transport coefficients on the restitution coefficient should be exactly the same. This is demonstrated to Navier-Stokes order in this section.
Return now to the eigenvalue problem (43) and consider the case of $`k<<1`$. Look for solutions with the expansion (a more complete characterization of the conditions for this perturbation expansion is given below),
$$\mathrm{\Psi }_i(๐ค)=\mathrm{\Psi }_i(0)+k\mathrm{\Psi }_i^{(1)}+k^2\mathrm{\Psi }_i^{(2)}+\mathrm{},$$
(59)
$$\lambda _i(๐ค)=\lambda _i(0)+k\lambda _i^{(1)}+k^2\lambda _i^{(2)}+\mathrm{}.$$
(60)
The reference eigenfunctions $`\mathrm{\Psi }_i(0)`$ and eigenvalues $`\lambda _i(0)`$ are taken to be the long wavelength hydrodynamic results of Eqs. (53) and (54). As already indicated, there is a $`d`$-fold degeneracy for the eigenvalue $`\lambda (0)=\zeta _0^{}/2`$. However, the $`d`$-dimensional subspace spanned by its eigenvectors was naturally partitioned by symmetry into the longitudinal and transverse components. The eigenvectors for the transverse modes decouple from the remaining three modes even at finite wavevector for the same symmetry reasons, so there are no complications of degenerate perturbation theory. In the longitudinal subspace, all the eigenvalues are distinct, except in the elastic limit $`\alpha =1`$, where there is a three-fold degeneracy in this subspace. Thus the two cases of unperturbed reference states with $`\alpha =1`$ and $`\alpha <1`$ must be distinguished. In the former case, the eigenvalues behave as
$$\lambda _i(๐,\alpha )=\left[\lambda _i(0)+k\lambda _i^{(1)}+k^2\lambda _i^{(2)}+\mathrm{}\right]_{\alpha =1}+\text{terms of order }(1\alpha ),$$
(61)
and are regular in $`(1\alpha )`$. This occurs when the degree of dissipation because of inelasticity is small relative to the effects of the spatial variation. The corresponding eigenvalues are then similar to those of a normal gas, with $`d1`$ shear diffusion modes, two sound modes, and a heat diffusion mode. Here attention is restricted to the more interesting and relevant second case of fixed $`\alpha <1`$ with small spatial perturbations. It will be seen that the modes are now qualitatively different, consistently with the results reported in Sec. II, since the degeneracy of the elastic limit is lifted at the outset by the finite dissipation.
To set up the perturbation expansion, projection operators $`๐ซ_i`$ for the biorthogonal set $`\{\mathrm{\Phi }_i,\mathrm{\Psi }_i(0)\}`$ are defined by
$$๐ซ_iX=\mathrm{\Psi }_i(0)(\mathrm{\Phi }_i,X),$$
(62)
for an arbitrary element $`X`$ in the Hilbert space. The eigenvalue problem (43) can be rearranged as
$$\left[^{}\lambda _i\left(0\right)\right]\mathrm{\Psi }_i\left(๐\right)=\left[\lambda _i\left(k\right)\lambda _i\left(0\right)i๐๐^{}\right]\mathrm{\Psi }_i\left(๐\right).$$
(63)
Then, operating on both sides of this equation with $`๐ฌ_i1๐ซ_i`$, and using the property $`\left[^{}\lambda _i\left(0\right)\right]๐ซ_i=0`$, this becomes
$$๐ฌ_i\left[^{}\lambda _i\left(0\right)\right]๐ฌ_i\mathrm{\Psi }_i\left(๐\right)=๐ฌ_i\left[\lambda _i\left(k\right)\lambda _i\left(0\right)i๐๐^{}\right]\mathrm{\Psi }_i\left(๐\right).$$
(64)
By construction, the right side of the above equation is orthogonal to the null space for the adjoint of $`\left[^{}\lambda _i\left(0\right)\right]`$ , and the Fredholm alternative assures solutions to this equation Fred . The eigenvalue problem for $`\mathrm{\Psi }_i\left(๐\right)`$ is determined only up to an overall scale factor, amounting to the choice of normalization. It is convenient to choose
$$(\mathrm{\Phi }_i,\mathrm{\Psi }_i\left(๐ค\right))=1,$$
(65)
implying $`(\mathrm{\Phi }_i,\mathrm{\Psi }_i^{(n)})=0`$ for $`n1`$. With this, Eq. (64) gives two sets of equations for the eigenvalues and eigenvectors,
$$\lambda _i\left(๐\right)=\lambda _i\left(0\right)+(\mathrm{\Phi }_i,i๐๐^{}\mathrm{\Psi }_i\left(๐ค\right))+(\mathrm{\Phi }_i,\left[^{}\lambda _i\left(0\right)\right]๐ฌ_i\mathrm{\Psi }_i\left(๐ค\right)),$$
(66)
$$๐ฌ_i\mathrm{\Psi }_i\left(๐\right)=\left\{๐ฌ_i\left[^{}\lambda _i\left(\mathrm{๐}\right)\right]๐ฌ_i\right\}^1๐ฌ_i\left[\lambda _i\left(k\right)\lambda _i\left(0\right)i๐๐^{}\right]\mathrm{\Psi }_i\left(๐\right).$$
(67)
To zeroth order in $`k`$, these equations give the hydrodynamic modes of the last section in the long wavelength limit, consistently. To first order in $`k`$, the eigenvectors are
$`\mathrm{\Psi }_i^{(1)}(๐)`$ $`=`$ $`๐ฌ_i\mathrm{\Psi }_i^{(1)}(๐)=\left\{๐ฌ_i\left[^{}\lambda _i\left(0\right)\right]\right\}^1๐ฌ_i\left(\lambda _i^{(1)}i\widehat{๐}๐^{}\right)\mathrm{\Psi }_i\left(0\right)`$ (68)
$`=`$ $`\left\{๐ฌ_i\left[^{}\lambda _i\left(0\right)\right]\right\}^1๐ฌ_ii\widehat{๐}๐^{}\mathrm{\Psi }_i\left(0\right).`$
The first equality in the above transformations is a consequence of the normalization condition (65). For the first order eigenvalues it is found:
$$\lambda _i^{(1)}=(\mathrm{\Phi }_i,i\widehat{๐}๐^{}\mathrm{\Psi }_i\left(0\right))+(\mathrm{\Phi }_i,^{}๐ฌ_i\mathrm{\Psi }_i^{(1)})=0.$$
(69)
Now the last equality follows from the fact that each $`\mathrm{\Phi }_i`$ and $`\mathrm{\Psi }_i\left(0\right)`$ have the same definite parity under the change $`๐^{}๐^{}`$ and the distribution function of the HCS, $`\varphi (v^{})`$, defining the scalar product is invariant under this change. Thus the first term on the right side vanishes. The second term also vanishes for similar reasons, since $`\mathrm{\Phi }_i`$ and $`\mathrm{\Psi }_i^{(1)}`$ have opposite parity and $`^{}`$ is invariant under the change in sign of the velocity.
To second order in $`k`$ the eigenvectors and eigenvalues are given by
$$๐ฌ_i\mathrm{\Psi }_i^{(2)}=\left\{๐ฌ_i\left[^{}\lambda _i\left(0\right)\right]\right\}^1๐ฌ_ii\widehat{๐}๐^{}\mathrm{\Psi }_i^{(1)},$$
(70)
$$\lambda _i^{(2)}=(\mathrm{\Phi }_i,i๐๐^{}\mathrm{\Psi }_i^{(1)})+(\mathrm{\Phi }_i,^{}๐ฌ_i\mathrm{\Psi }_i^{(2)}).$$
(71)
The above expressions for the eigenvalues are evaluated in Appendix C. The results have the same forms as given in Eqs. (27). Furthermore, the reduced transport coefficients $`\eta ^{},`$ $`\kappa ^{}`$, and $`\zeta _1^{}`$ are determined from the same integral equations as following from the Chapman-Enskog solution summarized in Appendix A. This confirms that the hydrodynamic modes determined from the spectrum of the linearized Boltzmann equation are consistently determined to Navier-Stokes order by both methods.
## VI โAgeingโ to Hydrodynamics
The existence of hydrodynamic excitations only assures that there is a hydrodynamic contribution to the dynamics of small perturbations of the HCS. To establish a description in terms of these hydrodynamic excitations alone, it is necessary to characterize the rest of the excitations in the spectrum. For gases with elastic collisions, it has been shown that for sufficiently small $`k`$ the hydrodynamic excitations are smaller in magnitude than all other excitations, and bounded away from them Mc89 ; Scharf . Consequently, there is a time scale beyond which only the hydrodynamic excitations persist. It is on this space and time scales that hydrodynamics in the usual sense applies. Typically, the conditions are wavelengths larger than the mean free path and times later than the mean free time. This leaves a large domain of macroscopic space and time scales for hydrodynamics.
The extension of this concept of โageing to hydrodynamicsโ for granular gases is expected, but its verification is not so straightforward. The mathematical analysis for elastic collisions does not transfer to the granular gas due to the significant differences in the linear collision operator. There are qualitative differences in the hydrodynamic modes. For example, the fact that energy is not conserved means that the hydrodynamic excitations cannot be made arbitrarily small simply by making $`k`$ small. The fastest decaying hydrodynamic modes is that with eigenvalue $`\zeta _0^{}/2`$ at long wavelengths. This has its maximum value at large dissipation, and the question arises as to whether the time scale for this mode can become comparable to or exceed those of the non-hydrodynamic modes at strong dissipation.
### VI.1 Model Kinetic Equation
Current analysis of the spectrum of the linearized Boltzmann operator for granular gases appears to be limited to the hydrodynamic excitations discussed here, with no characterization of the rest of the spectrum as yet. Consequently, in the remainder of this presentation these questions are addressed in the context of a model kinetic equation . This model BMyD96 is an extension of the familiar Bhatnager, Gross, Krook (BGK) single relaxation time model for normal gases Ce75 . The Boltzmann equation can be formally written in the form
$$\left(\frac{}{t}+๐ฏ\right)f=\nu \left(fg\right).$$
(72)
There are two significant differences of the model considered here with respect to the original Boltzmann equation. First, the collision frequency is replaced by a velocity independent function of the local density and temperature, $`\nu =\nu (n,T)`$. Second, the gain term of the Boltzmann equation is replaced by $`\nu g`$, where $`g`$ is taken to be a Gaussian function of the velocity,
$$g(๐,๐,t)=n\left[\frac{b(T)}{\pi }\right]^{d/2}e^{b(T)V^2}.$$
(73)
Here $`๐ฝ`$ is the peculiar velocity defined bellow Eq. (10). The parameter $`b`$ is chosen to be function of the local density and temperature so as to enforce the moment conditions (5) above. This leads to the identification
$$b(T)=\frac{m}{2T\left(1\zeta /\nu \right)},$$
(74)
In this way, it is assured that the exact macroscopic balance equations (6)-(10) are preserved by the model. Then the Chapman-Enskog method leads to the same Navier-Stokes hydrodynamic equations as for the Boltzmann equation, with only the transport coefficients being different. In the following, it will be considered that the expression for the collision frequency $`\nu `$ is chosen as scaling with $`nT^{1/2}`$, in order to mimic the hard sphere behavior. Dimensional analysis then implies the same scaling for the cooling rate $`\zeta `$. Moreover, note that consistency of the model kinetic equation requires that $`\zeta <\nu `$ for all values of $`\alpha `$. A possible choice for the cooling rate is to be the same as obtained from the Boltzmann equation by using a local Maxwellian for the distribution function. In the same spirit, the collision frequency can be fixed by fitting one of the transport coefficients of the model to that obtained from the Boltzamnn equation by the Chapman-Enskog procedure in the first Sonine approximation. If the shear viscosity $`\eta `$ is used, the above choices lead to Baskaran04
$$\nu (๐,t)=\frac{(33\alpha +2d)(1+\alpha )}{4d}\nu _0(๐,t),$$
(75)
$$\zeta (๐,t)=\frac{(2+d)(1\alpha ^2)}{4d}\nu _0(๐,t),$$
(76)
where $`\nu _0(๐,t)`$ is an average local collision frequency,
$$\nu _0(๐,t)=\frac{8\pi ^{(d1)/2}n\sigma ^{d1}}{(2+d)\mathrm{\Gamma }(d/2)}\left(\frac{T}{m}\right)^{1/2}.$$
(77)
The above expressions yield
$$\frac{\zeta (๐,t)}{\nu (๐,t)}=\frac{(2+d)(1\alpha )}{33\alpha +2d}\frac{2+d}{3+2d}<1,$$
(78)
in agreement with the model consistency requirement.
This kinetic model reduces to the BGK model for normal gases at $`\alpha =1`$ Ce75 . Otherwise it reproduces all of the qualitative features of the granular Boltzmann equation, including a nontrivial HCS and the same hydrodynamic excitations discussed in the sections above. In some respects, the model kinetic equation is more complex than the Boltzmann equation since the collision operator is a nonlinear functional of $`f`$ through the dependence of $`g`$ on $`T`$ and $`๐ฎ`$. However, the linearized model kinetic equation for small perturbations of the HCS is considerably simpler than that for the Boltzmann equation, as shown below.
### VI.2 Model HCS and Linear Model Kinetic Equation
The HCS equation (31) for this model becomes in dimensionless form
$$\frac{\zeta _0^{}}{2}\frac{}{๐^{}}\left(๐^{}\varphi \right)+\nu _0^{}\varphi =\nu _0^{}\left[\pi \left(1\zeta _0^{}/\nu _0^{}\right)\right]^{d/2}\mathrm{exp}\left(\frac{v^2}{1\zeta _0^{}/\nu _0^{}}\right),$$
(79)
where $`\zeta _0^{}\mathrm{}\zeta _{hcs}/v_{hcs}(t)`$, as in the preceding sections and, consistently, $`\nu _0^{}\mathrm{}\nu _{hcs}/v_{hcs}(t)`$. The solution to the above equation is
$$\varphi (v^{})=\nu _0^{}\left[\pi \left(1\zeta _0^{}/\nu _0^{}\right)\right]^{d/2}_0^{\mathrm{}}๐s\mathrm{exp}\left[\left(\frac{d\zeta _0^{}}{2}+\nu _0^{}\right)s\right]\mathrm{exp}\left(\frac{e^{\zeta _0^{}s}v^2}{1\zeta _0^{}/\nu _0^{}}\right).$$
(80)
It is easily verified that this distribution exhibits algebraic decay for large velocities,
$$\varphi (v^{})\frac{p\pi ^{d/2}}{2}\left(\frac{p}{p2}\right)^{p/2}\mathrm{\Gamma }\left(\frac{p+d}{2}\right)v^{\left(p+d\right)},$$
(81)
with $`p=2\nu _0^{}/\zeta _0^{}`$. Therefore, moments of degree $`p`$ or greater do not exist.
The linearized kinetic model equation for small perturbations of the HCS is obtained in Appendix D, with the result, in the dimensionless variables of the previous sections,
$$\left(_s+i๐๐^{}+_m^{}\right)\stackrel{~}{\mathrm{\Delta }}(๐,๐^{},s)=0.$$
(82)
The linear collision operator in this case is
$$_m^{}=\underset{i}{}\lambda _i(0)๐ซ_i+๐ฌ\left(\nu _{hcs}^{}+\frac{\zeta _0^{}}{2}\varphi ^1\frac{}{๐^{}}๐^{}\varphi \right)๐ฌ.$$
(83)
The projection operators $`๐ซ_i`$ are the same as defined in Eq. (62) and
$$๐ซ\underset{i}{}๐ซ_i,๐ฌ1๐ซ.$$
(84)
In the above expressions, the summations are over the hydrodynamic modes. The first term on the right hand side of Eq. (83) is the projection onto the hydrodynamic eigenfunctions, while the second one is orthogonal to this subspace. Consequently, the spectrum of $`_m^{}`$ for the model kinetic equation has the same $`k=0`$ hydrodynamic eigenfunctions and eigenvalues as the Boltzmann equation, i.e. it is
$$_m^{}\mathrm{\Psi }_i(0)=\lambda _i\left(0\right)\mathrm{\Psi }_i(0),i=1,\mathrm{},d+2.$$
(85)
Furthermore, the structure of $`_m^{}`$ decomposes into operators defined in the hydrodynamic subspace and its orthogonal complement. This allows more detailed analysis of the non-hydrodynamic spectrum. The associated eigenfunctions lie in the orthogonal complement $`๐ฌ`$. Consider the general form for $`\mathrm{\Psi }_Q=๐ฌ\mathrm{\Psi }`$
$`(\mathrm{\Psi }_Q,[^{}{\displaystyle \underset{i}{}}\lambda _i(0)๐ซ_i]\mathrm{\Psi }_Q)`$ $`=`$ $`(\mathrm{\Psi }_Q,\left(\nu _0^{}+{\displaystyle \frac{\zeta _0^{}}{2}}\varphi ^1{\displaystyle \frac{}{๐^{}}}๐^{}\varphi \right)\mathrm{\Psi }_Q)`$ (86)
$`=`$ $`(\left(\nu _0^{}{\displaystyle \frac{\zeta _0^{}}{2}}๐^{}{\displaystyle \frac{}{๐^{}}}\right)\mathrm{\Psi }_Q,\mathrm{\Psi }_Q).`$
The second term in the scalar product above can be simplified as
$`{\displaystyle \frac{\zeta _0^{}}{2}}(๐^{}{\displaystyle \frac{}{๐^{}}}\mathrm{\Psi }_Q,\mathrm{\Psi }_Q)`$ $`=`$ $`{\displaystyle \frac{\zeta _0^{}}{4}}(1,๐^{}{\displaystyle \frac{}{๐^{}}}\mathrm{\Psi }_Q^2)`$ (87)
$`=`$ $`{\displaystyle \frac{\zeta _0^{}}{4}}(\varphi ^1{\displaystyle \frac{}{๐^{}}}\left(๐^{}\varphi \right),\mathrm{\Psi }_Q^2)`$
$`=`$ $`{\displaystyle \frac{\nu _0^{}}{2}}(\mathrm{\Psi }_Q,\mathrm{\Psi }_Q){\displaystyle \frac{\nu _0^{}}{2}}(\varphi ^1g_0^{},\mathrm{\Psi }_Q^2),`$
where
$$g_0^{}(v^{})=n_{hcs}^1v^d(t)g_{hcs}(v,t)0.$$
(88)
Then Eq. (86) gives the desired inequality
$$\frac{(\mathrm{\Psi }_Q,\left[_m^{}_i\lambda _i(0)๐ซ_i\right]\mathrm{\Psi }_Q)}{(\mathrm{\Psi }_Q,\mathrm{\Psi }_Q)}=\frac{\nu _0^{}}{2}\left[1+\frac{(\mathrm{\Psi }_Q,\varphi ^1g_0^{}\mathrm{\Psi }_Q)}{(\mathrm{\Psi }_Q,\mathrm{\Psi }_Q)}\right]>\frac{\nu _0^{}}{2}.$$
(89)
Therefore, the non-hydrodynamic spectrum of $`_m^{}`$ consists of points (or continuum) with real parts larger than $`\nu _0^{}/2`$. The fastest decaying hydrodynamic eigenvalue is at $`\zeta _0^{}/2`$. The hydrodynamic excitations are isolated from the rest of the spectrum for $`\zeta _0^{}<\nu _0^{}`$, as it is always the case. Assuming analyticity in $`k`$, the hydrodynamic modes at finite wavelength also will be isolated from the rest of the spectrum for sufficiently small $`k`$.
In summary, the kinetic model considered here illustrates the expected behavior for the Boltzmann equation. The linearized kinetic equation for small perturbations of the HCS characterizes the complete complex response. Among the excitations there are $`d+2`$ hydrodynamic modes. At long wavelengths, these modes are isolated from and have smaller eigenvalues than the rest of the spectrum. Hence, there is a sufficiently long time scale on which only the hydrodynamic excitations persist.
This kinetic model can be used to explore the role of hydrodynamics in great detail as the linear equation can be solved exactly Baskaran04 . For example, it has been shown that the hydrodynamic modes extend to very short wavelengths far beyond the validity of the Navier-Stokes approximation.
## VII Discussion
The objective here has been to explore the role of hydrodynamics for a granular gas by a direct analysis of the spectrum of the linear inelastic Boltzmann equation. This analysis has been shown to lead to results equivalent to those obtained previously based on the Chapman-Enskog method to solve the Boltzmann equation. However, the current method is more straightforward and less susceptible to subjective questions about the applicability of the Chapman-Enskog method to granular gases. In addition, the formulation of the problem in terms of the spectrum of the linearized collision operator is the proper setting to explore the context for dominance of a hydrodynamic description.
The eigenvalue problem posed in Sec. IV, together with the macroscopic balance equations leads to a precise definition of the hydrodynamic modes. Here, it has been established that they exist for sufficiently long wavelengths. It remains to explore the form and extent to which they are meaningful at shorter wavelengths, and other methods to solve the eigenvalue problem are available to complement the simple perturbation theory applied here. This has been done in ref. Baskaran04 for the kinetic model introduced in Sec. V, demonstrating the extension of the hydrodynamic modes to wavelengths an order of magnitude shorter than those required for the Navier-Stokes approximation .
A second relevant result of the analysis here is the identification of the hydrodynamic eigenfunctions. This allows calculation of the hydrodynamic component for properties of interest, and provides the means to explore dynamical mechanisms beyond the Boltzmann description based on hydrodynamics (e.g., mode coupling phenomena for fluctuations). For normal gases, these eigenfunctions are the summational invariants $`(1,v^2,๐)`$ in the long wavelength limit. For a granular gas, they are replaced by derivatives of the logarithm of the HCS distribution function, which behave quite differently for large velocities, as illustrated in Figs. 1 and 2.
The hydrodynamic excitations are usually interpreted as a set of $`d+2`$ modes associated with the density, temperature, and flow field of the macroscopic balance equations. Here, they are simply $`d+2`$ points associated with the spectrum of the linearized inelastic Boltzmann equation for response to small perturbations. It is often claimed that the temperature should not be included in the set of hydrodynamic fields, as the energy is not conserved. However, the analysis here shows that in the proper reduced variables all but one eigenvalue is non-zero in the long wavelength limit. Thus, the exclusion of one field from the description will not recover a smaller set of eigenvalues clustered around zero at long wavelengths. Instead, the $`d+2`$ modes have the following qualitative behavior. They are all clustered near zero for weak dissipation, but become increasingly separated at stronger dissipation. This means that within the hydrodynamic description there are time scales set by both the wavelength and the cooling rate and these can be quite different.
A second question is the isolation of these hydrodynamic modes from the rest of the spectrum. At weak dissipation, it is expected that this is the case since the hydrodynamic eigenvalues are all small and the non-hydrodynamic spectrum is expected to be of the order of the collision frequency. At strong dissipation, there is a hydrodynamic mode of the order of $`\zeta _0^{}/2`$, which can be of the order of the collision frequency. It remains an open question regarding the size of this eigenvalue relative to the non-hydrodynamic spectrum. It is possible that the isolation of the hydrodynamic spectrum places some restriction on the degree of dissipation. However, the analysis based on the kinetic model of Sec. VI suggests this may not be the case. For the kinetic model, all of the hydrodynamic spectra remain isolated from all of the non-hydrodynamic spectra for any degree of dissipation. In this case, the $`d+2`$ hydrodynamic modes define a dominant set on a sufficiently large time scale, even when the separation of times among the hydrodynamic modes is significant.
The idealized model for a granular gas discussed here can be made more realistic by considering a more complex binary collision rule. In particular, a velocity dependent restitution coefficient (which approaches unity as the relative velocity goes to zero) and tangential friction are two additional qualitative features of real granular fluids. However, the general form of the hydrodynamic equations (macroscopic balance equations) is unchanged in that case; only the detailed values of the transport coefficients, pressure, and cooling rates differ from those of the present model. Similarly, the collision operator of the Boltzmann equation becomes more complex but its properties relevant for the macroscopic balance equations are the same. Most of the analysis given here depends only on those properties rather than the detailed form of the collision operator. Consequently, it is expected that the implications of our simple model extend to more realistic models as well.
The analysis of the linear Boltzmann collision operator for normal, elastic gases is quite complete Mc89 . It is hoped that the beginning provided here for granular gases will provoke the new mathematical analysis required for this case as well.
## VIII Acknowledgements
J.W.D thanks A. Baskaran for helpful discussions and both authors thank M.J. Ruiz-Montero for providing Figs. 1 and 2. The research of J.W.D. was supported in part by Department of Energy Grant DE-FG02ER54677. The research of J.J.B. was supported by the Ministerio de Educaciรณn y Ciencia (Spain) through Grant No. BFM2005-01398 (partially financed by FEDER funds).
## Appendix A Chapman-Enskog Results
The Chapman-Enskog method constructs a solution to the Boltzmann equation whose space and time dependence occurs entirely through the hydrodynamic fields and their gradients. For small spatial gradients, i.e. small relative variation of the hydrodynamic fields over a mean free path, the solution reads (we use the same notation as in ref. BDKyS98 )
$$f(๐,๐,t)=f_{hcs}^{(0)}(๐,๐,t)+๐\mathbf{}\mathrm{ln}T+\mathbf{}\mathrm{ln}n+๐ข:\mathbf{}๐+^2T+๐ฉ^2n.$$
(90)
Consider the contributions to the heat flux and the momentum flux from the terms of first order in the gradients of the above expression. Since these fluxes appear in the macroscopic balance equations under a gradient, the resulting contributions are of second order in the gradients (Navier-Stokes order). For consistency, the cooling rate which does not occur under a gradient in the balance equations must be calculated to second order. Therefore, the last two terms on the right hand side of Eq. (90) give contributions to the cooling rate at this order, but lead to higher order (Burnett) terms in the heat and momentum fluxes. Additional non-linear terms of second order in the gradients coming from the cooling rate have been omitted in Eq. (90) as only the linear hydrodynamic equations are considered here. The reference distribution function $`f_{hcs}^{(0)}(๐,๐,t)`$ has the same functional form as the distribution function of the HCS discussed at the beginning of Sec. III, but scaled with respect to the local exact hydrodynamic fields at time $`t`$ for the generally non-uniform, non-stationary state. Therefore, it can be considered as the local HCS distribution function.
The dissipative part of the pressure tensor and the heat flux are given by Eqs. (11) and (12), respectively. The transport coefficients in these equations are determined from the functions $`๐`$, $``$, and $`๐ข`$ appearing in Eq. (90 ) through
$$\eta =\frac{1}{d^2+d2}๐๐๐ฃ(๐):๐ข(๐),\kappa =\frac{1}{dT}๐๐๐บ(๐)๐(๐),$$
$$\mu =\frac{1}{dn}๐๐๐บ(๐)(๐),$$
(91)
where the functions $`๐ฃ(๐)`$ and $`๐บ(๐)`$ are defined by
$$๐ฃm\left(๐๐\frac{v^2}{d}๐จ\right),๐บ=\left(\frac{mv^2}{2}\frac{d+2}{2}T\right)๐.$$
(92)
The functions $`๐(๐,n,T)`$, $`(๐,n,T)`$, and $`๐ข(๐,n,t)`$ are solutions of the integral equations
$$\left(\zeta _{hcs}^{(0)}T\frac{}{T}+L^{(0)}\frac{\zeta _{hcs}^{(0)}}{2}\right)๐=๐จ,$$
(93)
$$\left(\zeta _{hcs}^{(0)}T\frac{}{T}+L^{(0)}\right)=๐ฉ+\zeta _{hcs}^{(0)}๐,$$
(94)
$$\left(\zeta _{hcs}^{(0)}T\frac{}{T}+L^{(0)}\right)๐ข=๐ฆ,$$
(95)
with the definitions
$$L^{(0)}(t)X(๐)=J[f_{hcs}^{(0)},X]J[X,f_{hcs}^{(0)}],$$
(96)
$$๐จ=\frac{๐}{2}\frac{}{๐}\left(๐f_{hcs}^{(0)}\right)\frac{T}{m}\frac{f_{hcs}^{(0)}}{๐},๐ฉ=๐f_{hcs}^{(0)}\frac{T}{m}\frac{f_{hcs}^{(0)}}{๐},$$
$$๐ฆ=\frac{}{๐}\left(๐f_{hcs}^{(0)}\right)\frac{๐จ}{d}\frac{}{๐}\left(๐f_{hcs}^{(0)}\right).$$
(97)
Moreover, $`\zeta _{hcs}^{(0)}`$ is the local form of $`\zeta _{hcs}`$.
The transport coefficients arising from the cooling rate are best described by first making the expression for the cooling rate more explicit. In general, it is given by a bilinear functional of the distribution function,
$$\zeta =\zeta [f,f],$$
(98)
where
$`\zeta [X,Y]`$ $`=`$ $`{\displaystyle \frac{m}{dnT}}{\displaystyle ๐๐v^2J[X,Y]}`$ (99)
$`=`$ $`\left(1\alpha ^2\right){\displaystyle \frac{m\pi ^{\left(d1\right)/2}\sigma ^{d1}}{4dnT\mathrm{\Gamma }\left(\frac{d+3}{2}\right)}}{\displaystyle ๐๐๐๐_1g^3X(๐)Y(๐_1)}=\zeta [Y,X].`$
The linear contributions from the cooling rate to second order order in the gradients are then
$$\zeta _L^{(2)}=\zeta _1^2T+\zeta _2^2n,$$
(100)
with
$$\zeta _1=2\zeta [,f_{hcs}^{(0)}],\zeta _2=2\zeta [๐ฉ,f_{hcs}^{(0)}].$$
(101)
The integral equations for $``$ and $`๐ฉ`$ are found to be
$$\left(\zeta _{hcs}^{(0)}T\frac{}{T}\frac{3}{2}\zeta _{hcs}^{(0)}+L^{(0)}\right)=T\zeta _1\frac{f_{hcs}^{(0)}}{T}\left(\frac{2\kappa }{dn}\frac{f_{hcs}^{(0)}}{T}+\frac{1}{dT}๐๐\right),$$
(102)
$`\left(\zeta _{hcs}^{(0)}T{\displaystyle \frac{}{T}}+L^{(0)}\right)๐ฉ`$ $`=`$ $`T\zeta _2{\displaystyle \frac{f_{hcs}^{(0)}}{T}}+{\displaystyle \frac{T\zeta _{hcs}^{(0)}}{n}}`$ (103)
$`\left({\displaystyle \frac{2\mu }{dn}}{\displaystyle \frac{f_{hcs}^{(0)}}{T}}+{\displaystyle \frac{1}{dn}}๐\right).`$
The above results can be easily transformed to be expressed in terms of the reduced units and quantities introduced in Secs. II and III. Dimensionless transport coefficients are again defined by Eqs. (23) and (24), but replacing $`n_{hcs}`$ and $`T_{hcs}`$ by their local values $`n`$ and $`T`$. Of course, now it is $`\mathrm{}=(n\sigma ^{d1})^1`$ and $`v_{hcs}(t)`$ is replaced by $`v(t)=(2T/m)^{1/2}`$. In this way, it is obtained:
$$\eta ^{}=\frac{1}{d^2+d2}\underset{i,j}{}(D_{ij}^{},C_{ij}^{}),\kappa ^{}=\frac{1}{d^2}\underset{i}{}(S_i^{},๐_i^{}),$$
$$\mu ^{}=\frac{1}{d^2}\underset{i}{}(S_i^{},_i^{}),$$
(104)
with the dimensionless functions $`๐ฃ^{}`$ and $`๐บ^{}`$ given by
$$๐ฃ^{}=๐^{}๐^{}\frac{v^2}{d}๐จ,๐บ^{}=\left(v^2\frac{d+2}{2}\right)๐^{}.$$
(105)
Moreover, $`๐^{}`$, $`^{}`$, and $`๐ข^{}`$ are now defined through the equations
$$\left(^{}\frac{\zeta _0^{}}{2}\right)๐^{}=๐จ^{},$$
(106)
$$^{}^{}=๐ฉ^{}+\zeta _0^{}๐^{},$$
(107)
$$\left(^{}+\frac{\zeta _0^{}}{2}\right)๐ข^{}=๐ฆ^{},$$
(108)
where the operator $`^{}`$ is defined in Eq. (37) and
$$๐จ^{}=๐^{}\mathrm{\Psi }_2(0)๐ฟ_3(0),๐ฉ^{}=2๐^{}๐ฟ_3(0),$$
$$๐ฆ^{}=๐^{}๐ฟ_3(0)\frac{๐จ}{d}๐^{}๐ฟ_3(0).$$
(109)
The expressions of the reduced transport coefficients $`\zeta _1^{}`$ and $`\zeta _2^{}`$ are:
$$\zeta _1^{}=\frac{1}{d}(a,^{}),\zeta _2^{}=\frac{1}{d}(a,๐ฉ^{}),$$
(110)
where
$$a(๐^{})=\frac{(1\alpha ^2)\pi ^{(d1)/2}}{2\mathrm{\Gamma }\left(\frac{d+3}{2}\right)}๐๐_1^{}\varphi (v_1^{})g^3.$$
(111)
The integral equations obeyed by $`^{}`$ and $`๐ฉ^{}`$ are
$$\left(^{}\frac{\zeta _0^{}}{2}\right)^{}=\left(\zeta _1^{}\kappa ^{}\right)\mathrm{\Psi }_2(0)+\frac{1}{d}๐^{}๐^{},$$
(112)
$$^{}๐ฉ^{}=\zeta _0^{}^{}+\left(\zeta _2^{}\mu ^{}\right)\mathrm{\Psi }_2(0)+\frac{1}{d}^{}๐^{}.$$
(113)
For later use, it is convenient to elaborate more the above expression for $`\zeta _1^{}`$. By construction, the velocity integrals of $``$ times $`1`$, $`๐`$, and $`v^2`$ vanish. This is equivalent to say that $`^{}`$ is orthogonal to the set of functions $`\mathrm{\Phi }_i`$ defined in Eq. (58), and in particular it verifies $`^{}=๐ฌ_2^{}`$, with $`๐ฌ_2=1๐ซ_2`$, $`๐ซ_i`$ being the projection operator defined in Eq. (62). Then, acting with $`๐ฌ_2`$ on both sides of Eq. (112) it is obtained:
$$^{}=๐ฌ_2^{}=\frac{1}{d}\left[๐ฌ_2\left(^{}\frac{\zeta _0^{}}{2}\right)\right]^1๐ฌ_2๐^{}๐^{}.$$
(114)
For the same reason, Eq. (106) yields
$`๐^{}`$ $`=`$ $`๐ฌ_2๐^{}=\left[๐ฌ_2\left(^{}{\displaystyle \frac{\zeta _0^{}}{2}}\right)\right]^1๐ฌ_2๐จ^{}`$ (115)
$`=`$ $`\left[๐ฌ_2\left(^{}{\displaystyle \frac{\zeta _0^{}}{2}}\right)\right]^1๐ฌ_2๐^{}\mathrm{\Psi }_2(0)\left[๐ฌ_2\left(^{}{\displaystyle \frac{\zeta _0^{}}{2}}\right)\right]^1๐ฟ_3(0)`$
$`=`$ $`\left[๐ฌ_2\left(^{}{\displaystyle \frac{\zeta _0^{}}{2}}\right)\right]^1๐ฌ_2๐^{}\mathrm{\Psi }_2(0)+{\displaystyle \frac{1}{\zeta _0^{}}}๐ฟ_3(0).`$
Substitution of this expression into Eq. (114) after some algebra leads to
$$^{}=\frac{1}{d}\left[๐ฌ_2\left(^{}\frac{\zeta _0^{}}{2}\right)\right]^1๐ฌ_2๐^{}\left[๐ฌ_2\left(^{}\frac{\zeta _0^{}}{2}\right)\right]^1๐ฌ_2๐^{}\mathrm{\Psi }_2(0)\frac{2}{\zeta _0^2}\mathrm{\Psi }_1(0),$$
(116)
and use of this into Eq. (110) gives the result
$$\zeta _1^{}=\frac{1}{\zeta _0^{}}+\frac{1}{d}(a,_1^{}),$$
(117)
with
$$_1^{}=\frac{1}{d}\left[๐ฌ_2\left(^{}\frac{\zeta _0^{}}{2}\right)\right]^1๐ฌ_2๐^{}\left[๐ฌ_2\left(^{}\frac{\zeta _0^{}}{2}\right)\right]^1๐ฌ_2๐^{}\mathrm{\Psi }_2(0).$$
(118)
Although the expression for $`\zeta _2^{}`$ can be written in a similar way, it will not be needed here.
## Appendix B Adjoint Linear Operator and Biorthogonal Set
The adjoint for $`^{}`$, $`^{}`$, is defined as usual by
$$(X,^{}Y)=(^{}X,Y),$$
(119)
for arbitrary $`X(๐^{})`$ and $`Y(๐^{})`$ belonging to the Hilbert space. The explicit form of $`^{}`$ is given in Eq. (37). From it, and using the above definition, the expression for $`^{}`$ is easily found,
$`^{}X(๐^{})`$ $`=`$ $`{\displaystyle ๐๐_1^{}\varphi (v_1^{})๐\widehat{๐}\mathrm{\Theta }(\widehat{๐}๐)\widehat{๐}๐\left[X(๐^{\prime \prime })+X(๐_1^{\prime \prime })X(๐^{})X(๐_1^{})\right]}`$ (120)
$`{\displaystyle \frac{\zeta _0^{}}{2}}๐^{}{\displaystyle \frac{}{๐^{}}}X(๐^{}),`$
where $`๐^{\prime \prime }`$ and $`๐_1^{\prime \prime }`$ are the postcollisional velocities corresponding to $`๐^{}`$ and $`๐_1^{}`$,
$`๐^{\prime \prime }`$ $`=`$ $`๐^{}{\displaystyle \frac{1+\alpha }{2}}(\widehat{๐}๐^{})\widehat{๐},`$
$`๐_1^{\prime \prime }`$ $`=`$ $`๐_1^{}+{\displaystyle \frac{1+\alpha }{2}}(\widehat{๐}๐^{})\widehat{๐}.`$ (121)
Equation (120) gives immediately
$$^{}1=0,^{}๐ฏ^{}=\frac{\zeta _0^{}}{2}๐^{},$$
(122)
so that $`1`$ and $`๐^{}`$ are eigenfunctions of the adjoint operator with eigenvalues $`0`$ and $`\zeta _0^{}/2,`$ respectively.
However, the kinetic energy is not an eigenfunction of the adjoint operator. Direct calculation gives
$$^{}v^2=\zeta _0^{}v^2+\frac{a(๐^{})}{2},$$
(123)
where $`a(๐)`$ is given by Eq. (111). Nevertheless, a biorthogonal set can be constructed from $`1,`$ $`๐^{}`$, and $`v^2`$ and it is given in Eq. (58). As noted in the main text, the choice of this set is not unique. The conditions of biorthogonality on $`\mathrm{\Phi }_2`$ are
$$(\mathrm{\Phi }_2,\mathrm{\Psi }_2(0))=(๐^{}\frac{\mathrm{\Phi }_2}{๐^{}},1)=1,$$
(124)
$$(\mathrm{\Phi }_2,\mathrm{\Psi }_1(0))=(\mathrm{\Phi }_2,1)1=0,$$
(125)
$$(\mathrm{\Phi }_2,๐ฟ_3(0))=(\frac{\mathrm{\Phi }_2}{๐^{}},1)=\mathrm{๐}.$$
(126)
A sufficient condition to guarantee that the above relations are verified, is that $`\mathrm{\Phi }_2`$ have the form
$$\mathrm{\Phi }_2(๐^{})=A+Bb(v^{}),$$
(127)
where $`b(v^{})`$ is an arbitrary scalar function of $`๐^{}`$, and the constants $`A`$ and $`B`$ are determined from
$$A=1(1,b)B,B=(๐^{}\frac{b}{๐^{}},1)^1.$$
(128)
Substitution of these expressions into Eq. (126) yields
$$\mathrm{\Phi }_2(๐^{})=1+\left[b(v^{})(1,b)\right](๐^{}\frac{b}{๐^{}},1)^1.$$
(129)
The optimal choice for $`b(v^{})`$ would be that implying that $`\mathrm{\Phi }_2(๐^{})`$ is an eigenfunction of $`^{}`$ corresponding to the eigenvalue $`\zeta _0^{}/2`$. This is accomplished if $`b(v^{})`$ is the solution to
$$^{}b(v^{})=\frac{\zeta _0^{}}{2}\left[b(v^{})+B^1\right].$$
(130)
The solution to this equation, if it exists, has not yet been found.
## Appendix C Evaluation of the Perturbation Theory Results
The expansion of the hydrodynamic eigenvalues of the linearized Boltzmann equation for small $`k`$ is given in Sec. V with the result
$$\lambda _i(๐ค)=\lambda _i(0)+k^2\lambda _i^{(2)}+\mathrm{},$$
(131)
where
$$\left\{\lambda _i(0)\right\}=\{0,\frac{\zeta _0^{}}{2},\frac{\zeta _0^{}}{2}\},$$
(132)
and
$`\lambda _i^{(2)}`$ $`=`$ $`(\mathrm{\Phi }_i,i\widehat{๐}๐^{}\mathrm{\Psi }_i^{(1)})+(\mathrm{\Phi }_i,^{}๐ฌ_i\mathrm{\Psi }_i^{(2)})`$ (133)
$`=`$ $`(\mathrm{\Phi }_i,\widehat{๐}๐^{}\left\{๐ฌ_i\left[^{}\lambda _i\left(0\right)\right]\right\}^1๐ฌ_i\widehat{๐}๐^{}\mathrm{\Psi }_i\left(0\right))+(\mathrm{\Phi }_i,^{}๐ฌ_i\mathrm{\Psi }_i^{(2)}).`$
The eigenvalue $`\zeta _0^{}/2`$ is $`d`$-fold degenerated and the convenient choice for the lowest order eigenfunctions has been discussed in Sec. IV. The formal expression for the second order eigenfunctions $`๐ฌ_i\mathrm{\Psi }_i^{(2)}`$ is given in Eq. (70). For $`i2`$, the second term on the right hand side of Eq. (133) vanishes since $`\mathrm{\Phi }_i`$ is an eigenvector in those cases and, therefore,
$$(\mathrm{\Phi }_i,^{}๐ฌ_i\mathrm{\Psi }_i^{(2)})=(^{}\mathrm{\Phi }_i,๐ฌ_i\mathrm{\Psi }_i^{(2)})(\mathrm{\Phi }_i,๐ฌ_i\mathrm{\Psi }_i^{(2)})=0.$$
(134)
Then, Eq. (133) can be rewritten as
$$\lambda _i^{(2)}=(\widehat{๐}๐^{}\mathrm{\Phi }_i,\left\{๐ฌ_i\left[^{}\lambda _i(0)\right]\right\}^1\widehat{๐}๐^{}\mathrm{\Psi }_i(0))+\delta _{i,2}\frac{1}{d}(^{}v^2,๐ฌ_2\mathrm{\Psi }_2^{(2)}),$$
(135)
where it has been used that
$$๐ซ_i\widehat{๐}๐^{}\mathrm{\Psi }_i(0)=\mathrm{\Psi }_i(0)(\mathrm{\Phi }_i,\widehat{๐}๐^{}\mathrm{\Psi }_i(0))=0.$$
(136)
This follows since $`\mathrm{\Phi }_i`$ and $`\mathrm{\Psi }_i(0)`$ have the same parity with respect to reflections of $`๐^{}`$.
For the first eigenvalue it is
$`\lambda _1^{(2)}`$ $`=`$ $`(\widehat{๐}๐^{},\left(๐ฌ_1^{}\right)^1\widehat{๐}๐^{}\mathrm{\Psi }_i(0))=(\widehat{๐}๐^{},^1\widehat{๐}๐^{}\mathrm{\Psi }_1(0))`$ (137)
$`=`$ $`(^1\widehat{๐}๐^{},\widehat{๐}๐^{}\mathrm{\Psi }_1(0))={\displaystyle \frac{2}{\zeta _0^{}}}((\widehat{๐}๐^{})^2,\mathrm{\Psi }_1(0))={\displaystyle \frac{1}{\zeta _0^{}}}.`$
In the first transformation, the property $`๐ซ_1^{}X=0`$, for arbitrary $`X(๐^{})`$, has been employed. Next, consider the eigenvalue associated to the longitudinal component of $`๐ฟ_3(0)`$ that we will denote by $`\lambda _{}^{(2)}`$,
$`\lambda _{}^{(2)}`$ $`=`$ $`(\left(\widehat{๐}๐^{}\right)^2,\left[๐ฌ_3\left(^{}+{\displaystyle \frac{\zeta _0^{}}{2}}\right)\right]^1\widehat{๐}๐^{}\widehat{๐}{\displaystyle \frac{\mathrm{ln}\varphi }{๐^{}}})`$ (138)
$`=`$ $`(v_i^2,\left[๐ฌ_3\left(^{}+{\displaystyle \frac{\zeta _0^{}}{2}}\right)\right]^1G_{ii}^{})`$
$`+(v_i^2,\left[๐ฌ_3\left(^{}+{\displaystyle \frac{\zeta _0^{}}{2}}\right)\right]^1\left[\mathrm{\Psi }_1(0)+{\displaystyle \frac{d+1}{d}}\mathrm{\Psi }_2(0)\right])`$
$`=`$ $`(v_i^2,\left(^{}+{\displaystyle \frac{\zeta _0^{}}{2}}\right)^1G_{ii}^{})+{\displaystyle \frac{2}{\zeta _0^{}}}(v_i^2,\mathrm{\Psi }_1(0))`$
$`+{\displaystyle \frac{d1}{d}}{\displaystyle \frac{1}{\zeta _0^{}}}(v_i^2,\mathrm{\Psi }_2(0)),`$
where it has been taken into account that $`๐ซ_3\left(^{}+\zeta _0^{}/2\right)`$=0. Using Eq. (108) in the first term on the right hand side of the above equation and evaluating explicitly the other two ones one gets
$$\lambda _{}^0=(v_i^2,C_{ii}^{})+\frac{1}{d\zeta _0^{}}=\frac{2(d1)}{d}\eta ^{}+\frac{1}{d\zeta _0^{}}.$$
(139)
The last equality follows from the fact that the expression of $`\eta ^{}`$ given in Eq. (104) is equivalent to
$$(D_{ij}^{},C_{ij}^{})=\eta ^{}\left(1+\frac{d2}{d}\delta _{i,j}\right),$$
(140)
because of the symmetry of the tensors $`๐ฃ`$ and $`๐ข`$. The calculation of the shear modes eigenvalues $`\lambda _{}^{(2)}`$ is straightforward:
$`\lambda _{}^{(2)}`$ $`=`$ $`(v_{}^{}v_i^{},\left(^{}+{\displaystyle \frac{\zeta _0^{}}{2}}\right)^1v_{}^{}\mathrm{\Psi }_{3,i})`$ (141)
$`=`$ $`(D_{,i}^{},\left(^{}+{\displaystyle \frac{\zeta _0^{}}{2}}\right)^1G_{,i}^{})`$
$`=`$ $`(D_{,i}^{},C_{,i}^{})=\eta ^{}.`$
Finally, the evaluation of $`\lambda _2^{(2)}`$ is somewhat more complicated. The contributions from each of the two terms in Eq. (135) will be computed separately. The first one is given by
$$(\widehat{๐}๐^{}\mathrm{\Phi }_2,\left[๐ฌ_2\left(^{}\frac{\zeta _0^{}}{2}\right)\right]^1\widehat{๐}๐^{}\mathrm{\Psi }_2(0))=(v_i^{}\mathrm{\Phi }_2,\left(\frac{\zeta _0^{}}{2}\right)^1v_i^{}\mathrm{\Psi }_2(0)),$$
(142)
that is equivalent to
$`{\displaystyle \frac{1}{d}}(S_i^{},\left({\displaystyle \frac{\zeta _0^{}}{2}}\right)^1v_i^{}\mathrm{\Psi }_2(0))+{\displaystyle \frac{d1}{d}}(v_i^{},\left({\displaystyle \frac{\zeta _0^{}}{2}}\right)^1v_i^{}\mathrm{\Psi }_2(0))`$
$`={\displaystyle \frac{1}{d}}(S_i^{},๐_i^{}){\displaystyle \frac{2}{d\zeta _0^{}}}(S_i^{},\mathrm{\Psi }_{3,i}(0)){\displaystyle \frac{d+1}{d}}{\displaystyle \frac{1}{\zeta _0^{}}}(v_i^{},v_i^{}\mathrm{\Psi }_2(0))`$
$`=\kappa ^{}{\displaystyle \frac{d+1}{d\zeta _0^{}}}.`$ (143)
The analysis of the second term on the right hand side of Eq. (135) is carried out in an analogous way,
$$\frac{1}{d}(^{}v^2,๐ฌ_2\mathrm{\Psi }_2^{(2)})=\frac{\zeta _0^{}}{d}(v^2,๐ฌ_2\mathrm{\Psi }_2^{(2)})+\frac{1}{d}(a,๐ฌ_2\mathrm{\Psi }_2^{(2)}).$$
(144)
It is
$`{\displaystyle \frac{\zeta _0^{}}{d}}(v^2,๐ฌ_2^2\mathrm{\Psi }_2^{(2)})`$ $`=`$ $`{\displaystyle \frac{\zeta _0^{}}{2}}(\mathrm{\Phi }_1,๐ฌ_2\mathrm{\Psi }_2^{(2)})`$ (145)
$`=`$ $`{\displaystyle \frac{1}{\zeta _0^{}}}(\mathrm{\Phi }_{3,i},v_i^{}\mathrm{\Psi }_2(0))={\displaystyle \frac{1}{\zeta _0^{}}}.`$
In the second equality above, the explicit expression of $`\mathrm{\Psi }_2^{(2)}`$ given in Eq. (70) as well as the properties of $`^{}`$ have been used. The second term on the right hand side of Eq. (144) can be rewritten as
$$\frac{1}{d}(a,๐ฌ_2\psi _2^{(2)})=\frac{1}{d}(a,_1^{}),$$
(146)
where $`_1^{}`$ is defined in Eq. (118). Substitution of Eqs. (145) and (146) into Eq. (144) yields
$$\frac{1}{d}(^{}v^2,๐ฌ_2\mathrm{\Psi }_2^{(2)})=\frac{1}{\zeta _0^{}}\frac{1}{d}(a,_1^{})=\zeta _1^{},$$
(147)
because of Eq. (117). Then, use of Eqs. (143) and (147) into Eq. (135) gives the final expression for $`\lambda _2^{(2)}`$,
$$\lambda _2^{(2)}=\kappa ^{}\zeta _1^{}\frac{d+1}{d\zeta _0^{}}.$$
(148)
Comparison of the results for the second order eigenvalues obtained in this Appendix with those reported in Sec. II shows that both agree, with the transport coefficients given by the same expressions in both cases.
## Appendix D Linearization of the Model Kinetic Equation
Solutions to the model kinetic equation (72) are sought of the form (32) and, consistently, $`g=g_{hcs}+\delta g`$ and $`\nu =\nu _{hcs}+\delta \nu `$, retaining only terms linear in $`\mathrm{\Delta }`$,
$$\left(\frac{}{t}+๐ฏ\right)(f_{hcs}\mathrm{\Delta })=\delta \nu \left(f_{hcs}g_{hcs}\right)\nu _{hcs}\left(f_{hcs}\mathrm{\Delta }\delta g\right).$$
(149)
Use has been made of the HCS model equation
$$\frac{f_{hcs}}{t}=\frac{\zeta _{hcs}}{2}\frac{}{๐}\left(๐f_{hcs}\right)=\nu _{hcs}\left(f_{hcs}g_{hcs}\right).$$
(150)
Next, from the choice of $`\nu `$ as corresponding to hard sphere behavior,
$$\delta \nu =\nu _{hcs}\left(\frac{\delta n}{n_{hcs}}+\frac{2\delta T}{T_{hcs}}\right).$$
(151)
Similarly, the definition of $`g`$ in Eq. (73) gives
$$\delta g=\frac{g_{hcs}}{n_{hcs}}\delta n+\frac{g_{hcs}}{T_{hcs}}\delta T\frac{g_{hcs}}{๐}๐.$$
(152)
The linear model kinetic equation (149) then becomes
$`\left({\displaystyle \frac{}{t}}+๐ฏ\right)(f_{hcs}\mathrm{\Delta })`$ $`=`$ $`\left({\displaystyle \frac{\delta n}{n_{hcs}}}+{\displaystyle \frac{\delta T}{2T_{hcs}}}\right)\nu _{hcs}\left(f_{hcs}g_{hcs}\right)`$ (153)
$`\nu _{hcs}\left(f_{hcs}\mathrm{\Delta }{\displaystyle \frac{g_{hcs}}{n_{hcs}}}\delta n{\displaystyle \frac{g_{hcs}}{T_{hcs}}}\delta T+{\displaystyle \frac{g_{hcs}}{๐}}๐\right).`$
It is convenient to eliminate $`g_{hcs}`$ in this result using the HCS equation (150) to get
$`\left({\displaystyle \frac{}{t}}+๐ฏ\right)(f_{hcs}\mathrm{\Delta })`$ $`=`$ $`\nu _{hcs}f_{hcs}\mathrm{\Delta }+\left({\displaystyle \frac{\delta n}{n_{hcs}}}+{\displaystyle \frac{\delta T}{2T_{hcs}}}\right){\displaystyle \frac{\zeta _{hcs}}{2}}{\displaystyle \frac{}{๐}}\left(๐f_{hcs}\right)`$ (154)
$`+{\displaystyle \frac{\delta n}{n_{hcs}}}\left[\nu _{hcs}f_{hcs}+{\displaystyle \frac{\zeta _{hcs}}{2}}{\displaystyle \frac{}{๐}}\left(๐f_{hcs}\right)\right]`$
$`+\delta T\left[\nu _{hcs}{\displaystyle \frac{f_{hcs}}{T_{hcs}}}+{\displaystyle \frac{\zeta _{hcs}}{2}}{\displaystyle \frac{}{๐}}\left(๐{\displaystyle \frac{f_{hcs}}{T_{hcs}}}\right)\right]`$
$`๐{\displaystyle \frac{}{๐}}\left[\nu _{hcs}f_{hcs}+{\displaystyle \frac{\zeta _{hcs}}{2}}{\displaystyle \frac{}{๐}}\left(๐f_{hcs}\right)\right].`$
The right hand side of this expression can be put in a more transparent form by noting that
$$\left\{(\mathrm{\Delta },\mathrm{\Phi }_i)\right\}=\{\frac{\delta n}{n_{hcs}},\frac{\delta T}{2T_{hcs}}+\frac{\delta n}{n_{hcs}},\frac{u_{}}{v_{hcs}},\frac{๐_{}}{v_{hcs}}\}.$$
(155)
where the scalar product $`(X,Y)`$ is defined in Eq. (40) and $`\left\{\mathrm{\Phi }_i\right\}`$ are the biorthogonal set given in Eq. (58). Moreover, the relation
$$T_{hcs}\frac{\mathrm{ln}f_{hcs}}{T_{hcs}}=\frac{1}{2\varphi }\frac{}{๐}\left(๐\varphi \right),$$
(156)
allows to write the dependence on $`f_{hcs}`$ on the right hand side in terms of the lowest order hydrodynamic eigenfunctions $`\{\mathrm{\Psi }_i(0)\}`$ of the linearized Boltzmann operator given in Eq. (54). Then, Eq. (154) can be rewritten as
$`\left({\displaystyle \frac{}{t}}+๐\right)(f_{hcs}\mathrm{\Delta })`$ $`=`$ $`\nu _{hcs}f_{hcs}\left[\mathrm{\Delta }{\displaystyle \underset{i}{}}\mathrm{\Psi }_i(0)(\mathrm{\Phi }_i,\mathrm{\Delta })\right]`$ (157)
$`+{\displaystyle \frac{\zeta _{hcs}}{2}}{\displaystyle \frac{}{๐}}\left[f_{hcs}๐{\displaystyle \underset{i}{}}\mathrm{\Psi }_i(0)(\mathrm{\Phi }_i,\mathrm{\Delta })\right]`$
$`f_{hcs}{\displaystyle \frac{v_{hcs}(t)}{\mathrm{}}}{\displaystyle \underset{i}{}}\lambda _i(0)\mathrm{\Psi }_i(0)(\mathrm{\Phi }_i,\mathrm{\Delta }).`$
Here $`\left\{\lambda _i(0)\right\}`$ are the eigenvalues of Eq. (53).
Finally, introducing the dimensionless variables of Secs. II and III, and the single Fourier component of Eq. (35) gives the linear kinetic equation for the model
$$\left(_s+i๐ค๐ฏ^{}+_m^{}\right)\stackrel{~}{\mathrm{\Delta }}(๐,๐^{},s)=0,$$
(158)
with
$$_m^{}=\underset{i}{}\lambda _i(0)๐ซ_i+\nu _0^{}๐ฌ+\frac{\zeta _0^{}}{2}\varphi ^1\frac{}{๐^{}}๐^{}\varphi ๐ฌ.$$
(159)
The projection operators $`๐ซ_i`$ are defined in Eq. (62), and
$$๐ซ\underset{i}{}๐ซ_i,๐ฌ1๐ซ.$$
(160)
It is directly verified that
$$๐ซ\varphi ^1\frac{}{๐^{}}๐^{}\varphi ๐ฌ=0,$$
(161)
so that $`_m^{}`$ can be written
$$_m^{}=\underset{i}{}\lambda _i(0)๐ซ_i+๐ฌ\left[\nu _0^{}+\frac{\zeta _0^{}}{2}\varphi ^1\frac{}{๐^{}}๐^{}\varphi \right]๐ฌ.$$
(162)
This is the expression for the generator of the linear dynamics used in main the text.
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# 1 Introduction
## 1 Introduction
The hierarchy problem manifests itself in the enormous difference between the standard model, gravity and in a wider sense the dark energy scales. It is assumed often that in the case of the standard model its solution requires some UV regulating physics. However, as it was suggested in -, the hierarchy problem can be addressed as the problem of a vacuum super-selection rule. The recent progress in the understanding of the string theory vacua landscape (see for example ) gives a hint on the different possibilities of vacua density distributions. This motivates the studying of an alternative mechanism as a possible solution to the hierarchy problem. The idea proposed in , based on an earlier work on cosmic attractors , is to put to work the multiplicity of vacua. The hierarchy problem is promoted into a problem of the super-selection rule among the infinite number of vacua, that are finely scanned by the Higgs mass. In this framework, the Higgs mass is promoted into a dynamical variable. An infinite number of vacua cluster around a certain point making it an attractor. On the resulting landscape in all but a measure zero set of vacua the Higgs mass has a common, hierarchically small value due to the attractor.
In this paper we will analyze a model which can be viewed as an effective theory obtained after integrating out the Higgs field in the model . As a result we get a setup in which branes adjust their charges according to the values of the field they produce. We consider the possibility of having an attractor for such system. Three-forms are sourced out by a number of membranes (two-branes) with charges $`q`$ that can be self-adjusted. In particular, we will consider the membrane charges being tuned as an arbitrary function of the field. This problem essentially can be reduced to that of field dependent charges in 1+1 electrodynamics. As we will show, the presence of an attractor is a very generic feature of such models.
## 2 Three-forms and two-branes
Letโs review the setup of the work . The spectrum of different string theories contains antisymmetric form fields, which after compactification to four dimensions give rise to three-forms, two-forms and one-forms. In particular, we are interested in three-forms $`C_{\alpha \beta \gamma }`$. The action for a three-form in four dimensions reads
$$S_C=_{3+1}\frac{1}{48}F_{\mu \alpha \beta \gamma }F^{\mu \alpha \beta \gamma },$$
(1)
where the four-form field strength
$$F_{\mu \alpha \beta \gamma }=d_{[\mu }C_{\alpha \beta \gamma ]}.$$
(2)
This action is gauge invariant, and this guarantees the decoupling of the time component. The gauge transformation is
$$C_{\alpha \beta \gamma }C_{\alpha \beta \gamma }+d_{[\alpha }\mathrm{\Omega }_{\beta \gamma ]},$$
(3)
where $`\mathrm{\Omega }_{\beta \gamma }`$ is some two-form depending on the coordinates and the square brackets denote anti-symmetrization. The form $`C`$ has no propagating degrees of freedom in four dimensions. The equations of motion stemming from the action (1) are
$$^\mu F_{\mu \nu \alpha \beta }=\mathrm{\hspace{0.17em}0}$$
(4)
and have a constant solution
$$F_{\mu \nu \alpha \beta }=F_0ฯต_{\mu \nu \alpha \beta };$$
(5)
here $`F_0=\mathrm{constant}`$ and $`ฯต_{\mu \nu \alpha \beta }`$ is the totally antisymmetric tensor. In the absence of interactions with other fields this constant changes the Lagrangian and contributes to the cosmological term. In the presence of interactions $`F_0`$ will contribute to those fields masses and to the couplings.
Three-forms couple to two-branes, e.g. membranes. The effective action is given by
$$S=\frac{q}{6}_{2+1}d^3\xi C_{\mu \nu \alpha }\left(\frac{Y^\mu }{\xi ^a}\frac{Y^\nu }{\xi ^b}\frac{Y^\alpha }{\xi ^c}\right)ฯต^{abc}_{3+1}\frac{1}{48}F^2.$$
(6)
where $`q`$ is the brane charge and $`Y^\mu (\xi )`$ describe the brane history as a function of its world volume coordinates $`\xi ^a,a=0,1,2`$ We can rewrite the interaction term as a four dimensional integral
$$d^4x\frac{1}{6}J^{\alpha \beta \gamma }C_{\alpha \beta \gamma }$$
(7)
where the brane current
$$J^{\alpha \beta \gamma }(x)=d^3\xi \delta ^4(xY(\xi ))q\left(\frac{Y^\alpha }{\xi ^a}\frac{Y^\beta }{\xi ^b}\frac{Y^\gamma }{\xi ^c}\right)ฯต^{abc}.$$
(8)
As long as $`q`$ is constant, the current is conserved. We end up with the equations of motion
$$_\mu F^{\mu \nu \alpha \beta }=qd^3\xi \delta ^4(xY(\xi ))\left(\frac{Y^\nu }{\xi ^a}\frac{Y^\alpha }{\xi ^b}\frac{Y^\beta }{\xi ^c}\right)ฯต^{abc}.$$
(9)
We consider the simple case of static and flat branes,
$$Y^\mu =\xi ^\mu ,\mu =0,1,2$$
(10)
$$Y^3=\mathrm{\hspace{0.17em}0}$$
(11)
We take $`x_3=z`$ as the coordinate transversal to the brane. Then the equations of motion reduce to
$$_\mu F^{\mu \nu \alpha \beta }=q\delta (z)ฯต^{\nu \alpha \beta z}.$$
(12)
The equations of motion show that the brane separates two vacua. In each of them the field strength is constant and the jump between the values of the field in different vacua is given by the brane charge $`q`$. This way, there is a solution with multiplicity of vacua and the vacua in this solution are labeled by an integer $`n`$,
$$\frac{1}{24}F_{\alpha \beta \gamma \mu }ฯต^{\alpha \beta \gamma \mu }=qn+F_0,$$
(13)
where $`F_0`$ is a constant which in the theory with an attractor mechanism will be fixed.
In the model , the lowest order parity and gauge-invariant Lagrangian describing a non-trivial interaction between the Higgs field $`\varphi `$ and the gauge field $`C`$ was suggested in the form
$$L=|_\mu \varphi |^2\frac{1}{48}F^2+|\varphi |^2\left(m^2+\frac{F^2}{48M^2}\right)\frac{\lambda }{2}|\varphi |^4+\mathrm{}$$
(14)
where $`\lambda `$ is the quartic coupling and $`m`$, $`M`$ are mass parameters. As a result, the gauge field determines the value of the effective mass and consequently the vacuum expectation value of the Higgs field. The Higgs field in turn readjusts the brane charges and gauge field closing a cycle. This can create an attractor depending on details of the interaction provided there is an additional symmetry forbidding higher loop corrections to the classical attractor.
## 3 Explicit gauge field dependence
The key idea is to consider a charge $`q(F)`$ being explicitly field dependent; this corresponds to effective โintegrating outโ the Higgs mass in the model . To simplify the derivations and make the physical content clearer, we will consider the 1+1 case and will call the field potential $`A`$. Thus we will have electrodynamics with self-adjusting charges. In the absence of the mentioned dependence the current is
$$J^\mu (x)=๐\xi q\delta ^2(xY(\xi ))\frac{Y^\mu (\xi )}{\xi },$$
(15)
where the charge $`q`$ acts as source for the gauge field. If we consider the charges being field depending sources, the field in turn readjusts the charge. The corresponding current can be written as
$$J^\mu (x)=๐\xi q\left(F(Y(\xi ))\right)\delta ^2(xY(\xi ))\frac{Y^\mu (\xi )}{\xi }.$$
(16)
This current is no longer conserved unless we rewrite the interaction so that the field couples only to the transverse part of the current. This can be guaranteed by an interaction term with a projection kernel
$$\mathrm{\Pi }_{\mu \nu }=g_{\mu \nu }\frac{_\mu _\nu }{^2}$$
(17)
so that the Lagrangian
$$=\frac{1}{4}F_{\mu \nu }F^{\mu \nu }+A^\mu \mathrm{\Pi }_{\mu \nu }J^\nu .$$
(18)
The potential $`A_\mu `$ couples only to the transverse part of the current $`J^\mu `$. For a single static charge current located at point $`x_1z=a`$
$$J^0(z)=q(F(z))\delta (za),J^1(z)=0$$
(19)
The interaction term
$$J^0\mathrm{\Pi }_{0\nu }A^\nu =J^0\mathrm{\Pi }_{00}A^0=J^0A_0J^0\frac{_0_0}{^2}A_0=J^0A_0$$
(20)
and the Lagrangian becomes
$$=\frac{1}{2}(_zA_0)^2+q(_zA_0)\delta (za)A_0.$$
(21)
The variation of the Lagrangian is
$$\delta =(_z^2A_0+q(_zA_0)\delta (za)_z(q^{}(dzA_0)\delta (za)A_0))\delta A_0$$
(22)
and gives the equation of motion
$`_z^2A_0=\delta (za)\left[_zA_0q^{}(_zA_0)q(_zA_0)+q^{\prime \prime }(_zA_0)_z^2A_0A_0\right]+`$
$`+_z\delta (za)q^{}(_zA_0)A_0.`$ (23)
Outside the brane the equation of motion reduces to
$$_z^2A_0=0.$$
(24)
The general solution for the field strength is a constant,
$$FF_{10}=_zA_0=\mathrm{const}.$$
(25)
Integrating in a small neighborhood near the brane leads to the boundary condition
$$F(a+0)F(a0)=$$
(26)
$$=\left[Fq^{}(F)q(F)+q^{\prime \prime }(F)_zFA_0_z(q^{}(_zA_0)A_0)\right]_{z=a}$$
(27)
or equivalently
$$F(a+0)F(a0)=q(F)|_{z=a}.$$
(28)
We can implement this boundary condition into the equation of motion as
$$_zF=q(F(a))\delta (za)$$
(29)
Letโs look at the behavior of the field when we add in succession $`N`$ charged branes at points $`z=a_n`$ (for convenience we take $`a_n>a_k`$ for $`n>k`$). In this case we get the following equation
$$_zF=\underset{k=1}{\overset{N}{}}q(F(a_k))\delta (za_k)$$
(30)
where the integer $`k=1,\mathrm{},N`$ labels the branes. The general solution to this equation is
$$F(z)=\underset{k=1}{\overset{N}{}}q(F(a_k))\theta (za_k)$$
(31)
with the boundary condition to be applied. These boundary conditions will fix the constants $`q(F(a_k))`$. The equations of motion lead to the following recursion relation
$$F_k+\frac{1}{2}q(F_k)=F_{k1}\frac{1}{2}q(F_{k1})$$
(32)
for the values of the fields between the branes. As $`N\mathrm{}`$, if there is an attractor point, the following limit should exist
$$\underset{N\mathrm{}}{lim}F_NF_A.$$
(33)
From the recursion one can see that when the limit exists, then $`q(F_A)=0`$. It means that the attractor point candidate $`F_A`$ should be a root of the equation
$$q(F)=0.$$
(34)
In Figures 1 and 2 we show how the addition of branes at points $`a_1`$, $`a_2`$,โฆ leads to an increasing number of vacua near the attractor point $`F=F_A`$. A sufficient condition for the existence of the above limit, i.e. the attractor point at the root $`F_A`$, is $`q^{}(F_A)>0`$.
## 4 Number Density near the Attractor Point
We will evaluate the vacuum number density near the attractor point. In the case of a charge depending linearly on the field $`F`$,
$$q(F)=c(FF_A)$$
(35)
with $`c`$ some positive constant, the recursion (32) with an initial value $`F=F_1`$ has the solution
$$F_k=\left(\frac{2c}{2+c}\right)^{k1}\left(F_1F_A\right)+F_A.$$
(36)
We can express the number of vacuum states outside the interval of fields $`FF_A`$
$$k=1+\frac{\mathrm{ln}\frac{FF_A}{F_1F_A}}{\mathrm{ln}\frac{2c}{2+c}}.$$
(37)
Correspondingly, the number density of the vacuum states in the linear case is given by
$$n\frac{dk}{dF}=\frac{1}{\mathrm{ln}\frac{2c}{2+c}}\frac{1}{|FF_A|}$$
(38)
and is divergent at the attractor point $`F=F_A`$.
We would like to calculate the number density for arbitrary self-interaction but we cannot explicitly write the number of states $`k`$ in terms of the field range for a generic function $`q(F)`$. Nevertheless we can estimate the derivative via the recursion relation
$$n\frac{\mathrm{\Delta }k}{\mathrm{\Delta }F}=\frac{1}{F_kF_{k1}}\frac{1}{q(F)}$$
(39)
This number of states is also divergent at the attractor point. In Figure (3) we have depicted the typical, divergent behavior of the number density of vacua near the attractor point.
Letโs find the domain of convergence to the attractor. The answer comes from the sufficient condition of the attractor existence at the point $`F=F_A`$
$$q(F_A)=0,q^{}(F_A)>0.$$
(40)
It follows from the recursion relation that in the vicinity of the attractor point $`F_A`$ to guarantee the convergence we should continuously satisfy
$$q^{}(F)>0,$$
(41)
$$q^{}(F)<4+q^{}(F_A).$$
(42)
In the linear example
$$q(F)=c(FF_A)$$
(43)
the fields converge everywhere for $`c>0`$.
For a quadratic dependence
$$q(F)=(FF_A)(FF_B)$$
(44)
the attractor will be located at $`F_A`$ for
$$F_A>F_B.$$
(45)
The range of the fields $`F`$ which converge to the attractor point is defined by
$$0<q^{}(F)<4+q^{}(F_A).$$
(46)
This restricts $`F`$ to the range
$$(F_A+F_B)/2<F<2+F_A.$$
(47)
## 5 Conclusions
We have shown that the self-adjusting charges or membranes have an attractor point. This implies that the number density of vacua within a small range around the attractor point blows up. From the physical point of view, the attractor adjusts to the point with a minimal self-interaction, creating an enormous number of vacua with close values. The sufficient condition for the attractor existence is vanishing of the charge at the attractor point and the charge being an increasing function of the field. Even if the interaction never reaches zero, this point still will be like an attractor. However, the number density of vacua will have a finite maximum sharp pick.
We hope that the suggested attractor model can be implemented for a natural explanation of large hierarchies, like scales of electroweak theory, gravity and dark energy.
Acknowledgments: it is pleasure to thank Gia Dvali for raising this problem and for useful discussions.
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# Optimal phase space projection for noise reduction
## I Introduction
Due to its simplicity in implementation and efficiency in computation, noise reduction based on phase space projection has been widely studied in previous literature. For example, Broomhead and King broomhead extracting (2) advocated that, in case of white noise, via singular value decomposition (SVD), one could extract qualitative dynamics from experimental (noisy) time series by removing the empirical orthogonal functions (EOFs) vautard singular (13) of the trajectory matrix which correspond to the noise components. To deal with the case of colored noise, Allen and Smith allen optimal (1) proposed a more general method, which would statistically pre-whiten colored noise by introducing a transformation to the covariance matrix of noise. In general, phase space projection based on these methods would not operate on the EOFs that span the signal-plus-noise subspace, therefore those operations could achieve a lowest possible distortion for the clean signal, but at the price of a highest possible residual noise level ephraim signal (4). To obtain an optimal tradeoff between signal distortion and residual noise so as to minimize the overall distortion, Ephraim and Trees proposed the time domain constraint (TDC) projector ephraim signal (4), which improves the performance of the existing methods by imposing a constraint on the residual noise, and which also includes the existing methods as its subcases. As a generalization, some authors also extended the TDC projector to the cases with colored noise doclo multimicrophone (3, 7).
Usually, these authors will make two assumptions concerning the experimental time series. The first assumption is that the time series is stationary and ergodic, and the second one is that the noise components are independent of the clean signal. In this communication we will re-examine the idea of the TDC projector and deduce a more universal version. We will also show that, with the first assumption, the second is not necessary in general.
The remainder of this article will go as follows: In the second section we will introduce the idea of the TDC projector. Based on the assumption that the noisy time series is stationary and ergodic, we will obtain the optimal TDC projector for a trajectory matrix in the sense of minimizing signal distortion subject to a permissible noise level. In the third section we will apply the optimal TDC projector to simulated data from the Rรถssler system and experimental speech data. We will also compare the performance of the projectors under different TDCs. Finally, a conclusion is available to summarize the whole article.
## II Mathematical deduction
Given a noisy time series $`s=\{s_i\}_{i=1}^M`$, we suppose that the corresponding clean signal and the additive noise components are $`d=\{d_i\}_{i=1}^M`$ and $`n=\{n_i\}_{i=1}^M`$ respectively, thus for each noisy data point $`s_i`$, we have $`s_i=d_i+n_i`$. In addition, we assume $`\{s_i\}_{i=1}^M`$ are (weakly) stationary and ergodic so that its expectation exists and its variance is finite, while its (auto)covariances only depend on the time difference between the subsets.
Following the definition in broomhead extracting (2), we could construct a $`(Mm+1)\times m`$ trajectory matrix $`๐`$ from $`\{s_i\}_{i=1}^M`$ by letting
$$๐=\left(\begin{array}{cccc}s_1& s_2& \mathrm{}& s_m\\ s_2& s_3& \mathrm{}& s_{m+1}\\ \mathrm{}& & \mathrm{}& \\ s_{Mm+1}& s_{Mm+2}& \mathrm{}& s_M\end{array}\right)_{(Mm+1)\times m}$$
with $`Mm+1>m`$. Similarly, we could also obtain the corresponding trajectory matrices $`๐`$ and $`๐`$ for components $`\{d_i\}_{i=1}^M`$ and $`\{n_i\}_{i=1}^M`$ respectively, and we have $`๐=๐+๐`$.
For the purpose of noise reduction, we introduce a projection operator $`๐`$ on the trajectory matrix $`๐`$ of noisy signal, through which we could obtain a matrix $`๐=\mathrm{๐๐}`$. We define $`๐_0=`$ $`๐๐=๐(๐๐_m)+\mathrm{๐๐}`$ as the matrix of residual signal, where the term $`๐(๐๐_m)`$ means signal distortion and the term $`\mathrm{๐๐}`$ is residual noise. With the intention of data augmentation, we would require to achieve as small signal distortion as possible. Thus $`๐=๐_m`$ would be an intuitive choice. However, in situations such as speech communication, one would also require a permissible residual noise level of the noisy signal, and the objective becomes to minimize signal distortion subject to achieving a permissible residual noise level. Thus if the initial data does not fulfil this requirement, one has to reduce the initial noise level at the price of introducing possible signal distortion. Similar to the idea proposed in ephraim signal (4), here we impose a time domain constraint (TDC) $`\mu `$ on the term of residual noise $`\mathrm{๐๐}`$ and treat $`๐=๐(๐๐_m)+\mu \mathrm{๐๐}`$ as the part that requires a minimal distortion, where $`\mu ^2[0,+\mathrm{})`$ is the Lagrange multiplier determined by the permissible noise level from the practical demand (see Eq. (33) of ephraim signal (4) and the related discussions therein). Thus our objective will be to minimize the average energy $`\mathrm{\Xi }=\left(\underset{i=1}{\overset{M}{}}r_i^2\right)/M`$ of the data set $`r=\{r_i\}_{i=1}^M`$ that (approximately) corresponds to the matrix $`๐`$. If $`M`$ $`m`$, then
$$\mathrm{\Xi }\frac{1}{(Mm+1)m}tr(๐^T๐)\text{,}$$
(1)
where $`tr()`$ means the trace of a square matrix, $`๐^T`$ denotes the transpose of the matrix $`๐`$.
Discarding the constant term $`tr(๐^T๐)`$ in $`tr(๐^T๐)`$, we have
$`tr(๐^T๐)`$ $`=`$ $`tr(๐^T(๐+\mu ๐)^T(๐+\mu ๐)๐)`$ (2)
$`2tr(๐^T(๐+\mu ๐)^T๐).`$
Taking $`m`$ as a constant note0 (14), for the minimization problem, by requiring $`tr(๐^T๐)/๐=\mathrm{๐}`$, we would have $`(๐+\mu ๐)^T(๐+\mu ๐)๐(๐+\mu ๐)^T๐=\mathrm{๐}`$ according to the differential rules in, for example, (Lutkepohl introduction, 10, p. 472). Therefore the optimal projector
$$๐_{\mathrm{min}}=\left\{(๐+\mu ๐)^T(๐+\mu ๐)\right\}^1(๐+\mu ๐)^T๐\text{.}$$
(3)
With the noise components, $`tr(๐^T๐)/๐^2=2(๐+\mu ๐)^T(๐+\mu ๐)`$ is positive definite, which confirms that the extremum taken at $`๐_{\mathrm{min}}`$ is a minimum. The corresponding minimal value
$$tr_{\mathrm{min}}(๐^T๐)=tr(๐^T๐)tr(๐^T(๐+\mu ๐)๐_{\mathrm{min}})\text{.}$$
(4)
But note that $`(๐+\mu ๐)^T`$ is not a square matrix, its (ordinary) inverse matrix usually is not defined, thus we could not cancel the terms of $`(๐+\mu ๐)^T`$ in Eq. (3).
Since $`๐=๐+๐`$, we could also write Eq. (3) in the form of
$`๐_{\mathrm{min}}`$ $`=`$ $`\left\{(๐+(\mu 1)๐)^T(๐+(\mu 1)๐)\right\}^1`$ (5)
$`\times (๐+(\mu 1)๐)^T(๐๐)\text{.}`$
If we assume the clean signal and the noise components are independent, statistically we have $`๐^T๐=๐^T๐=\mathrm{๐}`$ as $`M\mathrm{}`$, hence $`๐^T๐=๐^T๐+๐^T๐`$, and Eq. (5) reduces to
$$๐_{\mathrm{min}}=\left\{๐^T๐+(\mu ^21)๐^T๐\right\}^1(๐^T๐๐^๐๐)\text{.}$$
(6)
Let $`๐_S`$, $`๐_N`$ denote the covariance matrices of $`\{s_i\}_{i=1}^M`$ and $`\{n_i\}_{i=1}^M`$ respectively, by assuming the expectation values $`E(s)=E(n)=0`$, we have $`๐_S=๐^T๐/(Mm+1)`$ and $`๐_N=๐^T๐/(Mm+1)`$ as $`M\mathrm{}`$. Thus Eq. (6) would be expressed as
$$๐_{\mathrm{min}}=\left\{๐_S+(\mu ^21)๐_N\right\}^1(๐_S๐_๐)\text{,}$$
(7)
which is consistent with the result in, for example, Eq. (3) of hu generalized (7). But note that here we use $`\mu ^2`$ to substitute for the multiplier $`\mu `$ in Eq. (3) of hu generalized (7). Also note that $`๐_{\mathrm{min}}`$ in our work is the transpose of that in Eq. (3) of hu generalized (7), this is because the trajectory matrices in our work are essentially the transpose of those in doclo multimicrophone (3, 4, 7).
In many situations, although the noise components are theoretically uncorrelated to the clean signal, numerical calculations often indicate that the assumption $`๐^T๐=๐^T๐=\mathrm{๐}`$ does not hold strictly for finite data sets. As a more rigorous form, Eq. (5) needs no independence assumption between the noise components and the clean signal. Thus this expression is a further generalization of previous studies.
## III Numerical results
We note that the trajectory matrices previously introduced are all Hankel matrices. Take trajectory matrix $`๐`$ of the noisy signal as an example, its entries satisfy $`๐(i,j)=๐(k,l)`$ if $`i+j=k+l`$, where $`๐(i,j)`$ denote the element of matrix $`๐`$ on $`i`$-th row and $`j`$-th column. However, matrix $`๐=\mathrm{๐๐}`$ usually will not be a Hankel matrix, and we may have many ways to obtain the filtered (or projected) signal $`\{z_i\}_{i=1}^M`$. In our work we use the method of secondary diagonal averaging to extract signal from the matrix $`๐`$, which takes the average of the elements along the secondary diagonals of matrix $`๐`$ as the filtered signal $`\{z_i\}_{i=1}^M`$ (for details, see (golyandina analysis, 5, p. 24)), and thus can form a new trajectory (Hankel) matrix $`๐^H`$ from $`\{z_i\}_{i=1}^M`$. Golyandina et al. prove that this method is optimal among all Hankelization procedures in the sense that the matrix difference $`๐^H๐`$ has minimal Frobenius norm (golyandina analysis, 5, p. 24, p. 266).
We adopt the signal-to-noise ratio (SNR) as the metric to evaluate the performance of our noise reduction scheme, which is defined (in dB) as ephraim signal (4, 8)
$$SNR=10\mathrm{log}_{10}\frac{d^2}{zd^2}\text{,}$$
(8)
where $`d^2=\underset{i=1}{\overset{M}{}}d_i^2`$ and $`zd^2=`$ $`\underset{i=1}{\overset{M}{}}(z_id_i)^2`$.
We first apply our algorithm to a simulated data set, which is generated from the $`x`$ component of the Rรถssler system
$$\{\begin{array}{c}\dot{x}=(y+z)\hfill \\ \dot{y}=x+ay\hfill \\ \dot{z}=b+(xc)z\hfill \end{array}$$
(9)
with parameter $`a=0.15`$, $`b=0.2`$ and $`c=10`$. The data is evenly sampled for every $`0.1`$ time units. We generate $`10,000`$ data points and discard the first 1000 to avoid transition. To construct the trajectory matrices, we will set the window size $`m=20`$ .
Let $`\left\{s_i\right\}_{i=1}^M`$ and $`\left\{d_i\right\}_{i=1}^M`$ again denote the noisy and clean signals respectively. We consider adding three types of noise contamination to the clean data. The first one is additive white noise $`\left\{\xi _i\right\}_{i=1}^M`$ (so that $`s_i=d_i+\xi _i`$), which follows the normal Gaussian distribution $`N(0,1)`$. The second one is additive colored noise $`\left\{\eta _i\right\}_{i=1}^M`$ (so that $`s_i=d_i+\eta _i`$), which, as an example, is produced from a third order autoregressive ($`AR(3)`$) process in the form of $`\eta _i=0.8\eta _{i1}0.5\eta _{i2}+0.6\eta _{i3}+\xi _i`$, where variable $`\xi `$ follows the normal distribution $`N(0,1)`$. The last one is multiplicative noise $`\left\{\zeta _id_i\right\}_{i=1}^M`$ (so that $`s_i=(1+\zeta _i)d_i`$). As an example, we let $`\zeta _i=\eta _i^2`$, where $`\left\{\eta _i\right\}_{i=1}^M`$ is from the previous $`AR(3)`$ process, then the noise component $`\left\{\zeta _id_i\right\}_{i=1}^M`$ is correlated to the clean data $`\left\{d_i\right\}_{i=1}^M`$.
By varying the magnitude of the introduced noise, we have the initial noise level be $`20`$ dB, $`10`$ dB, $`0`$ dB respectively, and for each noise level, we will include $`10`$ different noise samples from the same process in calculation. We will also study the performance of the projectors under different constraints. As examples, we let TDC $`\mu =0`$, $`0.5`$ and $`1`$ separately. TDC $`\mu =0`$ will lead to the least-squares (LS) projector based on the SVD technique that appeared in, for example, allen optimal (1, 2, 13). We would need to specify the dimension of the signal-plus-noise subspace so as to group the EOFs and eigenvalues that correspond to the noisy signal and remove the complementary noise subspace, which is essentially related to the problem of choosing the embedding dimension for embedding reconstruction from a scalar time series (see the discussion in johnson generalized (8)). Thus here we adopt the criterion of false nearest neighbor Kennel (9), a method proposed for selection of appropriate embedding dimensions. To apply this criterion in calculation, we utilized the codes implemented in the TISEAN package Tisean (6) and found that the proper dimension size $`K`$ of the signal-plus-noise subspace is $`5`$ in our cases. For $`\mu =1`$, we will obtain the well-know linear minimum mean-squared-error (LMMSE) projector (detailed introductions available in, e.g., trees detection (12)). After all of the calculations, we finally list the performance of these TDC projectors in Table 1. For better comprehension of the presented results, we provide the waveforms of all of the data listed in Table 1 as the supplementary material supplementary (15). To keep our presentation concise, here we only take out the raw data contaminated with $`0`$ dB additive white noise as an example and depict its waveform of in panel $`(a)`$ of Fig. (1). For comparison, we also plot the augmented data with TDC$`=0,0.5`$ and $`1`$ in panel $`(b)`$, $`(c)`$ and $`(d)`$, whose mean noise levels are $`5.88,`$ $`9.71`$ and $`10.10`$ dB correspondingly.
From Table 1, we see that for the Rรถssler system, our algorithm works for all of the three types of contamination. But the data augmentation for additive colored noise is not as obvious as those for additive white noise and multiplicative noise (the possible explanation is explored in the appendix). We also see that, in general, the LMMSE projector has better performance than that of the LS projector in the sense that it can achieve better SNR as defined in Eq. (8).
We then apply our algorithm to a very noisy speech (vowel) data (with $`8,000`$ data points), which is sampled at $`44`$ kHz and quantized to $`16`$ bits. In this case we only know the background noise measured in the period without the signal. It would be preferred if we could produce a set of samples that mimic the behavior of the underlying noise. Here we adopt the pseudo-periodic surrogate (PPS) algorithm small surrogate (11) to generate $`9`$ surrogates based on the original background noise. With these data sets, the initial SNR of the speech data is estimated to be $`0.32\pm 0.18`$ dB via Eq.(8). To introduce phase space projection to the speech data, we let the window size $`m=30`$ and set the dimension size of signal-plus-noise subspace to be $`K=8`$, and then apply the TDC projectors $`๐`$ to its trajectory matrix. For the LS projector ($`\mu =0`$), the augmented SNR$`=4.36\pm 0.41`$ dB. While for TDC $`\mu =0.5`$ and $`1`$, the corresponding SNRs increase to $`6.28\pm 0.61`$ dB and $`6.97\pm 0.66`$ dB respectively. As an illustration, we plot the waveforms of the original speech record and three projected data under different TDCs in Fig. (2), from which we can see that, the LMMSE projector ($`\mu =1`$) would lead to a smoother speech waveform (panel (d)) than that of the LS projector (panel (b)). Although the speech data output from the LMMSE projector has lower (signal) magnitudes than those of the speech record from the LS projector, it is still preferred to its rival in speech communication since a smoother data will usually bring better communication quality.
## IV Conclusion
In this communication we re-examined the noise reduction technique based on phase space projection. By imposing a constraint on the residual noise, we deduced the optimal time domain constrained projector in the sense of minimizing signal distortion subject to a permissible noise level. We also showed that, in general we need not assume independence between clean signal and noise components as was previously done. This viewpoint was confirmed by our numerical results (see the third column of the calculation results in Table 1).
## Appendix
Here let us examine the metric of signal-to-noise ratio (SNR) in more detail. According to the definition in Eq. (8), $`SNR=10\mathrm{log}_{10}d^2/zd^2,`$where $`d^2=\underset{i=1}{\overset{M}{}}d_i^2`$ and $`zd^2=`$ $`\underset{i=1}{\overset{M}{}}(z_id_i)^2`$. Note that $`d^2=tr(๐^T๐)/m`$ and $`zd^2=tr((๐๐)^T(๐๐))/m`$ as $`M\mathrm{}`$, thus
$$SNR=10\mathrm{log}_{10}tr(๐^T๐)10\mathrm{log}_{10}tr((๐๐)^T(๐๐)).$$
Since $`๐=\mathrm{๐๐}`$, we have $`tr((๐๐)^T(๐๐))=tr(๐^T๐^T\mathrm{๐๐})2tr(๐^T\mathrm{๐๐})+tr(๐^T๐)`$. For the case that the noise and the clean signal are independent. substituting the optimal projector $`๐_{\mathrm{min}}`$ into the expression, it can be shown that $`tr_{\mathrm{min}}((๐๐)^T(๐๐))=tr(๐^T๐)tr(๐_{\mathrm{min}}๐^T๐)`$. For simplicity, we assume the expectation values $`E(d)=E(n)=0`$, then $`๐_D=๐^T๐/(Mm+1)`$ and $`๐_N=๐^T๐/(Mm+1)`$ as $`M\mathrm{}`$, where $`๐_D`$ and $`๐_N`$ are the covariance matrix of the clean signal and the noise respectively, and $`๐_{\mathrm{min}}`$ can be expressed in the form of Eq. (7), or equivalently, $`๐_{\mathrm{min}}=\{๐_D+\mu ^2๐_N\}^1๐_D`$. Therefore in this case, we have $`tr_{\mathrm{min}}((๐๐)^T(๐๐))=tr(๐_D)tr(๐_{\mathrm{min}}๐_D)`$, thus the maximal SNR can be expressed by
$`SNR_{\mathrm{max}}`$ $`=`$ $`10\mathrm{log}_{10}tr(๐_D)10\mathrm{log}_{10}(tr(๐_D)`$ (10)
$`tr(\{๐_D+\mu ^2๐_N\}^1๐_D^2)).`$
Through the SVD technique broomhead extracting (2), $`๐_D`$ can be written as $`๐_D=๐_D๐ฒ_D๐_D^T`$ , where $`๐_D`$ is the normalized eigenvector matrix of $`๐_D`$, and $`๐ฒ_D`$ is a diagonal matrix whose non-zero elements are the eigenvalues of $`๐_D`$ (in fact $`๐_D^T๐_D=๐_m`$ and $`๐_D๐_D=๐_D๐ฒ_D`$). Similarly, we have $`๐_N=๐_N๐ฒ_N๐_N^T`$. Let $`๐_N=๐_D๐_{DN}`$ (for better comprehension, $`๐_{DN}`$ can be thought as a kind of projection from $`๐_N`$ on $`๐_D`$), then $`๐_N=๐_D๐_{DN}๐ฒ_N๐_{DN}^T๐_D^T`$. Substitute it into Eq. (10), we have
$`SNR_{\mathrm{max}}`$ $`=`$ $`10\mathrm{log}_{10}tr(๐ฒ_D)10\mathrm{log}_{10}(tr(๐ฒ_D)`$
$`tr(\{๐ฒ_D+\mu ^2๐_{DN}๐ฒ_N๐_{DN}^T\}^1๐ฒ_D^2)).`$
If the noise components are white, we have $`๐ฒ_N=\sigma ^2๐_m`$ (with $`\sigma `$ being the standard deviation of the noise process) and $`๐_N=๐_D`$ (i.e. $`๐_{DN}=๐_m`$) allen optimal (1). However, for the case of colored noise, usually $`๐_{DN}๐_m`$. Instead it is possible that the absolute value of the elements in $`๐_{DN}`$ are relatively small. Thus even for the same clean signal $`\{d_i\}_i^M`$, the $`SNR_{\mathrm{max}}`$ performance of the colored noise might be much worse than that of the white noise. This fact might explain the observation that the results in Table. I are not that promising for the additive colored noise.
## Acknowledgement
This research was supported by Hong Kong University Grants Council Competitive Earmarked Research Grant (CERG) No. PolyU 5216/04E.
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# BeppoSAX/PDS serendipitous detections at high galactic latitudes
## 1 Introduction
The hard Xโray sky is still poorly explored at high galactic latitudes and the only truly all sky survey performed so far above $``$15 keV dates back to the 1980โs (Levine et al. (1984)). This pioneering work was performed with the *HEAO 1* A4 experiment and detected $``$70 sources in the 13โ80 keV band down to a flux limit of $`23\times 10^{10}`$ erg cm<sup>-2</sup> s<sup>-1</sup> with an angular resolution of $``$3; of these sources sixteen were located at high galactic latitudes ($`|b|13^{}`$) and only seven were of extragalactic nature.
A step forward will be provided by the imager onโboard *INTEGRAL* which is surveying a great fraction of the sky with a sensitivity better than a few mCrab <sup>1</sup><sup>1</sup>1For a Crabโlike spectrum 1 mCrab corresponds to $`2\times 10^{11}`$ erg cm<sup>-2</sup> s<sup>-1</sup> in the 20โ100 keV energy band in the 20โ100 keV energy range and an angular resolution of a few arcmin (Bassani et al. 2004, Bird et al. 2004); this exploratory work will be followed by the *SWIFT* mission which is expected to survey the hard Xโray sky with a sensitivity of $``$1 mCrab at high galactic latitudes (Gehrels et al. (2004)). In the meantime, the *Beppo*SAX archive can be used to probe the extragalactic Xโray sky in the same 20โ100 keV energy band. Pointed observations of *Beppo*SAX/PDS have unveiled many hard Xโray emitting AGN and provided the best yet spectroscopy of this type of objects above 10 keV. However, observations were sometimes limited by the lack of imaging capability: contaminating sources were found inside the target field of view as well as in the offset fields used for background measurements. These data have often been neglected by the original observers although they provide another powerful tool with which to study the extragalactic sky above 10 keV: new sources can be found and spectroscopically measured, while known objects can be reโobserved and compared to previous measurements in a search for variability, which is a dimension so far poorly explored. Furthermore, the argument put forward by Fabian (2001) that two of the three nearest AGN have very high column densities ($`>10^{24}`$ cm<sup>-2</sup>) indicates the need to constrain the statistics of highly absorbed AGN: increasing the number of such objects known is therefore an important task. In view of these facts, we have carried out a program to search systematically in the *Beppo*SAX/PDS archive (which at the moment is being reanalysed using the *XAS* software package) for observations which indicate the presence of a contaminating source either in the pointed field of view (search mode 1) or in the offset fields used for background measurements (search mode 2). Results related to search mode 1 are presented in this work, while a future work will be devoted to results found in search mode 2.
The paper is organized as follows: in Sec. 2 the selected sample is described, in Sec. 3 the observations, data reduction techniques and analysis procedures are presented, while Sec. 4 is devoted to the discussion of the results. Finally, conclusions and future work are summarized in Sec. 5.
## 2 Sample selection and contaminating source search
Although *Beppo*SAX was not designed to perform an Xโray survey, a systematic search in the entire PDS data archive could allow the discovery of new sources, down to a flux limit of $``$$`10^{11}`$ erg cm<sup>-2</sup> s<sup>-1</sup>. In particular, the simultaneous monitoring of three $`1^{}.3`$ sky regions (target field + two offset fields) allowed the PDS to survey a substantial fraction of the sky over the 15โ100 keV energy range.
The PDS instrument has a hexagonal field of view of $`1^{}.3\times 1^{}.3`$ FWHM and no imaging capability; its positional uncertainty can be approximated by an error box circle of $`1^{}.3`$ in radius. The MECS has instead a field of view of $`30^{}`$ radius and so covers about 25$`\%`$ of the PDS area. It is therefore likely that a serendipitous source can be undetected by the MECS, but still observed by the PDS.
In this paper we concentrate on the study of those sources discovered in the pointed fields of view or detected in search mode 1. The search in this mode was performed in the following way: first, we visually inspected in the *Beppo*SAX archive (www.asdc.asi.it) all observations performed above $`13^{}`$ in galactic latitude and available to the public as of October 2001 (635); from this set of data we extracted a sample containing those sources which clearly show a mismatch between the MECS and PDS spectra in the standard archive analysis: this mismatch was taken as strong evidence for the presence of a contaminating source in the PDS field of view. From this preliminary list we excluded all those sources that were likely to be Compton thick on the basis of various considerations (i.e. large iron line equivalent width, Xโray to \[$`O_{III}`$\] ratio less than 1, extreme absorption): in the spectrum of these type of sources a MECS/PDS spectral mismatch is simply due to a more complex spectral shape than used for the quick look analysis. Then, we performed a cut in the signal to noise ratio, accepting a source only if it had at least 3$`\sigma `$ detection in the PDS. We then confirmed with our own analysis the MECS/PDS mismatch by fitting the MECS and PDS data with a simple model, generally a power law either unabsorbed or absorbed: only when the crossโcalibration constant between these two instruments was significantly outside the nominal range of 0.75โ0.95 (Fiore et al. 1999a ), i.e. greater than 2, was a source maintained. Overall, this analysis has provided a sample of twelve regions which are discussed in the present paper. Although this sample is not complete (not all observations performed by *Beppo*SAX were screened) nevertheless this search gives an idea of the extragalactic sources that populate the hard Xโray sky; in this sense the sources serendipitously found in this work can be โlooselyโ taken as representative of the AGN population in the 20โ100 keV band.
An important step in the search described above, is related to the reduction of the PDS spectra which were extracted using the *XAS* v2.1 package (Chiappetti & Dal Fiume (1997)): this provides slightly different results than *SAXDAS*, i.e. the package usually used to perform the archive analysis. A preliminary comparison between PDS spectra extracted by means of these two software packages is under way<sup>2</sup><sup>2</sup>2On behalf of the PDS group, see ftp://ftp.tesre.bo.cnr.it in the directory /pub/sax/doc/software\_docs/xas\_vs\_saxdas.ps. and will be presented in a future work. The analysis performed on sources of different intensity show that the spectral parameters do not change when computed with the different packages, but the associated errors are smaller when using *XAS*. It is also important to underline that a significant improvement in the signal to noise ratio is obtained by means of *XAS*. Furthermore, the *XAS* package allows a more reliable check of the background fields by taking advantage of the rocking technique (Frontera et al. 1997b ); this is rather important when the source is faint in the PDS as is often our case. When one collimator is pointing ON source, the other collimator is pointing in one of the two OFF positions. The standard stay time in each position of either collimator was 96 s. At each cycle the two collimators were swapped: the one pointing to the source was moved to monitor the background and vice versa. In this way we can obtain, in addition to the target observation, two independently accumulated spectra of the two +OFF and โOFF fields which are offset by $`210^{}`$ with respect to the main pointing; these offsets are used as background in the computation of the spectrum of the target source. The comparison between the spectra of these two offset fields and, in particular, the difference between the +OFF and โOFF spectra in count rate, is a good diagnostic tool to investigate the presence of contamination.
If no contamination is present in either of the offset fields, we expect the difference of their count rates to be compatible with zero; on the contrary, a positive excess of counts in the difference indicates the presence of contamination in the +OFF field; vice versa a negative count rate provides evidence for contamination in the โOFF field.
Applying this method to our sample, we found the presence of a contaminating source in the +OFF field of the first observation of AD Leonis (excess: $`0.189\pm 0.070`$ counts s<sup>-1</sup>, 2.7$`\sigma `$) and in part of the first observation of VW Cephei (excess: $`0.139\pm 0.043`$ counts s<sup>-1</sup>, 3.2$`\sigma `$), while in part of the observations of MKN 1073 and NGC 7552 the โOFF fields show an excess of $`0.159\pm 0.060`$ counts s<sup>-1</sup> (2.7$`\sigma `$) and $`0.203\pm \mathrm{0.0.056}`$ counts s<sup>-1</sup> (3.6$`\sigma `$), respectively. To give an idea of the effect of contamination in an offset field on the PDS data, we have considered two extreme cases: a $``$3$`\sigma `$ excess reduces the source count rate of $``$5$`\%`$, while a 10$`\sigma `$ contamination provides a reduction at around 15$`\%`$.
In any case, in order to extract the uncontaminated source spectra, we excluded the contaminated fields and considered only the uncontaminated one in the computation of the background for these particular sources.
After assembling the sample, we search for the likely contaminating sources, adopting the following strategy. First, we considered all the sources present in the MECS image in addition to the target source and checked their consistency with the PDS data, i.e. if the crossโcalibration constant fell within the nominal interval. If no sources in the MECS field of view matched the PDS data, then, we searched for likely high energy emitters located inside the PDS field of view but not observed by the MECS given the significantly different fields of view. In particular, when adopting this second approach, we focused our search on bright sources in the 2โ10 keV band as these are the most likely to contaminate the PDS observation and are expected to appear in the HEASARC Xโray archives. In Table 1 we list the twelve cases we found, reporting in each case the *Beppo*SAX observation target with its relative observation date and object type, the PDS/MECS crossโcalibration constant obtained by fitting the data with a simple power law and finally the name and type of the contaminating source found. It is evident from the Table that the crossโcalibration constant is always outside the nominal range, confirming the presence of one or more contaminating objects in the PDS field of view. It is worth noting that extrapolation of a more complex model from the MECS to the PDS energy range, still provides in all cases a high cross calibration constant. We found only four fields where the contaminating source is so close ($``$$`25^{}`$) to the*Beppo*SAX target that it is also detected by the MECS instrument.
## 3 Observations and data analysis
In this work we made use of data from three of the *Narrow Field Instruments* (NFIs) onโboard the ItalianโDutch satellite *Beppo*SAX (Boella et al. 1997a ): the Low Energy Concentrator Spectrometer (LECS, 0.1โ10 keV, Parmar et al. (1997)), the Medium Energy Concentrator Spectrometer (MECS, 1.3โ10 keV, Boella et al. 1997b ), and the Phoswich Detection System (PDS, 15โ300 keV, Frontera et al. 1997b ).
For all sources, the LECS and MECS spectra were downloaded from the *Beppo*SAX archive; the analysis of the onโaxis source was standard. For all the offโaxis sources detected in the MECS field of view, contaminating or not, the MECS spectra were extracted from a region centered on the source and having a radius chosen according to the criteria suggested by Fiore et al. (1999a). For these sources, because of the lack of the appropriate ancillary response files, the LECS spectral analysis could not be performed. The background subtraction for the onโaxis sources was performed using blank sky spectra extracted from the same region of the source, while for the offโaxis sources we used a local background spectrum (extracted from a region with a radius equal to the source extraction radius) to account for possible contaminating effects inside the MECS field of view.
For the PDS data reduction, source visibility windows were selected following the criteria of no Earth occultation and high voltage stability during the exposure. In addition, the observations closest to the South Atlantic Anomaly were discarded from the analysis.
The LECS and MECS spectra were rebinned in order to sample the energy resolution of the detectors with an accuracy proportional to the count rate. The PDS data were instead rebinned so as to have logarithmically equal energy intervals. The data rebinning also required that at least 20 counts was in each bin so that the $`\chi ^2`$ statistic could reliably be used.
The energy bands used for spectral fitting were limited to those where the response functions are well known, i.e. 0.1โ4.5 keV, 1.5โ10.5 keV, and 15โ100 keV, for the LECS MECS, and PDS, respectively.
The spectral analysis was performed using the XSPEC v11.2.0 software package (Arnaud (1996)) and the instrument response matrices released by the *Beppo*SAX ASI Science Data Center. For the offโaxis sources, we used the appropriate MECS ancillary response files to correct for the effects of vignetting due to the mirrors. In addition, we introduced a flux correction factor in the PDS band for each contaminating source in order to estimate the real flux at the source. This correction is simply a function of the distance of the source from the main target and is related to the PDS response: the reduction in sensitivity is of a factor of $``$2 at $`38^{}`$ from the main pointing coordinates while at $`78^{}`$ the response is zero (Frontera et al. 1997a ).
Normalization constants have also been introduced to allow for known differences in the absolute crossโcalibration between the detectors. The LECS/MECS crossโcalibration constant was allowed to lie within the nominal range 0.7โ1.0 (Fiore et al. 1999a ), while the PDS/MECS crossโcalibration was left free to vary in order to search for the presence of contaminating sources and then to estimate their broad band spectrum. The absorption of Xโrays due to our galaxy in the direction of each object (Dickey & Lockman (1990)) is added in all models of the spectral analysis. All quoted errors correspond to $`90\%`$ confidence interval for one interesting parameter ($`\mathrm{\Delta }\chi ^2=2.71`$).
The main aim of this work is to estimate the PDS spectrum of each serendipitous source found. The method we used to achieve this was to fit simultaneously the LECS (not always available) and MECS data of the target source to determine the bestโfit model to these data; afterwards, we performed an extrapolation of this bestโfit model to the PDS energy band in order to estimate the contribution of the target source to the total PDS flux. Hence, by subtracting from the total PDS spectrum the contribution due to the onโaxis source, we were able to estimate the high energy spectrum of the serendipitous/contaminating object.
Four sources in our sample (NGC 1553, 1E 1839.6+8002, VW Cephei and AD Leonis), have been observed by *Beppo*SAX at different epochs (see Table 1). As a first step, our approach was to analyse each single pointing, for each source, individually. Although in some cases we found evidence for flux variability in the target and/or the serendipitous sources, the spectral parameters were always consistent within the respective uncertainties; therefore, we decided to sum, in all cases, all available pointings to improve the statistical quality of the data, reporting only the averaged spectral parameters.
## 4 Results
In the following, we briefly describe the main results obtained for each field. We discuss separately the four fields in which the contaminating object is detected in the MECS instrument because for these objects we can perform a broad band spectral analysis (1.5โ100 keV), while in the remaining eight cases we can only present and discuss the PDS spectrum.
### 4.1 Sources detected in the MECS field of view
In this section, we describe the four fields where the object likely to be responsible for the PDS spectrum is also observed by the MECS instrument. In Table 2 (and also in Table 5 afterwards) we report for each of these fields, all sources detected by the MECS, their fluxes extrapolated from the MECS to the PDS energy range (15โ100 keV) assuming their best fit model, their relative contribution (in $`\%`$) to the PDS flux, their distance from the target source and the relative flux correction factor, $`R`$. The observed flux is divided by $`R`$ to reconstruct the flux at the source. It is evident from the values reported in column 4 (Contribution) which source is in each case responsible for most of the PDS flux. In Table 3 we list for each of these likely associations, the MECS and PDS exposures and count rates, the 2โ10 keV flux, the observed high energy flux (20โ100 keV) and the effective flux (20โ100 keV) at the source. For these four sources we jointly fitted the MECS and PDS data to provide a broad band (20โ100 keV) spectrum. In Table 4 we summarize the bestโfit parameters and also report the range of the crossโcalibration constant PDS/MECS (column 6). As is evident from the values of column 6, the normalization constants obtained are always compatible with the suggested value of 0.75โ0.95. This is a further indication that these sources are indeed responsible for the high energy emission detected by the PDS instrument.
#### 4.1.1 The ON 325 Field
Within the MECS field of view (see Fig. 1) of this pointing we find two other sources besides the target. The first, located at $`\alpha (2000)=12^h18^m54^s.4`$ and $`\delta (2000)=+29^{}58^{}10^{\prime \prime }.5`$, is identified with the Seyfert 1.9 galaxy 1AXG J121854+2957, which belongs to the *ASCA* Medium Sensitivity Survey (AMSS) (Ueda et al. (2001)) and to the *Beppo*SAX High Energy Large Area Survey (HELLAS) (Fiore et al. 1999b ). The second is located at $`\alpha (2000)=12^h18^m25^s.2`$ and $`\delta (2000)=+29^{}48^{}48^{\prime \prime }.9`$, $``$$`20^{}`$ from ON 325 and corresponds to the Seyfert 1 galaxy MKN 766.
Following Perri et al. (2003), the LECS and MECS data of the blazar ON 325 are best fitted with a broken power law, for which we obtain a $`\chi ^2/\nu =32.4/25`$ and a 2โ10 keV flux of $``$$`8\times 10^{13}`$ erg cm<sup>-2</sup> s<sup>-1</sup>. The two power laws intersect at $`E_B=4.12_{0.83}^{+1.10}`$ keV and have photon indices $`\mathrm{\Gamma }_1=2.41_{0.27}^{+0.18}`$ and $`\mathrm{\Gamma }_2=0.44_{0.85}^{+0.79}`$ in agreement with Perri et al. (2003). Adopting this model, we estimate a contribution of ON 325 to the PDS data of the order of 10$`\%`$ (see Table 2).
The MECS spectra of 1AXG J121854+2957 is instead well fitted ($`\chi ^2/\nu =10.8/14`$) with an absorbed power law having a photon index $`\mathrm{\Gamma }`$ $``$2.3, a column density N$`{}_{\mathrm{H}}{}^{}<12\times 10^{22}`$ cm<sup>-2</sup> and a 2โ10 keV flux of $`8.5\times 10^{13}`$ erg cm<sup>-2</sup> s<sup>-1</sup>. These values are in agreement with those found by Loaring et al. (2003) during a *XMMโNewton* observation of the source. Extrapolation of this spectrum to the PDS band indicates a small contribution to the high energy flux, while a simultaneous fit to the MECS/PDS data provides a PDS/MECS crossโcalibration constant in the range 2โ30, i.e. well outside the nominal interval. A good fit ($`\chi ^2/\nu =30/34`$) of the MECS data of MKN 766 is also described by a simple power law with a photon index $`\mathrm{\Gamma }=1.92\pm 0.13`$ and a 2โ10 keV flux of $`1.7\times 10^{11}`$ erg cm<sup>-2</sup> s<sup>-1</sup>. The high contribution to the PDS flux (see Table 2) found with this model indicates that MKN 766 is the best candidate to account for the PDS flux. In fact, the fit of the MECS spectrum of MKN 766 with the PDS data gives a PDS/MECS crossโcalibration constant well inside the nominal interval. The 2โ10 keV flux (see Table 3), found during this pointing, is slightly lower than the value of $`2.05\times 10^{11}`$ erg cm<sup>-2</sup> s<sup>-1</sup> found by Matt et al. (2000a) during a previous (May 1997) *Beppo*SAX dedicated observation. We find no evidence for an iron line at around 6.4 keV (see Fig. 2), consistent with the fact that this feature seems to be present only when MKN 766 is in a high state (Leighly et al. (1996)).
To confirm if MKN 766 is really the contaminating source in the ON 325 field, we also reanalysed the PDS data of the two pointings of MKN 766 performed by *Beppo*SAX on May 1997 and 2001. In the first observation the comparison between the +OFF and โOFF field count rates shows an excess of $`3\sigma `$ in the โOFF field, while in the second there is no evidence for contamination.
Searching in the HEASARC archive for possible objects which could contaminate the โOFF field, we found two sources belonging to the *ROSAT* Bright Source Catalog: a quasar (1RSX J121320.3+270841) and an unidentified object (1RSX J121258.3+272653), at a distance of $``$40 and $``$59 from the field centre. Fitting the โOFF data with a simple power law, we obtain $`\mathrm{\Gamma }2`$ and a 20โ100 keV flux of $`3.6\times 10^{11}`$ erg cm<sup>-2</sup> s<sup>-1</sup>. Since both sources have similar count rates in *ROSAT*, we can only state that either or both could be responsible for this excess.
In any case, this provides a good example of the capability of our analysis: contamination in the offset fields can be found and identified, while at the same time a more correct evaluation of the target spectrum is provided.
The PDS spectrum of MKN 766 in the May 1997 observation can be corrected and reanalysed together with that of May 2001. In both cases the photon index $`\mathrm{\Gamma }`$ is $``$2, in agreement with that found in December 1998 (i.e. during the ON 325 observation), while the 20โ100 keV flux is roughly $``$$`3\times 10^{11}`$ erg cm<sup>-2</sup> s<sup>-1</sup>, higher than during the ON 325 measurement; this finding indicates that the flux variability always seen in this source at soft Xโray energies (Matt et al. 2000a and references therein) extends to above $``$10 keV.
#### 4.1.2 The NGC 5793 Field
NGC 5793 has been detected at $``$$`5\sigma `$ level in the MECS instrument, while it is insignificant ($`<1\sigma `$) in the LECS energy range. Fitting the MECS data with a simple power law of photon index $`\mathrm{\Gamma }=1.9`$ we obtain a very low flux of $``$$`7\times 10^{14}`$ erg cm<sup>-2</sup> s<sup>-1</sup>. This model, extrapolated to the high energy band, gives a negligible contribution to the PDS data (see Table 2).
Inspection of Fig. 3 indicates the presence of two extra sources in the MECS field of view. The first located at $`\alpha (2000)=14^h58^m50^s.7`$ and $`\delta (2000)=16^{}51^{}50^{\prime \prime }`$.7, corresponds to a bright irregular spiral galaxy belonging to the 2 Micron All Sky Survey eXtended Source Catalog (i.e. 2MASX J14585116โ1652223). The second is located at $`\alpha (2000)=14^h59^m23^s.8`$ and $`\delta (2000)=16^{}37^{}13^{\prime \prime }.4`$ and is identified with the normal galaxy NGC 5796.
We extracted the MECS spectra of both sources and extrapolated them to high energies. We find that the MECS spectrum of NGC 5796 is well fitted ($`\chi ^2/\nu =15/20`$) with a simple power law having a photon index $`\mathrm{\Gamma }=1.95_{0.56}^{+0.67}`$ and a 2โ10 keV flux of $`1.5\times 10^{13}`$ erg cm<sup>-2</sup> s<sup>-1</sup>. The contribution of NGC 5796 to the PDS flux is only a few percent and the PDS/MECS crossโcalibration constant is extremely high (20โ103).
The MECS bestโfit model in the case of 2MASX J14585116โ1652223 is a simple power law with photon index $`\mathrm{\Gamma }=1.15_{0.18}^{+0.33}`$ and a 2โ10 keV flux of $`9.3\times 10^{13}`$ erg cm<sup>-2</sup> s<sup>-1</sup> ($`\chi ^2/\nu =15/20`$). The extrapolation of this model to high energies (see Table 2) suggests that 2MASX J14585116โ1652223 is likely to be the contaminating object. The simultaneous fit of the MECS/PDS data (see Table 3 and Fig. 4) shows a flat spectrum ($`\mathrm{\Gamma }=1.20_{0.23}^{+0.26}`$), a more suitable value (0.7โ3.5) for the PDS/MECS crossโcalibration constant and no evidence for extra absorption. Fixing $`\mathrm{\Gamma }=1.9`$ and adding absorption to the power law provides a similar fit and an upper limit to the column density of $``$$`4\times 10^{22}`$ cm<sup>-2</sup>. As displayed in Table 3, the 20โ100 keV flux turns out to be of $`1.11\times 10^{11}`$ erg cm<sup>-2</sup> s<sup>-1</sup>, one order of magnitude higher than the 2โ10 keV flux, confirming that 2MASX J14585116โ1652223 emits mostly in the high energy band. Very little is known of this object except for its near infrared properties: the source is fairly bright with total photometry of 14.4, 13.6 and 13 magnitudes in the J, H and K bands respectively and an extent of $`11^{\prime \prime }`$; the optical counterpart has magnitudes B $`=16`$ and R $`=15`$. No previous Xโray data are reported, including the lack of detection by *ROSAT*, which would be consistent with the presence of strong absorption. The RโK colour is $``$2, i.e. quite red which again point to an object with considerable extinction. No detection in the radio band was found in the literature or in the HEASARC archive.
#### 4.1.3 The 1E 1839.6+8002 Field
The M4 Ve star 1E 1839.6+8002 was observed by *Beppo*SAX twice (October 2000 and February 2001). This source was detected for the first time during the *Einstein* observation of the Broad Line Radio Galaxy 3C 390.3 in the 1980โs. We stress that a detailed study of the spectral behaviour of this source is beyond the aim of this paper. Nevertheless, following Pan et al. (1997), the LECS and MECS spectra of both observations are well described by a thermal model (RaymondโSmith model in XSPEC, Raymond & Smith (1977)), with a temperature $`kT`$ $``$1.8 keV and a 2โ10 keV flux of $`9.5\times 10^{14}`$ erg cm<sup>-2</sup> s<sup>-1</sup> (first observation) and $`7.2\times 10^{14}`$ erg cm<sup>-2</sup> s<sup>-1</sup> (second observation).
The source appears to be faint and variable ($``$25$`\%`$) in the 2โ10 keV range, with a negligible contribution to the high energy emission in both observations (in Table 2 only the averaged contribution of the source is reported).
Within the MECS field of view (see Fig. 5) of both measurements, a bright object located at $`\alpha (2000)=18^h41^m51^s.5`$ and $`\delta (2000)=+79^{}47^{}01^{\prime \prime }.3`$ is clearly visible: it corresponds to the well known radio galaxy 3C 390.3. The *Einstein* IPC data of 3C 390.3 revealed the presence of strong intrinsic absorption (Kruper et al. (1990)), while at higher energies *Ginga* data showed the presence of an iron line at 6.4 keV (Inda et al. (1994)) and a reflection component (Nandra & Pounds (1994)). *ASCA* data confirmed the presence of the iron K emission line but could not constrain the reflection hump, most probably because of the limited energy range (Eracleous, Halpern & Livio (1996); Leighly et al. (1997)). 3C 390.3 has also been detected by *OSSE* in the softโ$`\gamma `$ energy range, above 50 keV (Dermer & Gherels (1995)). All these results have been confirmed by Grandi et al. (1999), on the basis of a dedicated *Beppo*SAX observation performed in January 1997. The above authors also found a column density variability when comparing *Beppo*SAX results with previous measurements.
No evidence for variability is present in the MECS data of 3C 390.3 nor in the PDS analysed in this work. Hence, in order to improve the statistics, we summed the two data sets together to perform our spectral analysis. First, we analysed only the MECS data of 3C 390.3, finding as bestโfit model ($`\chi ^2/\nu =13.5/17`$) a power law with a slope of 1.7 plus a narrow gaussian line having a centroid energy of 6.4 keV and an equivalent width of $`EW=115`$ eV, compatible with K<sub>ฮฑ</sub> line emission from neutral iron located at the source redshift. The contribution of the 3C 390.3 flux to the total PDS flux is significant (see Table 2) and clearly indicates that this source is indeed responsible for the contamination.
The joint fit of the MECS data of this source with the PDS is also well reproduced by a power law plus a narrow gaussian line ($`\chi ^2/\nu =16.5/22`$) having parameters similar to those found fitting only the MECS data. As for the column density, we can only provide an upper limit of $``$$`8\times 10^{20}`$ cm<sup>-2</sup>, compatible with the values found and reported by Grandi et al. (1999). The broad band spectrum of the source is shown in Fig. 6. The flux turns out to be $`2.2\times 10^{11}`$ erg cm<sup>-2</sup> s<sup>-1</sup> and $`4.3\times 10^{11}`$ erg cm<sup>-2</sup> s<sup>-1</sup> in the 2โ10 and 20โ100 keV band respectively, again in agreement with that found by Grandi et al. (1999). The PDS/MECS crossโcalibration constant turns out to be consistent (see Table 4) with the suggested values, confirming that 3C 390.3 is the serendipitous source contaminating the PDS emission during the observations of the 1E 1839.6+8002 field of view (see Table 2). There is no evidence for variability between the present (2000โ2001) *Beppo*SAX observations and that reported by Grandi et al. (1999) in the 1997 observation.
#### 4.1.4 The SCG 2353โ6101 Field
In the case of the cluster of galaxies SCG 2353โ6101 (Abell 4067) there is another source in the MECS field of view, PKS 2356โ611, located at $`\alpha (2000)=23^h59^m04^s.1`$ and $`\delta (2000)=60^{}55^{}01^{\prime \prime }.5`$ (see Fig. 7).
PKS2356โ611 ($`z=0.096`$, B $`=12.3`$, R $`=12.8`$, as reported in the USNOโB1 Catalog (Monet et al. (2003))) is one of the strongest southern FR II sources, with a total radio power $`P_{1.4\mathrm{GHz}}`$ $``$$`10^{25.8}`$ W Hz<sup>-1</sup> (Koekemoer et al. (1998)), while optically it shows strong highโexcitation narrow line emission with \[O<sub>III</sub>\]$`\lambda `$5007/H<sub>ฮฒ</sub> $`>`$ 10. Lipovetsky et al. (1988) have cataloged this source as a Seyfert galaxy of type 2. It is also listed in the 2 Micron All Sky Survey eXtended Source Catalog as 2MASX J23590436โ6054594 (with magnitudes J $`=13.7`$, H $`=12.7`$, K $`=12.9`$) and was not detected by *ROSAT*.
The LECS and MECS spectra of SCG 2353โ6101 are well fitted ($`\chi ^2/\nu =32/37`$) with a thermal bremsstrahlung model, giving a temperature of $`4.74_{2.26}^{+1.70}`$ keV. The extrapolation of this model to the PDS energy band provides negligible flux. (see Table 2). This, combined with the high value (8โ28) of the PDS/MECS crossโcalibration, indicates the presence of a high energy emitting source and PKS2356โ611 is a good candidate.
As shown in Fig. 7, PKS2356โ611 is located near one of the two <sup>55</sup>Fe calibration sources of the MECS instrument. In this case a careful choice of the background region is required to avoid contamination by 5.95 keV photons produced by the calibration source: this is done taking the background region as near as possible to the calibration sources itself.
Fitting the MECS data of PKS 2356โ611 and the PDS points with a single power law plus intrinsic absorption provides an acceptable fit ($`\chi ^2/\nu =60/43`$), a steep spectrum ($`\mathrm{\Gamma }=1.75`$) and a column density of $`N_\mathrm{H}`$ $``$$`1\times 10^{23}`$ cm<sup>-2</sup> (Fig. 8). Both spectral parameters are typical of Seyfert 2 galaxies, and seem to confirm the source classification made by Lipovetsky et al. (1988).
The crossโcalibration constant in this case turns out to be lower (0.2โ0.5) than generally observed. This could be due to imperfect background correction in the MECS because of the location of the source near the calibrator. In any case, the spectral parameters do not change significantly when this constant is constrained to vary within the nominal range of values.
An inspection of the broad band spectrum of PKS 2356โ611 (see Fig. 8) suggests the possible presence of an excess around $``$6 keV. Although the addition of an extra component in the form of a narrow gaussian line at around 6.4 keV with $`EW`$ $``$500 eV provides an improvement in the fit ($`\mathrm{\Delta }\chi ^2=11.5`$ for two additional parameters), its reality is questioned by the location of the source near the <sup>55</sup>Fe calibrator: residuals at line energy could still be due to a non perfect subtraction of the calibration line. Also in this case the crossโcalibration constant turns out to be low ($`<`$ 0.5) and the $`\chi ^2`$ values as well as the spectral parameters do not change significantly if we vary the crossโcalibration constant within the allowed range. The absorption measured would produce an iron line of 100โ200 eV in transmission, while reflection in the torus would increase the $`EW`$ to about 400 eV (Turner et al. (1997)) possibly indicating the presence of a reflection component. However, the addition in the fit of this component via the pexrav model in XSPEC (Magdziarz & Zdziarski (1995)) is not required by the data and does not improve the crossโcalibration constant, which turns out to be very small ($`<0.4`$). A large $`EW`$ is also possible if the source is Compton thick; this type of object is generally characterized by a low $`F_\mathrm{X}`$/$`F_{\mathrm{OIII}}`$ ratio (Bassani et al. (1999)). In PKS 2356โ611 this ratio is $``$67, when $`F_{\mathrm{OIII}}`$ is corrected for the reddening in the galaxy using the observed value of $`H_\alpha `$/$`H_\beta =4.32`$ (Koekemoer et al. (1998)). This result suggests that PKS 2356โ611 is a Compton thin Seyfert 2 galaxy rather than a Compton thick one and so the large $`EW`$ is probably not due to strong absorption ($`N_\mathrm{H}>1.5\times 10^{24}`$ cm<sup>-2</sup>).
### 4.2 Sources not detected in the MECS field of view
In this section we describe the remaining eight fields where the source responsible for the PDS spectrum is not observed by the MECS. In Table 5 we report for each of these fields, the source/s observed in the MECS, their flux extrapolated to the 20โ100 keV band and their contribution to the PDS flux. In Table 6 instead we list for each field, the PDS exposure and count rate as well as the photon index $`\mathrm{\Gamma }`$ and observed 20โ100 keV flux obtained assuming a simple power law fit. We also report the effective flux at the source, estimated from that observed by applying the correction factor $`R`$.
#### 4.2.1 The IRAS 01025โ6423 Field
In this field only the target source, IRAS 01025โ6423, a Seyfert 2 galaxy is barely detected by MECS at a $``$2$`\sigma `$ level, while the LECS data are not available. Due to the low signal to noise ratio, we only attempted to put an upper limit on the flux assuming a power law model with photon index $`\mathrm{\Gamma }=2`$. The source appears to be faint in the 2โ10 keV range with a flux $`3.1\times 10^{14}`$ erg cm<sup>-2</sup> s<sup>-1</sup> and its contribution to the PDS signal is obviously null (see Table 5). Searching in the HEASARC archive for possible contaminating sources within the PDS field of view, we found two possible candidates belonging to the *ROSAT* Bright Source Catalog: 1RXS J011316.2โ641142 (CPDโ64 120), classified as a K1 Ve star, and 1RXS J010333.7โ643925 (PKS 0101โ649), identified as a radio source, at distances from the target source of $``$58 and $``$33, respectively. PKS 0101โ649, which is the closest to the target source, is more likely to be the object responsible for the high energy emission, as a star is not expected to emit above 10 keV. PKS0101โ649 belongs to the 2 Micron All Sky Survey eXtended object (i.e. 2MASX J01033376โ6439079 in NED), with total photometry of J $`=15`$, H $`=14`$ and K $`=13.4`$ magnitudes; the USNOโA2 R magnitude is 16.5 (Monet et al. (1999)) thus providing an RโK value of 3.1 which suggests a quite red object; at radio frequencies it is fairly bright with a 6 cm flux of $``$270 mJy.
The PDS spectrum is characterized by a photon index $`\mathrm{\Gamma }`$ $``$1.5 (see Table 6) compatible, within uncertainties, with an AGN spectrum; the 20โ100 keV flux is $`2.3\times 10^{11}`$ erg cm<sup>-2</sup> s<sup>-1</sup>, making PKS 0101โ649 a bright source at these energies. If we extrapolate this power law to soft Xโray energies, we find a 0.1โ2.4 keV flux of $`0.22\times 10^{11}`$ erg cm<sup>-2</sup> s<sup>-1</sup> which is consistent with the *ROSAT* value of PKS 0101โ649. Alternatively, the Xโray source spectrum could be steeper but highly absorbed even above 10 keV as observed in Compton thick objects (Matt et al. 2000b ). If we add absorption to the power law, fixing the photon index to 1.9, we find an upper limit to the column density $`<2\times 10^{25}`$ cm<sup>-2</sup>. Obviously, an optical spectrum would be highly desirable to provide more information and establish the true nature of this source.
#### 4.2.2 The MKN 1073 Field
The Seyfert 2 galaxy MKN 1073 was not clearly detected by the LECS and the MECS instruments. In both detections there is evidence for diffuse emission probably associated with the Perseus cluster, which is at $`62^{}`$ from the centre of the MECS pointing. The signal in the PDS is high ($`>10\sigma `$ level), with an exceptionally steep spectrum ($`\mathrm{\Gamma }`$ $``$3.3, see Table 6). We ascribe the PDS emission to the Perseus cluster, which has also been targetted by *Beppo*SAX on September 1996. The spectral analysis of these PDS data also shows a very steep spectrum ($`\mathrm{\Gamma }=3.80_{0.18}^{+0.17}`$) and a 20โ100 keV flux of $``$$`5\times 10^{11}`$ erg cm<sup>-2</sup> s<sup>-1</sup>, lower than the flux seen during the MKN 1073 observation (see Table 6). Since flux variability in the cluster is unlikely, the change in the emission could be due to the Seyfert 2 galaxy NGC 1275, often taken as responsible for the high energy ($`>`$ 20 keV) emission observed from this region; comparison of data from high energy instruments over years shows a large variation in the flux above 20 keV (Osako et al. (1994)). It is also possible that both cluster and AGN contribute to the high energy emission, but it is difficult at present to disentangle one contribution from the other.
#### 4.2.3 The NGC 1553 Field
*Chandra* observations of the Xโray faint SO galaxy NGC 1553 (Blanton et al. (2001)) has spatially and spectrally resolved the source of the Xโray emission. A significant fraction of this ($``$70$`\%`$) is detected as diffuse flux while the remainder is due to 49 objects. The strongest source in the field is located at the centre of NGC 1553 to within $`0^{\prime \prime }`$.5 and shows a hard spectrum, typical of an AGN. NGC 1553 was observed twice by *Beppo*SAX (January 1997 and November 1997). Trinchieri et al. (2000) performed the spectral analysis of the LECS and MECS data finding as bestโfit a RaymondโSmith model ($`kT`$ $``$0.26 keV) for the soft component and a thermal bremsstrahlung ($`kT`$ $``$4.8 keV) for the hard component. Our analysis of the LECS and MECS data from both observations is in agreement with that reported by these authors; the extrapolation of this model to high energies provides a very low contribution to the PDS flux (see Table 5). Searching in past Xโray mission archives, we find that the Seyfert 1 galaxy NGC 1566 is the best candidate for the PDS emission as also suggested by Trinchieri et al. (2000). This source is at a distance of $`60^{}`$.2 from the target object and is therefore just inside the PDS field of view. NGC 1566 ($`z=0.005`$, B $`=9.9`$, R $`=9.6`$) is a very bright Xโray source known since the *HEAO 1* Xโray Source Catalog (Wood et al. (1984)). The source belongs also to the 2 Micron All Sky Survey eXtended Source Catalog (2MASX J04200041โ5456161), with magnitudes J $`=7.8`$, B $`=7.2`$ and K $`=6.9`$, it is reported by *ROSAT* in the Bright Source Catalog and it is a radio source of $``$100 mJy at 6 cm. The 2โ10 keV flux reported in the literature ranges from $`<`$ 0.5 to $`1.7\times 10^{11}`$ erg cm<sup>-2</sup> s<sup>-1</sup> (Halpern (1982)) indicating variability by around a factor of 3. The PDS emission shows a variation of $``$70$`\%`$ during our two observations, but due to the poor statistics in both measurements (the signal to noise ratio in the PDS is around $`3\sigma `$ level in each measurement), we summed the observations to improve the statistics and performed an averaged spectral analysis. As shown in Table 6, this analysis indicates a photon index $`\mathrm{\Gamma }`$ of $``$2.5 and an average 20โ100 keV flux, corrected for the offset, of $``$$`8\times 10^{11}`$ erg cm<sup>-2</sup> s<sup>-1</sup>. By extrapolating the flux down to the 2โ10 keV range, we find a value of $``$$`2\times 10^{10}`$ erg cm<sup>-2</sup> s<sup>-1</sup>, much higher with respect to previous measurements; however, if the photon index is restricted to values more appropriate for a Seyfert 1 galaxy (i.e. 1.7โ1.9), then the 2โ10 keV flux reduces to $``$$`6\times 10^{11}`$ erg cm<sup>-2</sup> s<sup>-1</sup>, only a factor of 3 higher. This result indicates that NGC 1566 could indeed dominate the PDS emission and that flux variability is likely.
#### 4.2.4 The AD Leonis Field
The Mโdwarf Ad Leonis has been the subject of an observational campaign performed by *Beppo*SAX (April 1997, 1, 8 and 12 May 1999). We adopt the multiโtemperature model (mekal code in XSPEC, Mewe et al. (1985)) proposed by van den Besselaar et al. (2003) for the *XMMโNewton* and *Chandra* observations, and find that the LECS and MECS data of the four observations are best described by a twoโtemperature model, in which the metal abundances are left free to vary. The 2โ10 keV flux varies in the range $`1.42.3\times 10^{12}`$ erg cm<sup>-2</sup> s<sup>-1</sup>, corresponding to a flux variation of $``$60$`\%`$, while the spectral parameters do not show significant variation. In view of these indications, we summed together the data from different pointings and estimated the average spectral parameters. The twoโtemperature model remains the best fit ($`\chi ^2/\nu =97/82`$), giving $`kT_1=0.65_{0.02}^{+0.03}`$ keV and $`kT_2=2.04_{0.17}^{+0.24}`$ keV metal abundances of $`0.14\pm 0.02`$ and $`0.56_{0.21}^{+0.25}`$ respectively and an averaged 2โ10 keV flux of $``$$`2\times 10^{12}`$ erg cm<sup>-2</sup> s<sup>-1</sup>; this spectrum if extrapolated to high energies provides almost no contribution to the PDS flux (see Table 5).
At a distance of $``$55 from AD Leonis, we find the Seyfert 1 galaxy NGC 3227, which could be the contaminating source we are searching for. NGC 3227, first detected in the *Ariel V* all sky survey, is one of the bright sources in the Piccinotti sample of AGN (Piccinotti et al. (1982)). NGC 3227 is known to be variable (on time scales of hoursโdays) and characterized by a spectrum with a photon index $`\mathrm{\Gamma }`$ $``$1.6, flatter than that typically observed in Seyfert galaxies (George et al. (1998)). *Ginga* observations performed in 1988 indicate a 2โ10 keV flux of $``$$`4\times 10^{11}`$ erg cm<sup>-2</sup> s<sup>-1</sup> (Pounds et al. (1989)), while *ASCA* observations performed during 1993 and 1995 provide 2โ10 keV flux measurements in the range $`2.42.6\times 10^{11}`$ erg cm<sup>-2</sup> s<sup>-1</sup>. Recently, Lamer et al. (2003), analyzing the *Rossi Xโray Timing Explorer (RXTE)* data of NGC 3227, confirmed the source variability, explaining it in terms of transient absorption by a gas cloud (neutral or weakly ionized) of column density of $``$$`3\times 10^{23}`$ cm<sup>-2</sup> moving across the line of sight to the Xโray source. This interpretation has been confirmed also by *XMMโNewton* observations (Gondoin et al. (2003)). NGC 3227 has been detected at higher energies by *OSSE* and *BATSE* onโboard the *Compton Gamma Ray Observatory (CGRO)*. The *OSSE* data when fitted in the 50โ200 keV range, show a spectral index $`\mathrm{\Gamma }=1.86_{0.38}^{+0.36}`$, softer with respect to that found in the 2โ10 keV band and a photon flux of ($`3.40\pm 0.55)\times 10^4`$ cm<sup>-2</sup> s<sup>-1</sup> in the 50โ150 keV band (Zdziarski et al. (2000)). Based on these measurements, we estimate a 20โ100 keV flux measurement in the range $`5.07.8\times 10^{11}`$ erg cm<sup>-2</sup> s<sup>-1</sup>. In addition, the *BATSE* data (Malizia et al. (1999)) show a 20โ100 keV flux in the range $`9.712.8\times 10^{11}`$ erg cm<sup>-2</sup> s<sup>-1</sup>. The spectral analysis of our PDS data indicate that, during the four observations of AD Leonis, the 20โ100 keV flux measurements of NGC 3227 ranged in the interval $`6.314.7\times 10^{11}`$ erg cm<sup>-2</sup> s<sup>-1</sup>, with an averaged value of $`9.8\times 10^{11}`$ erg cm<sup>-2</sup> s<sup>-1</sup> and a photon index $`\mathrm{\Gamma }`$ $``$1.85 (see Table 6). These results turn out to be in perfect agreement with both *OSSE* and *BATSE* data. We want to stress that our measurement contains the only spectral information about NGC 3227 obtained by *Beppo*SAX (the source has never been observed by the satellite) and so represent an interesting addition to the list of Piccinotti sample sources for which *Beppo*SAX high energy data are available.
#### 4.2.5 The H1846โ786 Field
H1846โ786 is a Seyfert 1 galaxy belonging to the Piccinotti sample (Piccinotti et al. (1982)). Despite its Xโray brightness, the source has been poorly studied. A simple power law, with a photon index $`\mathrm{\Gamma }`$ $``$1.95 is the bestโfit model to the LECS and MECS data ($`\chi ^2/\nu =125/100`$) as obtained by Quadrelli et al. (2003). This result indicates that no extra absorption and/or an iron line are present in the spectrum. The extrapolation of this bestโfit model to the PDS energy range shows a significant contribution from H1846โ786 (see Table 5). In fact, as can be seen in Table 1, this case represents the lowest range for the PDS/MECS normalization constant, indicating that both the target (H1846โ786) and another object, located within the PDS field of view, contribute to the PDS flux. This excess emission could be due to the near ($`z=0.0285`$, B $`=8.8`$, R $`=10.4`$) bright Seyfert 1 galaxy ESO 025โG002, which has a *ROSAT* flux only 5 times lower than the target source. Also this source belongs to the 2 Micron All Sky Survey eXtended Source Catalog (i.e. 2MASX J18544039โ7853544), with magnitudes J $`=11.4`$, H $`=10.6`$ and K $`=10.3`$, and shows weak radio emission ($``$10 mJy at 36 cm). The spectral analysis of the PDS data yields a contribution from this object of $``$60$`\%`$. The spectrum is steep ($`\mathrm{\Gamma }`$ $``$3) but still compatible with the canonical AGN value of 2 and the 20โ100 keV flux turns out to be of the order of $`2.0\times 10^{11}`$ erg cm<sup>-2</sup> s<sup>-1</sup> (see Table 6 for more details). Thus, ESO 025โG002 shows at high energies a flux of the same order of that of H1846โ786, while at soft energies is 5 time lower. This can be explained if the source is absorbed at soft energies; in fact, fixing the PDS power law index to a value more appropriate to an AGN and allowing for absorption, we find an upper limit to the column density of a few $`10^{22}`$ cm<sup>-2</sup>.
Finally, we can conclude that in this particular case both target and serendipitous source are likely to give roughly the same contribution to the PDS emission.
#### 4.2.6 The VW Cephei Field
VW Cephei, a W UMaโtype binary system, has been observed by *Beppo*SAX twice (May 1998 and October 1998). Following the investigations performed by Gondoin (2004), we fit separately the LECS and MECS data of both data sets with the mekal optically thin plasma model. In each observation, a threeโcomponent model, having different temperatures but the same metal abundance, gives the best representation of the data. The 2โ10 keV flux turns out to be $`1.6\times 10^{12}`$ erg cm<sup>-2</sup> s<sup>-1</sup> and $`0.8\times 10^{12}`$ erg cm<sup>-2</sup> s<sup>-1</sup> in the first and second pointings, respectively. Although the flux varies ($``$50$`\%`$) from the first to the second observation, the spectral shape shows, within the uncertainties, no evidence of variability; therefore, we summed the two data sets together. The threeโtemperature model applied to the combined data still provides a good fit ($`\chi ^2/\nu =37/65`$) and gives temperature values of 2.4, 0.69 and 0.13 keV; extrapolation of this model to the PDS band indicates that the VW Cephei contribution to the PDS data is negligible (see Table 5).
As a possible contaminating source, we find at around $`34^{}`$ from VW Cephei an object identified with 4C +74.26, a radioโloud active galaxy. This is a particularly interesting object because its *ASCA* Xโray spectrum (Brinkmann et al. (1998); Sambruna et al. (1999)) shows features that are typical of Seyfert galaxies more than of Broad Line Radio Galaxies: a warm absorber and a significant Compton reflection hump. It is also the only quasar in the collection of Sambruna et al. (1999) with a detectable Fe line. 4C +74.26 was targetted by *Beppo*SAX on May 1999. Analysis of these data by Hasenkopf et al. (2002) indicated as the bestโfit model a power law continuum modified by Compton reflection at high energies and by absorption at low energies. An iron line was also detected, but its energy could not be tightly constrained, falling between 6.4 and 6.9 keV. Our spectral analysis of the average PDS emission obtained by combining the two observations of VW Cephei yields a photon index $`\mathrm{\Gamma }`$ $``$2 and a 20โ100 keV flux, corrected for the offโaxis effect, of $`4.6\times 10^{11}`$ erg cm<sup>-2</sup> s<sup>-1</sup> (see Table 6). If we fit the PDS spectrum with the pexrav model, fixing the inclination angle to the value indicated by Hasenkopf et al. (2002), we find a satisfactory fit ($`\chi ^2/\nu =1.2/2`$), but the reflection coefficient of $`R`$ $``$1.3 and the cutโoff energy $`E_C>148`$ keV are not well constrained, although in agreement with what reported by Hasenkopf et al. (2002). In order to check the selfโconsistency of our spectral analysis, we reanalysed the PDS data of the May 1999 observation with a simple power law and obtained spectral parameters for this epoch (a photon index $`\mathrm{\Gamma }`$ $``$2 and a 20โ100 keV flux of $`4.0\times 10^{11}`$ erg cm<sup>-2</sup> s<sup>-1</sup>) consistent with those obtained during the VW Cephei pointings.
Overall, the data indicate that 4C +74.26 is the most likely contaminating source present in the PDS field of view of VW Cephei and further suggest no variability in the high energy flux.
#### 4.2.7 The NGC 7331 Field
A good fit to the MECS and LECS data of the LINER NGC 7331 is provided by a single power law having a photon index $`\mathrm{\Gamma }=1.95_{0.18}^{+0.12}`$ and 2โ10 keV flux of $`5.1\times 10^{13}`$ erg cm<sup>-2</sup> s<sup>-1</sup>. This result confirms the nonโthermal Xโray spectrum of NGC 7331 reported by Stockdale et al. (1998) on the basis of both *ROSAT* and radio observations. The above model however gives a small contribution (see Table 5) to the PDS emission when extrapolated to high energies.
Searching for potential high energy emitters in the HEASARC archives, we find a source located at about $`30^{}`$ from the centre of the MECS pointing and identified with the Seyfert 2 galaxy NGC 7319, which belongs to the Stephanโs Quintet, a compact group of galaxies. Recent *Chandra* data of NGC 7319 (Trinchieri et al. (2003)) describe the source emission as due to the superposition of a strong and heavy absorbed nuclear source plus diffuse softer emission. Modeling the nuclear source with a combination of absorbed and unabsorbed power law components having the same photon index, plus a narrow 6.4 keV emission line, the above authors found a photon index $`\mathrm{\Gamma }`$ $``$1.7, a column density $`N_\mathrm{H}=4\times 10^{23}`$ cm<sup>-2</sup>, an $`EW`$ $``$110 eV and a 2โ10 keV flux of $`8.2\times 10^{12}`$ erg cm<sup>-2</sup> s<sup>-1</sup>. The spectral fit was similar to that derived from a previous *ASCA* observation (Awaki et al. (1997)), except for the smaller $`EW`$ and higher 2โ10 keV flux. The Xโray spectral parameters as well as the $`F_\mathrm{X}`$/$`F_{\mathrm{OIII}}`$ ratio (Bassani et al. (1999)) indicate that the source is likely to be Compton thin. NGC 7319 ($`z=0.0225`$, B $`=14`$, R $`=12.7`$) belongs to the 2 Micron All Sky Survey eXtended Source Catalog (2MASX J22360355+3358327), with magnitudes J $`=11.1`$, H $`=10.3`$ and K $`=10.1`$, and is also detected in radio ($``$7 mJy at 6 cm).
The PDS data are well described by a power law with a photon index $`\mathrm{\Gamma }`$ $``$2.1 (see Table 6) and a 20โ100 keV flux of $`2.4\times 10^{11}`$ erg cm<sup>-2</sup> s<sup>-1</sup> at the source. By extrapolating the PDS spectrum to the 2โ10 keV band, we find a flux of $`1.7\times 10^{11}`$ erg cm<sup>-2</sup> s<sup>-1</sup>, which is higher by a factor of $``$9 and $``$2 with respect to the *ASCA* and *Chandra* measurements, respectively. However, if we fit the PDS spectrum with an absorbed power law, fixing the column density $`N_\mathrm{H}`$ and the photon index $`\mathrm{\Gamma }`$ to the bestโfit values found by Trinchieri et al. (2003) and extrapolate to the 2โ10 keV energy range, we find a 2โ10 keV flux of $``$$`7\times 10^{12}`$ erg cm<sup>-2</sup> s<sup>-1</sup> slightly lower than the *Chandra* value, but still compatible with it. This result strongly indicates that NGC 7319 is a good candidate to explain the PDS emission in the NGC 7331 observation.
#### 4.2.8 The NGC 7552 Field
In this field, although two extra sources are present within the MECS image (see Fig. 9), neither seems to be responsible for the high energy emission seen in the PDS.
The LECS and MECS data of the starburst galaxy NGC 7552 are well described by a twoโcomponent model, consisting of a thermal part (RaymondโSmith model in XSPEC) with $`kT=0.92_{0.34}^{+0.28}`$ keV and a power law with a steep spectrum ($`\mathrm{\Gamma }=2.74_{0.75}^{+0.91}`$). The 2โ10 keV flux is $`1.6\times 10^{13}`$ erg cm<sup>-2</sup> s<sup>-1</sup>, and the extrapolation of this model to high energies shows an almost null ($`<`$ 1$`\%`$) contribution to the PDS emission.
As displayed in Fig. 9, within the MECS field of view we find two other sources. The first, located at $`\alpha (2000)=23^h17^m28^s.8`$ and $`\delta (2000)=42^{}47^{}36^{\prime \prime }.2`$ is associated to the Seyfert 1 galaxy FRL 1041, while the second at $`\alpha (2000)=23^h14^m00^s.7`$ and $`\delta (2000)=42^{}43^{}17^{\prime \prime }.3`$ is identified with the cluster of galaxies Sersic 159โ03. Fitting the MECS data of FRL 1041 with a simple power law, provides a satisfactory fit ($`\chi ^2/\nu =14.9/15`$), a photon index $`\mathrm{\Gamma }=1.63_{0.60}^{+0.67}`$, compatible with values found for Seyfert galaxies, and a 2โ10 keV flux of $`5.0\times 10^{13}`$ erg cm<sup>-2</sup> s<sup>-1</sup>. The extrapolation of this model to high energies gives a low (see Table 5) contribution and a joint fit of the MECS/PDS data gives a crossโcalibration constant (37โ82) well outside the nominal range. Concerning the galaxy cluster, Sersic 159โ03, we find that the bestโfit model ($`\chi ^2/\nu =26.4/32`$), for the MECS data only, is a thermal bremsstrahlung having a temperature value ($`kT=2.54_{0.22}^{+0.24}`$ keV), compatible with that found by *XMMโNewton* observations (Kaastra et al. (2001)), plus a narrow gaussian line centered at $`6.47_{0.21}^{+0.24}`$ keV with an $`EW`$ of $`552_{327}^{+229}`$ eV.
The 2โ10 keV flux is $`7.0\times 10^{12}`$ erg cm<sup>-2</sup> s<sup>-1</sup>. The quality of the fit is rather worse when we fit simultaneously the MECS and PDS data; in fact, we cannot put any constraint on the spectral parameters and, in addition, the PDS/MECS crossโcalibration constant varies within the interval 91โ303. The above results indicate that neither source is likely to be responsible for the high energy emission (see also Table 5), even in the case of the cluster which is strongly detected in the MECS instrument ($``$50$`\sigma `$ level), but has an extremely soft spectrum.
In view of these findings, we performed a further search in order to understand if there was any other nearby object which could produce the PDS result. Just outside the MECS field of view we find NGC 7582, a Seyfert 2 galaxy, belonging to the Piccinotti sample, i.e. a very bright source in hard Xโrays (Piccinotti et al. (1982); Malizia et al. (1999)).
NGC 7582 has already been observed by *Beppo*SAX in November 1998 (Turner et al. (2000)). In that case the broad spectrum (2โ100 keV) was described by a power law of photon index $`\mathrm{\Gamma }=1.95_{0.18}^{+0.09}`$ (steeper than the index found during a previous *ASCA* observation) transmitted through a dual absorber with column densities $`N_\mathrm{H}`$ $``$$`1.6\times 10^{24}`$ cm<sup>-2</sup> (covering 60$`\%`$ of the nucleus) and $``$$`1.44\times 10^{23}`$ cm<sup>-2</sup> (fully covering the source). The authors reported a 10โ100 keV flux of $``$1.2$`\times 10^{10}`$ erg cm<sup>-2</sup> s<sup>-1</sup> and the presence of variability on timeโscales down to a few thousand of seconds. In addition, Turner et al. (2000) found that NGC 7582 was significantly brighter than the average level sampled by *OSSE* ($``$$`4\times 10^{11}`$ erg cm<sup>-2</sup> s<sup>-1</sup> in the 50โ150 keV band) and *BATSE* ($``$$`8.9\times 10^{11}`$ erg cm<sup>-2</sup> s<sup>-1</sup> in the 20โ100 keV band) (Johnson et al. (1997); Malizia et al. (1999)), both estimated using a power law model.
By subtracting from the PDS the contribution due to each of the three sources present in the MECS image of Fig. 9 and fitting the remaining emission with a simple power law, we find a photon index of 1.72 (see Table 6) and a 20โ100 keV flux (50โ150 keV) at the source of $`8.13\times 10^{11}`$ erg cm<sup>-2</sup> s<sup>-1</sup> ($``$$`5\times 10^{11}`$ erg cm<sup>-2</sup> s<sup>-1</sup>), which is in perfect agreement with the values reported by *BATSE* and *OSSE*. If we fit our data, taking into account the double absorber proposed by Turner et al. (2000), we find a photon index $`\mathrm{\Gamma }`$ $``$1.9 and a 10โ100 keV flux of $`2.3\times 10^{10}`$ erg cm<sup>-2</sup> s<sup>-1</sup>, a factor of 2 higher than the value reported for the 1998 observation by Turner et al. (2000).
To conclude, NGC 7582 is very likely the serendipitous source responsible for the emission detected above 10 keV in the pointing of NGC 7552. Furthermore, we provide evidence for variability ($``$$`90\%`$) in the flux between the two *Beppo*SAX measurements taken one year apart.
## 5 Conclusions
In this work we have shown how a careful search in the *Beppo*SAX public archive can provide evidence for new hard Xโray emitting sources, most likely associated with active galaxies, and/or confirm the high energy emission of known objects. In particular, we report the detection of six new hard Xโray emitters (two type 1 Seyfert galaxies, two type 2, one quasar and a galaxy not yet classified as active); for two of these sources broad band spectra are presented, while in four cases only results above 10 keV are reported.
For the remaining 6 objects in our sample, emission above 10 keV was known before this work, mainly from previous dedicated *Beppo*SAX pointings. In the case of NGC 3227, an object known to emit at high energies from past *OSSE* and *BATSE* measurements, a spectrum is published here for the first time. Comparison of the present data with previous observations either with *Beppo*SAX and/or with other high energy missions indicates that flux variability is present in three (MKN 766, NGC 7582 and NGC 3227) possibly four (NGC 1275) objects.
As the AGN reported in this paper are โlooselyโ representative of the population of extragalactic objects with emission above 10 keV, it is worth examining their overall characteristics. The first consideration to make is that most of our objects (6 out of 10 with optical classification) are broad line emitting or unabsorbed AGN, contrary to the expectation of finding more absorbed than unabsorbed objects. The other four (possibly five if we include PKS 0101โ649) show intrinsic absorption which, however, is compatible with a Compton thin nature ($`N_H<10^{24}`$ cm<sup>-2</sup>). Only in the case of PKS 0101โ649 could the source be Compton thick. This is probably due to the limited capability of current hard Xโray detectors which allow just the brightest and nearest extragalactic objects to be found. In fact, all but two (Perseus and 4C +74.26) or even one (if NGC 1275 is responsible for the emission detected by the PDS) of our AGNs are nearby ($`z<0.1`$) and belong to the 2 Micron All Sky Survey eXtended Source Catalog (i.e. the surrounding galaxy is visible in the nearโinfrared). At least eight (possibly nine if we also consider NGC 1275) of our objects are radio loud, according to Terashima $`\&`$ Wilson (2003) definition.
Overall, we can conclude that our objects are bright, nearby, with a detectable emission in all wavebands and only a fraction ($`3040\%`$) have a column density in excess of $`10^{22}`$ cm<sup>-2</sup>.
As shown here it is also possible to encounter offset fields contaminated by the presence of serendipitous sources. We are now performing a systematic analysis of all PDS data present in the archive (above $`15^{}`$ in galactic latitudes) evaluating the relative offset fields; the aim of this work is to provide a list of positive detections in these background measurements. We are confident that this analysis will provide new high energy sources as highlighted here in three cases.
###### Acknowledgements.
This research has made use of SAXDAS linearized and cleaned event files of LECS and MECS produced at ASI Science Data Center; and of the High Energy Science Archive Research Center (HEASARC), provided by NASAโs Goddard Space Flight Center. We acknowledge the financial support of the Italian Space agency (ASI) through contract ASI/CNR I/R/073/02. We thank J.B. Stephen for a careful reading of the manuscript. We also thank the referee for useful comments and suggestions.
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# Avoided crossings in three coupled oscillators as a model system of acoustic bubbles
## I Introduction
Avoided crossings ref1 have been observed theoretically and experimentally in a large variety of physical systems involving eigenvalues (e.g., natural frequencies, eigenenergies) ref2 ; ref3 ; ref4 ; ref5 ; ref6 ; ref7 ; ref8 ; ref9 ; ref10 ; ref11 ; ref12 ; ref13 ; ref14 ; ref15 , and they have attracted much attention even in recent years because of their rich physics and practical importance in, for example, mechanical engineering ref11 ; ref12 ; ref13 ; ref14 and quantum physics ref2 ; ref3 ; ref4 ; ref5 . In the avoided crossing regions, eigenvalues of the system first approach each other as a system parameter is varied but then veer abruptly from each other without crossing. In those regions a drastic change of some characteristic of the system occurs along the eigenvalue loci. In Ref. ref12 , for example, Pierre illustrated that the mode shapes of a disordered chain of coupled pendulums change in the regions where avoided crossings of the eigenfrequencies of the system take place. In that study, disorders in the lengths of the pendulums were used as the system parameters. Also, in Ref. ref2 , Walkup et al. studied in detail avoided crossings observed in the energy levels of diamagnetic hydrogen as functions of the magnetic field strength or the angular momentum, which lead to the diabatic exchange of the states of the wave functions. A study of Bose-Einstein condensation (BEC) ref5 showed that in order to trap molecules created in an atomic BEC through a Feshbach resonance, an avoided crossing of two bound states of the molecules must be exploited, through which the vibrational quantum number and size of the trapped molecules change.
In the present paper, we show theoretically that avoided crossings can be observed in acoustically coupled bubbles, which has to the authorsโ knowledge not been stated in the literature. Furthermore, based on analyses of transition frequencies ref16 ; ref17 , we propose a way to detect a state exchange occurring in the avoided crossing region. The theoretical model used in this study, reviewed in Sec. II, is a forced coupled oscillator model that describes acoustic coupling of pulsating bubbles. Using the model, we show in Sec. III that if at least three bubbles exist, it is possible that the resonance frequencies of the bubbles exhibit an avoided crossing when they are plotted as functions of the separation distances between the bubbles. As has been demonstrated (e.g., Refs. ref16 ; ref18 ), in double-bubble systems, neither crossings nor avoided crossings of the resonance frequencies as functions of the separation distance are observed, since the higher of the two resonance frequencies of the systems increases and the lower one decreases as the bubbles approach each other. However, as shown in the present paper, by introducing one more bubble whose monopole (i.e., decoupled) resonance frequency crosses with one of the resonance frequencies of a double-bubble system, one can observe the avoided crossing of the resonance frequencies when all three bubbles are coupled.
In Sec. IV, we examine the phase properties of the three coupled bubbles to show that a state exchange actually occurred between the bubbles in the avoided crossing region. In this effort, the notion of a transition frequency plays an important role. The transition frequencies introduced in Refs. ref16 ; ref17 are characteristic frequencies of acoustically coupled bubbles, around which the oscillation phase of bubbles inverts, e.g., from in-phase to out-of-phase with the driving sound. It was proved in Ref. ref17 that a bubble in a $`N`$-bubble system has up to $`2N1`$ transition frequencies, only $`N`$ ones of which correspond to the resonance frequencies of the system. That is, observing the transition frequencies allows us to obtain richer insight into the phase properties than that obtained by only observing the resonance frequencies. This notion has already been exploited as a powerful tool to understand the sign reversal of the secondary Bjerknes force ref18 ; ref19 in which the oscillation phases play a crucial role. Using this notion and observing directly the oscillation phases, we show that the coupled bubbles exchange their oscillation states through the avoided crossing and the state exchange takes place at the separation distances where an avoided crossing resonance frequency crosses with a transition frequency that is not a resonance frequency. The present findings appear to reveal a taste of bubblesโ hidden complexity.
Section V summarizes this paper, and the Appendixes present additional remarks.
## II Coupled oscillator model, resonance frequency, and transition frequency
The theoretical model used in the present study is a forced oscillator model in which $`N`$ harmonic oscillators are coupled (ref16 ; ref17 and references therein):
$`\ddot{e}_i+\omega _{i0}^2e_i+\delta _i\dot{e}_i`$ $`=`$ $`{\displaystyle \frac{p_{\mathrm{ex}}}{\rho R_{i0}}}{\displaystyle \frac{1}{R_{i0}}}{\displaystyle \underset{j=1,ji}{\overset{N}{}}}{\displaystyle \frac{R_{j0}^2}{D_{ij}}}\ddot{e}_j`$ (1)
$`\text{for}i`$ $`=`$ $`1,2,\mathrm{},N,`$
where $`N`$ corresponds to the number of bubbles, $`R_{i0}`$ is the equilibrium radius of bubble $`i`$, $`e_i`$ is the deviation of radius assumed as $`\left|e_i\right|R_{i0}`$, $`\omega _{i0}`$ is the monopole (angular) resonance frequency of bubble $`i`$, defined as
$$\omega _{i0}=\sqrt{\frac{3\kappa _iP_0+(3\kappa _i1)2\sigma /R_{i0}}{\rho R_{i0}^2}},$$
(2)
$`\delta _i`$ is the damping factor, the overdots denote the time derivation, $`p_{\mathrm{ex}}`$ is the pressure of the external sound, $`\rho `$ is the density of the surrounding liquid, $`D_{ij}`$ ($`=D_{ji})`$ is the separation distance between the centers of bubbles $`i`$ and $`j`$, $`\kappa _i`$ is the polytropic exponent of the gas inside the bubbles, $`P_0`$ is the static pressure, and $`\sigma `$ is the surface tension. In this linear model, the following assumptions are made: the surrounding liquid is incompressible, the sound amplitude is sufficiently low, the separation distances are much larger than the bubblesโ radii, and the shape deformation of the bubbles is negligible. The last term of Eq. (1), representing the pressures of the sounds that the neighboring bubbles emit, describes the acoustic coupling between the bubbles. As in the double-bubble case ref20 , this model may be assumed to be of third order with respect to the inverse of the separation distances (i.e., the truncated terms are of fourth or higher order); see Appendix A.
Using this model with $`N=3`$, a matrix equation for determining the amplitudes and phases of the radial oscillations is derived. Assuming $`p_{\mathrm{ex}}=P_a\mathrm{exp}(\mathrm{i}\omega t)`$ and $`e_i=\beta _i\mathrm{exp}(\mathrm{i}\omega t)`$ with $`P_a`$ being a positive constant, $`\omega `$ being the driving (angular) frequency, and $`\beta _i`$ being a complex amplitude, we have
$$๐\left(\begin{array}{cccccccccccccccccccc}\beta _1& & & & & & & & & & & & & & & & & & & \\ \beta _2& & & & & & & & & & & & & & & & & & & \\ \beta _3& & & & & & & & & & & & & & & & & & & \end{array}\right)=\frac{P_a}{\rho }๐,$$
(3)
where $`๐`$ is a $`3\times 3`$ matrix whose elements, $`a_{i,j}`$ ($`i,j=1,2,3)`$, are defined as
$$a_{i,j}\{\begin{array}{cccccccccccccccccccc}R_{i0}[(X\omega _{i0}^2)\mathrm{i}\omega \delta _i]& \text{for}i=j,& & & & & & & & & & & & & & & & & & \\ \frac{R_{j0}^2}{D_{ij}}X& \text{otherwise,}& & & & & & & & & & & & & & & & & & \end{array}$$
(4)
with
$$X\omega ^2,$$
(5)
and $`๐=(1,1,1)^T`$. We should note here that essentially the same matrix equations can be found in previous papers (e.g. ref21 ; ref22 ; ref17 ). The solution of Eq. (3) is represented as
$`\left(\begin{array}{cccccccccccccccccccc}\beta _1& & & & & & & & & & & & & & & & & & & \\ \beta _2& & & & & & & & & & & & & & & & & & & \\ \beta _3& & & & & & & & & & & & & & & & & & & \end{array}\right)`$ $`=`$ $`{\displaystyle \frac{P_a}{\rho }}๐^1๐`$
$`=`$ $`{\displaystyle \frac{P_a}{\rho }}{\displaystyle \frac{|๐|^{}๐๐}{|๐|^{}|๐|}},`$
where $`|๐|`$ and $`๐`$ are the determinant and the cofactor matrix of $`๐`$, respectively, and $`|๐|^{}`$ is the complex conjugate of $`|๐|`$. We used here an expression in which the denominator is real.
The eigenfrequencies of the system are determined by
$$|๐|=0,$$
(10)
which is a cubic equation in terms of $`X`$. For $`\delta _i0`$, the roots of this equation are equivalent to the resonance frequencies of the system. The transition frequencies of bubble $`i`$, defined as the driving frequencies at which the phase difference between bubble $`i`$ and the driving sound is $`\pi /2`$ (or $`3\pi /2)`$ ref16 ; ref17 ; ref18 , are determined by
$$\mathrm{Re}(\tau _i)=0,$$
(11)
where
$$\left(\begin{array}{cccccccccccccccccccc}\tau _1& & & & & & & & & & & & & & & & & & & \\ \tau _2& & & & & & & & & & & & & & & & & & & \\ \tau _3& & & & & & & & & & & & & & & & & & & \end{array}\right)|๐|^{}๐๐.$$
(12)
(See Appendix B for the concrete forms of $`\left|๐\right|`$ and $`๐๐`$.) From the mathematical proof given in Ref. ref17 , one knows that Eq. (11) is a fifth-order polynomial in terms of $`X`$, meaning that the bubbles may have up to five transition frequencies.
The phase delay of bubble $`i`$, denoted by $`\varphi _i`$, measured from the phase of the driving sound is determined using the $`\mathrm{atan2}(y,x)`$ function in the C language, which returns $`\mathrm{tan}^1(y/x)[\pi ,\pi ]`$, as
$$\varphi _i=\{\begin{array}{cccccccccccccccccccc}\psi _i\hfill & \text{if }\psi _i0,\hfill & & & & & & & & & & & & & & & & & & \\ \psi _i+2\pi \hfill & \text{otherwise}\hfill & & & & & & & & & & & & & & & & & & \end{array}$$
with
$$\psi _i=\mathrm{atan2}(\mathrm{Im}(\tau _i),\mathrm{Re}(\tau _i)).$$
The next section shows that in certain cases an avoided crossing is observed in the solution of Eq. (10). In the discussion, to obtain real eigenfrequencies that correspond to the resonance frequencies of the triple-bubble system for weak damping, we for the moment assume $`\delta _i0`$ (but $`\delta _i0)`$. Under this assumption, one obtains
$$\mathrm{Im}(|๐|)\mathrm{๐},$$
(13)
$$|๐||๐|^{},$$
(14)
and
$$\tau _i\mathrm{Re}(\tau _i).$$
(15)
Influences of the damping effect on the phase properties are briefly discussed in Sec. IV.
## III Avoided crossings of resonance frequencies
To begin with, a double-bubble system is briefly reconsidered to confirm that no avoided crossings are observed in the resonance frequencies of the system as functions of the separation distance. The solid lines in Fig. 1 indicate the resonance frequencies of two coupled bubbles (bubbles 1 and 2) of $`(R_{10},R_{20})=(50\mu \mathrm{m},51\mu \mathrm{m})`$ as functions of $`l_{12}=D_{12}/(R_{10}+R_{20})`$. The other parameters are set to $`\rho =1000`$ kg/m<sup>3</sup>, $`\kappa _i=1.4(i=1,2,3)`$, $`P_0=1`$ atm, and $`\sigma =0.0728`$ N/m. As has been proved theoretically ref23 ; ref24 ; ref16 ; ref20 , two resonance (or natural) frequencies appear in this system, each of which, for $`D_{12}\mathrm{}`$, converges to the monopole resonance frequency of a bubble. The higher resonance frequency increases and the lower decreases as the separation distance decreases. It is therefore obvious that avoided crossings cannot occur.
Here we introduce one more bubble into the system. The dashed line displayed in Fig. 1 denotes the monopole resonance frequency of the introduced bubble, bubble 3, whose radius $`R_{30}=51.5\mu \mathrm{m}`$. Note that this resonance frequency crosses with a resonance frequency of the double-bubble system. This crossing, as shown immediately, triggers an avoided crossing when the third bubble is coupled with the double-bubble system.
Figure 2 shows the resonance frequencies in the case where all three bubbles are coupled. The separation distances are set to $`D_{12}=l_{12}(R_{10}+R_{20})`$, $`D_{23}=l_{23}(R_{20}+R_{30})`$, and $`D_{31}=D_{12}+D_{23}`$; that is, the bubbles are arranged in line (see Fig. 3(a)). Here the nondimensional quantities $`l_{12}`$ and $`l_{23}`$ are used as the system parameters. Figures 2(a), 2(b), and 2(c) show the results for $`l_{23}=100`$ ref25 , $`50`$, and $`20`$, respectively. In the figures, an avoided crossing is clearly seen that takes place around the point at which the two decoupled resonance frequencies cross. The line of the resonance frequency originating with bubble 3 is divided into two parts, and each of them connects smoothly, like blending, with the curve of a resonance frequency of the double-bubble system, also divided into two parts. As bubble 3 comes closer to the others, the avoided crossing becomes broader and the origin of each resonance frequency becomes increasingly unclear.
An avoided crossing is also observed when bubble 3 is smaller than the others. Figure 4 shows the resonance frequencies when $`R_{30}=49.5\mu \mathrm{m}`$. Here the bubbles are arranged as illustrated in Fig. 3(b). If bubble 3 is so large or so small that its monopole resonance frequency does not cross with a resonance frequency of the double-bubble system, no distinct avoided crossing is observed, though this situation is not shown here.
## IV State exchange in the avoided crossing region
To manifest a state exchange like that which the coupled bubbles experience through the avoided crossing, we examined the oscillation phases of the bubbles. In bubble dynamics, the phase of radial oscillation plays important roles in many situations, including acoustic levitation ref26 ; ref27 ; ref28 , bubble-bubble interaction ref18 , and multibubble sonoluminescence ref29 , and hence an accurate understanding of it is crucial. In fact, by carefully examining the oscillation phases of two coupled bubbles for weak driving, we have recently succeeded in presenting a novel interpretation, which may be more accurate than previous ones, of the sign reversal of the secondary Bjerknes force ref18 ; ref19 , a paradoxical phenomenon that is considered to be the cause of the stable structure formation of bubbles in a weak acoustic field ref24 ; ref30 . In that discussion, it was suggested that the transition frequencies seem to be essential components for gaining an accurate understanding of the phenomenon, since the sign reversal takes place at the transition frequencies that cannot be obtained by resonance-frequency analysis. In the present paper, we show by examining the oscillation phases that the bubbles exchange their oscillation states through the avoided crossing. As shown later, the point at which the state exchange occurs can be clearly detected by observing the transition frequencies.
Figure 5 shows the transition frequencies for $`(R_{10},R_{20},R_{30})=(50\mu \mathrm{m},51\mu \mathrm{m},51.5\mu \mathrm{m})`$ with $`l_{23}=20`$. The thick lines denote the transition frequencies that correspond to the resonance frequencies already shown in Fig. 2(c). As expected from the mathematical proof presented in ref17 , the bubbles have up to five transition frequencies, all of which invert the oscillation phase of the corresponding bubble. It is worth noting that in each panel of Fig. 5 the second-highest resonance frequency (denoted below by $`\omega _{2\mathrm{n}\mathrm{d}})`$ crosses once with a transition frequency in the avoided crossing region. Such crossings have not been found in double-bubble systems ref16 ; ref18 . In the following discussion, we focus our attention on the phase properties of the bubbles in this region to elucidate what happens around the intersecting points.
The phase delays $`\varphi _i`$ for different $`l_{12}`$ as functions of $`\omega /\omega _{10}`$ are shown by the solid lines in Fig. 6. In the computation of $`\varphi _i`$, we used very small but nonzero $`\delta _i`$ to obtain continuous results. Figures 6(a, b) and 6(c, d), respectively, show $`\varphi _i`$ for $`l_{12}`$ smaller and larger than the intersecting point $`l_{12}=l_{\mathrm{int}}`$ ($``$8.89). Here, we only displayed $`\varphi _i`$ in the frequency range around the two avoided crossing resonance frequencies. The vertical dotted lines indicate the two lowest resonance frequencies. As in double-bubble cases ref18 ; ref19 , at the resonance frequencies the phase delays of all bubbles shift simultaneously by $`+\pi `$, whereas at the remaining transition frequencies only one phase delay shifts by $`\pi `$.
The $`\varphi _i`$-curves, as can be clearly seen in the figures, have different convexities on different sides of the intersecting point. For $`l_{12}`$ smaller than $`l_{\mathrm{int}}`$, at $`\omega _{2\mathrm{n}\mathrm{d}}`$, $`\varphi _1`$ and $`\varphi _2`$ shift from $`\pi `$ to $`2\pi `$ but $`\varphi _3`$ shifts from $`0`$ to $`\pi `$ as $`\omega `$ increases. For $`l_{12}`$ larger than $`l_{\mathrm{int}}`$, on the other hand, an opposite tendency is seen; $`\varphi _1`$ and $`\varphi _2`$ shift from $`0`$ to $`\pi `$ but $`\varphi _3`$ shifts from $`\pi `$ to $`2\pi `$. That is, a kind of state exchange takes place between bubble 3 and the other two bubbles at the intersecting point.
Regarding the relationship between the state exchange and the phase properties, in the frequency range around $`\omega _{2\mathrm{n}\mathrm{d}}`$, bubble 3 oscillates out-of-phase with the other bubbles regardless of whether $`l_{12}<l_{\mathrm{int}}`$ or $`l_{12}>l_{\mathrm{int}}`$, although the individual phase delays experience rapid shifts at $`\omega _{2\mathrm{n}\mathrm{d}}`$. This means that the state exchange cannot be perceived accurately by observing whether the bubbles oscillate in-phase or out-of-phase with each other or by observing the sign of the secondary Bjerknes force, which is determined by the cosine of the phase difference between two bubbles ref31 ; ref32 . Just the individual phase delays (or transition frequencies) should be examined.
In the $`\varphi _i`$-curves, we can find several similarities with double-bubble cases. Bubbles 1 and 2, or bubble 3, have a phase delay greater than $`\pi `$ in the frequency range from $`\omega _{2\mathrm{n}\mathrm{d}}`$ to a certain higher frequency (equal to the next-higher transition frequency of the corresponding bubble). A similar observation can be found for double-bubble systems ref18 ; ref19 . In Ref. ref18 we discovered and elucidated that such a large phase delay can appear when two bubbles interact with each other through sound. In the double-bubble case, the larger one of the two bubbles has a phase delay greater than $`\pi `$ in the frequency range between the higher of two resonance frequencies and the highest of the transition frequencies of the bubble. We can, for a wider frequency range, also find a similarity between the double- and triple-bubble cases. In the frequency range $`\omega /\omega _{10}<0.995`$, the profiles of $`\varphi _1`$ and $`\varphi _2`$ for $`l_{12}<l_{\mathrm{int}}`$ and that of $`\varphi _3`$ for $`l_{12}>l_{\mathrm{int}}`$ are very similar to the profile of the phase delay of the larger bubble in a double-bubble system; those phase delays first exhibit two sharp rises and then one sharp fall as $`\omega `$ increases. Also, the profiles of the remaining phase delays are very similar to that of the phase delay of the smaller bubble in a double-bubble system, exhibiting one sharp rise, one sharp fall, and then one sharp rise. This seems to indicate that in the frequency range considered, for $`l_{12}<l_{\mathrm{int}}`$ bubbles 1 and 2 act as โlarger bubblesโ while bubble 3 acts as a โsmaller bubble,โ but for $`l_{12}>l_{\mathrm{int}}`$ each bubble acts in the opposite way; that is, the physical roles that the bubbles play are exchanged through the avoided crossing. This observation could also be interpreted as a result of the change of a physical meaning of $`\omega _{2\mathrm{n}\mathrm{d}}`$. As illustrated in Fig. 2, $`\omega _{2\mathrm{n}\mathrm{d}}`$ is a hybrid of two resonance frequencies having different origins. We assume here that the origin, or the principal origin, of each avoided crossing resonance frequency is switched at $`l_{\mathrm{int}}`$. This assumption allows us to consider that $`\omega _{2\mathrm{n}\mathrm{d}}`$ for $`l_{12}<l_{\mathrm{int}}`$, for example, is the resonance frequency whose principal origin is bubble 3. This suggestion is consistent, not only with the observation for large $`l_{23}`$ where the origin of each resonance frequency is relatively clear, but also with the above speculation that bubbles 3 acts as a โsmaller bubbleโ for $`l_{12}<l_{\mathrm{int}}`$, because $`\omega _{2\mathrm{n}\mathrm{d}}`$ is higher than the lowest resonance frequency that is one of the two avoided crossing resonance frequencies. The observation for $`l_{12}>l_{\mathrm{int}}`$ can be interpreted in a similar manner.
Lastly, we briefly examine how the damping affects the state exchange. For the damping coefficient, we use the value for viscous damping,
$$\delta _i=\frac{4\mu }{\rho R_{i0}^2}$$
(16)
with viscosity $`\mu =1.002\times 10^3`$ kg/(m s). The dashed curves in Fig. 6 show the phase delays in the damped case. The viscous effect smoothes the phase profiles, but the convexity of the curves is not altered from that for $`\delta _i0`$, as in the double-bubble cases ref18 ; ref19 . The state exchange is clearly detected even in the present case. The qualitative tendencies of the phase delays are not changed by the viscous damping.
## V Conclusion
We have shown theoretically that avoided crossings can be observed in the resonance frequencies of acoustically coupled gas bubbles plotted as functions of the separation distances. A state exchange taking place between the bubbles in the avoided crossing region has been clearly exhibited by examining the oscillation phases and transition frequencies of the coupled bubbles. We have clarified that the state exchange is perceived by observing the individual oscillation phases of the bubbles, not by observing whether the bubbles oscillate in-phase or out-of-phase with each other. Since the individual phase (or more properly, the phase difference between a bubble and the external sound) determines the sign of the primarily Bjerknes force ref26 ; ref27 ; ref28 acting on the corresponding bubble, this state exchange should play a role in, e.g., acoustic levitation using the force. The results of this study suggest that the transition frequencies introduced in Ref. ref16 can be a useful tool for detecting the state exchange, which takes place at the separation distance where an avoided crossing resonance frequency crosses with a transition frequency that is not a resonance frequency. Though we only considered triple-bubble systems in a linear arrangement, extensions to systems containing a larger number of bubbles and in different arrangements may be straightforward. Also, nonlinear effects on the avoided crossings and oscillation phases could be examined using nonlinear models ref20 ; ref21 ; ref29 ; ref33 . As with other physical systems, the avoided crossings in acoustically coupled bubbles might be real.
###### Acknowledgements.
The author thanks Dr. Akemi Nishida for valuable comments. This work was supported by a Grant-in-Aid for Young Scientists (B) (17760151) from the Ministry of Education, Culture, Sports, Science, and Technology of Japan.
## Appendix A
High-order nonlinear models for $`N`$ pulsating bubbles in a liquid have been proposed, in which terms proportional to $`D_{ij}^k`$ ($`k2)`$ appear that involve the translational velocities of the bubbles ref21 ; ref33 . In Ref. ref33 , for example, Doinikov derived the following model equation for $`N`$ spherical bubbles:
$$R_i\ddot{R}_i+\frac{3}{2}\dot{R}_i^2\frac{P_i}{\rho }=\frac{\dot{๐ฉ}_i^2}{4}\underset{j=1,ji}{\overset{N}{}}\left\{\frac{R_j^2\ddot{R}_j+2R_j\dot{R}_j^2}{D_{ij}}+H_{ij}\right\},$$
(17)
$`{\displaystyle \frac{1}{3}}R_i\ddot{๐ฉ}_i+\dot{R}_i\dot{๐ฉ}_i={\displaystyle \frac{๐
_i}{2\pi \rho R_i^2}}`$ $`+`$ $`{\displaystyle \underset{j=1,ji}{\overset{N}{}}}\{{\displaystyle \frac{1}{D_{ij}^2}}(R_iR_j^2\ddot{R}_j+2R_iR_j\dot{R}_j^2+\dot{R}_i\dot{R}_jR_j^2)๐ญ_{ij}`$ (18)
$`{\displaystyle \frac{R_j^2}{2D_{ij}^3}}[R_iR_j\ddot{๐ฉ}_j+(\dot{R}_iR_j+5R_i\dot{R}_j)\dot{๐ฉ}_j]`$
$`+{\displaystyle \frac{3R_j^2}{2D_{ij}^3}}\{๐ญ_{ij}[R_iR_j\ddot{๐ฉ}_j+(\dot{R}_iR_j+5R_i\dot{R}_j)\dot{๐ฉ}_j]\}๐ญ_{ij}\},`$
with
$`H_{ij}`$ $``$ $`{\displaystyle \frac{R_j^2}{2D_{ij}^2}}(R_j\ddot{๐ฉ}_j+\dot{R}_j\dot{๐ฉ}_i+5\dot{R}_j\dot{๐ฉ}_j)๐ญ_{ij}`$ (19)
$`{\displaystyle \frac{R_j^3}{4D_{ij}^3}}\left[\dot{๐ฉ}_j(\dot{๐ฉ}_i+2\dot{๐ฉ}_j)3(\dot{๐ฉ}_j๐ญ_{ij})[๐ญ_{ij}(\dot{๐ฉ}_i+2\dot{๐ฉ}_j)]\right],`$
$$๐ญ_{ij}\frac{๐ฉ_j๐ฉ_i}{D_{ij}},$$
$$P_i\left(P_0+\frac{2\sigma }{R_{i0}}\right)\left(\frac{R_{i0}}{R_i}\right)^{3\gamma }\frac{2\sigma }{R_i}\frac{4\mu \dot{R}_i}{R_i}P_0p_{\mathrm{ex}},$$
(20)
where $`R_i`$ and $`๐ฉ_i`$ are the instantaneous radius and position vector, respectively, of bubble $`i`$, $`๐
_i`$ denotes external forces on bubble $`i`$, $`๐ญ_{ij}`$ is a unit vector, $`\gamma `$ is the specific heat ratio of the gas inside the bubbles, and $`\mu `$ is the viscosity. Here we showed only the incompressible version, though Doinikov also derived a model for bubbles in a compressible liquid. Equations (17) and (18) represent the volume oscillation of bubble $`i`$ and its translational motion, respectively. The linear coupled oscillator model used in the present study is recovered from Eq. (17) by truncating the high-order terms $`H_{ij}`$ and assuming weak driving and $`\gamma =\kappa _i`$.
Since the velocity field forming around a pulsating sphere is proportional to $`1/r^2`$, where $`r`$ is the distance measured from the center of the sphere, the truncated terms $`H_{ij}`$, which are composed of the translational velocities $`\dot{๐ฉ}_i`$, might be considered to be of fourth, or higher, order with respect to the inverse of the separation distances. This speculation is consistent with the suggestion by Harkin et al. for double-bubble systems ref20 .
Equation (17) further suggests that under the assumption of $`\dot{๐ฉ}_i\mathrm{๐}`$ one cannot construct a linear model that has higher-order accuracy than that of Eq. (1), since this assumption makes the high-order terms inaccurate.
## Appendix B
For the convenience of readers, we show the concrete forms of $`\left|๐\right|`$ and $`๐๐`$:
$`{\displaystyle \frac{|๐|}{R_{10}R_{20}R_{30}}}`$ $`=`$ $`L_1L_2L_3+s_{21}s_{32}s_{13}+s_{12}s_{23}s_{31}`$ (21)
$`L_1(M_2M_3+s_{23}s_{32})L_2(M_3M_1+s_{31}s_{13})L_3(M_1M_2+s_{12}s_{21})`$
$`+\mathrm{i}[M_1M_2M_3M_1(L_2L_3s_{23}s_{32})M_2(L_3L_1s_{31}s_{13})`$
$`M_3(L_1L_2s_{12}s_{21})],`$
$$๐๐=(c_1,c_2,c_3)^T,$$
$`{\displaystyle \frac{c_i}{R_{j0}R_{k0}}}`$ $`=`$ $`(L_js_{ij})(L_ks_{ik})+(s_{ij}s_{kj})(s_{jk}s_{ik})M_jM_k`$ (22)
$`+\mathrm{i}\left[M_j(s_{ik}L_k)+M_k(s_{ij}L_j)\right]`$
$`\text{for }(i,j,k)=(1,2,3),(2,3,1),\text{ or }(3,1,2),`$
where
$$L_iX\omega _{i0}^2,$$
$$M_i\omega \delta _i,$$
$$s_{ij}\frac{R_{j0}}{D_{ij}}X.$$
For $`\delta _i0`$, Eqs. (21) and (22) reduce, respectively, to
$`{\displaystyle \frac{|๐|}{R_{10}R_{20}R_{30}}}`$ $``$ $`L_1L_2L_3+s_{21}s_{32}s_{13}+s_{12}s_{23}s_{31}`$ (23)
$`L_1s_{23}s_{32}L_2s_{31}s_{13}L_3s_{12}s_{21},`$
$$\frac{c_i}{R_{j0}R_{k0}}(L_js_{ij})(L_ks_{ik})+(s_{ij}s_{kj})(s_{jk}s_{ik}).$$
(24)
Equation (23) and the real part of Eq. (21) are cubic functions and Eq. (24) and the real part of Eq. (22) are quadratic functions in terms of $`X`$. The imaginary parts of Eqs. (21) and (22) can be written in a form of $`\omega f(X)`$, where $`f`$ is quadratic in Eq. (21) and linear in Eq. (22). (As proved theoretically in Ref. ref17 , the imaginary parts are composed of terms of odd orders with respect to $`M`$ that are proportional to $`\omega X^n`$ with $`n`$ being a positive integer.)
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# Modeling a coronal loop heated by MHD-turbulence nanoflares
## 1 Introduction
Nanoflares (Parker 1988) are among the best candidates to explain the heating of the solar corona, and, in particular, of the coronal loops (e.g. Peres et al. 1993, Cargill 1993, Kopp & Poletto 1993, Shimizu 1995, Judge et al. 1998, Mitra-Kraev & Benz 2001, Katsukawa & Tsuneta 2001, Warren et al. 2002, 2003, Spadaro et al. 2003, Cargill & Klimchuk 1997, 2004, Mรผller et al. 2004, Testa et al. 2004).
Although the evidence of nanoflares appears to be well established, it is still unclear whether, and to what extent, they really can provide enough energy to heat the whole corona (e.g. Aschwanden 1999). More recently, models of nanoflares with a prescribed random time distribution of the pulses deposited at the footpoints of multi-stranded loops have been proposed (Warren et al. 2002, Warren et al. 2003), and have been shown to describe several observed features.
According to some models, nanoflares are the result of dissipation in an MHD turbulence, generated inside closed magnetic structures in the corona, and due to nonlinear interactions among fluctuations generated by photospheric motions. Possible evidence of turbulent motions has been detected from line broadenings in coronal loops (Saba & Strong 1991). Most of these models include direct numerical solution of MHD equations in two or three dimensions (Einaudi et al., 1996; Hendrix & Van Hoven, 1996; Dmitruk & Gรณmez, 1997; Dmitruk et al., 1998; Dmitruk & Gรณmez, 1999; Buchlin et al., 2003) using relatively low Reynolds/Lundquist numbers. Recently Nigro et al. 2004 (hereafter NMCV04) have related coronal nanoflares to intermittent dissipative events in the MHD turbulence produced in a coronal magnetic structure by footpoint motions. The injected energy is stored in the loop up to significant levels in the form of magnetic and velocity fluctuations and released intermittently through nonlinear interactions which process these fluctuations and generate cascades toward smaller scales where energy is dissipated. The derived probability distribution functions of the peak maximum power, peak duration time, energy dissipated in a burst and waiting time between bursts are in good agreement with those obtained from the analysis of coronal impulsive events (Datlowe et al. 1974, Lin et al. 1984, Dennis 1985, Crosby et al. 1993, Shimizu & Tsuneta 1997, Krucker & Benz 1998, Boffetta et al. 1999, Parnell & Jupp 2000, Aschwanden et al. 2000a, b). This heating model does not need any ad hoc hypothesis, once the loop length and the characteristic Alfven speed, i.e. the strength of the ambient magnetic field (if the density does not change much), are fixed.
In the present work we model the plasma confined in a coronal loop heated according to the events dissipation rate and distribution described in NMCV04. We will compute the evolution of the distributions of the density, temperature and velocity of the loop plasma by means of the time-dependent thermo-hydrodynamic Palermo-Harvard (Peres et al. 1982, Betta et al. 1997) loop model assuming the output of the hybrid shell model illustrated in NMCV04 as the basis of the heating function.
In Section 2 we describe the set up of the loop model with the MHD-turbulence dissipation rate as input heating; in Sec. 3 we show relevant results and discuss them in Sec. 4.
## 2 The loop model
Our purpose here is to model the evolution of the plasma confined in a coronal loop under the effect of the energy dissipation predicted in NMCV04. According to their settings, we model a magnetic loop, with a total length of 30,000 km. The plasma is described as a compressible fluid moving and transporting energy only along the magnetic field lines, i.e. along the loop itself. Thus, the magnetic field has only the role of confining the plasma. The loop model assumes constant loop cross-section.
We use the Palermo-Harvard code (Peres et al. 1982, Betta et al. 1997), a 1-D hydrodynamic code that consistently solves the time-dependent density, momentum and energy equations for the plasma confined by the magnetic field:
$$\frac{dn}{dt}=n\frac{v}{s},$$
(1)
$$nm_\mathrm{H}\frac{dv}{dt}=\frac{p}{s}+nm_\mathrm{H}g+\frac{}{s}(\mu \frac{v}{s}),$$
(2)
$$\frac{dฯต}{dt}+(p+ฯต)\frac{v}{s}=Hn^2\beta P(T)+\mu (\frac{v}{s})^2+\frac{}{s}(\kappa T^{5/2}\frac{T}{s}),$$
(3)
with $`p`$ and $`ฯต`$ defined by:
$$p=(1+\beta )nK_\mathrm{B}Tฯต=\frac{3}{2}p+n\beta \chi ,$$
(4)
where $`n`$ is the hydrogen number density, $`s`$ the spatial coordinate along the loop, $`v`$ the plasma velocity, $`m_\mathrm{H}`$ the mass of hydrogen atom, $`\mu `$ the effective plasma viscosity, $`P(T)`$ the radiative losses function per unit emission measure, $`\beta `$ the fractional ionization, i.e. $`n_\mathrm{e}/n_\mathrm{H}`$, $`\kappa `$ the thermal conductivity Spitzer (1962), $`K_\mathrm{B}`$ the Boltzmann constant, and $`\chi `$ the hydrogen ionization potential. $`H(s,t)`$ is a function of both space and time which describes the heat input in the loop. This function will be described in detail in Sec. 2.1. The numerical code uses an adaptive spatial grid to follow adequately the evolving profiles of the physical quantities, which can vary dramatically in the transition region and under the effect of the evolution. The loop is not symmetric, the apex is at half the numerical grid and there is a chromosphere on each side. The boundary conditions at the loop footpoints are the same as in Reale et al. (2000).
### 2.1 The heating function
The original version of the Palermo-Harvard hydrodynamic code includes a space- and time-dependent heating function, which describes the input of external energy triggering transient events (Peres et al. 1987). Several formulations are possible and the code can be easily adapted. For the present work, the heating function is given by the output dissipation rate of NMCV04 (in the form of a numerical table).
The model developed in NMCV04 has been derived within the framework of the Reduced Magnetohydrodynamics (RMHD) (Strauss 1976, Zank & Matthaeus, 1992), with the assumptions that: (i) the plasma is permeated by a strong uniform magnetic field $`๐_0`$ in the longitudinal direction; (ii) there is low thermal to magnetic pressure ratio $`\beta _P=8\pi p/B^21`$; (iii) the longitudinal scale $`ล_{||}`$ of transverse velocity $`v_{}`$ and magnetic field $`B_{}`$ fluctuations is much larger than the transverse scale $`l_{}`$; indeed, the MHD turbulence is anisotropic (e.g., Carbone & Veltri, 1990), the energy cascade being more efficient perpendicularly to $`๐_0`$. (iv) Small amplitude perturbations $`B_{}/B_0=v_{}/c_{A0}<l_{}/l_{||}1`$, where $`c_{A0}`$ is the background Alfvรฉn velocity, commonly assumed of the order of $`c_{A0}10^8`$ cm s<sup>-1</sup>, while the fluctuating velocity can be estimated using nonthermal broadening of coronal spectral lines: $`v_{}3\times 10^61.5\times 10^7`$ cm s<sup>-1</sup>. Under the above assumptions the set of the RMHD equations can be derived; they describe the evolution of magnetic and velocity fluctuations in terms of two distinct effects: (a) wave propagation in the longitudinal direction, at the Alfvรฉn velocity; (b) nonlinear couplings, which generate a turbulent cascade perpendicularly to $`๐_0`$. The model proposed by NMCV04 (hybrid shell model) includes both these dynamical mechanisms, but nonlinear effects are described in a simplified way by using a shell technique (Boffetta et al., 1999): a Fourier expansion is carried out in the perpendicular directions and the resulting spectral space is divided into concentric shells of exponentially increasing radius. In each shell velocity and magnetic field fluctuations are represented by complex scalar quantities. Nonlinear effects are reproduced by quadratic terms representing the interactions between nearest and next nearest neighbor shells; the coefficients are chosen so as to conserve 2D quadratic invariants: total energy, cross helicity and squared magnetic potential. The equation of the hybrid shell model is written as:
$`\left({\displaystyle \frac{}{t}}\sigma {\displaystyle \frac{}{s}}\right)Z_n^\sigma (x,t)=\chi k_n^2Z_n^\sigma (s,t)+`$ (5)
$`\mathrm{i}k_n({\displaystyle \frac{13}{24}}Z_{n+2}^\sigma Z_{n+1}^\sigma +{\displaystyle \frac{11}{24}}Z_{n+2}^\sigma Z_{n+1}^\sigma {\displaystyle \frac{19}{48}}Z_{n+1}^\sigma Z_{n1}^\sigma `$
$`{\displaystyle \frac{11}{48}}Z_{n+1}^\sigma Z_{n1}^\sigma +{\displaystyle \frac{19}{96}}Z_{n1}^\sigma Z_{n2}^\sigma {\displaystyle \frac{13}{96}}Z_{n1}^\sigma Z_{n2}^\sigma ))^{}`$
where $`Z_n^\sigma (s,t)=v_n(s,t)+\sigma b_n(s,t)`$ (with $`n=0,1,\mathrm{},n_{max}`$ and $`\sigma =\pm 1`$) are the Elsรคsser variables; $`k_n=k_02^n`$ the transverse wavenumber, with $`k_0=2\pi (L/L_{})`$; $`\chi =\lambda /(c_{A0}L)`$, where the magnetic diffusivity $`\lambda `$ has been assumed equal to the transverse kinematic viscosity; the asterisk means complex conjugate. Lengths are normalized to the loop length $`L`$, and time to the Alfvรฉn transit time $`t_A=L/c_{A0}`$; the velocity $`v_n`$ and magnetic field $`b_n`$ fluctuations are normalized to $`c_{A0}`$ and $`B_0`$, respectively.
The shell technique allows us to describe the turbulence at high Reynolds/Lundquist numbers with a relatively small number of degrees of freedom. In particular, we used a number of shells $`n_{max}=11`$, with a very small dissipation coefficient $`\chi =10^7`$. Since the longitudinal spatial dependence is retained, the hybrid shell model can describe effects of longitudinal resonance. Moreover, it is possible to implement boundary conditions to describe the effects of transverse motions at the loop bases. In particular, the system is excited through the boundary at $`s=0`$, by imposing a given velocity perturbation at large transverse scales, simulating photospheric motions. This boundary perturbation amounts to $`10^5`$ cm s<sup>-1</sup>, is gaussian distributed and has a correlation time $`t_c=300`$ s. At the other boundary $`s=1`$ total reflection conditions are imposed. The equations (5) are numerically solved using second order finite difference schemes, both in space and in time.
During the evolution fluctuating energy enters or exits the driven boundary, so the total energy content in the loop fluctuates erratically in time. At the same time nonlinear effects transfer energy to smaller transverse scales, thus building a turbulence spectrum. Dissipation takes place mainly at the smallest scales. Occasionally, the velocity imposed at the lower boundary drives the loop near to one longitudinal resonance: then, the velocity fluctuations increase at the driven large scale shells, enhancing the energy cascade process towards small dissipative scales. This process results in a spike of dissipated energy, converted to heat. The dissipated power at time $`t`$ and position $`s`$ along the loop is calculated as:
$$H(s,t)=\frac{\chi }{2}\underset{\sigma ,n}{}k_n^2|Z_n^\sigma (s,t)|^2$$
(6)
and is the heating input in the loop plasma model (Eq. 3). The hybrid shell model yields the energy distribution along the loop integrated in the transverse direction, and provides therefore the heat input for the one-dimensional loop model. The power in the whole loop is:
$$W(t)=_0^1H(s,t)๐s$$
(7)
The profile of $`W(t)`$ contains a sequence of spikes of different amplitudes and durations. The space and time profile of the heating function results from the interplay between the external driver (photospheric motions), the loop resonance and the nonlinear turbulent cascade.
The heat spatial distribution is sampled every 0.1 Alfven time. For an Alfven speed of $`2\times 10^8`$ cm/s, one Alfven transit time is 15 s (NMCV04). The numerical table yields the heat distribution per unit time and volume along the loop (sampled every 37.5 km) and span a total time of 307.5 ks, i.e. 3.56 days. We assume a circular cross section and an aspect ratio d/L=0.2, where d is the cross-section diameter; the cross-section area is $`A=2.83\times 10^{17}`$ cm<sup>2</sup>. Fig. 1 shows a few selected segments of the evolution of the average loop heating rate $`W(t)/(AL)`$; they are essentially zooms of the dissipation power shown in Fig. 1 in NMCV04. The heating per unit volume is negligible in the first 1000 s. After this (relatively short) transient, the heating is steadily above $`10^6`$ erg cm<sup>-3</sup> s<sup>-1</sup>. The evolution of the average heating rate is highly irregular, with sharp pulses whose duration spans all time scales from few seconds to a few ks. Some pulses resemble flares. Also the pulses intensity is highly irregular. Most of them are entirely below $`10^4`$ erg cm<sup>-3</sup> s<sup>-1</sup>. A few of them are higher (although mostly below $`10^3`$ erg cm<sup>-3</sup> s<sup>-1</sup>); in fact, eleven heating pulses reach values well above $`3\times 10^4`$ erg cm<sup>-3</sup> s<sup>-1</sup> and occur around 10.5, 22, 25, 57.5, 69.5, 78, 90, 99, 121, 182, 249 ks, as shown in Fig. 1. The most intense pulse is the seventh one (90 ks) and is higher than $`10^3`$ erg cm<sup>-3</sup> s<sup>-1</sup>. The high pulses are noticeably less frequent in the second half of the heating time interval: nine of them occur in the first 150 ks. Most of these pulses last $`0.31`$ ks and are rather peaked.
The heating rate per unit volume averaged over the whole heating duration is $`3\times 10^5`$ erg cm<sup>-3</sup> s<sup>-1</sup>. According to the loop scaling laws (Rosner et al. 1978), for the prescribed length this is the heating rate (per unit volume) of a loop at an equilibrium base pressure of $`0.025`$ dyne cm<sup>-2</sup> and a maximum temperature of $`5\times 10^5`$ K.
Fig. 2 shows distributions of the heating rate per unit volume along the loop sampled during the fourth segment in Fig. 1 (from 22.5 to 27 ks, hereafter segment Ref1). For each time, a couple of distributions are shown, one at 1.5 s from the other. The heating distribution is quite uniform for low heating. During the high intensity phase of the heating, the distribution becomes less uniform, with large peaks propagating back and forth along the loop and extending over $`1/5`$ of the loop.
### 2.2 The initial conditions
Since our scope is to investigate the structure, stability, and observable properties of the simulated loop both in time and on the average, the initial conditions ought to be moderately important: we should start with an initially cool and empty loop, thereafter entirely governed by the new time-dependent heating. For technical reasons, our choice has been to set up this condition by letting an initially hotter loop relax to a much cooler condition. The initial loop is obtained from the model of Serio et al. (1981) with a uniform steady heating and a base pressure 0.03 dyne cm<sup>-2</sup>, corresponding to a loop maximum temperature of $`5\times 10^5`$ K, i.e. the expected average condition of the nano-flare heated loop. In order to let this loop relax, we made a preliminary time-dependent simulation assuming zero coronal heating in the loop (but keeping the chromospheric heating on, to have stable footpoints). The simulation followed the loop evolution for 2000 s, i.e. approximately 2.5 loop thermal decay times (Serio et al. 1991). At the end of the simulation, the loop maximum temperature decreased to $`60,000`$ K, and the pressure to $`1.5\times 10^4`$ dyne cm<sup>-2</sup>. A residual velocity field was present in the loop, with speeds not larger than 6 km/s, an amply subsonic (Mach 0.2) value. We took this final status as the initial condition for the simulations with the nanoflare heating.
## 3 Results
Our main purpose here is to explore how the dissipation rate described in NMCV04 can bring a loop to coronal conditions and maintain it. In this perspective we will describe in detail the solution obtained in a segment containing a heat pulse of medium intensity, specifically the fourth segment (named Ref1, between 22.5 and 26.3 ks) in Fig. 1. We will also discuss the segment including the highest heat pulse, i.e. the eighth segment (which we will call RefH). The solutions in the other segments do not differ much from those that we are going to illustrate.
### 3.1 Medium pulse
Fig. 3 shows the evolution of the temperature, particle density, pressure and velocity distributions along the loop obtained from the loop simulations during segment Ref1. The temperature is steadily below 0.2 MK until the pulse at $`t24.5`$ ks. Then it gradually increases due to the enhanced heating. Fig. 3 clearly shows that the effects of the spatial heating structure (Fig. 2) are smoothed by the efficient thermal conduction. The pulse drives also plasma evaporation from the chromosphere, visible in the density, pressure and velocity distributions (the negative velocity peaks indicate plasma moving upwards from the far footpoint). The density distributions shows more significant fluctuations traveling along the loop.
For more quantitative information, Fig. 4 shows selected distributions of temperature, particle density, velocity, pressure along the loop around the times marked in Fig. 2. Each column of the figure shows the distributions along the loop at the exact time, as well as 100 s before and after this time. In the low heating state (left column), the temperature is steadily between 0.2 and 0.3 MK along most of the loop with a profile very similar to that of a static loop. Also the density does not change much along the loop and is always below $`10^8`$ cm<sup>-3</sup> in most of the loop. The distribution of plasma velocity shows fluctuations with amplitude $`10`$ km/s propagating back and forth along the loop. During the heat pulse, the temperature increases to about 1 MK (in $`100`$ s). The distribution at the time of the temperature maximum appears to be more peaked than in the cool state and the position of the maximum slightly oscillates around the loop apex. At later times ($`t>25`$ ks), the temperature slowly decreases and its distribution flattens (right panel). Asymmetric fronts of plasma evaporation develop as the heating increases (center panel, solid line) and the density starts to increase. The density continues to increase even after the temperature maximum (right panel), staying above $`2\times 10^8`$ cm<sup>-3</sup> for a long time. During the heat pulse, the plasma evaporation fronts are clearly visible also in the velocity profiles: two similar strong fronts rise from both footpoints after t=24.8 ks, reaching a speed of about 50 km/s at intermediate positions along the loop. Then the plasma noticeably becomes less dynamic. During the heating decay, the loop slowly returns to a cool average state around 0.4 MK. The plasma velocity continues to decrease until the plasma becomes practically static around t=25.5 ks. Then the velocity distribution gets inverted: plasma begins to drain along the loop, at very low speed (lower than 10 km/s). The pressure distribution along the loop is quite stable in the cool state. When the heating increases, the pressure increases as well (together with the temperature and the density). The pressure distribution then settles to a very flat distribution during the pulse decay at about 0.04 dyne cm<sup>-2</sup>.
Fig. 5 shows the evolution of the loop maximum temperature, the loop minimum density and pressure, and of the maximum velocity. The first three quantities are typical of the upper region of the loop, close to the apex, the last midway between the apex and the footpoint of the loop. The evolution of the loop maximum temperature is globally similar to that of the average heating (Fig. 1), but much less noisy. Consequently, it is similar also to the evolution of the maximum temperature expected from the evolution of the average loop heating through the loop scaling laws (Rosner et al. 1978). The former temperature is slightly higher ($`10`$ %) and decays more slowly than the latter one. The peak temperature is different because scaling laws assume a constant and uniform heating, while the actual heating function in the simulation is variable and non-uniform along the loop. The slower decay is due to the fact that the plasma response to heating decrease is not instantaneous, and the cooling processes have their own characteristic times. The density enhancement due to the heat pulse of this segment is significantly delayed ($`300`$ s) with respect to the temperature increase, as typical of loop plasma evaporation. For comparison, Fig. 5 shows the equilibrium loop density values as expected from the loop scaling laws. The comparison clearly shows the delay mentioned above, but emphasizes as well that during the pulse rise the loop is significantly underdense, and becomes overdense in the later decay phase. This is expected in dynamically heated loops: while the heating is on, the loop is filling with plasma and therefore below the density equilibrium conditions; when the heating stops, the loop cools down but the plasma drains even more slowly. The maximum pressure has an evolution in between that of the density and of the temperature, and explains why the plasma dynamics is time-shifted with respect to the plasma thermal evolution. Fig. 5 shows that the plasma velocity is constantly below 20 km/s except during the heat pulse, when it grows to about 50 km/s. These values are well subsonic.
From the output results of the hydrodynamic simulations, i.e. distributions of temperature, density and velocity along the loop sampled at regular time intervals, it is possible to compute the UV and X-ray emission from the confined plasma. Fig. 6 shows the emission along the loop in three representative XUV lines, i.e. Ca X 558 ร
, Mg IX 368 ร
, Mg X 625 ร
, peaking at $`\mathrm{log}T=5.9,6.0`$ and 6.1, respectively, at the same times as the distributions shown in the left two columns in Fig. 4. Since the line emission is sensitive both to the temperature and to the square of the density, the emission distributions are less uniform and fluctuate more. This may be a distinctive signature of this model in loop observations. In these lines the loop is visible for a limited time during this segment. In the hottest line (Mg X 625 ร
) it decays very rapidly.
### 3.2 High pulse
In the course of the whole sequence of heating evolution, the most intense heat pulse โ which we will label RefH โ occurs little after time t=90 ks (eighth panel in Fig. 1). Fig. 7 shows the evolution of the loop maximum temperature, the loop minimum density and pressure, and of the maximum velocity, to be compared with the evolution obtained in segment Ref1 (Fig. 5). The loop maximum temperature reaches 1.5 MK around time t=90.5 ks. Then it decays below 1 MK, but stays above 0.5 MK for the rest of the segment because of the occurrence of other minor heat pulses. The density at the apex reaches about $`4\times 10^8`$ cm<sup>-3</sup> and a pressure of 0.1 dyne cm<sup>-2</sup> around time t=91 ks, about 500 s later than the temperature peak. The velocity gets above 60 km/s, always amply subsonic.
Fig. 8 shows the light curves integrated along the whole loop during segment RefH in the 171 A and in the 195 filter bands of the Transition Region and Coronal Explorer (TRACE, Handy et al. 1999). The light curve in the 171 A filter band resembles the evolution of the heat pulses (although much smoother). In the 195 A filter band, only the first pulse is significant, and only in its initial phase the emission is significant, giving the impression of an anticipated evolution. This evolution resembles more closely the evolution of the maximum temperature shown in Fig. 7.
## 4 Discussion and conclusions
This work is devoted to exploring the effect of nanoflares due to the magnetic energy dissipation through MHD-turbulence on the dynamic and thermal evolution of the plasma in a coronal loop. The parameters considered in NMCV04, i.e. and Alfven speed of 2000 km/s corresponding to a magnetic field of about 10 G in corona, lead to a loop with a typical maximum temperature of $`5\times 10^5`$ K. Since coronal loops are typically observed at higher temperatures, $`1`$ MK, here we focus on the effects produced by the most intense heat pulses predicted in NMCV04. We compute in detail the hydrodynamics and thermodynamics of the loop plasma during the pulses and analyze the results.
Although the spatial distribution of the heating has significant fluctuations traveling along the loop and also rapid fluctuations in time, we find that the plasma is not so fast to react and smoothes out the fluctuations both in space and in time. We find that, under the effect of a medium heat pulse, the loop plasma reaches $`T1MK`$ and density $`0.2\times 10^9`$ cm<sup>-3</sup>. The efficient thermal conduction makes the plasma respond promptly to the heating deposition but also smooths the heating fluctuations. The plasma rapidly reaches the equilibrium temperature (according to the loop scaling laws) and then cools following the decay of the heat pulse. The same evolution occurs for a higher heat pulse, which produces a higher peak temperature of 1.5 MK and a higher density of $`0.5\times 10^9`$ cm<sup>-3</sup>. The density (and pressure) of the plasma shows more significant fluctuations traveling along the loop but globally responds on longer time scales. The heat pulses do not last long enough to let the plasma reach the thermo/hydrostatic equilibrium: the plasma is underdense during the heat pulse and overdense after the pulse with respect to thermal equilibrium. This density evolution is a consequence of the impulsive heating (Winebarger et al. 2003a, Warren et al. 2003). The speed of the plasma driven by the heat pulse is relatively small, largely subsonic, and speeds of few tens of km/s occur only for very few minutes. The emission distribution in relevant spectral lines may be relatively more sensitive to fluctuations due to the turbulent heating and may be used to diagnose this model. For the highest heat pulse, our model also predicts the light curves in two relevant TRACE filter bands to be โout of phaseโ one from the other. This phase difference is in qualitative agreement with observations (Winebarger et al. 2003b) but also predicted by other different loop models (Warren et al. 2003).
The heating model used here has very few free parameters (essentially the magnetic field strength and the loop length) and depends on basic physical effects. The shell model does not yield a detailed description of turbulence, and cannot reproduce the energy distribution, in the direction transverse to the magnetic field. However, it should be adequate to describe the behaviour of the loop integrated in the transverse direction and the detailed energy dissipation along the loop, matching the scope of the Palermo-Harvard loop model.
A series of questions are opened by this work. First, characterizing features of the proposed heating are the disturbances traveling along the loop. We have shown that observations in single spectral lines may be sensitive to disturbances in the loop, but detecting such effects may not be trivial with present day instruments. Also, one may wonder on the effect of changing the magnetic field strength; can a stronger field lead to hotter active region loops or even major flares? Even if the heating function may be modified with a simple scaling, this question requires anyhow additional detailed loop modeling, since the loop plasma evolves non-linearly under the effect of the heating, coupled with the dynamics and the cooling processes.
As further issue to investigate, we note that the heating function is modified by the local plasma conditions, e.g. the density stratification and its time variation (the Alfven speed depends on the density). Including self-consistently a feedback of the loop plasma conditions on the energy dissipation may easily modify some characteristics of the heating function, such as the pulse duration, and thus influence the results. Tackling this question requires to couple the hybrid MHD turbulence model with the loop time-dependent hydrodynamic model, a task planned for future work.
This first work paves the path to future works along several lines, such as the time decomposition analysis of results and the coupling of the heating and loop models, the comparison with observations, encompassing the selection (or acquisition) and analysis of observations made of long and regularly sampled image sequences.
FR and GP acknowledge support for this work from Ministero dellโIstruzione, Universitร e Ricerca. G.N., F.M. and P.V. acknowledge partial support by the MIUR (Ministero dell Istruzione, dellโUniversitร e della Ricerca) through a National Project Fund (cofin 2002) and by the European Community within the Research Training Network Turbulence in Space Plasmas, Theory, Observation and Simulation. RMHD numerical calculations were performed in the framework of HPCC (Center for High Performance Computing) of the University of Calabria.
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# 1 Introduction
## 1 Introduction
It is well known that the standard theory of fiber bundles is designed to deal with the situations when one has ambiguously defined functions โ sections of fiber bundles over a space $`X`$. In the modern language this means that, knowing the value of a section of a fiber bundle at a point $`x`$, we can find its value at any other point $`y`$. But the result will depend on the path taken from $`x`$ to $`y`$ and from the choice of the oneโform connection on the bundle. The path ordered exponent of the connection gives the map between the fibers $`V_x`$ and $`V_y`$ over the points $`x`$ and $`y`$.
This construction assumes that the oneโform connection is unambiguously defined up to the gauge transformations โ rotations in the fibers. But what if the oneโform connection was ambiguous?( The simplest example of such a situation, which we know, is the WuโYang monopole.) To describe the situation with the ambiguous oneโform connections one has to have a โconnectionโ for the connections โ the one to be ordered over the surface inside $`X`$. The latter orderings give 2D holonomies which describe maps between the sequences of fibers over different curves, along which the original oneโform was ordered. One can consult e.g. , on the mathematics behind such a construction.
One badly needs a local field theory description of the 2D holonomies in terms of twoโtensor fields. In and here we propose such a theory. To get it one has to organize a triangulationโindependent areaโordering of twoโtensor gauge field carrying color indices. The triangulationโindependence is very crucial for the consistency of the theory. It means the following fact. To construct the areaโordering we should approximate the Riemann surfaces by their triangulated versions. For the discretized case the area ordering is obtained as follows. We glue the exponents of the twoโtensor field over the whole simplicial surface by putting them at different simplices (triangles), which are sitting at different points inside the base space $`X`$. Now it is easy to see that the twoโtensor connection should carry three color indices along with two spacial ones : $`B_{\mu \nu }^{ijk}`$, $`\mu ,\nu =1,\mathrm{},\mathrm{dim}X`$ and $`i,j,k=1,\mathrm{},\mathrm{dim}V`$. This is due to the three wedges of each 2D simplex. In we have found a way to exponentiate cubic matrices, which we refer to as the surface exponent. With such an exponent we can take the continuum limit from the simplicial surfaces to the smooth ones. It is important that with the use of the surface exponent the result for the limit will not depend on the way it is taken! This is what we refer to as triangulationโindependence.
The organization of the paper is as follows. In the next section we briefly review the paper and concisely formulate its statements. In the section 3 we find the differential equation for the surface exponent and define a new way of the exponentiation of the quadratic matrices. We establish the relation between such an exponent of the quadratic matrices and the surface exponent for the cubic ones. At the end of the section 3 we discuss a possible matrix integral formulation of the surface exponent. In the section 4 we derive the gauge transformation rules and curvature for the twoโtensor gauge fields. As well in that section we discuss the meaning of the theory of the nonโAbelian tensor fields. In the section 5 we establish links between the latter theory and the lattice integrable models and String Field Theory. We conclude with summary in the section 6.
## 2 Exponentiation of cubic matrices and the areaโordering
### 2.1 The surface exponent
In the paper we have defined the surface exponent of a cubic matrix $`\widehat{B}=B_{ijk}`$, $`i,j,k=1,\mathrm{},N`$ as follows<sup>2</sup><sup>2</sup>2The reason for the subscripts on the LHS of this equation will be explained below.:
$`\mathrm{Tr}\left(E_{g,I,\kappa }^{\widehat{B}}\right)\underset{M\mathrm{}}{lim}{\displaystyle \underset{\mathrm{graph},g}{\overset{M}{}}}\left(\widehat{I}+{\displaystyle \frac{\widehat{B}}{M}}\right),`$ (1)
where the limit is taken over a sequence of closed (because on the LHS we take Tr โ no any free indices), connected threeโvalent graphs<sup>3</sup><sup>3</sup>3Threeโvalent graphs are those in which three wedges terminate in each their vertex. Note that in we have considered the triangulation graphs โ those which triangulate Riemann surfaces and are dual to the fat three-valent graphs (see fig. 1). with $`M`$ vertices. At a fixed $`M`$ the product on the RHS of eq.(1) is taken over an $`M`$โvertex threeโvalent graph: At each vertex of the graph we put matrix $`I_{ijk}+B_{ijk}/M`$ and we glue the indices of these matrices with the use of a bilinear form $`\kappa ^{ij}`$ (as it is shown in the fig. 2). It seem that to take the limit $`M\mathrm{}`$ we have to choose a sequence of graphs. But in we have proved that the limit (1) does NOT depend on the choice of the sequence of graphs under the following conditions:
* One should consider fat threeโvalent graphs so that it is obvious on which minimal genus Riemann surface this graph can be mapped. The genus $`g`$ of the graphs in the sequence should be fixed: I.e. at every $`M`$ the same topology graphs should be taken.
* The matrices $`\widehat{I}`$ and $`\widehat{\kappa }`$ should obey the following conditions:
$`I_{ijk}=I_{kij}=I_{jki}\mathrm{cyclic}\mathrm{symmetry},`$
$`{\displaystyle \underset{j,k=1}{\overset{N}{}}}I_{ij}^{}{}_{}{}^{k}I_{lk}^{}{}_{}{}^{j}=\kappa _{il}\mathrm{normalization},`$
$`{\displaystyle \underset{n=1}{\overset{N}{}}}I_{inl}I_{}^{n}{}_{jk}{}^{}={\displaystyle \underset{n=1}{\overset{N}{}}}I_{ijn}I_{}^{n}{}_{kl}{}^{}\mathrm{fusion}.`$ (2)
and we higher the indices of $`I_{ijk}`$ with the use of $`\kappa ^{ij}`$. One can find the graphical representation of the last two equations in the fig. 3.
* In the sequence we confine to the graphs of the following kind: As the limit $`M\mathrm{}`$ is taken the number of wedges of each face of the graph should be suppressed in comparison with $`M`$. The geometric meaning of this condition is explained in .
In we consider one extra condition. However, this condition is not necessary for the limit (1) to be independent of the choice of the sequence of the graphs. It just simplifies the considerations of .
In general at given $`N`$ there are many solutions to the equations for $`\widehat{I}`$ and $`\widehat{\kappa }`$ and the exponent depends on the choice of $`\widehat{I}`$, $`\widehat{\kappa }`$ and $`g`$. This explains the subscripts in eq.(1) . The result of the limit is:
$`\mathrm{Tr}\left(E_{g,I,\kappa }^{\widehat{B}}\right)={\displaystyle \underset{F=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{F!}}B^{j_1^{(1)}j_2^{(1)}j_3^{(1)}}\mathrm{}B^{j_1^{(F)}j_2^{(F)}j_3^{(F)}}I_{j_3^{(1)}j_2^{(1)}j_1^{(1)}\left|\mathrm{}\right|j_3^{(F)}j_2^{(F)}j_1^{(F)}}^g,`$ (3)
where we use the same letter $`I`$ to denote another matrix $`I_{j_3^{(1)}j_2^{(1)}j_1^{(1)}\left|\mathrm{}\right|j_3^{(F)}j_2^{(F)}j_1^{(F)}}^g`$ with $`3F`$ indices, which is obtained via multiplication of $`I_{ijk}`$ over any oriented threeโvalent graph with $`g`$ handles, $`F`$ holes and 3 external wedges at each hole (see fig. 4). At the same time $`I^g(F=0)`$ is just a number obtained via the multiplication of $`I_{ijk}`$ over any closed genus $`g`$ graph. For example, $`I^0=_{i,k}\kappa _i^k\kappa _k^i=_i\delta _i^i=N`$. The reason for the division of the indices of $`I_{j_3^{(1)}j_2^{(1)}j_1^{(1)}\left|\mathrm{}\right|j_3^{(F)}j_2^{(F)}j_1^{(F)}}^g`$ into groups separated by vertical lines is as follows: These matrices are only cyclicly symmetric under the exchange of their indices inside each triple, but are completely symmetric under any exchange of the triples between themselves . All that follows from the topology of the underlaying graphs (see fig. 4) and the condition in fig. 3.
Once we have defined the trace of the surface exponent, we can define the exponent with any number of indices or, better to say, for any twoโdimensional topology โ for a given number of handles $`g`$, holes $`L`$ and distribution of external indices over the holes. To begin with, let us define the exponent for the disc topology and three external wedges at the boundary. To do that in eq.(1) we have to use open graph with one hole and three wedges at the boundary. The result for the corresponding limit is:
$`\left(E_{I,\kappa }^{\widehat{B}}\right)_{m_1m_2m_3}={\displaystyle \underset{F=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{F!}}B^{j_1^{(1)}j_2^{(1)}j_3^{(1)}}\mathrm{}B^{j_1^{(F)}j_2^{(F)}j_3^{(F)}}I_{m_1m_2m_3\left|j_3^{(1)}j_2^{(1)}j_1^{(1)}\right|\mathrm{}|j_3^{(F)}j_2^{(F)}j_1^{(F)}}.`$ (4)
The exponent in eq.(4) is the building block for the construction of a generic surface exponent and for the areaโordering below. In fact, we can glue a twoโdimensional surface of any topology with the use of triangles โ dual to the threeโvalent vertices. Now we have to put $`\left(E_{I,\kappa }^{\widehat{B}}\right)_{m_1m_2m_3}`$ instead of $`(I+B/M)_{ijk}`$ or $`I_{ijk}`$ in the vertices. For example, eq.(3) is obtained via multiplication of eq.(4) over any genus $`g`$ closed triangulated Riemann surface: To obtain Tr$`E^{\widehat{B}}`$ we have to take $`\widehat{B}/J`$ in eq.(4) , where $`J`$ is the total number of triangles out of which the closed surface in question is constructed. This is true due to the main property of the surface exponent which is discussed in the section 3.
### 2.2 Triangulationโindependent areaโordering
With the use of the surface exponent, in we have constructed the triangulationโindependent ordering of nonโAbelian twoโtensor fields over twoโdimensional surfaces. Let us repeat that construction here. It is natural to consider, within this context, a triangulated approximation $`\stackrel{~}{\mathrm{\Sigma }}(\stackrel{~}{\gamma _1},\mathrm{},\stackrel{~}{\gamma _L})`$ of an oriented Riemann surface $`\mathrm{\Sigma }(\gamma _1,\mathrm{},\gamma _L)`$ in a target space $`X`$. This surface has $`L`$ boundary closed loops $`\gamma `$โs which are approximated by the closed broken lines $`\stackrel{~}{\gamma }`$โs (see fig. 5). At the end we take the continuum limit and the result does not depend on the way the limit is taken: It does not depend on the choice of the sequence of triangulated surfaces $`\stackrel{~}{\mathrm{\Sigma }}`$ which are approaching the smooth surface $`\mathrm{\Sigma }`$ in the limit in question! It is this fact which we refer to as the triangulation independence.
To get the triangulationโindependent areaโordering, we have to assign the matrix
$$U_{ijk}(x,\mathrm{\Delta }\sigma ^{\mu \nu })=\left(E_{I,\kappa }^{\widehat{B}_{\mu \nu }(x)\mathrm{\Delta }x^\mu \mathrm{\Delta }x^\nu }\right)_{ijk}$$
to each simplex (triangle) sitting at the point $`x`$ of the image of the triangulated surface in $`X`$. Here $`\mathrm{\Delta }\sigma ^{\mu \nu }=\mathrm{\Delta }x^\mu \mathrm{\Delta }x^\nu `$ is the oriented area element associated with the triangle in question. The indices $`i,j`$ and $`k`$ are assigned to the three wedges of the corresponding triangle. Hence, $`\widehat{U}:V^3`$ for an $`N`$โdimensional vector space $`V`$ โ the โfiberโ of our โfiber bundleโ, while $`X`$ is the base of the bundle.
The areaโordering is obtained by gluing, via the use of the biโlinear form $`\kappa ^{ij}:V^2`$, the matrices $`U_{ijk}`$ on each triangle over the whole simplicial surface (see fig. 6):
$`U\left[\stackrel{~}{\mathrm{\Sigma }}(\stackrel{~}{\gamma _1},\mathrm{},\stackrel{~}{\gamma _L})\right]_{j_1^{(1)}\mathrm{}j_{n_1}^{(1)}\left|j_1^{(2)}\mathrm{}j_{n_2}^{(2)}\right|\mathrm{}|j_1^{(L)}\mathrm{}j_{n_L}^{(L)}}`$
$`{\displaystyle \underset{k_1,k_2,\mathrm{},k_w}{}}U_{j_1^{(1)}k_1k_2}(x_1^{},\mathrm{\Delta }\sigma _1)U_{}^{k_1k_3}{}_{k_4}{}^{}(x_2^{},\mathrm{\Delta }\sigma _2)U_{}^{k_3}{}_{j_2^{(1)}j_3^{(1)}}{}^{}(x_3^{},\mathrm{\Delta }\sigma _3)\mathrm{}`$ (5)
where $`w`$ is the total number of internal wedges of the graph; $`j^{(l)}`$ are the indices corresponding to the wedges of the $`l`$-th broken line $`\stackrel{~}{\gamma }_l`$ and $`n_l`$ is the total number of the wedges of this broken line. We higher (lower) the indices via the use of the aforementioned bilinear form $`\kappa ^{ij}`$ ($`\kappa _{ij}\kappa ^{jk}=\delta _i^k`$).
In the continuum limit the discrete indices $`1,\mathrm{},n_l`$ are converted into the continuum ones $`s_l[0,\mathrm{\hspace{0.17em}2}\pi )`$ and we obtain:
$`U\left[\stackrel{~}{\mathrm{\Sigma }}(\stackrel{~}{\gamma _1},\mathrm{},\stackrel{~}{\gamma _L})\right]_{j_1^{(1)}\mathrm{}j_{n_1}^{(1)}\left|j_1^{(2)}\mathrm{}j_{n_2}^{(2)}\right|\mathrm{}|j_1^{(L)}\mathrm{}j_{n_L}^{(L)}}U\left[\mathrm{\Sigma }^{}(\gamma _1,\mathrm{},\gamma _L)\right]_{j^{(1)}(s_1)\left|j^{(2)}(s_2)\right|\mathrm{}|j^{(L)}(s_L)},`$
$`\mathrm{where}\widehat{U}(\mathrm{\Sigma }):V^{\mathrm{}}(1)\mathrm{}V^{\mathrm{}}(L).`$ (6)
Here $`j^{(l)}(s_l)`$ is the โcolorโ index assigned to the continuous number of points enumerated by $`s_l`$ โ a parametrization of the $`l`$-th loop $`\gamma _l`$. The expression (6) is not so unfamiliar for the string theoreticians as it could seem from the first sight. In fact, if we substitute $`j^{(l)}(s_l)`$ by $`x^{(l)}(s_l)`$, then $`U\left[\mathrm{\Sigma }(\gamma _1,\mathrm{},\gamma _L)\right]`$ can be considered as a kind of string amplitude whose endโloops $`\gamma `$โs are mapped by $`x(s)`$โs. This is the subject of the section 5.
Using the surface exponent (4), we can write the explicit expression for the areaโordered exponent (โ$`AE`$โ) (5)โ(6). For example, for the case of the disc $`D`$ we obtain:
$`U_{j(s)}(D,\widehat{B})=\left(AE_{g,I,\kappa }^{{\scriptscriptstyle _D\widehat{B}_{\mu \nu }๐x_\mu ๐x_\nu }}\right)_{j(s)}{\displaystyle \underset{F=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{F!}}I_{j(s)\left|j_3^{(1)}j_2^{(1)}j_1^{(1)}\right|\mathrm{}|j_3^{(F)}j_2^{(F)}j_1^{(F)}}^g\times `$
$`\times {\displaystyle }{\displaystyle _D}B_{\mu _1\nu _1}^{j_1^{(1)}j_2^{(1)}j_3^{(1)}}dx_{\mu _1}dx_{\nu _1}\mathrm{}{\displaystyle }{\displaystyle _D}B_{\mu _F\nu _F}^{j_1^{(F)}j_2^{(F)}j_3^{(F)}}dx_{\mu _F}dx_{\nu _F}=\left(E_{g,I,\kappa }^{{\scriptscriptstyle _D\widehat{B}_{\mu \nu }๐x_\mu ๐x_\nu }}\right)_{j(s)},`$ (7)
where $`s`$ is the parametrization of the boundary and $`j(s)`$ is the index function at the boundary $`D`$. The last equality holds because of the symmetry properties of the matrices $`I_{j(s)\left|j_3^{(1)}j_2^{(1)}j_1^{(1)}\right|\mathrm{}|j_3^{(F)}j_2^{(F)}j_1^{(F)}}`$. These properties are discussed below the eq.(3) . As the result, such an areaโordering is rather trivial because there is no need to order anything. In fact, we can easily interchange the order of integrals over $`B`$ in eq.(7) . This should be the generic property of all โhigher dimensionalโ exponents (for multi-index matrices) due to the triviality of the corresponding homotopy groups.
Let us clarify this important subject. In the oneโdimensional ordering for oneโform gauge fields the order is important for the following reason: If one throws away a point (insertion of a local observable) from an open curve (along which the ordering is done) it becomes disconnected. Then it is important what is on the left and what is on the right from the point in question. On the other hand, in the case of twoโdimensional orderings, throwing away a point is not sufficient to make the surface disconnected. Hence, the order is unimportant. As the result there are no commutators of local fields which describe twoโdimensional areaโorderings. We come back to this point again at the end of the section 4.
However, the commutative nature of the areaโordering does not mean that we have obtained something trivial! In fact, if we consider a surface $`\mathrm{\Sigma }`$ with one designated point $`x`$ and two external wedges at this point the corresponding $`U(\mathrm{\Sigma }_x)_{ij}`$ (or $`U(\mathrm{\Sigma }_x)_{i}^{}{}_{}{}^{j}=U(\mathrm{\Sigma }_x)_{ik}\kappa ^{kj}`$) gives a nontrivial (nonโdiagonal matrix) map. All that goes without saying that more complicated matrices $`U(\mathrm{\Sigma }_{x,y,\mathrm{}})_{ijm\mathrm{}}^{}{}_{}{}^{kl\mathrm{}}`$ give completely nonโtrivial nonโlinear maps. Somehow the latter objects intrinsically include interactions because the $`B`$โfield and especially the background $`\widehat{I}`$ give nonโlinear maps ($`\widehat{B},\widehat{I}:V^3`$). As the result $`U`$โs give nonโtrivial nonโlinear maps between the endโloops of the corresponding surfaces.
In any case, to our mind this is a deep fact which is yet to be understood to establish the relation of the subject in question to the String Field Theory and to the standard theory of fiber bundles. Note that in we describe how to construct modified surface โexponentโ, which obeys less trivial relations, and we think that the surface exponent can appear in the applications in the latter form.
## 3 Relevant properties of the exponent
### 3.1 Differential equation for the exponent
Let us derive the differential equation for the surface exponent. To do that recall that the exponent of a quadratic matrix $`A_{i}^{}{}_{}{}^{j}`$ has the following main property (which is at the basis of the probability theory and all evolution phenomena):
$`\left(e^{t_1\widehat{A}}\right)_{i}^{}{}_{}{}^{j}\left(e^{t_2\widehat{A}}\right)_{j}^{}{}_{}{}^{k}=\left(e^{\left(t_1+t_2\right)\widehat{A}}\right)_{i}^{}{}_{}{}^{k}`$ (8)
for any two numbers $`t_1`$ and $`t_2`$. This equation defines the exponent unambiguously. In fact, from eq.(8) one can derive the differential equation for the exponent: Choose $`t_1=t`$ and $`t_2=dtt`$. Then, expanding over $`dt`$ both sides of eq.(8) , we obtain:
$`{\displaystyle \frac{d}{dt}}\left(e^{t\widehat{A}}\right)_i^j=A_i^k\left(e^{t\widehat{A}}\right)_k^j.`$ (9)
The surface exponent can have any number of external indices โ not only zero or two like the exponent of a quadratic matrix. As the result it obeys many different identities following from the conditions of the triangulation independence. But all these conditions originate from the basic one<sup>4</sup><sup>4</sup>4Note that according to the discussion at the end of the section 2 we actually have the relation as: $`\left(E_{I,\kappa }^{\widehat{B}_1}\right)_{j_1j_2}^{}{}_{}{}^{k_1}\left(E_{I,\kappa }^{\widehat{B}_2}\right)_{j_3j_4k_1}=\left(E_{I,\kappa }^{\widehat{B}_1+\widehat{B}_2}\right)_{j_1j_2j_3j_4}`$ for any cubic matrices $`\widehat{B}_1`$ and $`\widehat{B}_2`$.:
$`\left(E_{I,\kappa }^{t_1\widehat{B}}\right)_{j_1j_2}^{}{}_{}{}^{k_1}\left(E_{I,\kappa }^{t_2\widehat{B}}\right)_{j_3j_4k_1}=\left(E_{I,\kappa }^{\left(t_1+t_2\right)\widehat{B}}\right)_{j_1j_2j_3j_4},`$ (10)
which is shown graphically in the fig. 7. This equation, however, does not define unambiguously the surface exponent, because $`\widehat{I}`$ and $`\widehat{\kappa }`$ do not explicitly present in this equation.
However, if $`t_1=t`$ and $`t_2=dtt`$, we can obtain the following differential equation:
$`{\displaystyle \frac{d}{dt}}\left(E_{I,\kappa }^{t\widehat{B}}\right)_{j_1j_2}^{}{}_{}{}^{k_1}I_{j_3j_4k_1}=\left(E_{I,\kappa }^{t\widehat{B}}\right)_{j_1j_2}^{}{}_{}{}^{k_1}B_{j_3j_4k_1},`$ (11)
which fixes the exponent unambiguously for a particular choice of $`\widehat{I}`$ and $`\widehat{\kappa }`$ matrices, because now they are explicitly present in the equation.
Similarly one can write the variational equation for the areaโordered exponent (7):
$`{\displaystyle \frac{\delta ^2}{\delta \sigma _{\mu \nu }^2(s)}}\left(E_{g,I,\kappa }^{{\scriptscriptstyle _D\widehat{B}_{\mu \nu }๐x_\mu ๐x_\nu }}\right)_{\mathrm{}j_s}I^{j_skl}=\left(E_{g,I,\kappa }^{{\scriptscriptstyle _D\widehat{B}_{\mu \nu }๐x_\mu ๐x_\nu }}\right)_{\mathrm{}j_s}B_{\mu \nu }^{j_skl}\left[x(s)\right],`$ (12)
where $`\delta ^2/\delta \sigma _{\mu \nu }^2(s)`$ is the standard variation with respect to the addition of a small area $`(\mathrm{\Delta }\sigma ^{\mu \nu })`$ to the disc $`D`$ and $`j_s`$ means single index $`j`$ at the point $`s`$.
### 3.2 New way to exponentiate quadratic matrices and its relation to the surface exponent
In this subsection we will define nonโstandard way of exponentiation of quadratic matrices. It is inspired by the surface exponent. In fact, we can put in the vertices of the graphs in eq.(1) the matrix $`I_{ijk}`$ rather than $`I_{ijk}+B_{ijk}/M`$, but glue their indices with the use of $`\kappa ^{ij}+B^{ij}/W`$ rather than just with the use of $`\kappa ^{ij}`$, i.e.:
$`\mathrm{Tr}\left(\stackrel{~}{E}_{g,I,\kappa }^{\widehat{B}}\right)\underset{W\mathrm{}}{lim}{\displaystyle \underset{\mathrm{graph}(g)}{\overset{W}{}}}\left(\widehat{\kappa }+{\displaystyle \frac{\widehat{B}}{W}}\right),`$ (13)
where under โgraphโ we assume the same kind of threeโvalent genus $`g`$ closed graph as in eq.(1) , which has $`W`$ wedges. In the light of the above discussion the limit (13) does not depend on the choice of the sequence of graphs. The result for the limit is:
$`\mathrm{Tr}\left(\stackrel{~}{E}_{g,I,\kappa }^{\widehat{B}}\right)={\displaystyle \underset{F=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{F!}}B^{j_1^{(1)}j_2^{(1)}}\mathrm{}B^{j_1^{(F)}j_2^{(F)}}I_{j_1^{(1)}j_2^{(1)}\left|\mathrm{}\right|j_1^{(F)}j_2^{(F)}}^g,`$ (14)
where $`I_{j_1^{(1)}j_2^{(1)}\left|\mathrm{}\right|j_1^{(F)}j_2^{(F)}}^g`$ is the matrix with $`2F`$ indices obtained via multiplication of $`I_{ijk}`$ over any oriented graph with $`g`$ handles, $`F`$ holes and 2 external wedges at each hole. Similarly we can define exponents with any distribution of indices.
Now we are going to show that this exponent is equivalent to the original one. In fact, it is easy to see that:
$`\underset{W\mathrm{}}{lim}{\displaystyle \underset{\mathrm{graph}(g)}{\overset{W}{}}}\left(\widehat{\kappa }+{\displaystyle \frac{\widehat{B}}{W}}\right)=\underset{W\mathrm{}}{lim}{\displaystyle \underset{\mathrm{graph}(g)}{\overset{W}{}}}\left(\widehat{\kappa }+{\displaystyle \frac{\widehat{B}}{2W}}\right)^2=\underset{M\mathrm{}}{lim}{\displaystyle \underset{\mathrm{graph}(g)}{\overset{M}{}}}\left(1+{\displaystyle \frac{\widehat{B}}{3M}}\right)^3\widehat{I},`$ (15)
where $`M`$ is the number of vertices of the same graph. (Note that $`2W=3M`$ for the threeโvalent graphs.) In the last expression of eq.(15) we put in the vertices the following matrix:
$`\left(1+{\displaystyle \frac{\widehat{B}}{3M}}\right)^3\widehat{I}\left(1+{\displaystyle \frac{B}{3M}}\right)_i^i^{}\left(1+{\displaystyle \frac{B}{3M}}\right)_j^j^{}\left(1+{\displaystyle \frac{B}{3M}}\right)_k^k^{}I_{i^{}j^{}k^{}}=`$
$`=I_{ijk}+{\displaystyle \frac{B_i^i^{}}{3M}}I_{i^{}jk}+{\displaystyle \frac{B_j^j^{}}{3M}}I_{ij^{}k}+{\displaystyle \frac{B_k^k^{}}{3M}}I_{ijk^{}}+๐ช\left({\displaystyle \frac{1}{M^2}}\right).`$ (16)
We higher the indices in this expression with the use of $`\kappa ^{ij}`$ and use in eq.(15) and eq.(16) the fact that $`\kappa _{ij}\kappa ^{jk}=\delta _i^k`$. In the limit $`M\mathrm{}`$ the terms $`๐ช(1/M^2)`$ in the eq.(16) do not survive and we obtain the relation between the cubic and quadratic matrices:
$`\stackrel{~}{B}_{ijk}={\displaystyle \frac{B_i^i^{}}{3}}I_{i^{}jk}+{\displaystyle \frac{B_j^j^{}}{3}}I_{ij^{}k}+{\displaystyle \frac{B_k^k^{}}{3}}I_{ijk^{}}`$ (17)
for which the following equality holds:
$`\stackrel{~}{E}_{g,I,\kappa }^{\widehat{B}}=E_{g,I,\kappa }^{\widehat{\stackrel{~}{B}}}.`$ (18)
Thus, basically we have the unique โtwoโdimensionalโ exponent. It is worth pointing out now that all the relevant properties of the surface exponent are due to the fact that it is the background $`I_{ijk}`$ who carries three color indices rather than the matrix $`B_{ijk}`$ under the exponent. In any case, we are going to use the new way of exponentiation of the quadratic matrices in the section 4.
### 3.3 Towards Matrix model representation for the exponent
In this subsection we discuss our attempts to represent the surface exponent via the use of a matrix integral. In fact, the mentioned above matrices $`I_{j_3^{(1)}j_2^{(1)}j_1^{(1)}\left|\mathrm{}\right|j_3^{(F)}j_2^{(F)}j_1^{(F)}}^g`$ are 2D topological invariants in the sense discussed say in . This is true due to the conditions (2) imposed on $`\widehat{I}`$ and $`\widehat{\kappa }`$: The matrix $`I_{j_3^{(1)}j_2^{(1)}j_1^{(1)}\left|\mathrm{}\right|j_3^{(F)}j_2^{(F)}j_1^{(F)}}^g`$ depends only on the topology of the discretized Riemann surface<sup>5</sup><sup>5</sup>5By the โtopologyโ in this case we mean a given number of handles $`g`$, number of holes $`L`$ and the distribution of external wedges over the holes., but does not depend on its concrete triangulation. We would like to write an explicit representation for this matrix through the correlation function in a 2D topological theory.
To do that let us consider the matrix integral, which is a generalization of the integral considered in :
$`Z(K,N,\kappa ,I)={\displaystyle \underset{i=1}{\overset{N}{}}d\widehat{\mathrm{\Phi }}^i\mathrm{exp}\left\{K\mathrm{Tr}\left[\widehat{\mathrm{\Phi }}^i\widehat{\mathrm{\Lambda }}\widehat{\mathrm{\Phi }}^j\kappa _{ij}+\frac{\mathrm{i}}{3}I_{ijk}\widehat{\mathrm{\Phi }}^i\widehat{\mathrm{\Phi }}^j\widehat{\mathrm{\Phi }}^k\right]\right\}},`$ (19)
where from now on we put $`\widehat{\mathrm{\Lambda }}=1`$. The integral is taken over the unitary matrices $`\widehat{\mathrm{\Phi }}_i=\mathrm{\Phi }_i^{ab}`$, where $`i=1,\mathrm{},N`$ and $`a=1,\mathrm{},K`$. Taylor expanding this integral over the powers of $`I_{ijk}`$ and using the decoupling of correlations for the Gaussian integral, usually referred to as Wickโs theorem, we obtain the Feynman diagram expressions for the contributions to the integral. Then, it is easy to see that:
$`{\displaystyle \frac{^{\chi (g)}}{K^{\chi (g)}}}Z(K,N,\kappa ,I)|_{K=0}I^g,g1,`$ (20)
because by taking this derivative we extract the genus $`g`$, threeโvalent, closed graph (Feynman diagram), where in the vertices the matrix $`I_{ijk}`$ is standing, while in the wedges one puts the โpropagatorโ $`\kappa ^{ij}`$ ($`\kappa ^{ij}\kappa _{jk}=\delta _k^i`$). Note that the integral (19) is convergent and explicitly calculable:
$`\mathrm{log}Z(K,N,\kappa ,I)={\displaystyle \underset{g=0}{\overset{\mathrm{}}{}}}C_gK^{\chi (g)}I^g,`$ (21)
where $`C_g`$ are some combinatorial coefficients counting the numbers (with alternating signs) of threeโvalent connected graphs at the given $`g`$.
It seems that by considering the $`g`$ handle contribution to the correlation function:
$`{\displaystyle \underset{i=1}{\overset{N}{}}d\widehat{\mathrm{\Phi }}^i}:\widehat{\mathrm{\Phi }}_{j_3^{(1)}}\widehat{\mathrm{\Phi }}_{j_2^{(1)}}\widehat{\mathrm{\Phi }}_{j_1^{(1)}}:\mathrm{}:\widehat{\mathrm{\Phi }}_{j_3^{(F)}}\widehat{\mathrm{\Phi }}_{j_2^{(F)}}\widehat{\mathrm{\Phi }}_{j_1^{(F)}}:\mathrm{exp}\left\{K\mathrm{Tr}\left[\widehat{\mathrm{\Phi }}^i\widehat{\mathrm{\Phi }}^j\kappa _{ij}+{\displaystyle \frac{\mathrm{i}}{3}}I_{ijk}\widehat{\mathrm{\Phi }}^i\widehat{\mathrm{\Phi }}^j\widehat{\mathrm{\Phi }}^k\right]\right\}`$ (22)
will give us the desired matrix $`I_{j_3^{(1)}j_2^{(1)}j_1^{(1)}\left|\mathrm{}\right|j_3^{(F)}j_2^{(F)}j_1^{(F)}}^g`$ if we do not Wick contract the $`\widehat{\mathrm{\Phi }}`$โs inside the โnormal orderingโ โ $`:\mathrm{\Phi }\mathrm{\Phi }\mathrm{\Phi }:`$. In fact, this correlation function gives a 2D topological invariant in the sense of . However this topological invariant does not coincide with $`I_{j_3^{(1)}j_2^{(1)}j_1^{(1)}\left|\mathrm{}\right|j_3^{(F)}j_2^{(F)}j_1^{(F)}}^g`$ because we can not separate the contact terms from the desired ones: In the result of the calculation of the genus $`g`$ contribution to the eq.(22) (if $`F=2`$), along with the matrix $`I_{j_3^{(1)}j_2^{(1)}j_1^{(1)}|j_3^{(2)}j_2^{(2)}j_1^{(2)}}^g`$ the following kind of contribution appears: $`I_{j_3^{(1)}j_2^{(1)}j_2^{(2)}j_1^{(2)}}^g\kappa _{j_1^{(1)}j_3^{(2)}}`$ โ one of the contact terms.
At this stage we do not know how to separate this two kinds of contributions to the correlation functions (22). However, we would like to find a way to do that, because it will establish a closer relation of our considerations to the subject of the String Field Theory.
## 4 Theory of nonโAbelian twoโtensor fields
In this section we discuss gauge transformations and curvature for the twoโtensor gauge connection. To explain the idea of our argument and to set the notations let us present here the derivation of the gauge transformation and curvature of an ordinary gauge field. Consider a holonomy matrix
$$\widehat{U}(x,\mathrm{\Delta }x)=e^{\mathrm{i}\widehat{A}_\mu (x)\mathrm{\Delta }x^\mu }$$
for a small $`\mathrm{\Delta }x`$. This matrix transforms, under the rotations in the fibers, as:
$$\widehat{\stackrel{~}{U}}(x,\mathrm{\Delta }x)=\widehat{g}^1(x)\widehat{U}(x,\mathrm{\Delta }x)\widehat{g}(x+\mathrm{\Delta }x).$$
Then, expanding in powers of $`\mathrm{\Delta }x`$ both sides of this expression, we obtain:
$$\widehat{\stackrel{~}{A}}_\mu (x)=\widehat{g}^1(x)\left[_\mu +\mathrm{i}\widehat{A}_\mu (x)\right]\widehat{g}(x).$$
Which is the gauge transformation for the gauge field.
As well, by considering the holonomy matrix for a small square loop and expanding in powers of its size, we obtain:
$`\widehat{U}(x,\mathrm{\Delta }x^\mu )\widehat{U}(x+\mathrm{\Delta }x^\mu ,\mathrm{\Delta }x^\nu )\widehat{U}^+(x+\mathrm{\Delta }x^\nu ,\mathrm{\Delta }x^\mu )\widehat{U}^+(x,\mathrm{\Delta }x^\nu )=1+\mathrm{i}\widehat{F}_{\mu \nu }\mathrm{\Delta }x^\mu \mathrm{\Delta }x^\nu +\mathrm{},`$ (23)
where $`\widehat{F}_{\mu \nu }=_{[\mu }\widehat{A}_{\nu ]}\mathrm{i}[\widehat{A}_\mu ,\widehat{A}_\nu ]`$.
Obviously for the twoโform gauge connection everything works in a similar way if we substitute the line holonomy for quadratic matrices with the surface holonomy for cubic ones. For example, it seems to be natural to postulate the following transformation for the areaโordered exponent over the disc $`D`$:
$`\left(E_{I,\kappa }^{{\scriptscriptstyle _D\widehat{\stackrel{~}{B}}_{\mu \nu }๐x_\mu ๐x_\nu }}\right)_{j(s)}=\left(E_{I,\kappa }^{{\scriptscriptstyle _D\widehat{B}_{\mu \nu }๐x_\mu ๐x_\nu }}\right)_{k(s)}\left(e^{\mathrm{i}_D๐s\dot{x}^\mu \widehat{A}_\mu }\right)_{}^{k(s)}{}_{j(s)}{}^{},`$ (24)
where $`s`$ is the parametrization of the boundary and $`j(s)`$ and $`k(s)`$ are the index functions at the boundary $`D`$. Here
$`\left(e^{\mathrm{i}_D๐s\dot{x}^\mu \widehat{A}_\mu }\right)_{}^{k(s)}{}_{j(s)}{}^{}=\underset{|\mathrm{\Delta }x|0,L\mathrm{}}{lim}{\displaystyle \underset{a=1}{\overset{L}{}}}\left(e^{\mathrm{i}\widehat{A}_\mu (x_a)\mathrm{\Delta }x^\mu }\right)_{}^{k_a}{}_{j_a}{}^{}`$ (25)
is a kind of a hedgehog line โ โWilson lineโ with all free (without contraction) indices of all $`\widehat{A}`$โs at each $`s[0,2\pi )`$. The exponent here is just the standard one for quadratic matrices.
However the $`B`$โfield can not transform in this way. In fact, consider such a transformation for a small square with four external indices (see fig. 8):
$`\left(E_{I,\kappa }^{\widehat{\stackrel{~}{B}}_{\mu \nu }[x]\mathrm{\Delta }x^\mu \mathrm{\Delta }x^\nu }\right)_{k_1k_2n}\left(E_{I,\kappa }^{\widehat{\stackrel{~}{B}}_{\mu \nu }\left[x\right]\mathrm{\Delta }x^\mu \mathrm{\Delta }x^\nu }\right)_{}^{n}{}_{k_3k_4}{}^{}=\left(E_{I,\kappa }^{\widehat{B}_{\mu \nu }[x]\mathrm{\Delta }x^\mu \mathrm{\Delta }x^\nu }\right)_{j_1j_2n}\times `$
$`\times \left(E_{I,\kappa }^{\widehat{B}_{\mu \nu }\left[x\right]\mathrm{\Delta }x^\mu \mathrm{\Delta }x^\nu }\right)_{}^{n}{}_{j_3j_4}{}^{}\left(e^{\mathrm{i}\widehat{A}_\mu \mathrm{\Delta }x^\mu }\right)_{}^{j_1}{}_{k_1}{}^{}\left(e^{\mathrm{i}\widehat{A}_\nu \mathrm{\Delta }x^\nu }\right)_{}^{j_2}{}_{k_2}{}^{}\left(e^{\mathrm{i}\widehat{A}_\mu \mathrm{\Delta }x^\mu }\right)_{}^{j_3}{}_{k_3}{}^{}\left(e^{\mathrm{i}\widehat{A}_\nu \mathrm{\Delta }x^\nu }\right)_{}^{j_4}{}_{k_4}{}^{}.`$ (26)
Expanding this expression in powers of $`\mathrm{\Delta }x`$, we obtain:
$`2I_{k_1k_2k_3k_4|m_3m_2m_1}\stackrel{~}{B}_{\mu \nu }^{m_1m_2m_3}\mathrm{\Delta }x^\mu \mathrm{\Delta }x^\nu =2I_{k_1k_2k_3k_4|m_3m_2m_1}B_{\mu \nu }^{m_1m_2m_3}\mathrm{\Delta }x^\mu \mathrm{\Delta }x^\nu +`$
$`+\mathrm{i}I_{j_1j_2j_3j_4}\left[\left(A_\mu \mathrm{\Delta }x^\mu \right)_{k_1}^{j_1}\delta _{k_2}^{j_2}\delta _{k_3}^{j_3}\delta _{k_4}^{j_4}\delta _{k_1}^{j_1}\delta _{k_2}^{j_2}\left(A_\mu \mathrm{\Delta }x^\mu \right)_{k_3}^{j_3}\delta _{k_4}^{j_4}\right]+`$
$`+\mathrm{i}I_{j_1j_2j_3j_4}\left[\delta _{k_1}^{j_1}\left(A_\nu \mathrm{\Delta }x^\nu \right)_{k_2}^{j_2}\delta _{k_3}^{j_3}\delta _{k_4}^{j_4}\delta _{k_1}^{j_1}\delta _{k_2}^{j_2}\delta _{k_3}^{j_3}\left(A_\nu \mathrm{\Delta }x^\nu \right)_{k_4}^{j_4}\right]+\mathrm{}`$ (27)
We do not spell here the higher terms in $`A`$ because they are not relevant to observe that there is no way to transform the $`B`$โfield to compensate the linear terms in $`\mathrm{\Delta }x`$.
Then we propose that the areaโordered exponent transforms with the use of the surface exponent:
$`\left(E_{I,\kappa }^{{\scriptscriptstyle _D\widehat{\stackrel{~}{B}}_{\mu \nu }๐x_\mu ๐x_\nu }}\right)_{j(s)}=\left(E_{I,\kappa }^{{\scriptscriptstyle _D\widehat{B}_{\mu \nu }๐x_\mu ๐x_\nu }}\right)_{k(s)}\left(\stackrel{~}{E}_{I,\kappa }^{\mathrm{i}_D๐s\dot{x}^\mu \widehat{A}_\mu }\right)_{}^{k(s)}{}_{j(s)}{}^{},`$ (28)
where we consider it more natural to integrate the quadratic matrix $`\widehat{A}_\mu =A_\mu ^{ab}`$ (rather than a cubic one) over the boundary of the disc $`D`$. Hence in this expression we are using the exponent defined in the subsection 3.2.
Now if we take the disc as in the fig. 8 and expand both sides of the corresponding expression (analogous to eq.(26) ) in powers of $`\mathrm{\Delta }x`$, we observe that all the linear terms (as in eq.(27) ) do cancel out and no other dangerous terms appear. As the result the $`B`$โfield transforms as:
$`\stackrel{~}{B}_{\mu \nu }^{ijk}=B_{\mu \nu }^{ijk}+\mathrm{i}{\displaystyle \frac{_{[\mu }A_{\nu ]}^{}{}_{l}{}^{i}}{3}}I^{ljk}+\mathrm{i}{\displaystyle \frac{_{[\mu }A_{\nu ]}^{}{}_{l}{}^{j}}{3}}I^{ilk}+\mathrm{i}{\displaystyle \frac{_{[\mu }A_{\nu ]}^{}{}_{l}{}^{k}}{3}}I^{ijl},`$ (29)
i.e. as a collection of the Abelian stringy $`B`$โfields for each group of indices $`i,j,k`$. This is not a surprise for us once we understood the main properties of the surface exponent, which are discussed at the end of the section 2.
We can calculate as well the curvature of the โnonโAbelianโ $`B`$โfield. To do that we have to consider a small closed surface having the geometry of the cube (see fig. 9). Then the corresponding holonomy matrix has the following expansion:
$`\mathrm{Tr}[\widehat{U}(x,\mathrm{\Delta }x^\mu \mathrm{\Delta }x^\nu )\widehat{U}(x+\mathrm{\Delta }x^\nu ,\mathrm{\Delta }x^\mu \mathrm{\Delta }x^\alpha )\widehat{U}(x,\mathrm{\Delta }x^\nu \mathrm{\Delta }x^\alpha )\times `$
$`\times \widehat{U}^1(x,\mathrm{\Delta }x^\nu \mathrm{\Delta }x^\mu )\widehat{U}^1(x+\mathrm{\Delta }x^\mu ,\mathrm{\Delta }x^\alpha \mathrm{\Delta }x^\nu )\widehat{U}^1(x+\mathrm{\Delta }x^\alpha ,\mathrm{\Delta }x^\nu \mathrm{\Delta }x^\alpha )]=`$
$`=N+I_{ijk}_{[\mu }B_{\nu \alpha ]}^{kji}\mathrm{\Delta }x^\mu \mathrm{\Delta }x^\nu \mathrm{\Delta }x^\alpha +\mathrm{},`$ (30)
where each of the $`\widehat{U}`$ matrices corresponds to the rectangle and represents the following product of its triangular constituents:
$`U_{j_1j_2j_3j_4}(x,\mathrm{\Delta }x^\mu \mathrm{\Delta }x^\nu )=\left(E_{I,\kappa }^{\widehat{B}_{\mu \nu }[x]\mathrm{\Delta }x^\mu \mathrm{\Delta }x^\nu }\right)_{j_1j_2n}\left(E_{I,\kappa }^{\widehat{B}_{\mu \nu }\left[x\right]\mathrm{\Delta }x^\mu \mathrm{\Delta }x^\nu }\right)_{}^{n}{}_{j_3j_4}{}^{}`$ (31)
Thus, we have the following curvature for the twoโtensor field:
$`F_{\mu \nu \alpha }^{ijk}=_{[\mu }B_{\nu \alpha ]}^{ijk}.`$ (32)
As the result, the corresponding local field theory is free. Again this is not a surprise for us. Somehow once we consider the local fields as nonโlinear variables ($`\widehat{B}:V^3`$), the theory for them itself becomes linear! As we have already mentioned, however, the 2D holonomy matrices
$$\left(E_{I,\kappa }^{{\scriptscriptstyle _\mathrm{\Sigma }\widehat{B}_{\mu \nu }๐x_\mu ๐x_\nu }}\right)_{\mathrm{}}^{\mathrm{}}$$
include interactions due to the nonโlinear nature of the $`B`$โfield and especially of the background $`\widehat{I}`$. It is worth mentioning at this point that in the twoโtensor field was considered to carry two color indices, i.e. to be a linear field ($`B_{\mu \nu }^{ij}\mathrm{\Delta }x^\mu \mathrm{\Delta }x^\nu :VV`$). Then to organize the triangulationโindependent areaโordering with the use of the standard exponent one had to include the oneโform gauge field ($`A`$) as well. As the result the theory for such a couple of fields ($`B,A`$) becomes interacting due to the presence of commutators.
Let us clarify these observations. All this idea of the surface exponent and the triangulationโindependent areaโordering was invented just to make a local description of the theory of nonโAbelian twoโtensor fields<sup>6</sup><sup>6</sup>6By โnonโAbelianโ we mean just those fields which carry color indices.. But for the local fields everything becomes commutative if we make them (and especially the background $`\widehat{I}`$) to carry three color indices. As the result there are no commutators in this case as we pointed out at the end of the section 2. And once there are no commutators, then there are no interactions!
It seems that to obtain a nonโtrivial ordering one has to throw away a line out of the surface. The corresponding line observable can be sensitive to the nonโtrivial ordering. But such a nonโlocal observable should not be constructed only from the local $`B`$โfield. As the result, it is nonโlocal theory which can be nonโcommutative, while that for local fields is commutative . Thus, at this stage the problem is not that we do not know just how to write the interacting local theory for nonโAbelian twoโtensor fields, rather we even do not know the nature of their possible interactions. What can we say?.. Probably this is just the way the nature is and we should be happy that it is so, because nonโlinear fields are free and easy to deal with! The question is whether one can construct anything nonโtrivial with the use of such free fields. We are going to discuss this point right now.
## 5 Future directions
### 5.1 Relation to the YangโBaxter equation
In this subsection we would like to point out a curious relation of the above considerations to the YangโBaxter equations in the theory of integrable lattice models. The relation is as follows (see e.g. ). Consider the surface holonomy matrix for the cube (see fig. 9). The condition of the zero curvature establishes that the cube holonomy matrix is trivial:
$`\widehat{U}(x,\mathrm{\Delta }x^\mu \mathrm{\Delta }x^\nu )\widehat{U}(x+\mathrm{\Delta }x^\nu ,\mathrm{\Delta }x^\mu \mathrm{\Delta }x^\alpha )\widehat{U}(x,\mathrm{\Delta }x^\nu \mathrm{\Delta }x^\alpha )\times `$
$`\times \widehat{U}^1(x,\mathrm{\Delta }x^\nu \mathrm{\Delta }x^\mu )\widehat{U}^1(x+\mathrm{\Delta }x^\mu ,\mathrm{\Delta }x^\alpha \mathrm{\Delta }x^\nu )\widehat{U}^1(x+\mathrm{\Delta }x^\alpha ,\mathrm{\Delta }x^\nu \mathrm{\Delta }x^\alpha )=1.`$ (33)
This equation is just the YangโBaxter equation where the matrices $`U_{j_1j_2j_3j_4}(x,\mathrm{\Delta }x^\mu \mathrm{\Delta }x^\nu )`$ in eq.(30) play the role of the Rโmatrix.
In this context it is interesting to see whether one can clasify all possible Rโmatrices according to all possible choices of $`\left(E_{I,\kappa }^{\widehat{B}}\right)_{j_1j_2j_3j_4}`$. We would like to remind at this point about the starโtriangle equality (which is well in the spirit of the considerations of our paper) and its relation to the YangโBaxter equation (see e.g. ).
Apart from other things, these considerations are important to establish Hamiltonian formalism for multiโdirectional evolution. This is the subject of the next subsection.
### 5.2 Relation to the String Field Theory
Consider for the time being the case when the target space $`X`$ is just $`R^D`$. As well let us take $`V=`$, where $``$ is the Hilbert space. Then the indices $`i,j`$ and $`k`$ take continuous values and we take them to be $`x,y`$ and $`z`$ โ the coordinates in $`X`$. Let us take as the base of our fiber bundle the worldโsheet (space of $`\sigma `$ and $`\tau `$) rather than the target space $`X`$. I.e. in this case the $`B_{\mu \nu }`$โfield is just the density $`\widehat{B}_{\sigma \tau }(\sigma ,\tau )=\widehat{B}(\sigma ,\tau )`$ with three color indices: $`\widehat{B}:^3`$.
As the result the observables of the theory are as follows:
$`U\left[\mathrm{\Sigma }^{}(\gamma _1,\mathrm{},\gamma _L)\right]_{x^{(1)}(s_1)\left|x^{(2)}(s_2)\right|\mathrm{}|x^{(n)}(s_n)}^{y^{(1)}(t_1)\left|y^{(2)}(t_2)\right|\mathrm{}|y^{(m)}(t_m)}=`$
$`=y^{(1)}(t_1)\left|y^{(2)}(t_2)\left|\mathrm{}y^{(m)}(t_m)\left|E_{g,I,\kappa }^{{\scriptscriptstyle _\mathrm{\Sigma }๐\sigma ๐\tau \widehat{B}(\sigma ,\tau )}}\right|x^{(1)}(s_1)\right|x^{(2)}(s_2)\mathrm{}\right|x^{(n)}(s_n).`$ (34)
Here $`|x(s)=_s|x_s`$ and $`|x_s`$ is the standard coherent state in quantum mechanics. The observables (34) represent interacting string theory amplitudes for the proper choice of $`\widehat{I}`$ and $`\widehat{\kappa }`$ and $`\widehat{B}(\sigma ,\tau )`$. This is a much better situation with respect to the standard String Field Theory Hamiltonian quantization , where one starts with the free string theory amplitude:
$`U(\mathrm{cylinder})=x(\sigma )\left|e^{HT}\right|y(\sigma ),\mathrm{where}H={\displaystyle _0^{2\pi }}๐\sigma \left({\displaystyle \frac{\delta ^2}{\delta x^2(\sigma )}}+{\displaystyle \frac{\left[_\sigma x(\sigma )\right]^2}{2\pi }}\right).`$ (35)
Only after that one includes interactions as perturbative expansion in the cubic interaction over the quadratic background<sup>7</sup><sup>7</sup>7We do not pay attention to the constraints and ghosts at this stage.. The latter formalism suffers from the well known background dependence, apart from many other things. We see, however, that this problem does not appear in our formalism which is intrinsically cubic, i.e. does not have to be expanded over a quadratic part! To make this statement completely rigorous one has to establish the following relation explicitly:
$`x(\sigma )\left|E_{I,\kappa }^{{\scriptscriptstyle _\mathrm{\Sigma }๐\sigma ๐\tau \widehat{B}(\sigma ,\tau )}}\right|y(\sigma )=x(\sigma )\left|e^{HT}\right|y(\sigma ).`$ (36)
I.e. for given $`H`$ one has to find explicit values for $`\widehat{I}`$, $`\widehat{\kappa }`$ and $`\widehat{B}(\sigma ,\tau )`$ in this equation for the arbitrary choice of $`x(\sigma )`$ and $`y(\sigma )`$. Which is going to be done in a separate publication.
It is worth mentioning at this point that the considerations of this section along with those of the subsection 3.3. bring us very close to the subject of simplicial strings and open/closed string duality in the spirit of ,, and . As well we think that these considerations will elaborate on the nature of the integrability in the large $`N_c`$ YangโMills theory .
## 6 Conclusions and Acknowledgments
Thus, we have obtained local field theory description of the nonโAbelian twoโtensor fields. Due to the triviality of the corresponding homotopy groups in more than oneโdimension, there are no commutator terms for the local twoโtensor fields. Once there are no commutators, then there are no interactions, i.e. we obtain a free field theory for the local nonโAbelian twoโtensor fields. (We call them nonโAbelian because they carry color indices.) However, this does not mean that theory for them does not have any nonโtrivial content. In fact, because the background $`I_{ijk}`$ and the field $`B_{\mu \nu }^{ijk}`$ carry three color indices the 2D holonomy matrices in this theory describe nonโtrivial maps between fibers of the bundles over loop spaces.
Due to the latter fact we find a promising link of our constructions to the background independent formulation of String Field Theory. At the same time our theory is linked with the Rโmatrix and the integrability in the lattice models. These observations all together form a knot which remains to be disentangled for the better understanding of all mentioned subjects.
Apart from the future directions listed in the main body of the text we would like to bring attention to another important problem. We would like to find a reduction of the nonโAbelian twoโtensor fields to the nonโAbelian oneโform gauge fields. If the surface, over which the areaโordering is done, is of the cylindrical topology, it can be degenerated into a curve in the target space. In this case we can expect to reproduce somehow the path ordering of a oneโform gauge field over the curve. However, this problem is much more complicated in comparison with the one presented in the subsection 5.2. In fact, if we deal with a oneโform gauge field $`A_\mu `$, $`\mu =1,\mathrm{},D`$ for $`D2`$ the ordering over the curve is nonโtrivial. And the question is how to reproduce such an ordering from the trivial one for the twoโtensor fields. This is a problem for a separate study.
I would like to acknowledge valuable discussions with U.Schreiber, T.Tada, T.Pilling, N.Amburg, D.Vasiliev, A.Losev, A.Rosly and especially with A.Gerasimov, G.Sharigin, V.Dolotin, S.Loktev and A.Morozov. I would like to thank Tada san and Ishimoto san for very worm hospitality at RIKEN, Tokyo, where this work was finished. This work was done under the partial support of grants RFBR 04-01-00646, INTAS 03โ51โ5460 and the Grant from the President of Russian Federation MKโ2097.2004.2.
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# Appendix
## Appendix
We describe here the remaining GQSs for the Lie superalgebra $`sl(m|n)`$ with odd CAOs only. According to , there are two classes. For the first class, $`l`$ can be any index between $`1`$ and $`m1`$, so assume that $`l`$ is fixed ($`1l<m`$). The CAOs are then described by the root vectors of \[1, eq. (3.9)\], but in order to deduce solutions for the CCs we need to multiply them by some overall constant. This gives, for $`k=1,\mathrm{},m`$ and $`r=1,\mathrm{},n`$:
$$x_{rk}^+=\{\begin{array}{ccc}\sqrt{|2mn2l|}e_{m+r,k}\hfill & \text{for}& kl\\ \sqrt{|n2l|}e_{k,m+r}\hfill & \text{for}& k>l\end{array}$$
(21)
and
$$x_{rk}^{}=\{\begin{array}{ccc}\sqrt{|2mn2l|}e_{k,m+r}\hfill & \text{for}& kl\\ ฯต\sqrt{|n2l|}e_{m+r,k}\hfill & \text{for}& k>l\end{array}$$
(22)
where $`ฯต=\text{sgn}((n2l)(2mn2l))`$. Of course, we have to assume that $`l`$ is such that these factors do not vanish, i.e. $`(n2l)(2mn2l)0`$. Then, one can deduce that
$$\underset{r=1}{\overset{n}{}}\underset{k=1}{\overset{m}{}}[\{x_{rk}^+,x_{rk}^{}\},x_{sj}^\pm ]=\nu n(mn)x_{sj}^\pm $$
(23)
where $`\nu =\text{sgn}(2mn2l)`$. Clearly, for $`mn`$ such systems provide solutions for the CCs for the $`N`$-particle $`D`$-dimensional oscillator whenever $`mn=DN`$.
For the second class, $`l`$ can be any index between $`1`$ and $`n1`$. Now the CAOs are described by the root vectors of \[1, eq. (3.8)\], again multiplied by some appropriate constant. This gives, for $`k=1,\mathrm{},m`$ and $`r=1,\mathrm{},n`$:
$$x_{rk}^+=\{\begin{array}{ccc}\sqrt{|2nm2l|}e_{m+r,k}\hfill & \text{for}& rl\\ \sqrt{|m2l|}e_{k,m+r}\hfill & \text{for}& r>l\end{array}$$
(24)
and
$$x_{rk}^{}=\{\begin{array}{ccc}\sqrt{|2nm2l|}e_{k,m+r}\hfill & \text{for}& rl\\ ฯต\sqrt{|m2l|}e_{m+r,k}\hfill & \text{for}& r>l\end{array}$$
(25)
where $`ฯต=\text{sgn}((m2l)(2nm2l))`$. Again we assume that $`l`$ is such that these factors do not vanish, i.e. $`(m2l)(2nm2l)0`$. Now one can deduce that
$$\underset{r=1}{\overset{n}{}}\underset{k=1}{\overset{m}{}}[\{x_{rk}^+,x_{rk}^{}\},x_{sj}^\pm ]=\nu m(nm)x_{sj}^\pm $$
(26)
where $`\nu =\text{sgn}(2nm2l)`$. For $`mn`$ such systems provide another class of solutions for the CCs for the $`N`$-particle $`D`$-dimensional oscillator whenever $`mn=DN`$.
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# On the NP-Completeness of Some Graph Cluster Measures
## 1 Introduction
Clustering is an important issue in the analysis and exploration of data. There is a wide area of applications in data mining, VLSI design, parallel computing, web searching, software engineering, computer graphics, gene analysis, etc. See also for an overview. Intuitively clustering consists in discovering natural groups (clusters) of similar elements in data set. An important variant of data clustering is graph clustering where the similarity relation is expressed by a graph. In this paper, we restrict to unweighted, undirected graphs with no self-loops.
We first recall some basic definitions from graph theory. Let $`G=(V,E)`$ be an *undirected* graph and denote by $`E(S)=\{\{u,v\}E;u,vS\}`$ the set of edges in a *subgraph* $`G(S)=(S,E(S))`$ *induced* by a subset of vertices $`SV`$. We say that $`SV`$ creates a *clique* of *size* $`|S|`$ if edges in $`E(S)=\{\{u,v\};u,vS,uv\}`$ join every two different vertices in $`S`$. Further denote by $`d_G(v)=|\{uV;`$ $`\{u,v\}E\}|`$ the *degree* of vertex $`vV`$ in $`G`$. We say that graph $`G`$ is a *cubic* graph if $`d_G(v)=3`$ for every $`vV`$. Moreover, any subset of vertices $`AV`$ creates a *cut* of $`G`$, that is a partition of $`V`$ into disjoint sets $`A`$ and $`VA`$. The *size* of cut $`A`$ is defined as
$$c_G(A)=\left|\{\{u,v\}E;uA,vVA\}\right|,$$
(1)
and
$$d_G(S)=\underset{vS}{}d_G(v)$$
(2)
denotes the sum of degrees in cut $`SV`$.
A canonical definition of a graph cluster does not exist, but it is commonly agreed that a cluster should be a connected subgraph induced by a vertex set $`S`$ with many internal edges $`E(S)`$ and few edges to outside vertices in $`VS`$ . In this paper we consider several locally computable fitness functions that are used for measuring the quality of a cluster within the graph. The prominent position among graph cluster measures is occupied by the *conductance* which is defined for any cut $`\mathrm{}SV`$ in graph $`G`$ as follows
$$\mathrm{\Phi }_G(S)=\frac{c_G(S)}{\mathrm{min}(d_G(S),d_G(VS))}.$$
(3)
Furthermore, the *local density* $`\delta _G(S)`$ (cf. the *average degree* ) of a subset $`\mathrm{}SV`$ in graph $`G`$ is the ratio of the number of edges in subgraph $`G(S)`$ induced by $`S`$ over the number of edges in a clique of size $`|S|`$ vertices, that is
$$\delta _G(S)=\frac{|E(S)|}{\left(\genfrac{}{}{0pt}{}{|S|}{2}\right)}=\frac{2|E(S)|}{|S|(|S|1)}$$
(4)
for $`S`$ containing at least two vertices whereas define $`\delta _G(S)=0`$ for $`|S|=1`$. Similarly, we define the *relative density* of cut $`\mathrm{}SV`$ as follows
$$\varrho _G(S)=\frac{|E(S)|}{|E(S)|+c_G(S)}.$$
(5)
Yet another graph cluster measure which we call *single cluster editing* (cf. ) of a subset $`SV`$ counts the number of edge operations (both additions and deletions) needed to transform $`S`$ into an isolated clique:
$$\epsilon _G(S)=\left(\genfrac{}{}{0pt}{}{|S|}{2}\right)|E(S)|+c_G(S).$$
(6)
Proposed clustering algorithms usually search for clusters that are optimal with respect to the above-mentioned fitness measures. Therefore the underlying optimization problems of finding the clusters that minimize the conductance or maximize the densities or that need a small single cluster editing are of special interest. In this paper we will formally prove that the associated decision problems for the conductance (Section 2), local and relative densities (Section 3), and single cluster editing (Section 4) are NP-complete. These complexity results appear to be well-known or at least intuitively credible, but not properly documented in the literature.
## 2 Conductance
Finding a subset of vertices that has the minimum conductance in a given graph has been often stated to be an NP-complete problem in the literature . However, we could not find an explicit proof anywhere. For example, the NP-completeness proof due to Papadimitrou for the problem of finding the minimum *normalized cut* which is in fact the conductance of a weighted graph does not imply the hardness in the unweighted case. Thus we provide the proof in this section. The decision version for the conductance problem is formulated as follows:
Minimum Conductance (Conductance)
*Instance:* An undirected graph $`G=(V,E)`$ and positive integer $`\varphi `$.
*Question:* Is there a cut $`SV`$ such that $`\mathrm{\Phi }_G(S)\varphi `$?
###### Theorem 1
Conductance is NP-complete.
Proof: Clearly, Conductance belongs to NP since a nondeterministic algorithm can guess a cut $`SV`$ and verify $`\mathrm{\Phi }_G(S)\varphi `$ in polynomial time. For the NP-hardness proof the following maximum cut problem on cubic graphs will be reduced to Conductance in polynomial time.
Maximum Cut for Cubic Graphs (Max Cutโ3)
*Instance:* A cubic graph $`G=(V,E)`$ and positive integer $`a`$.
*Question:* Is there a cut $`AV`$ such that $`c_G(A)a`$?
The Max Cutโ3 problem was first stated to be NP-complete in which became a widely used reference although an explicit proof cannot be found there and we were unable to reconstruct the argument from the sketch. Nevertheless, the NP-completeness of Max Cutโ3 follows from its APX-completeness presented in . The following reduction to Conductance is adapted from that used for the minimum edge expansion problem .
Given a Max Cutโ3 instance, i.e. a cubic graph $`G=(V,E)`$ with $`n=|V|`$ vertices, and positive integer $`a`$, a corresponding undirected graph $`G^{}=(V^{},E^{})`$ for Conductance is composed of two fully connected copies of the complement of $`G`$, that is $`V^{}=V_1V_2`$ where $`V_i=\{v^i;vV\}`$ for $`i=1,2`$, and $`E^{}=E_1E_2E_3`$ where $`E_i=\{\{u^i,v^i\};u,vV,uv,\{u,v\}E\}`$ for $`i=1,2`$, and $`E_3=\{\{u^1,v^2\};u,vV\}`$. In addition, define the required conductance bound
$$\varphi =\frac{1}{2n4}\left(n\frac{2a}{n}\right).$$
(7)
The number of vertices in $`G^{}`$ is $`|V^{}|=2n`$ and the number of edges $`|E^{}|=(2n4)n`$ since
$$d_G^{}(v)=2n4\text{for every }vV^{}$$
(8)
due to $`G`$ is a cubic graph. It follows that $`G^{}`$ can be constructed in polynomial time.
For a cut $`\mathrm{}SV^{}`$ in $`G^{}`$ with $`k=|S|2n`$ vertices denote by
$$S_i=\{vV;v^iS\}\text{for }i=1,2$$
(9)
the cuts in $`G`$ that are projections of $`S`$ to $`V_1`$ and $`V_2`$, respectively. Since $`c_G^{}(S)=c_G^{}(V^{}S)`$ it holds $`\mathrm{\Phi }_G^{}(S)=\mathrm{\Phi }_G^{}(V^{}S)`$ according to definition (3). Hence, $`kn`$ can be assumed without loss of generality when computing the conductance in $`G^{}`$. Thus,
$$\mathrm{\Phi }_G^{}(S)=\frac{|S||V^{}S|c_G(S_1)c_G(S_2)}{(2n4)|S|}$$
(10)
follows from condition (8) and the fact that $`G^{}`$ is composed of two fully connected complements of $`G`$, which can be rewritten as
$$\mathrm{\Phi }_G^{}(S)=\frac{1}{2n4}\left(2nk\frac{c_G(S_1)+c_G(S_2)}{k}\right).$$
(11)
Now we verify the correctness of the reduction by proving that the Max Cutโ3 instance has a solution if and only if the corresponding Conductance instance is solvable. First assume that a cut $`AV`$ exists in $`G`$ whose size satisfies
$$c_G(A)a.$$
(12)
Denote by
$$S^A=\{v^1V_1;vA\}\{v^2V_2;vVA\}V^{}$$
(13)
the cut in $`G^{}`$ whose projections (9) to $`V_1`$ and $`V_2`$ are $`S_1^A=A`$ and $`S_2^A=VA`$, respectively. Since $`|S^A|=n`$ and $`c_G(A)=c_G(VA)`$ the conductance of $`S^A`$ can be upper bounded as
$$\mathrm{\Phi }_G^{}\left(S^A\right)=\frac{1}{2n4}\left(n\frac{2c_G(A)}{n}\right)\frac{1}{2n4}\left(n\frac{2a}{n}\right)=\varphi $$
(14)
according to equations (11), (12), and (7), which shows that $`S^A`$ is a solution of the Conductance instance.
For the converse, assume that the conductance of cut $`\mathrm{}SV^{}`$ in $`G^{}`$ meets
$$\mathrm{\Phi }_G^{}(S)\varphi .$$
(15)
Let $`AV`$ be the maximum cut in $`G`$. For cut $`S^A`$ defined according to (13) we prove that
$$\mathrm{\Phi }_G^{}\left(S^A\right)\mathrm{\Phi }_G^{}(S)$$
(16)
which is rewritten to
$$\frac{1}{2n4}\left(n\frac{2c_G(A)}{n}\right)\frac{1}{2n4}\left(2nk\frac{c_G(S_1)+c_G(S_2)}{k}\right)$$
(17)
according to (14) and (11) where $`k=|S|n`$ and $`S_1,S_2`$ are defined in (9). Since $`2c_G(A)c_G(S_1)+c_G(S_2)`$ due to $`A`$ is the maximum cut in $`G`$, it suffices to show
$$nk+\left(\frac{1}{n}\frac{1}{k}\right)(c_G(S_1)+c_G(S_2))0$$
(18)
which follows from $`\frac{1}{n}\frac{1}{k}0`$ and $`c_G(S_1)+c_G(S_2)|S_1|n+|S_2|n=kn`$. Thus,
$$\frac{1}{2n4}\left(n\frac{2c_G(A)}{n}\right)=\mathrm{\Phi }_G^{}\left(S^A\right)\mathrm{\Phi }_G^{}(S)\varphi =\frac{1}{2n4}\left(n\frac{2a}{n}\right)$$
(19)
holds according to (14), (16), (15), and (7), which implies $`c_G(A)a`$. Hence, $`A`$ solves the MAX CUT-3 instance. $`\mathrm{}`$
## 3 Local and Relative Density
The decision version of the maximum density problem is formulated as follows:
Maximum Density (Density)
*Instance:* An undirected graph $`G=(V,E)`$, positive integer $`k|V|`$, and a rational number $`0r1`$.
*Question:* Is there a subset $`SV`$ such that $`|S|=k`$ and the density of $`S`$ in $`G`$ is at least $`r`$?
We distinguish between Local Density and Relative Density problems according to the particular density measure used which is the local density (4) and the relative density (5), respectively. Clearly, Local Density is NP-complete since this problem for $`r=1`$ coincides with the NP-complete Clique problem . Also the NP-completeness of Relative Density can easily be achieved:
###### Theorem 2
Relative Density is NP-complete.
Proof: Obviously, Relative Density belongs to NP since a nondeterministic algorithm can guess a cut $`SV`$ of cardinality $`|S|=k`$ and verify $`\varrho _G(S)r`$ in polynomial time. For the NP-hardness proof the following minimum bisection problem on cubic graphs which is known to be NP-complete will be reduced to Relative Density in polynomial time.
Minimum Bisection for Cubic Graphs (Min Bisectionโ3)
*Instance:* A cubic graph $`G=(V,E)`$ with $`n=|V|`$ vertices and positive integer $`a`$.
*Question:* Is there a cut $`SV`$ such that $`|S|=\frac{n}{2}`$ and $`c_G(S)a`$?
Given a Min Bisectionโ3 instance, i.e. a cubic graph $`G=(V,E)`$ with $`n=|V|`$ vertices, and positive integer $`a`$, a corresponding Relative Density instance consists of the same graph $`G`$, parameters $`k=\frac{n}{2}`$ and
$$r=\frac{3n2a}{3n+2a}.$$
(20)
Now for any subset $`SV`$ such that $`|S|=k=\frac{n}{2}`$ it holds
$$|E(S)|=\frac{3|S|c_G(S)}{2}=\frac{3n2c_G(S)}{4}$$
(21)
due to $`G`$ is a cubic graph, which gives
$$\varrho _G(S)=\frac{3n2c_G(S)}{3n+2c_G(S)}$$
(22)
according to (5). It follows from (20) and (22) that $`\varrho _G(S)r`$ iff $`c_G(S)a`$$`\mathrm{}`$
## 4 Single Cluster Editing
The problem of deciding whether a given graph can be transformed into a collection of cliques using at most $`m`$ edge operations (both additions and deletions) which is called Cluster Editing is known to be NP-complete . When the desired solution must contain exactly $`p`$ cliques, the so called pโCluster Editing problem remains NP-complete for every $`p2`$. Here we study the issue of whether a given graph contains a subset $`S`$ of exactly $`k`$ vertices such that at most $`m`$ edge additions and deletions suffice altogether to turn $`S`$ into an isolated clique:
Minimum Single Cluster Editing (1โCluster Editing)
*Instance:* An undirected graph $`G=(V,E)`$, positive integers $`k|V|`$ and $`m`$.
*Question:* Is there a subset $`SV`$ such that $`|S|=k`$ and $`\epsilon _G(S)m`$?
###### Theorem 3
1โCluster Editing is NP-complete.
Proof: Obviously, 1โCluster Editing belongs to NP since a nondeterministic algorithm can guess a subset $`SV`$ of cardinality $`|S|=k`$ and verify $`\epsilon _G(S)m`$ in polynomial time. For the NP-hardness proof the Min Bisectionโ3 problem is used again (cf. the proof of Theorem 2) which will be reduced to 1โCluster Editing in polynomial time.
Given a Min Bisectionโ3 instance, i.e. a cubic graph $`G=(V,E)`$ with $`n=|V|`$ vertices, and positive integer $`a`$, a corresponding 1โCluster Editing instance consists of the same graph $`G`$, parameters $`k=\frac{n}{2}`$ and
$$m=\frac{12a+n(n8)}{8}.$$
(23)
Now for any subset $`SV`$ such that $`|S|=k=\frac{n}{2}`$ it holds
$$\epsilon _G(S)=\frac{|S|(|S|1)}{2}\frac{3|S|c_G(S)}{2}+c_G(S)=\frac{12c_G(S)+n(n8)}{8}$$
(24)
according to (6) and (21). It follows from (23) and (24) that $`\epsilon _G(S)m`$ iff $`c_G(S)a`$$`\mathrm{}`$
## 5 Conclusion
In this paper we have presented the explicit NP-completeness proofs for the decision problems associated with the optimization of four possible graph cluster measures; namely the conductance, the local and relative densities, and single cluster editing. In clustering algorithms, combinations of fitness measures are often preferred as only optimizing one may result in anomalies such as selecting small cliques or connected components as clusters. An open problem is the complexity of minimizing the product of the local and relative densities (e.g. their sum is closely related to the edge operation count for the single cluster editing problem). Another important area for further research is the complexity of finding related approximation solutions .
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# Marginally trapped tubes and dynamical horizons
## I Motivation
Over the past decade, new definitions of black hole horizons have emerged which provide powerful tools for studying the behavior of black holes in the strongly dynamical regime. These ideas share the common philosophy that black holes should be thought of as physical objects in a spacetime that may be identified by local measurements. By contrast, traditional black holes and event horizons are globally defined properties of the causal structure of a spacetime hawkingellis72 ; wald .
Though there has always been a certain amount of interest in the dynamics of apparent horizons and their relation to black hole physics (see for example collins ), Hayward began this line of research in earnest with his definition of trapping horizons. These were initially used to formulate dynamic versions of the laws of black hole mechanics in a quasi-local context hayward94a . Since then however, they have found a variety of applications including, for example, studying interactions between black holes and gravitational waves haywardPert1 ; haywardGravWav . A bit later, the isolated horizons of Ashtekar et al. ashtekar ; ashtekar02a were developed to provide a quasi-local characterization of the equilibrium states of black holes. In those works, it was shown that these objects obey a phase space formulation of the zeroth and first laws of black hole mechanics. Further, loop quantum gravity has made important use of them as boundary conditions in calculations of black hole entropy entropy .
Closely related to both trapping and isolated horizons are dynamical horizons ashtekar02b ; ashtekar03a , which characterize the dynamical phase of smooth black hole evolutions. In particular, it has been shown that flux laws can be formulated for these horizons that measure the growth of such quantities as energy and entropy. These laws contain terms that may be identified with fluxes of particular physical quantities such as matter and gravitational waves. A similar law exists for trapping horizons conservation and it has been shown that the energy expressions involved agree with those derived in a Hamiltonian analysis of horizons as spacetime boundaries thebeast05 .
In other developments, studies have been made of the perturbative, โalmost isolatedโ regime haywardPert1 ; haywardPert2 ; booth04a and isolated horizons have been given a convenient characterization in terms of multipole moments, which can be extended to a definition for the multipole moments for dynamical and trapping horizons multipoles . Mathematical investigations have been made into such properties as existence and uniqueness abhaygreg ; ams05 and recently these notions have also begun to find application in the physical interpretation of numerical results num1 ; num2 ; num3 .
Thus, isolated, dynamical and trapping horizons constitute an increasingly well-developed quasi-local framework for analytical as well as numerical studies of black hole dynamics in the strong field regime. However, in order to make effective use of them it is important to have clear intuition about how they behave. Unfortunately, there is a paucity of good, analytical examples of spacetimes containing dynamical and/or trapping horizons and this has recently given rise to some confusion about the generic behavior of dynamical and trapping horizons. In this paper we will attempt to clarify matters by presenting several analytic and numerical examples of spacetimes containing these horizons.
Before considering these in more detail, let us recall the relevant definitions. From Hayward hayward94a we have:
Definition 1. A trapping horizon $`H`$ is a hypersurface in a 4-dimensional spacetime that is foliated by 2-surfaces (which we will take to be of spherical topology) such that $`\theta _{(\mathrm{})}|_H=0`$, $`\theta _{(n)}|_H0`$ and $`_n\theta _{(\mathrm{})}|_H0`$. A trapping horizon is called outer if $`_n\theta _{(\mathrm{})}|_H<0`$, inner if $`_n\theta _{(\mathrm{})}|_H>0`$, future if $`\theta _{(n)}|_H<0`$ and past if $`\theta _{(n)}|_H>0`$.
In this definition, and what follows, $`\mathrm{}^a`$ and $`n^a`$ are respectively the future-directed outgoing and ingoing null normals to a leaf of the foliation while $`\theta _{(\mathrm{})}`$ and $`\theta _{(n)}`$ are the expansion of the congruences of curves generated by those vector fields. Further, it is assumed that $`n_a`$ has been extended so that it is surface generating (ie. $`ndn=0`$) in some neighbourhood of $`H`$. Then, from the definition, outer trapping horizons have trapped surfaces โjust insideโ them while inner trapping horizons have trapped surfaces โjust outsideโ.
We will mainly be interested in a slight generalization of future trapping horizons that was recently introduced by Ashtekar and Galloway abhaygreg :
Definition 2. A marginally trapped tube (MTT), $`T`$ is a hypersurface in a 4-dimensional spacetime that is foliated by two-surfaces (again assumed to be of spherical topology) such that $`\theta _{(n)}|_T<0`$ and $`\theta _{(\mathrm{})}|_T=0`$.
We refer to the leaves of the foliation as marginally trapped surfaces.<sup>1</sup><sup>1</sup>1This follows the terminology used in abhaygreg . Note however that the exact definition of marginally trapped surfaces varies somewhat in the literature. For example, Wald wald defines a marginally trapped surface to be a compact, spacelike, two-surface for which both null-expansions are non-positive. MTTs have no restriction on their signature, which is allowed to vary over the hypersurface. However, if an MTT is everywhere spacelike it is referred to as a *dynamical horizon* ashtekar02a ; ashtekar03a , if it is everywhere timelike then it called a *timelike membrane* (TLM) ashtekar03a ; abhaygreg , and if it is everywhere null and non-expanding then we have an *isolated horizon*.<sup>2</sup><sup>2</sup>2More precisely, this is a non-expanding horizon of which isolated horizons are a special case. However, we will follow the common usage of isolated horizon in its general rather than specific meaning throughout this paper.
Note that the distinction between dynamical horizons and timelike membranes is more than just a technical difference of signature. For example, it is clear that since dynamical horizons are spacelike they may only be crossed in one direction by causal curves; this is a key characteristic of black hole horizons. By contrast, a timelike membrane, being timelike, obviously does not share this property. Further, it may be shown that while dynamical horizons always expand as they evolve, timelike membranes always shrink. Finally, in contrast to the dynamical horizon flux laws of ashtekar02a ; ashtekar03a , both the geometrical and matter contributions have indefinite signatures in the corresponding equations on timelike membranes. In particular, the geometric term no longer has a natural interpretation as a flux of gravitational radiation energy.
Simple explicit examples of trapping horizons/MTTs are provided by the Vaidya spacetimes ashtekar03a . These are spherically symmetric solutions to Einsteinโs equations describing the formation/growth of a black hole as null dust falls in from infinity. In the absence of a cosmological constant, the growth of the resulting black hole is always characterized by a dynamical/future outer trapping horizon. In the presence of a cosmological constant the situation is a little more complicated as a second MTT appears โ a timelike membrane which can be associated with the cosmological horizon. Still, even in this case, the MTT associated with the black hole remains a dynamical horizon.
Heuristic arguments presented in ashtekar03a and hayward00a suggest that under physically reasonable circumstances the MTTs associated with black hole formation and growth will always be spacelike and so future *outer* trapping horizons are generic in this context. *This intuition needs to be amended*. In Oppenheimer-Snyder spacetimes oppenheimer39 , which describe the gravitational collapse of spherical dust clouds, it has been shown that timelike membranes, rather than dynamical horizons, appear during the formation of black holes. These spacetimes are constructed out of a piece of the closed FRW universe (which forms a homogeneous and isotropic dust ball) surgically inserted into a Schwarzschild spacetime. As this โstarโ is allowed to collapse the horizon structure develops in the following way. At first there are no horizons. Then, as the dust reaches a critical density an isolated horizon forms that cuts off the continuing collapse from the rest of the universe (in this case it is also an event horizon). Coincidentally, a timelike membrane appears and contracts inside the isolated horizon until it reaches zero area.
In this case, the timelike membrane is clearly associated with the formation of the black hole and cannot be dismissed as โcosmologicalโ. One might argue that the OS spacetime is not very physical. Nevertheless, the question remains as to what is the generic behavior of MTTs during the formation or growth of a black hole. When are they spacelike and when are they timelike? Equivalently, when do they grow in the way that one would intuitively expect and when do they shrink?
In this paper, we present several new examples of MTT-containing spacetimes. From these, we see that in most circumstances horizons are either isolated or dynamical, spacelike, and expanding. However, under certain conditions more interesting behaviours are seen. There are two ways in which new horizons may form where they did not exist before. In the first, an MTT may appear out of a singularity (such as a point of infinite density). In the second, large (though not infinite) concentrations of matter may force the creation of a dynamical horizon-timelike membrane pair. Such formations are closely related to the well-known phenomena of โjumpingโ apparent horizons hawkingellis72 .
Horizons may also disappear. Timelike membrane-dynamical horizon pairs can annihilate each other or alternatively MTTs can vanish into singularities. Figure 1 schematically displays some of these behaviours, and also suggests another interpretation. One can think of a single MTT that winds its way forwards and backwards through time rather than multiple dynamical horizons, timelike membranes, and isolated horizons that appear and disappear.
For simplicity we will usually only consider spherically symmetric spacetimes. Manifolds will be smooth unless stated otherwise. We use units such that $`G=c=1`$, and spacetimes have signature $`(,+,+,+)`$. We will assume that the Einstein equations hold, with matter satisfying the null energy condition.
The paper is structured as follows. In section II we explain the basic properties of MTTs. In particular, we discuss how the signature of the metric induced on an MTT by the spacetime metric depends on the matter present and then examine several different types of matter. In section III we illustrate the analytic evolution of MTTs in various situations for the case of pressureless dust. In section IV we consider a couple of numerical examples of MTTs in the presence of a massive scalar field. Conclusions are presented in section V.
## II Marginally trapped tubes
### II.1 General properties
As we have seen, the definition of MTTs is weaker than that of both dynamical horizons and timelike membranes (there is no restriction on the signature of the metric on $`T`$), and of trapping horizons (there is no restriction that $`_n\theta _{(\mathrm{})}0`$). On an MTT, both the induced metric signature and $`_n\theta _{(\mathrm{})}`$ are allowed to vary.
There are simple procedures which may be used to determine this signature. Here we will restrict ourselves to the case of spherically symmetric spacetimes; the general case is closely related (see, for example, hayward94a or ashtekar03a ). With this condition, it is natural to restrict our attention to similarly symmetric MTTs. <sup>3</sup><sup>3</sup>3Non-spherically symmetric MTTs will also exist in these spacetimes. However, they are more complicated to locate and so will not be considered in this first analysis. See booth and references therein for a discussion of what is know about the allowed behaviours of such MTTs. Further, it will be convenient to use the symmetry to foliate the spacetime into spacelike two-spheres and so extend the definitions of $`\mathrm{}^a`$, $`n^a`$ and $`\theta _{(\mathrm{})}`$ off of $`T`$. We require that all such quantities share the symmetry and for simplicity will also impose the standard condition that $`\mathrm{}n=1`$.
Next, following the conventions of booth04a ; thebeast05 , we let $`๐ฑ^a`$ be a vector field which is: 1) tangential to $`T`$, 2) everywhere orthogonal to the foliation by marginally trapped surfaces, and 3) generates a flow which preserves the foliation. Thus, if $`v`$ is a foliation label, $`_๐ฑv`$ is a function of $`v`$ only โ it is independent of the exact position on a leaf. Then it is always possible to find a function $`C`$ and normalization of $`\mathrm{}^a`$ such that $`๐ฑ^a=\mathrm{}^aCn^a`$. Moreover, the definition of $`๐ฑ^a`$ implies that $`_๐ฑ\theta _{(\mathrm{})}=0`$, which gives us an expression for $`C`$:
$$C=\frac{_{\mathrm{}}\theta _{(\mathrm{})}}{_n\theta _{(\mathrm{})}}.$$
(1)
Note that $`๐ฑ_a๐ฑ^a=2C`$, so that the sign of $`C`$ determines the signature of $`T`$: if $`C>0`$ it is spacelike, if $`C=0`$ ($`_{\mathrm{}}\theta _{(\mathrm{})}=0`$) or becomes undefined ($`_{\mathrm{}}\theta _{(\mathrm{})}0`$ while $`_n\theta _{(\mathrm{})}=0`$) it is null, and if $`C<0`$ it is timelike. In addition, the sign of $`C`$ determines whether $`T`$ is expanding, contracting, or unchanging in area. Indeed, denoting the area element of the two-sphere cross-sections by $`\stackrel{~}{ฯต}`$, one has
$$_๐ฑ\stackrel{~}{ฯต}=C\theta _{(n)}\stackrel{~}{ฯต}.$$
(2)
This means that the expansion or contraction of an MTT is linked to its signature (since $`\theta _{(n)}<0`$). In particular, when $`T`$ is spacelike it expands while when it is timelike it contracts.
Further (1) explains the close relationship between the various trapping horizons and MTTs. If the null energy condition holds, then the numerator is non-positive (by the Raychaudhuri equation). Thus, the sign of $`C`$ is determined by the sign of $`_n\theta _{(\mathrm{})}`$. Away from isolation and cases where $`_n\theta _{(\mathrm{})}=0`$,<sup>4</sup><sup>4</sup>4To understand why the latter implies an important distinction, see the example in senovilla2 . an MTT is a dynamical horizon if and only if it is a future outer trapping horizon. Similarly it is a timelike membrane if and only it is a future inner trapping horizon.
To understand the generic behavior of these spherically symmetric MTTs we focus on this function $`C`$. Then, keeping in mind that $`\theta _{(\mathrm{})}=0`$, from the Raychaudhuri equation it is easy to see that $`_{\mathrm{}}\theta _{(\mathrm{})}=4\pi T_{ab}\mathrm{}^a\mathrm{}^b`$ while from $`G_{ab}\mathrm{}^an^b=8\pi T_{ab}\mathrm{}^an^b`$ one can show that $`_n\theta _{(\mathrm{})}=\stackrel{~}{}/4+4\pi T_{ab}\mathrm{}^an^b`$, where $`\stackrel{~}{}`$ is the scalar curvature of the two-sphere cross-section of the MTT foliation (the easiest way to see this is to consult a table of the Newman-Penrose equations, for example in stewart ; chandra ). Then, using eq. (1),
$$C=\frac{T_{ab}\mathrm{}^a\mathrm{}^b}{1/(2A)T_{ab}\mathrm{}^an^b},$$
(3)
where $`A`$ is the area of a two-sphere cross-section of the MTT. As noted above, if the null energy condition holds then the numerator is non-negative. Thus, the sign of $`C`$ depends on the relative magnitude of $`1/(2A)`$ and $`T_{ab}\mathrm{}^an^b`$.
The MTT behaviours implied by (1) and (3) are closely related to Theorem 2 of ams05 , which may be summarized as follows. Suppose a (not necessarily spherically symmetric) spacetime is foliated by a family $`\mathrm{\Sigma }_t`$ of spacelike hypersurfaces. Let $`S\mathrm{\Sigma }_o`$ be a marginally outer trapped surface, i.e., $`\theta _{(\mathrm{})}=0`$ for an outgoing null normal $`\mathrm{}^a`$ while the ingoing null expansion is unrestricted. Suppose $`S`$ is also โstrictly stably outermostโ, which roughly means that if $`S`$ is deformed outward, the corresponding deformation of the outgoing null expansion is non-negative and positive somewhere. Then first of all, $`S`$ is contained in a horizon $`H`$ foliated by marginally outer trapped leaves that lie in $`\mathrm{\Sigma }_t`$, which exists at least as long as these leaves remain strictly stably outermost. Moreover, if the null energy condition holds, $`H`$ is achronal. If $`G_{ab}\mathrm{}^a\mathrm{}^b>0`$ somewhere on $`S`$, then $`H`$ is spacelike everywhere near $`\mathrm{\Sigma }_o`$.
The condition that $`S`$ be strictly stably outermost is equivalent to the requirement that a certain operator $`L_{\mathrm{\Sigma }_o}`$ acting on functions $`\psi `$ on $`S`$ has a strictly positive principal eigenvalue. Restricting ourselves to spherical symmetry again and using the Einstein equations, this operator reduces to
$$L_{\mathrm{\Sigma }_o}\psi =\stackrel{~}{\mathrm{\Delta }}\psi +\left(\frac{1}{2}\stackrel{~}{}8\pi T_{ab}\mathrm{}^an^b\right)\psi ,$$
(4)
where $`\stackrel{~}{\mathrm{\Delta }}`$ is the Laplacian operator on the round 2-sphere $`S`$ and $`\stackrel{~}{}=8\pi /A`$ the scalar curvature of $`S`$. The principal eigenvalue (corresponding to $`\psi =\text{constant}`$) is then
$$\lambda =8\pi \left[1/(2A)T_{ab}\mathrm{}^an^b\right].$$
(5)
According to the theorem in ams05 , if the null energy condition holds then $`\lambda >0`$ implies local achronality of the horizon, consistent with our considerations above. If, moreover, $`G_{ab}\mathrm{}^a\mathrm{}^b=8\pi T_{ab}\mathrm{}^a\mathrm{}^b>0`$ on $`S`$ then the horizon must be spacelike, which is again as we found. Thus, the main results of ams05 applied to the case of spherical symmetry are neatly encapsulated in the expression (3) for the function $`C`$, which, however, will also tell us when we are dealing with a timelike membrane.
We now consider eq. (3) for some particular types of matter.
### II.2 Behavior for some matter sources
#### II.2.1 Timelike perfect fluid
For a perfect fluid that moves along timelike worldlines with unit tangent $`u^a`$, the stress-energy tensor takes the form
$$T_{ab}=(\rho +P)u_au_b+Pg_{ab},$$
(6)
where $`\rho `$ is the matter density of the fluid and $`P`$ is the pressure. Writing $`u^a=\xi \mathrm{}^a+(2\xi )^1n^a`$ for some function $`\xi `$, we find that
$$C=\frac{1}{2\xi ^2}\frac{\rho +P}{(1/A)+P\rho },$$
(7)
and see that a priori, the MTT could have spacelike, null, or timelike evolution, the deciding factor being the magnitude of $`(\rho P)`$ relative to $`1/A`$.
Simple examples of the various behaviours may be found in Robertson-Walker spacetimes bendov04a ; senovilla1 . In the case where these cosmological models are collapsing, one can find spherical MTTs through all points in $`M`$. Picking a single MTT for definiteness (or equivalently selecting a โcentreโ for the universe), and assuming an equation of state of the form $`P=\sigma \rho `$, where $`\sigma `$ is some constant, these surfaces are: i) timelike and contracting if $`\sigma <1/3`$ (this includes a dust-filled universe, $`\sigma =0`$, and so the timelike contractions seen in Oppenheimer-Snyder collapse), ii) null and contracting if $`\sigma =1/3`$ (a radiation filled universe with divergent $`C`$), and iii) spacelike and expanding if $`\sigma >1/3`$.<sup>5</sup><sup>5</sup>5Of course, these potential behaviours are not confined to barotropic equations of state. In general, given the dominant energy condition and assuming that $`\rho P`$, the signature of the MTT is determined by $`\rho 3P`$. For now, we are not too concerned with the physical interpretation of such models and evolutions, but instead just note that these are all allowed mathematical possibilities.
In the case of pressureless dust, $`P=0`$, one has the following heuristic picture deciding how the MTT will behave at a given two-sphere cross-section. Foliate spacetime by the level surfaces of dust particle proper time. Fix such a slice, let $`r`$ be some radial coordinate, and define the โmass enclosed by a sphere with coordinate radius $`r`$โ to be $`m(r)=m_oN(r)`$, with $`m_o`$ the mass of a single particle and $`N(r)`$ the number of particles with radial coordinate smaller than $`r`$. Then eq. (7) can be rewritten as
$$C=\frac{f}{\stackrel{~}{\rho }_T\stackrel{~}{\rho }},$$
(8)
where $`f=R\rho /(4\xi ^2)0`$ with $`R=(A/4\pi )^{1/2}`$ the areal radius, and $`\stackrel{~}{\rho }_T`$ and $`\stackrel{~}{\rho }`$ are two quantities with dimensions of surface density. $`\stackrel{~}{\rho }=dm/dA`$ is the change of enclosed mass with area, while $`\stackrel{~}{\rho }_T`$ is a purely geometric quantity associated with the MTT:
$$\stackrel{~}{\rho }_T=\frac{1}{8\pi }M_T\stackrel{~}{},$$
(9)
with $`M_T=R/2`$ the instantaneous mass<sup>6</sup><sup>6</sup>6See ashtekar03a for a detailed discussion in the context of dynamical horizons; similar considerations hold for MTTs. of the MTT and $`\stackrel{~}{}=8\pi /A`$ the scalar curvature of the 2-sphere cross-section. $`\stackrel{~}{\rho }_T`$ is a straightforward generalization of the *mass aspect* for isolated and dynamical horizons multipoles : it determines the mass multipoles of the MTT and as such its definition and physical meaning are not restricted to spherical symmetry. Thus, the relative magnitude of the โmatter surface densityโ $`\stackrel{~}{\rho }`$ and the โgeometric surface densityโ $`\stackrel{~}{\rho }_T`$ is what governs the behavior of the MTT. Given a 2-sphere cross-section of the MTT, if the matter surface density is larger than the geometric surface density then the MTT will contract because of eq. (2), and in that case it must be timelike. If on the other hand the matter surface density is smaller than the geometric one, the MTT will expand, in which case it must be spacelike. When the two balance each other the MTT is null (though not isolated, as in this case $`C\mathrm{}`$ rather than $`0`$).
#### II.2.2 Null fluids
We next consider a null fluid that moves inwards from infinity towards some centre with tangent vector $`n^a`$. Then, the stress-energy tensor is
$$T_{ab}=(\rho +P)n_an_b,$$
(10)
and so one quickly finds that
$$C=2A(\rho +P)0.$$
(11)
Hence, for matter of this type, the MTT can only be either isolated (if $`\rho +P=0`$) or spacelike and expanding (otherwise). Indeed, this is the type of matter in the Vaidya spacetime where it is well known that MTTs demonstrate only these behaviours ashtekar03a .
We note in passing that the Vaidya spacetimes can be generalized to include both outgoing and ingoing null dust, plus a distribution of energy whose rest frame is stationary. In such cases, the MTT can again be timelike, as follows from the discussion in israel1 . We suspect that with a careful choice of the parameters in these models, one could also generate examples of MTTs that are partly timelike, partly null, and partly spacelike. However, we will not discuss these examples further here, choosing instead to focus our attention on the more astrophysically relevant timelike dust spacetimes of section III.
#### II.2.3 Scalar fields
The final matter model that we will consider is that of a scalar field $`\varphi `$. This has stress-energy tensor
$$T_{ab}=\frac{1}{4\pi }\left(_a\varphi _b\varphi \frac{1}{2}_c\varphi ^c\varphi g_{ab}V(\varphi )g_{ab}\right),$$
(12)
where $`V(\varphi )`$ is a potential function. For arbitrary potentials, $`C`$ takes the form
$$C=\frac{2(_{\mathrm{}}\varphi )^2}{(1/A)2V(\varphi )}.$$
(13)
If $`V=0`$ (a massless scalar field), then $`_{\mathrm{}}\varphi 0`$ implies that the MTT must be spacelike. However, unless $`V`$ is negative definite, a non-zero potential a priori allows for spacelike, null, and timelike evolutions. In particular, this is the case for a massive Klein-Gordon field where $`V(\varphi )=m_o\varphi ^2/2`$ for some $`m_o>0`$.
## III Pressureless dust
Examples of all of the MTT behaviours discussed above can be seen within the Tolman-Bondi family of solutions. In this section we will review this surprisingly rich set of solutions and then discuss specific members which display the various behaviours.
### III.1 Tolman-Bondi solutions
These solutions describe the gravitational collapse of spherically symmetric dust clouds. They are very easy to work with and allow us to trace the evolution of a spacetime from specified initial conditions. A nice discussion can be found in the early part of gonc and with minor changes, we will follow that description below.
Initial conditions are given on a spherically symmetric, spacelike, three-surface $`\mathrm{\Sigma }_o`$. On that surface, we may specify: i) the dust density $`\rho _o`$, so that on that surface $`T_{ab}=\rho _ou_au_b`$ where $`u^a`$ is the forward-in-time pointing timelike unit normal to $`\mathrm{\Sigma }_o`$, and ii) the initial areal velocity $`v_o=\frac{dr}{d\tau }`$ of the dust, where $`r=\sqrt{A/4\pi }`$ is the areal radius and $`\tau `$ is the proper time as measured by observers comoving with the dust. Then, taking the areal radius $`r`$ as a coordinate on $`\mathrm{\Sigma }`$ along with the usual spherical coordinates $`(\theta ,\varphi )`$, these two functions are sufficient to specify both the intrinsic metric $`h_{ab}`$ and extrinsic curvature $`K_{ab}`$ of $`\mathrm{\Sigma }_o`$. Defining
$`m(r)`$ $`=`$ $`4\pi {\displaystyle _0^r}\rho _o(\stackrel{~}{r})\stackrel{~}{r}^2๐\stackrel{~}{r}\text{and}`$ (14)
$`k(r)`$ $`=`$ $`{\displaystyle \frac{2m(r)}{r}}v_o^2(r),`$ (15)
we have
$$ds^2=\frac{dr^2}{1k(r)}+r^2(d\theta ^2+\mathrm{sin}^2\theta d\varphi ^2),$$
(16)
and
$$K_{ab}=\frac{dv_o}{dr}\widehat{r}_a\widehat{r}_b+\frac{v_o}{r}\mathrm{\Omega }_{ab},$$
(17)
where $`\mathrm{\Omega }_{ab}=[d\theta ]_a[d\theta ]_b+\mathrm{sin}^2\theta [d\varphi ]_a[d\varphi ]_b`$ and $`\widehat{r}^a`$ is the spacelike, unit normal, outward pointing radial vector. With the restriction that $`\mathrm{\Sigma }_o`$ be spacelike, we note that initial conditions must be chosen so that $`k(r)<1`$.
The functions $`m(r)`$ and $`k(r)`$ have well-defined physical interpretations: $`m(r)`$ is a mass function which measures the amount of matter contained within a sphere of areal radius $`r`$ while $`k(r)`$ determines whether or not the system is gravitationally bound. We restrict our attention to gravitationally bound systems ($`k(r)>0`$) and so the allowed values of $`k(r)`$ are $`0<k(r)<1`$. Note too that if $`v_o(r)=0`$, $`\mathrm{\Sigma }_o`$ will represent an instant of time symmetry with $`K_{ab}=0`$. Then $`k(r)=\frac{2m(r)}{r}`$ and so $`02m(r)<r`$.
If we further restrict our attention to initial conditions for which all matter is initially either stationary or infalling ($`v_o(r)0`$), it is not hard to use the Einstein equations to show that with $`\tau `$ as the proper time measured by observers comoving with the dust,
$$\dot{R}(\tau ,r)\frac{dR(\tau ,r)}{d\tau }=\sqrt{\frac{2m(r)}{R(\tau ,r)}k(r)},$$
(18)
where $`R(\tau ,r)`$ is defined as the areal radius at time $`\tau `$ of the dust shell that had initial areal radius $`r`$ on $`\mathrm{\Sigma }_o`$.
Then, in Gaussian normal coordinates, these initial conditions can be evolved to give us a four-dimensional metric:
$$ds^2=d\tau ^2+\frac{(R^{}(\tau ,r))^2}{1k(r)}dr^2+R^2(\tau ,r)d\mathrm{\Omega }^2,$$
(19)
where $`R^{}=R/r`$ and the stress-energy tensor takes the form $`T_{ab}=\rho (\tau ,r)_a\tau _b\tau `$ with:
$$\rho (\tau ,r)=\frac{1}{4\pi R^2(\tau ,r)}\frac{dm}{dR}=\frac{r^2\rho _o(r)}{R^2(\tau ,r)R^{}(\tau ,r)}.$$
(20)
Now, there is an exact, parametric, solution to equation (18). Specifically, for $`0\eta <\pi `$ and initial areal radius $`r`$:
$`\tau (\eta ,r)`$ $`=`$ $`\tau _o(r)+{\displaystyle \frac{m(r)}{k^{3/2}(r)}}(\eta +\mathrm{sin}\eta )\text{ and}`$ (21)
$`R(\eta ,r)`$ $`=`$ $`{\displaystyle \frac{2m(r)}{k(r)}}\mathrm{cos}^2\left({\displaystyle \frac{\eta }{2}}\right),`$ (22)
where
$$\tau _o(r)=\frac{rv_o(r)}{2k(r)}+\frac{m(r)}{k^{3/2}(r)}\mathrm{arccos}\sqrt{1\frac{rv_o^2(r)}{2m(r)}}.$$
(23)
In the special case where $`v_o(r)=0`$, then $`\tau _o(r)=0`$ and $`\tau =0`$ corresponds to $`\eta =0`$. For simplicity, our explicit examples will be restricted to such evolutions; it turns out that these are sufficient to demonstrate all of the potential MTT evolutions. The equations (21) and (22) then reduce to:
$`\tau (\eta ,r)`$ $`=`$ $`\left({\displaystyle \frac{r^3}{8m(r)}}\right)^{1/2}(\eta +\mathrm{sin}\eta ),`$ (24)
$`R(\eta ,r)`$ $`=`$ $`r\mathrm{cos}^2\left({\displaystyle \frac{\eta }{2}}\right).`$ (25)
In constructing our examples, we will sometimes find it convenient to excise the interior or exterior part of a dust spacetime and replace it by a Schwarzschild region. Then, the Einstein equations are satisfied at the three-dimensional junction surface if and only if its intrinsic and extrinsic curvatures are the same whether measured on the interior or the exterior of the surface. For our purposes, it will be sufficient to only consider excisions along the comoving surfaces of constant $`r`$. Specifically, suppose we wish to take the junction surface as $`r=\widehat{r}`$, making the spacetime Schwarzschild either for $`r<\widehat{r}`$ or $`r>\widehat{r}`$. Then assuming $`v_o=0`$, the induced metric from the dust side on the corresponding 3-dimensional timelike junction surface is
$$d\widehat{s}^2=\frac{\widehat{r}^3}{8m(\widehat{r})}(1+\mathrm{cos}\eta )^2d\eta ^2+\widehat{r}^2\mathrm{cos}^4\left(\frac{\eta }{2}\right)(d\theta ^2+\mathrm{sin}^2\theta d\varphi ^2).$$
(26)
The induced metric from the Schwarzschild side is the same with $`m(\widehat{r})`$ replaced by the Schwarzschild mass $`M`$. The non-zero components of the extrinsic curvature from the dust side are
$$K_{\theta \theta }=\frac{K_{\varphi \varphi }}{\mathrm{sin}^2\theta }=\widehat{r}\mathrm{cos}^2\left(\frac{\eta }{2}\right)\sqrt{1\frac{2m(\widehat{r})}{\widehat{r}}},$$
(27)
and from the Schwarzschild side we again get the same with $`m(\widehat{r})`$ replaced by $`M`$. Thus the resulting restrictions are quite simple. If it is the *exterior* that is being replaced by Schwarzschild beyond some coordinate radius $`\widehat{r}`$, then the matching of intrinsic and extrinsic curvatures of the boundary requires that the Schwarzschild mass of the exterior geometry be $`M=m(\widehat{r})`$. If the *interior* is being replaced by a Schwarzschild geometry with mass $`M`$ for $`r<\widehat{r}`$, then the relationship between mass function and initial density should be
$$m(r)=M+4\pi _{\widehat{r}}^r\overline{r}^2\rho _o(\overline{r})๐\overline{r}$$
(28)
for $`r\widehat{r}`$.
Thus, given initial conditions, we can analytically calculate the full, four-dimensional metric as long as the Gaussian normal coordinate system remains valid. As can be inferred from the preceding discussion, the reason for this simplicity is that any shell of radius $`r`$ effectively evolves along the geodesics of a Schwarzschild solution of mass $`m(r)`$ โ outer shells do not affect the evolution of inner shells and inner shells only affect outer shells through the total mass function $`m(r)`$. This makes these spacetimes especially easy to work with but it also points to a potential problem. If the shells of constant $`r`$ do not maintain their original ordering, then the mass contained within a shell can change with time, and in this case the evolution will no longer be described by the solution discussed above. Further, apart from this physical problem, shell-crossings will also cause the Gaussian normal coordinates to break down as those surfaces of constant $`r`$ were also used as coordinates.
There is a large body of literature on shell crossings (see for example the references in gonc ), and the Tolman-Bondi spacetimes are often used to study this phenomenon. In this paper, however, we are only interested in using these spacetimes to provide concrete examples of MTT evolutions. As such we will studiously avoid such complications. A sufficient condition to guarantee that no shell-crossings occur is easily seen to be
$$R^{}(\tau ,r)0,$$
(29)
as this will ensure a physical separation of $`r=\text{constant}`$ surfaces. If $`v_o(r)=0`$ it is not hard to see that this reduces to
$$\frac{d\tau _c}{dr}0,$$
(30)
where
$$\tau _c(r)=\pi \sqrt{\frac{r^3}{8m(r)}},$$
(31)
is the time for the shell of initial areal radius $`r`$ to collapse to zero area.
### III.2 MTTs in Tolman-Bondi spacetimes
The above considerations define the Tolman-Bondi spacetimes. We now consider the location of *spherically symmetric* marginally trapped surfaces within those spacetimes. For definiteness we choose
$$\mathrm{}_a=[d\tau ]_a+\frac{R^{}(\tau ,r)}{\sqrt{1k(r)}}[dr]_a\text{ and}n_a=\frac{1}{2}[d\tau ]_a\frac{R^{}(\tau ,r)}{2\sqrt{1k(r)}}[dr]_a,$$
(32)
as the forward-in-time outward and inward null vectors satisfying $`\mathrm{}n=1`$. Note too that these expressions will always be well-defined since our requirement that $`\mathrm{\Sigma }_o`$ be purely spacelike fixes $`k(r)<1`$.
Of course $`\mathrm{}`$ and $`n`$ can be rescaled, but for the purposes of this paper that scaling is irrelevant (we only care about whether quantities defined with respect to them are zero, negative or positive) and so we can choose to work with this convenient choice of vectors. Then, keeping in mind that $`\theta _{(\mathrm{})}=(g^{ab}+\mathrm{}^an^b+n^a\mathrm{}^b)_a\mathrm{}_b`$ (and similarly for $`\theta _{(n)}`$):
$$\theta _{(\mathrm{})}=\frac{2(\dot{R}(\tau ,r)+\sqrt{1k(r)})}{R(\tau ,r)}\text{and}\theta _{(n)}=\frac{\dot{R}(\tau ,r)\sqrt{1k(r)}}{R(\tau ,r)}.$$
(33)
Solving $`\theta _{(\mathrm{})}=0`$ with a bit of help from evolution equation (18), it is not hard to see that marginally trapped surfaces occur whenever
$$R(\tau ,r)=2m(r).$$
(34)
This is the expected result if one recalls that shells of constant $`r`$ essentially move along the geodesics of a Schwarzschild spacetime with mass $`m(r)`$. Furthermore, it is clear that on such a surface,
$$\theta _{(n)}=\frac{\sqrt{1k(r)}}{m(r)}<0.$$
(35)
Thus, in these spacetimes, any two-sphere on which $`\theta _{(\mathrm{})}=0`$ is part of an MTT.
Finally we can calculate the expansion parameter $`C`$. Using condition (34), definition (1), the evolution equation (18), and definition (14) one can directly calculate
$$C=\frac{2m^{}(r)}{R^{}(\tau ,r)m^{}(r)}=\frac{2\rho (\tau ,r)}{1/A\rho (\tau ,r)},$$
(36)
where $`A=16\pi m^2(r)`$. As would be expected this result agrees with the earlier, more general, perfect fluid results (7) and (8).
With this background established, generating examples of MTT spacetimes is simply a matter of picking initial conditions and using (21) and (22) to generate the evolution. The expression for $`C`$, eq. (36) can then be used to determine the signature of the tube. Following these simple procedures we generate the examples below, all of which will start from an instant of time symmetry ($`v_o=0`$). Often we will consider situations where the dust accretes onto a pre-existing hole. In those cases, excisions will be performed inside of some shell so that a black hole may be inserted into the spacetime. When this is done we will use the Schwarzschild mass of the interior as a reference scale for masses, lengths, and times. In other cases the scale will be set according to the physics of the particular situation.
### III.3 Dust ball collapse
#### III.3.1 Collapse of a dust ball with Gaussian initial density
We begin with the example of a non-uniform dust ball collapsing to form a black hole. The initial density distribution is taken to be radially Gaussian so that
$$\rho _o=\frac{m_o}{\pi ^{3/2}r_o^3}e^{r^2/r_o^2}.$$
(37)
In this expression $`m_o`$ is the total mass of the cloud as read from the corresponding mass function $`m(r)`$ (equation 14) and also corresponds to its ADM mass. The parameter $`r_o`$ determines how much the cloud is initially โspread outโ. Such a distribution is pictured in figure 2a), where we have taken $`r_o=100m_o`$. $`\rho `$ is plotted in units of $`m_o^2`$ and $`r`$ in units of $`m_o`$.
For this choice of parameters $`r2m(r)>0`$ everywhere and so there are no marginally trapped surfaces in the initial time slice. Further, for any distribution of this kind condition (30) is met and therefore there are no shell crossings. In this case,
$$\tau _c(0)=\frac{\pi ^{5/4}\sqrt{3}}{4\sqrt{2}}\sqrt{\frac{r_o^3}{m_o}}1281m_o,$$
(38)
and the first $`r=\text{constant}`$ shells collapse to zero area at that time. As this happens, the central density $`\rho (\tau _c,0)`$ diverges to infinity (equation 20). Referring to 2c) we see that an MTT is born out of this divergence, while 2b) shows that it is everywhere spacelike and so is a dynamical horizon.
Note that in figure 2c) (as well as in future evolution graphs), the MTT is shown as the thick black line. The rest of the lines are $`r=\text{constant}`$ surfaces and correspond to the timelike geodesics that trace the evolution of individual dust shells. In this figure, most of them are very nearly horizontal as they have fallen in from a long way out and are moving very quickly (relative to the constant $`\tau `$ foliation) by the time that they approach the horizon.
Returning to the evolution graph 2c) we note that as $`\tau \mathrm{}`$, $`R2m_o`$ and $`C0`$ as the last bits of matter fall through the horizon and it asymptotes towards a null and isolated state.
Many of these features of the evolution may also be seen in the Penrose-Carter diagram for the spacetime which is given in figure 3. Thus, we again see the MTT created out of the central spacelike singularity, evolving in a spacelike fashion, and finishing its evolution by asymptoting to the null event horizon. Note that for simplicity of presentation, this diagram shows a sharp cut-off of the dust distribution at some finite $`r`$. In our example this does not occur and instead the density asymptotes to zero. Thus, properly one should view the cut-off as as marking, say, the $`r`$ which contains $`99.9\%`$ of the mass. This diagram also nicely shows that while the MTT appears at the same time as the singularity, nothing in particular is happening at $`r=0`$ when the event horizon โappearsโ.
#### III.3.2 Smooth versions of Oppenheimer-Snyder collapse
We now consider a family of spacetimes parametrized by a real number $`\sigma `$ where the initial density profiles are also not homogeneous but uniformly converge to a step function as $`\sigma \mathrm{}`$. These may then be viewed as interpolating between the collapse of a highly non-homogeneous dust ball like in the previous example, and Oppenheimer-Snyder collapse. As one would expect, in these examples we will encounter marginally trapped tubes that are partially spacelike and partially timelike.
In particular, consider initial densities of the form<sup>7</sup><sup>7</sup>7Similar initial density profiles can more easily be constructed by means of $`\mathrm{tanh}`$ functions. However, the associated mass functions $`m(r)`$ are awkward expressions in terms of polylogarithms, which are multi-branched and need to be treated with care. By defining initial density distributions in terms of error functions we avoid such complications.
$$\rho _o=\frac{m_oF(\sigma )}{r_o^3}\left(1\text{erf}\left[\sigma \left(\frac{r}{r_o}1\right)\right]\right)$$
(39)
where $`F(\sigma )`$ is a complicated function<sup>8</sup><sup>8</sup>8For those who are interested: $`F(\sigma )={\displaystyle \frac{3\sigma ^3}{2\pi \sigma (2\sigma ^2+3)(1+\text{erf}\sigma )+4\sqrt{\pi }e^{\sigma ^2}(1+\sigma ^2).}}`$ chosen so that $`m_o=lim_r\mathrm{}m(r)`$ is the total dust mass, $`r_o`$ is the location on the โstepโ where $`d\rho /dr`$ is a maximum, $`\sigma `$ characterizes the steepness of that step, and $`\text{erf}(x)`$ is the usual error function.
In the examples of Figure 4 we take $`m_o`$ as our length scale, choose $`r_o=2m_o`$, and consider a variety of values of $`\sigma `$. Then graph a) shows how the initial densities converge towards a step function with increasing $`\sigma `$. In the meantime graph b) shows that for smaller values of $`\sigma `$ the expansion parameter $`C`$ is always positive, while for the larger values it starts out negative with $`C6`$ while the density is approximately constant. This corresponds to the expected value of $`C`$ in the corresponding OS spacetime. It then diverges to $`\mathrm{}`$, switches to being positive around $`r=2m_o`$ (when the main part of the step has fallen through), and then asymptotes to $`0`$ as the density of dust falling through the horizon also goes to zero. As expected, the switch in sign corresponds to the corresponding switch of $`_n\theta _{(\mathrm{})}`$ โ that is when the MTT becomes instantaneously tangent to $`n^a`$.
In c) we also see that in all cases a dynamical horizon asymptotes to $`R=2m_o`$ as the last bits of dust fall through it. However, the MTT behaviour before that time varies greatly. For small values of $`\sigma `$ it continuously increases in area. For large values, we see that there are both increasing and decreasing regions. Somewhat confusingly it appears that for $`\sigma =2`$ there are both increasing and decreasing regions as well, even though b) shows that $`C>0`$ everywhere. We will return to examples such as this in section III.5.1, but for now we just keep in mind that spacelike surfaces can intersect in non-trivial ways. Here the dynamical horizon is demonstrating this as it intersects some of the $`\tau =\text{constant}`$ surfaces twice. The possibility of such foliation effects was also noted in ams05 .
Note that the timelike membranes in this example all go to zero areal radius around $`\tau =\pi m_o`$. This is not surprising as for large values of $`\sigma `$, $`\tau _c\pi m_o`$ (eq. (31)). Thus, the membranes vanish as the first dust shells also collapse to zero areal radius โ that is they disappear into the density singularity.
Finally, note that the spacetime diagrams for these spacetimes would be very similar to that shown in figure 3. The only significant difference would be that for the spacetimes with timelike membranes, the MTT would emerge from/vanish into the singularity as a timelike rather than spacelike surface โ the slope of the MTT would be greater than $`45^{}`$.
### III.4 Accretion onto a pre-existing hole
The next two examples study the accretion of dust shells onto a pre-existing black hole. In the first example a black hole will substantially increase its mass while the MTT remains null or spacelike everywhere through the evolution. In the second example we study a very large matter shell falling into a black hole; here we will again see timelike membranes.
#### III.4.1 Small dust shell falls into black hole
We begin with a dust shell of the form
$$\rho _o=\frac{m_oe^{\left(\frac{r}{r_o}\alpha \right)^2}}{2\pi ^{3/2}(1+2\alpha ^2)r_o^3},$$
(40)
where $`m_o`$ measures the total mass of the shell (ie. the asymptotic behaviour of the associated mass function), $`r_o`$ characterizes its โthicknessโ, and $`\alpha `$ gives its initial position in terms of $`r_o`$. Again this distribution is Gaussian, but this time it is a shell with peak density at $`\alpha r_o`$. We wish to study the accretion of such a shell onto a pre-existing black hole and so excise a small region in the interior of the shell spacetime and replace it by a Schwarzschild geometry with mass parameter $`M`$, putting the junction at some $`r=\widehat{r}>2M`$. From the discussion in subsection III.1, for $`r>\widehat{r}`$ the mass function must then be $`m(r)=M+4\pi _{\widehat{r}}^r\rho _o(\overline{r})\overline{r}^2๐\overline{r}`$.
An example of such a spacetime is shown in Figure 5, where we have chosen $`\widehat{r}=2.5M`$, $`m_o=M/2`$, $`\alpha =10`$, and $`r_o=10M`$. Note that the parameters cannot be chosen arbitrarily in this case as shell crossings can easily develop. In fact, a small increase of the mass parameter so that $`m_o0.7M`$ will be sufficient to cause these. That said, for the parameters that we have chosen this does not occur, and we see that the horizon is initially quiescent and begins to expand as the dust falls through it. The expansion is spacelike (from figureIIIb), $`C>0`$), peaks as the largest density of matter crosses the horizon, and then tails off along with the infalling dust. Asymptotically the horizon becomes null again with areal radius $`R=3M`$, so that the mass function tends to $`m=\frac{3}{2}M`$.<sup>9</sup><sup>9</sup>9The asymptotic values will in fact be slightly smaller than that; due to the excision at $`r=\widehat{r}`$, the parameter $`m_o`$ is not *exactly* equal to the mass of the shell. However, in this example the difference is negligible.
Figure 6 shows the corresponding spacetime diagram. From the spacetime surgery, inside the dust shell the spacetime is Schwarzschild with mass $`M`$. However, as the dust reaches and crosses the MTT, it begins to expand in a spacelike manner and continues to do so as long as the dust continues to fall in. Ultimately as that density goes to zero the MTT asymptotes to the null event horizon. As in figure 3 the MTT clearly reacts to physical events while the event horizon, being teleologically defined, does not. Note too that as in that earlier diagram, the dust distribution is again, for simplicity, shown with an edge.
#### III.4.2 Large shell of approximately constant density falls into black hole
In the previous example, the expansion was everywhere spacelike. It described a fairly dramatic situation, a significant expansion of a pre-existing black hole, however there was no sign of timelike membranes and the MTTs were isolated or dynamical horizons. We will now construct an example with timelike membranes. This will essentially be a smooth version of one of the more complicated Oppenheimer-Snyder examples presented in bendov04a . Thus, we build a spacetime in which a large amount of (initially) constant density dust falls into a black hole. The initial density is given by
$$\rho _o=\frac{3m_o\left(\text{erf}(\frac{rr_1}{M})\text{erf}(\frac{rr_2}{M})\right)}{4\pi (r_2r_1)(2r_1^2+2r_1r_2+2r_2^2+3M^2)},$$
(41)
where $`r_1`$ and $`r_2`$ mark the (approximate) start and end of the shell, $`m_o`$ is the total mass of the shell as read from the mass function. Excision is performed and the interior is replaced by a black hole with mass $`M`$, which joins onto the dust exterior at some $`r=\widehat{r}>2M`$. The full mass function for $`r>\widehat{r}`$ will then be $`m(r)=M+4\pi _{\widehat{r}}^r\rho _o(\overline{r})\overline{r}^2๐\overline{r}`$.
For this example we choose $`r_1=100M`$, $`r_2=2000M`$, and $`m_o=600M`$. Thus, the mass of the hole will increase dramatically during the evolution from $`m=M`$ to $`m=601M`$. This evolution is shown in figure 7. Note that despite the appearance of the density function, it is actually smooth for $`r>\widehat{r}`$ since the error functions themselves are smooth. Further, the density function will โspreadโ as it falls towards the hole and so the dust will not be of constant density as it crosses the horizon.
That said, we see that the horizon begins to expand in the usual way as the initial matter falls into it. Then however, something different happens. Considering evolution with respect to $`\tau `$ we see that around $`\tau =3000M`$, a new marginally trapped surface appears at $`R=1202M`$ and bifurcates into a dynamical horizon and a timelike membrane. The dynamical horizon quickly asymptotes to an isolated horizon while the timelike membrane contracts and eventually annihilates with the dynamical horizon that grows out to meet it.
Alternatively if we consider evolution with respect to the initial areal radius $`r`$, the MTT starts off almost-isolated, becomes dynamical as the first mass falls in, then goes null (with $`n^a`$ as a tangent) and becomes a timelike membrane that continues to expand as it travels backwards in time. Finally, as the last dust falls through, it again goes null (with $`n^a`$ as a tangent), dynamical, and then asymptotes back towards isolation.
The spacetime diagram for this evolution would be very similar to that shown in figure 6. The only difference would be that, during its active phase, the MTT would sometimes be timelike and so have a slope greater than $`45^{}`$.
### III.5 More complicated collapse
The previous examples display the basic behaviours of marginally trapped tubes โ spacelike expansions, creation from singularities, and timelike contractions/backwards-in-time-expansions (depending how one views the evolution). In this section, to get a feel for the possible range of evolutions, we will consider examples generated from more complicated initial conditions that combine several of these behaviours.
#### III.5.1 Spacelike expansion/multiple horizons
This example elaborates on a behaviour that we noted in section III.3.2. Namely it shows that the appearance of multiple horizons in a given leaf of a spacetime foliation does not always signal the existence of timelike membrane sections of the MTT. That is, a dynamical horizon as a spacelike surface can interweave with a foliation (of other spacelike surfaces) in highly non-trivial ways.
We will consider matter distributions of the form
$$\rho _o=\frac{\alpha \mu }{2\pi ^2r_or^2}\mathrm{sin}^2\left(\alpha \frac{r}{r_o}\right),$$
(42)
where $`r_o`$ is an arbitrary reference length scale. When this is integrated to a mass function, we see that for any positive integer $`N`$ there is an amount of mass $`\mu `$ between $`r=N\pi r_o/\alpha `$ and $`r=(N+1)\pi r_o/\alpha `$. Thus, physically this corresponds to a series of shells each with the same mass (though decreasing density). Again one has to choose parameters with care to avoid shell crossings and/or initial black holes. A particular choice that meets these criteria is $`\mu =(8\pi /5)r_o`$ and $`\alpha =1/4`$. This is the distribution whose initial configuration and evolution is shown in figure 8.
Then, from Figure 8c), it is clear that $`\tau =\text{constant}`$ surfaces in this spacetime will contain either one or three marginally trapped surfaces (except at turning points of the MTT, where they contain two). Further, these will expand and contract in apparently the same kinds of ways that we have seen previous example which include timelike membranes. However, an examination of Figure 8b), shows that despite this behaviour, the expansion parameter $`C`$ is always greater than zero and so the MTT is everywhere either spacelike or null/isolated. Then the apparent contractions/expansions arise simply because the MTT intersects the foliation in non-trivial ways.
In Figure 8c), the horizon goes vertical/null at intervals of $`R=2\mu 10r_o`$. Each of these corresponds to the density going to zero between each shell of mass $`\mu =(8\pi /5)r_o5r_o`$ and so the horizon becoming instantaneously isolated. Finally, note that with a careful choice of the parameters, the newly created MTT can absorb shells of this type for arbitrarily long periods of time, always alternating between being dynamical and isolated. Thus, a horizon can absorb an arbitrarily large amount of mass without ever going timelike. It is the rate of absorption, not the total amount that is significant in this regard.
#### III.5.2 Multiple timelike membranes
Our final example will demonstrate a spacetime that contains a MTT made up of multiple dynamical horizon and timelike membrane regions. The dust density will take the form:
$$\rho _0=\{\begin{array}{cc}\frac{\alpha }{r_o^2}\left[\pi \frac{1}{5}\frac{r}{r_o}\left(3+2\mathrm{cos}^2\left(5\frac{r}{r_o}\right)\right)\right]\hfill & 0r\pi r_o\hfill \\ & \\ 0\hfill & r>\pi r_o\hfill \end{array}$$
(43)
where $`\alpha `$ is a dimensionless constant. The exterior of $`\pi r_o`$ was excised to avoid negative density dust and so violations of the energy conditions. Thus, outside of that radius, the geometry will be Schwarzschild with mass $`M=m(\pi r_o)`$.
Taking $`\alpha =1/120`$ there are neither shell-crossings nor initial black holes in this spacetime. The evolution is then shown in the graphs of figure 9. The initial density is irregular and gives rise to an MTT which, if we think of evolution as parameterized by $`r`$, starts out as a dynamical horizon and then alternates back and forth between being timelike and spacelike with (non-isolated) null cross-sections separating these regions. From the point of view of evolution with respect to $`\tau `$, spacelike slices have anywhere between one and five marginally trapped surfaces in this example. However, all of these intricacies are contained within the outermost isolated horizon (and here it really is isolated since we cut the density distribution at $`r=\pi r_o)`$ and would not be visible to outside observers.
## IV Scalar fields
In the previous section we have seen that the potential MTT behaviours suggested by (7) and (8) are all achieved by the dust spacetimes. Thus, depending on the initial configuration, the MTT can be any of spacelike, null, or timelike. In this section we will attempt the same demonstration for the scalar fields. Thus, we will consider initial configurations of scalar fields, evolve them in time, and examine the behaviour of any MTT that forms with the help of the expansion parameter $`C`$.
Given that scalar fields are significantly more complicated than dust we will necessarily take a numerical rather than an analytic approach.<sup>10</sup><sup>10</sup>10 There are a few exact scalar field solutions in the literature, see for example viqar and the references listed therein. However, they are much more restricted than the Tolman-Bondi solutions (for instance, we are not aware of any that are asymptotically flat) and cannot produce the same variety of examples. Section IV.1 will introduce the numerical model and the subsequent sections will present the results for two different scalar fields configurations. These will be analogous to configurations considered in the last section. The first, in section IV.2, examines the evolution of a initially smooth step-like configuration that collapses to form a black hole. The second, in section IV.3, studies a shell of scalar field which falls into an existing black hole. Due to numerical complexities, these results will be less complete than those considered in the last section, but will still demonstrate some interesting and complementary behaviours.
### IV.1 Numerical approach
We will be interested in the scalar fields briefly considered in section II.2.3. The equations governing spacetimes containing these fields are generated by the Lagrangian:
$$=\sqrt{g}\left[\frac{1}{16\pi }R\frac{1}{2}_\alpha \varphi ^\alpha \varphi V(\varphi )\right].$$
(44)
This class of models includes the massive Klein-Gordon field with $`V=m_0\varphi ^2/2`$ and in the following we will restrict ourselves to this case.
To solve the resulting coupled Einstein-Klein-Gordon equations we work with the standard 3+1 approach based on the ADM equations ADM and in particular adopt the techniques introduced in Regular to perform the evolutions. Restricting to spherical symmetry, we can then study the general evolutions by considering the dynamics of spacetimes with metrics of the form
$$ds^2=\alpha ^2(r,t)dt^2+A(t,r)dr^2+r^2B(r,t)\left(d\theta ^2+\mathrm{sin}^2\theta d\phi ^2\right).$$
(45)
The spherical symmetry implies that all dynamical functions depend only on $`r`$ and $`t`$. $`\alpha `$ is the usual lapse function and for simplicity we impose a vanishing-shift gauge condition; equivalently the โtime-evolutionโ vector, $`/t`$ is everywhere orthogonal to the $`t=\text{constant}`$ โinstantaneousโ three-surfaces.
Focusing on this natural foliation with respect to $`t`$, we note that these hypersurfaces have intrinsic three-metric and extrinsic curvature :
$`h_{ab}`$ $`=`$ $`A[dr]_a[dr]_b+Br^2\mathrm{\Omega }_{ab}\text{and}`$ (46)
$`K_{ab}`$ $`=`$ $`{\displaystyle \frac{1}{2\alpha }}\left({\displaystyle \frac{A}{t}}\right)[dr]_a[dr]_b+{\displaystyle \frac{1}{2\alpha }}\left({\displaystyle \frac{B}{t}}\right)\mathrm{\Omega }_{ab},`$ (47)
respectively, where $`\mathrm{\Omega }_{ab}`$ was defined following equation (17). These, together with the value of the scalar field, are the variables whose evolution we will study.
Then, in the usual way we can rewrite the Einstein-Klein-Gordon equations in terms of the Hamiltonian and momentum constraints which restrict initial values of these quantities, in addition to evolution equations (which preserve the constraints). The actual implementation consists of picking an initial configuration for the scalar field, solving the constraints for $`h_{ab}`$ and $`K_{ab}`$ and finally using the evolution equation to obtain the dynamics of the corresponding spacetime. For brevity we will not include all of the details here; the interested reader is directed to Regular for a more detailed discussion.
Once we have a spacetime evolution, the next step is to search for an MTT. In this case it is natural to search for apparent horizons/marginally trapped surfaces within the $`t=\text{constant}`$ slices and this is what we will do. Since we are in spherical symmetry we need only consider the null expansions of the $`r=\text{constant}`$ two-surfaces. Then, with
$$u_a=\alpha [dt]_a\text{ and }s_a=\sqrt{A}[dr]_a,$$
(48)
as the forward-in-time and towards-infinity pointing unit normals to the two-surfaces, we define
$$\mathrm{}_a=u_a+s_a\text{ and }n_a=\frac{1}{2}(u_as_a)$$
(49)
as the outward and inward forward-in-time pointing null normals to those same surfaces. As in the dust examples, the exact โnormalizationโ is irrelevant as we are only interested in the signs of the ensuing quantities, not their magnitudes. Thus, we can make this specific, convenient, choice. Then, it is straightforward to see that:
$`\theta _{(\mathrm{})}`$ $`=`$ $`๐_as^a+KK_{ab}s^as^b={\displaystyle \frac{_tB}{\alpha B}}+{\displaystyle \frac{1}{A^{1/2}}}\left({\displaystyle \frac{2}{r}}+{\displaystyle \frac{_rB}{B}}\right)`$ (50)
$`\theta _{(n)}`$ $`=`$ $`{\displaystyle \frac{1}{2}}(๐_as^aK+K_{ab}s^as^b)={\displaystyle \frac{_tB}{2\alpha B}}{\displaystyle \frac{1}{2A^{1/2}}}\left({\displaystyle \frac{2}{r}}+{\displaystyle \frac{_rB}{B}}\right),`$ (51)
where $`๐_a`$ is the covariant derivative compatible with $`h_{ab}`$ and $`K=h^{ab}K_{ab}`$ (as a quick check that these expressions are reasonable, note that for Minkowski space where $`A=B=1`$, $`\theta _{(\mathrm{})}=2/r>0`$ and $`\theta _{(n)}=1/r<0`$). In the code we evaluate these expressions over the whole numerical grid and look for places where $`\theta _{(\mathrm{})}`$ changes signs and $`\theta _{(n)}<0`$. The corresponding two-spheres are then identified as the marginally trapped two-spheres which foliate the MTTs, and their area is
$$A_{AH}=4\pi r_{AH}^2B_{AH}$$
(52)
where $`B_{AH}`$ is the metric function $`B`$ evaluated at $`r=r_{AH}`$, the coordinate radius of these two-spheres.
Finally, the expansion parameter $`C`$ for this surface can be calculated using (13).
### IV.2 Scalar field โstepโ collapses to form a black hole
Our first example will be analogous to the dust ball collapses considered in section III.3. As for those dust examples, our initial slice will be a moment of time symmetry ($`K_{ab}=0`$) and further we will choose our coordinates so that the coordinate radius $`r`$ will also be the areal radius $`R`$ on that slice (thus $`B=1`$ initially). Then, we specify a step-like configuration of the scalar field as shown in figure 10a).
Solving the constraint equations to find the initial form of $`A`$, we can then integrate the data as discussed above to find how the geometry and fields evolve in time. The results are shown in figures 10b) and 10c). The first thing to note is that throughout this collapse $`C>0`$ and so the MTT is a dynamical horizon. Secondly, as in the earlier examples it asymptotes to null as the spacetime settles down to become Schwarzschild. Note from figure 10c) that initially the hypersurfaces do not contain an apparent horizon. However, during the collapse an apparent horizon appears and then keeps growing until it reaches a constant size.
### IV.3 Scalar field โshellโ accretes onto existing black hole
Our second scalar field example is analogous to the dust shell examples of section III.4. Thus we will consider an initial โshellโ of scalar field that accretes onto an existing black hole of mass $`M`$. After surgically inserting the black hole we again start out on a slice of time symmetry, though in this case for technical reasons do not start with $`r`$ as the areal radius. Instead, on the initial slice $`\stackrel{~}{B}=1`$, where $`B=\left(1+\frac{M}{2r}\right)^4\stackrel{~}{B}`$.
Then, as our initial scalar field conditions we consider a $`\varphi (0,r)`$ of the form shown in figure 11a). The corresponding initial intrinsic geometry (essentially only $`A`$ remains unknown) is then found by solving the constraint equations. The results for the corresponding evolution are then shown in figures 11b) and 11c). On those graphs we again see that $`C`$ remains everywhere non-negative and asymptotes to zero at late times. Note too the apparent horizon on the initial hypersurface which grows as the scalar field falls into the hole.
### IV.4 Outlook for scalar fields
Neither of the preceding examples included timelike membrane regions of the MTT. We believe that this is a reflection of the examples that our code has been able to integrate rather than a fundamental result. From (13) it is clear that a timelike membrane will only appear in situations where there is a sufficiently large concentration of the scalar field relative to the (inverse) area of the horizon. Our code had difficulty evolving such examples and so the lack of timelike membranes is not especially surprising. Physically, one would also expect it to be more difficult to obtain such examples for a massive scalar field which, in contrast to pressureless dust, resists compression. We expect that future investigations will find the appropriate combination of initial conditions needed to generate examples of timelike membranes in scalar field spacetimes. For now though, we simply note that our examples show timelike evolutions are not the rule for scalar field spacetimes.
## V Conclusions and Speculations
In this paper we have seen that dynamical horizons and timelike membranes characterize two possible modes of black hole expansion. In the first a black hole/MTT smoothly expands as matter falls into it. By contrast, the second occurs when matter densities are high enough to force the formation of a new horizon of non-zero area that encloses any already existing MTTs.
Further insight into these two possibilities can be gained by considering the magnitudes of the quantities involved. As seen in the discussion of II, assuming that the null energy condition holds, the signature of an MTT is determined by the relative magnitudes of $`1/A`$ and $`2T_{ab}\mathrm{}^an^b`$, or, specializing to (pressureless) dust, $`1/A`$ and $`\rho `$. Converting into physical units, it is straightforward to see that for a Schwarzschild black hole of mass $`M=\mu M_{}`$, the inverse horizon area corresponds to a density
$$\rho _A=\frac{c^6}{16\pi G^3M_{}^2}\frac{1}{\mu ^2},$$
(53)
where $`M_{}=2\times 10^{33}`$ g is the approximate mass of the Sun and $`c`$ and $`G`$ are respectively the speed of light and the gravitational constant. Then, for a solar mass black hole, $`\rho _A10^{16}\text{g}/\text{cm}^3`$ which is about an order of magnitude higher than the density of a neutron star. By contrast, for a supermassive black hole of mass $`10^8M_{}`$, one has $`\rho _A1\text{g}/\text{cm}^3`$, the density of water.
Though strictly speaking these results only apply in situations of spherical symmetry, it seems fairly safe to use them to draw some more general conclusions. Specifically, it is likely that both modes of expansion are not only mathematically possible as we have seen in this paper, but also both occur in physical situations. For small black holes, even if the lack of spherical symmetry changes these estimates by several orders of magnitude, it appears that dynamical horizon spacelike expansions are probably the dominant mode in all but the most extreme situations, such as black hole and/or neutron star collisions. Numerical studies support this contention as in such extreme situations the occurrence of multiple horizons appears to be generic num3 ; badridavid (though also relatively unstudied in a systematic way, since in most studies it is the exterior spacetime that is of interest and so the interior of the outermost apparent horizon is excised and thrown away). By contrast for supermassive black holes, it is likely that horizon jumps/TLMs are much more common.
These examples also suggest other (possible) properties of MTTs. First, dynamical horizons appear to only originate either out of singularities (the density singularities of III.3) or as part of a dynamical horizon/TLM pair (as in III.5.2). Equivalently, in all of our examples there is just one MTT associated with each black hole. This originates in a singularity and then weaves backwards and forwards in time, always expanding in area (relative to a foliation parameter that monotonically increases as one moves away from the singularity). This suggests that a similar result may be true away from spherical symmetry โ the multiple horizons/jumps seen in numerical studies of black hole collisions may actually all be part of a single MTT that weaves backwards and forwards in time.
In our examples it is also true that TLMs and dynamical horizons always occur in such a way that a causal signal originating from the MTT would never be detectable by sufficiently far-away observers. This is consistent with a gravitational confinement theorem due to Israel israel1 ; israel2 which states that if the weak energy condition holds (as is the case in all our examples), a trapped two-sphere can be extended to a spacelike three-cylinder foliated by trapped two-spheres of *constant area*. Assuming reasonably regular spacetimes, this three-cylinder will act as a permanent one-way membrane for causal effects. Even though an observer โinsideโ an MTT can escape that trapped region, he will not be able to send signals beyond the areal radius of the two-sphere on which the MTT was first crossed. We note that Israelโs confinement theorem is quite general; it does not assume spherical symmetry or even asymptotic flatness. Moreover, in appropriate asymptotically flat spacetimes, trapped surfaces must necessarily be contained within event horizons and so unable to send causal signals to null infinity hawkingellis72 ; wald . Our asymptotically flat examples are also consistent with the latter: the outermost parts of MTTs are dynamical horizons which asymptote to a null surface with some finite area. Consequently, in physically realistic spacetimes, any TLMs that may be present will be hidden from observers far from a black hole, also in the absence of spherical symmetry.
In summary, even though the examples and (dust) calculations that we have seen in this paper are quite simple we believe that they are extremely useful in forming a correct intuition about the behaviour of MTTs during general black hole evolution. In particular they provide a convenient testing ground for ideas about these evolutions which is much simpler than full numerical simulations of black hole collisions and yet still significantly richer (and more realistic) than the heretofore studied analytical examples (Vaidya and Oppenheimer-Snyder). As such they should be useful in, among other things, suggesting possible extensions of the recent mathematical investigations abhaygreg ; ams05 of MTT properties.
## Acknowledgements:
The authors would like to thank Abhay Ashtekar, Chris Beetle, Steve Fairhurst, Greg Galloway, Sean Hayward, Werner Israel, Badri Krishnan, Josรฉ Senovilla, the participants of Black Holes V and CCGRRA 11, and an anonymous referee who all made useful suggestions and comments on this work during its development. I. Booth was supported by NSERC. C. Van Den Broeck was supported in part by the Eberly Research Fund of Penn State, NSF grant PHY-00-90091, and the Edward M. Frymoyer Honors Scholarship Program. J.A. Gonzalez was supported by DFG grant โSFB Transregio 7: Gravitationswellenastronomieโ and by NSF grants PHY-02-18750 and PHY-02-44788.
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# Transitive and self-dual codes attaining the Tsfasman-Vladut-Zink bound
## 1. Introduction and Main Results
Let $`๐ฝ_q`$ be the finite field of cardinality $`q`$. In this paper we consider primarily linear $`[n,k,d]`$-codes $`C`$ over $`๐ฝ_q`$; i.e., the parameters $`n=n(C),k=k(C)`$ and $`d=d(C)`$ are the length, the dimension and the minimum distance of the code. The ratios $`R=R(C)=k(C)/n(C)`$ and $`\delta =\delta (C)=d(C)/n(C)`$ denote the information rate and the relative minimum distance, resp., of the code.
A crucial role in the asymptotic theory of codes plays the set $`U_q[0,1]\times [0,1]`$ which is defined as follows: a point $`(\delta ,R)^2`$ with $`0\delta 1`$ and $`0R1`$ belongs to $`U_q`$ if and only if there exists a sequence $`(C_i)_{i0}`$ of codes over $`๐ฝ_q`$ such that
$$n(C_i)\mathrm{},\delta (C_i)\delta \text{and}R(C_i)R,\text{as}i\mathrm{}.$$
One then defines the function $`\alpha _q:[0,1][0,1]`$ by
$$\alpha _q(\delta )=sup\{R;(\delta ,R)U_q\},\text{for}\delta [0,1].$$
The following facts are well-known (and easy to prove), see \[M\], \[T-V\]:
###### Proposition 1.1.
i) A point $`(\delta ,R)[0,1]\times [0,1]`$ belongs to the set $`U_q`$ if and only if $`0R\alpha _q(\delta )`$.
ii) The function $`\alpha _q`$ is continuous and non-increasing.
iii) $`\alpha _q(0)=1`$, and $`\alpha _q(\delta )=0`$ for $`1q^1\delta 1`$.
Many upper bounds for $`\alpha _q(\delta )`$ are known, see \[vL\], \[T-V\]. More interesting, however, are lower bounds for $`\alpha _q(\delta )`$, since any non-trivial lower bound assures the existence of arbritrarily long linear codes with good error correction parameters. The classical lower bound for $`\alpha _q(\delta )`$ is the asymptotic Gilbert-Varshamov bound, which says:
###### Proposition 1.2.
(see \[vL\]). For all $`\delta (0,1q^1)`$ one has
$$\alpha _q(\delta )1\delta \mathrm{log}_q(q1)+\delta \mathrm{log}_q(\delta )+(1\delta )\mathrm{log}_q(1\delta ).$$
For sufficiently large non-prime $`q`$ and for certain ranges of the variable $`\delta `$, the Gilbert-Varshamov bound is improved by the Tsfasman-Vladut-Zink bound as follows:
###### Proposition 1.3.
(see \[T-V-Z\], \[N-X\]). Let
$$A(q)=lim\underset{g\mathrm{}}{sup}N_q(g)/g,$$
where $`N_q(g)`$ denotes the maximum number of rational places that a function fields $`F`$ over $`๐ฝ_q`$ of genus $`g`$ can have. Then
$$\alpha _q(\delta )1\delta A(q)^1\text{for}0\delta 1.$$
It is well known that $`A(q)q^{1/2}1`$ (this is the Drinfeld-Vladut bound), and $`A(q)=q^{1/2}1`$ if $`q`$ is a square, see \[I\], \[T-V-Z\], \[G-S1\]. It then follows easily that the Tsfasman-Vladut-Zink bound in Proposition 1.3 improves the Gilbert-Varshamov bound for all squares $`q49`$. For non-linear codes over $`๐ฝ_q`$, the Tsfasman-Vladut-Zink bound was further improved recently, see \[X\], \[N-O1\], \[N-O2\], \[S-X\], \[E1\], \[E2\].
In order to prove the Gilbert-Varshamov and the Tsfasman-Vladut-Zink bound one constructs families of long codes with sufficiently good parameters. However, the proofs provide linear codes without any particular structure. For instance, one of the most challenging problems in coding theory is still open (see \[P-H\], \[M-W\]): Do there exist sequences $`(C_i)_{i0}`$ of cyclic codes $`C_i`$ over $`๐ฝ_q`$ with
$$n(C_i)\mathrm{},\underset{i\mathrm{}}{lim}R(C_i)>0\text{and}\underset{i\mathrm{}}{lim}\delta (C_i)>0\mathrm{?}$$
Cyclic codes can be understood as a special case of what we call in this paper transitive codes. Recall that a subgroup $`U`$ of the symmetric group $`S_n`$ is called transitive if for any pair $`(i,j)`$ with $`i,j\{1,\mathrm{},n\}`$ there is a permutation $`\pi U`$ such that $`\pi (i)=j`$. A permutation $`\pi S_n`$ is called an automorphism of the code $`C๐ฝ_q^n`$ if
$$(c_1,\mathrm{},c_n)C(c_{\pi (1)},\mathrm{},c_{\pi (n)})C$$
holds for all codewords $`(c_1,\mathrm{},c_n)C`$. The automorphism group $`\mathrm{Aut}(C)S_n`$ is the group of all automorphisms of the code $`C`$.
###### Definition 1.4.
A code $`C`$ over $`๐ฝ_q`$ of length $`n`$ is said to be transitive if its automorphism group $`\mathrm{Aut}(C)`$ is a transitive subgroup of $`S_n`$.
It is obvious that any cyclic code is transitive. We can now state our first result.
###### Theorem 1.5.
Let $`q=\mathrm{}^2`$ be a square. Then the class of transitive codes meets the Tsfasman-Vladut-Zink bound. More precisely, let $`R,\delta 0`$ be real numbers with $`R=1\delta 1/(\mathrm{}1)`$. Then there exists a sequence $`(C_j)_{j0}`$ of linear codes $`C_j`$ over $`๐ฝ_q`$ with parameters $`[n_j,k_j,d_j]`$ with the following properties:
1. All $`C_j`$ are transitive codes.
2. $`n_j\mathrm{}`$ as $`j\mathrm{}`$.
3. $`lim_j\mathrm{}k_j/n_jR`$ and $`lim_j\mathrm{}d_j/n_j\delta `$.
Other important classes of codes are the self-orthogonal codes and the self-dual codes. Recall that a linear code $`C`$ is called self-orthogonal if $`C`$ is contained in its dual code $`C^{}`$, and $`C`$ is called self-dual if $`C=C^{}`$. It is clear that the information rate of a self-orthogonal codes satisfies $`R(C)1/2`$; the information rate of self-dual codes is $`R(C)=1/2`$. It is well-known that self-dual codes reach the Gilbert-Varshamov bound, see \[MW-S\]. In this paper we shall prove:
###### Theorem 1.6.
Let $`q=\mathrm{}^2`$ be a square. Then the class of self-orthogonal codes and the class of self-dual codes meet the Tsfasman-Vladut-Zink bound. More precisely we have:
1. Let $`0R1/2`$ and $`\delta 0`$ with $`R=1\delta 1/(\mathrm{}1)`$. Then there is a sequence $`(C_j)_{j0}`$ of linear codes $`C_j`$ over $`๐ฝ_q`$ with parameters $`[n_j,k_j,d_j]`$ such that:
1. All $`C_j`$ are self-orthogonal codes.
2. $`n_j\mathrm{}`$ as $`j\mathrm{}`$.
3. $`lim_j\mathrm{}k_j/n_jR`$ and $`lim_j\mathrm{}d_j/n_j\delta `$.
2. There is a sequence $`(C_j)_{j0}`$ of self-dual codes $`C_j`$ over $`๐ฝ_q`$ with parameters $`[n_j,n_j/2,d_j]`$ such that $`n_j\mathrm{}`$ and
$$\underset{j\mathrm{}}{lim}d_j/n_j1/21/(\mathrm{}1).$$
Note that the bounds given in Theorem 1.5 and Theorem 1.6 are better than the Gilbert-Varshamov bound, for all squares $`q=\mathrm{}^249`$.
The main tool to prove Theorem 1.5 and Theorem 1.6 is a new asymptotically good tower of function fields over $`๐ฝ_q`$ which has particularly nice properties, see Theorem 1.7 below. Using that tower, we shall construct sequences of codes over $`๐ฝ_q`$ with the desired properties, analogously to the proof of Proposition 1.3 by Tsfasman-Vladut-Zink.
Before stating Theorem 1.7, we recall some notations from the theory of algebraic function fields, cf. \[S1\].
1. For a function field $`F/๐ฝ_q`$ we denote by $`g(F)`$ the genus and by $`N(F)`$ the number of rational places of $`F`$. For an element $`uF\{0\}`$, we denote by $`(u)^F`$, $`(u)_0^F`$ and $`(u)_{\mathrm{}}^F`$ the principal divisor, the zero divisor and the pole divisor, resp., of the element $`u`$. In particular we have $`(u)^F=(u)_0^F(u)_{\mathrm{}}^F`$. The divisor of a differential $`\mu 0`$ of $`F/๐ฝ_q`$ is denoted by $`(\mu )^F`$.
2. Let $`๐ฝ_q(x)`$ be a rational function field; then we denote, for $`\alpha ๐ฝ_q`$, by $`(x=\alpha )`$ the zero of the function $`(x\alpha )`$ and by $`(x=\mathrm{})`$ the pole of the function $`x`$ in $`๐ฝ_q(x)`$.
3. Let $`E/F`$ be an extension of function fields over $`๐ฝ_q`$. Let $`P`$ be a place of $`F`$ and let $`Q`$ be a place of $`E`$ lying above $`P`$. Then $`e(Q|P)`$ and $`d(Q|P)`$ denote the ramification index and the different exponent, resp., of $`Q|P`$. The different of $`E/F`$ (which is a divisor of the function field $`E`$) is denoted by $`\mathrm{Diff}(E/F)`$.
###### Theorem 1.7.
Let $`q=\mathrm{}^2`$ be a square. Then there exists an infinite tower $`=(E_0E_1E_2\mathrm{})`$ of function fields $`E_i/๐ฝ_q`$ with the following properties:
1. $`๐ฝ_q`$ is the full constant field of $`E_i`$, for all $`i0`$.
2. $`E_0=๐ฝ_q(z)`$ is the rational function field.
3. There exists an element $`wE_1`$ such that $`w^\mathrm{}1=z`$. So we have $`E_0=๐ฝ_q(z)๐ฝ_q(w)E_1`$, and the extension $`๐ฝ_q(w)/E_0`$ is cyclic of degree $`(\mathrm{}1)`$.
4. All extensions $`E_n/E_0`$ are Galois, and the degree of $`E_n/E_0`$ is
$$[E_n:E_0]=(\mathrm{}1)\mathrm{}^np^{t(n)},$$
where $`p=\mathrm{char}(๐ฝ_q)`$ is the characteristic of $`๐ฝ_q`$ and $`t(n)`$ is a non-negative integer.
5. The place $`(z=1)`$ of $`E_0`$ splits completely in all extensions $`E_n/E_0`$; i.e., there are $`[E_n:E_0]`$ distinct places of $`E_n`$ above the place $`(z=1)`$, and all of them are rational places of $`E_n`$. In particular we have that the number of rational places satisfies $`N(E_n)[E_n:E_0]=(\mathrm{}1)\mathrm{}^np^{t(n)}`$.
6. The principal divisor of the function $`w`$ (as in item c)) in the field $`E_n`$ has the form
$$(w)^{E_n}=e_0^{(n)}A^{(n)}e_{\mathrm{}}^{(n)}B^{(n)},$$
where $`A^{(n)}>0`$ and $`B^{(n)}>0`$ are positive divisors of the function field $`E_n`$. The ramification index $`e_0^{(n)}`$ of the place $`(w=0)`$ in $`E_n/๐ฝ_q(w)`$ has the form
$$e_0^{(n)}=\mathrm{}^{n1}p^{r(n)}\text{with}r(n)0,$$
and the ramification index $`e_{\mathrm{}}^{(n)}`$ of the place $`(w=\mathrm{})`$ in the extension $`E_n/๐ฝ_q(w)`$ has the form
$$e_{\mathrm{}}^{(n)}=\mathrm{}^np^{s(n)}\text{with}s(n)0.$$
7. The different of the extension $`E_n/๐ฝ_q(w)`$ is given by
$$\mathrm{Diff}(E_n/๐ฝ_q(w))=2(e_0^{(n)}1)A^{(n)}+2(e_{\mathrm{}}^{(n)}1)B^{(n)},$$
with $`e_0^{(n)},e_{\mathrm{}}^{(n)},A^{(n)}`$ and $`B^{(n)}`$ as in item f).
8. The genus $`g(E_n)`$ satisfies
$$g(E_n)=[E_n:๐ฝ_q(w)]+1(\mathrm{deg}A^{(n)}+\mathrm{deg}B^{(n)})[E_n:๐ฝ_q(w)],$$
with $`A^{(n)}`$ and $`B^{(n)}`$ as in item f).
9. The tower $``$ attains the Drinfeld-Vladut bound; i.e.,
$$\underset{n\mathrm{}}{lim}N(E_n)/g(E_n)=q^{1/2}1.$$
This paper is organized as follows: In Section 2 we prove Theorem 1.7 which is the basis for our code constructions. In Section 3 we deal with transitive codes and give the proof of Theorem 1.5. We also explain briefly that the method of proof of Theorem 1.5 yields an improvement of the Tsfasman-Vladut-Zink bound for transitive non-linear codes. Finally, in Section 4 we discuss self-orthogonal and self-dual codes and we prove Theorem 1.6.
## 2. An Asymptotically Optimal Galois Tower of Function Fields
For basic notations and facts in the theory of algebraic function fields we refer to \[S1\] and \[N-X\]. We will in particular use the notations introduced in Section 1 after Theorem 1.6.
A tower of function fields over $`๐ฝ_q`$ is an infinite sequence $`=(F_0,F_1,F_2,\mathrm{})`$ of function fields $`F_i`$ over $`๐ฝ_q`$ with the following properties:
1. $`F_0F_1F_2\mathrm{}`$, and all extensions $`F_{i+1}/F_i`$ are separable of degree $`[F_{i+1}:F_i]>1`$.
2. $`๐ฝ_q`$ is the full constant field of $`F_i`$, for all $`i0`$.
3. The genus $`g(F_i)`$ tends to infinity as $`i\mathrm{}`$.
Recall that $`N(F_i)`$ denotes the number of rational places of $`F_i`$ over $`๐ฝ_q`$. It is well-known that the limit of the tower $``$,
$$\lambda ():=\underset{i\mathrm{}}{lim}N(F_i)/g(F_i)$$
does exist (see \[G-S2\]). As follows from the Drinfeld-Vladut bound (see Section 1), one has that
$$0\lambda ()A(q)q^{1/2}1.$$
The tower $``$ is said to be asymptotically optimal if $`\lambda ()=A(q)`$. For $`q=\mathrm{}^2`$ a square number we have that $`A(q)=\mathrm{}1`$, see Section 1. Therefore a tower $``$ over $`๐ฝ_q`$ is asymptotically optimal if and only if $`\lambda ()=\mathrm{}1`$ (for $`q=\mathrm{}^2`$). The tower $`=(F_0,F_1,F_2,\mathrm{})`$ is called a Galois tower if all extensions $`F_i/F_0`$ are Galois.
From here on, $`q=\mathrm{}^2`$ is a square. We will construct an asymptotically optimal Galois tower $`=(E_0,E_1,E_2,\mathrm{})`$ over $`๐ฝ_q`$ with the properties stated in Theorem 1.7. The starting point is the asymptotically optimal tower $`=(F_0,F_1,F_2,\mathrm{})`$ over $`๐ฝ_q`$ which was introduced in \[G-S2\], see also \[G-S3\]. It is defined as follows:
1. $`F_0=๐ฝ_q(x_0)`$ is the rational function field.
2. For all $`i0`$ we have $`F_{i+1}=F_i(x_{i+1})`$ with
$$x_{i+1}^{\mathrm{}}+x_{i+1}=\frac{x_i^{\mathrm{}}}{x_i^\mathrm{}1+1}.$$
$`(2.1)`$
We will need the following properties (F1) - (F5) of this tower $``$; see \[G-S2, Sec. 3\] for the proof of (F1), (F2), (F3), (F5), and \[G-S3, Sec. 3\] for the proof of (F4).
1. All extensions $`F_{i+1}/F_i`$ are Galois of degree $`\mathrm{}`$.
2. The only places of $`F_0=๐ฝ_q(x_0)`$ which are ramified in the tower $``$, are the places $`(x_0=\alpha )`$ with $`\alpha ^{\mathrm{}}+\alpha =0`$ and the place $`(x_0=\mathrm{})`$.
3. The places $`(x_0=\mathrm{})`$ and $`(x_0=\alpha )`$ with $`\alpha ^\mathrm{}1+1=0`$ are totally ramified in all extensions $`F_n/F_0`$; i.e., their ramification index in $`F_n/F_0`$ is $`\mathrm{}^n`$.
4. One can refine the extensions $`F_{i+1}/F_i`$ to Galois steps of degree $`p=\mathrm{char}(๐ฝ_q)`$ as follows:
$$F_i=H_i^{(0)}H_i^{(1)}\mathrm{}H_i^{(a)}=F_{i+1}$$
with $`[H_i^{(j+1)}:H_i^{(j)}]=p`$. For any place $`P`$ of $`H_i^{(j)}`$ and $`Q`$ of $`H_i^{(j+1)}`$ lying above $`P`$, the different exponent $`d(Q|P)`$ satisfies
$$d(Q|P)=2(e(Q|P)1).$$
5. All places $`(x_0=\alpha )`$ of $`F_0`$ with $`\alpha ๐ฝ_q`$ and $`\alpha ^{\mathrm{}}+\alpha 0`$ split completely in the tower $``$; i.e., any of these places has $`\mathrm{}^n`$ extensions in $`F_n|F_0`$, and all of them are rational places of $`F_n`$.
We set
$$w:=x_0^{\mathrm{}}+x_0\text{and}z:=w^\mathrm{}1;$$
$`(2.2)`$
then
$$๐ฝ_q(z)๐ฝ_q(w)F_0=๐ฝ_q(x_0)F_1F_2\mathrm{}.$$
The extension $`๐ฝ_q(w)/๐ฝ_q(z)`$ is cyclic of degree $`(\mathrm{}1)`$, and the extension $`F_0/๐ฝ_q(w)`$ is Galois of degree $`\mathrm{}`$. In the extension $`F_0/๐ฝ_q(z)`$ we have the following ramification and splitting behaviour (which is easily checked):
1. The place $`(z=\mathrm{})`$ of $`๐ฝ_q(z)`$ is totally ramified in $`F_0/๐ฝ_q(z)`$; the only place of $`F_0`$ lying above $`(z=\mathrm{})`$ is the place $`(x_0=\mathrm{})`$.
2. Exactly $`\mathrm{}`$ places of $`F_0`$ lie above the place $`(z=0)`$, namely the places $`(x_0=\alpha )`$ with $`\alpha ^{\mathrm{}}+\alpha =0`$. Their ramification index in $`F_0/๐ฝ_q(z)`$ is $`\mathrm{}1`$.
3. No other places of $`๐ฝ_q(z)`$ are ramified in $`F_0`$.
4. One can refine the extension $`F_0/๐ฝ_q(w)`$ to Galois steps of degree $`p=\mathrm{char}(๐ฝ_q)`$ as follows:
$$๐ฝ_q(w)=H^{(0)}H^{(1)}\mathrm{}H^{(a)}=F_0$$
with $`[H^{(j+1)}:H^{(j)}]=p`$. For any place $`P`$ of $`H^{(j)}`$ and $`Q`$ of $`H^{(j+1)}`$ lying above $`P`$, the different exponent $`d(Q|P)`$ satisfies
$$d(Q|P)=2(e(Q|P)1).$$
5. The place $`(z=1)`$ splits completely in the extension $`F_0/๐ฝ_q(z)`$; the places of $`F_0`$ lying above $`(z=1)`$ are exactly the places $`(x_0=\alpha )`$ with $`\alpha ๐ฝ_q`$ and $`\alpha ^{\mathrm{}}+\alpha 0`$.
After these preparations we can now prove Theorem 1.7. We start with the tower $`=(F_0,F_1,F_2,\mathrm{})`$ as above; in particular we consider the elements $`w,zF_0`$ as defined in (2.2) above. Then we define the tower $`=(E_0,E_1,E_2,\mathrm{})`$ as follows: $`E_0=๐ฝ_q(z)`$ is the rational function field. For all $`n1`$,
$$E_n\text{is the Galois closure of field extension}F_{n1}/E_0.$$
We have then
$$E_0=๐ฝ_q(z)๐ฝ_q(w)๐ฝ_q(x_0)E_1E_2\mathrm{},$$
and items b), c) of Theorem 1.7 are clear. By Galois theory, the field $`E_n`$ is the composite of the fields
$$F_{n1},\tau (F_{n1}),\rho (F_{n1}),\mathrm{},$$
where $`\tau ,\rho ,\mathrm{}`$ run through all embeddings of the field $`F_{n1}`$ over $`E_0`$ into a fixed algebraically closed field $`\overline{E}E_0`$. The extension $`๐ฝ_q(w)/E_0`$ is Galois, hence the field $`๐ฝ_q(w)`$ is mapped onto itself by all such embeddings of $`F_{n1}/E_0`$. By items (F4) and (F9) above we can therefore obtain the field $`E_n`$ by iterated composites of $`F_{n1}`$ with Galois extensions of degree $`p=\mathrm{char}(๐ฝ_q)`$. It follows that the degree of $`E_n/F_q(w)`$ is a power of $`p`$. Since $`[F_{n1}:๐ฝ_q(w)]=\mathrm{}^n`$, item d) of Theorem 1.7 follows.
We consider now the place $`(z=1)`$ of the rational function field $`E_0=๐ฝ_q(z)`$. By items (F5) and (F10), this place splits completely in the extension $`F_{n1}/E_0`$; hence it splits completely also in $`\tau (F_{n1})/E_0`$ for all embeddings $`\tau `$ as above. As follows from ramification theory, the place $`(z=1)`$ then splits completely in the composite field of $`F_{n1},\tau (F_{n1}),\mathrm{}`$ (see \[S1, III.8.4\]). We have thus proved item e) of Theorem 1.7. An immediate consequence is that $`๐ฝ_q`$ is the full constant field of $`E_n`$; this is item a) of Theorem 1.7. Item f) of Theorem 1.7 follows easily from (F3) and (F6).
The core of the proof of Theorem 1.7 is item g). For its proof we need a result from \[G-S3\]:
###### Lemma 2.1.
Let $`F/๐ฝ_q`$ be a function field and let $`G_1/F`$ and $`G_2/F`$ be linear disjoint Galois extensions of $`F`$, both of degree $`p=\mathrm{char}(๐ฝ_q)`$. Denote by $`G=G_1G_2`$ the composite field of $`G_1`$ and $`G_2`$. Let $`Q`$ be a place of $`G`$ and denote by $`Q_1,Q_2`$ and $`P`$ its restrictions to the subfields $`G_1,G_2`$ and $`F`$. Suppose that we have
$$d(Q_i|P)=2(e(Q_i|P)1),\text{for}i=1,2.$$
Then $`d(Q|Q_i)=2(e(Q|Q_i)1)`$ holds for $`i=1,2`$.
###### Proof.
See \[G-S3, Lemma 1\]. โ
Now we prove item f) of Theorem 1.7. First of all, it follows from items (F2), (F6), (F7), (F8) that the places $`(w=0)`$ and $`(w=\mathrm{})`$ of $`๐ฝ_q(w)`$ are the only ramified places in $`F_{n1}/๐ฝ_q(w)`$ and hence in $`E_n/๐ฝ_q(w)`$. We consider now a place $`\stackrel{~}{Q}`$ of $`E_n`$ which is ramified in the extension $`E_n/E_0`$. By items (F2), (F6), (F7), (F8), $`\stackrel{~}{Q}`$ is either a zero or a pole of the function $`w`$; i.e., $`\stackrel{~}{Q}`$ is in the support of the divisor $`A^{(n)}`$ or $`B^{(n)}`$ (notation as in item f) of Theorem 1.7).
Let $`Q:=\stackrel{~}{Q}F_{n1}`$ be the restriction of $`\stackrel{~}{Q}`$ to the field $`F_{n1}`$. We refine the extension $`F_{n1}/๐ฝ_q(w)`$ to Galois steps of degree $`p`$:
$$๐ฝ_q(w)=K_0K_1\mathrm{}K_m=F_{n1}E_n,$$
$`(2.3)`$
with $`[K_{j+1}:K_j]=p`$. Let $`P_j:=QK_j`$ for $`j=0,\mathrm{},m`$. By items (F4), (F9), the different exponents $`d(P_{j+1}|P_j)`$ are given by
$$d(P_{j+1}|P_j)=2(e(P_{j+1}|P_j)1),\text{for}j=0,\mathrm{},m1.$$
$`(2.4)`$
The Galois closure $`E_n`$ of $`F_{n1}/E_0`$ is obtained by iterated composites of the chain
$$K_0K_1\mathrm{}K_m$$
with the chains
$$\tau (K_0)\tau (K_1)\mathrm{}\tau (K_m),$$
where $`\tau `$ runs through the embeddings of $`F_{n1}/E_0`$. So we can refine the chain in (2.3) to a chain
$$๐ฝ_q(w)=K_0K_1\mathrm{}K_m=F_{n1}K_{m+1}\mathrm{}K_r=E_n,$$
where all extensions $`K_{j+1}/K_j`$ are Galois of degree $`p`$ (for $`j=0,\mathrm{},r1)`$. We set $`P_j:=\stackrel{~}{Q}K_j`$ for $`j=m+1,\mathrm{},r`$. It then follows from (2.4) and Lemma 2.1 that the different exponents $`d(P_{j+1}|P_j)`$ satisfy $`d(P_{j+1}|P_j)=2(e(P_{j+1}|P_j)1)`$, for $`j=0,\mathrm{},r1`$. Using the transitivity of different exponents (cf. \[S1, III.4.11\]) we obtain that
$$d(\stackrel{~}{Q}|P_0)=2(e(\stackrel{~}{Q}|P_0)1).$$
This finishes the proof of item g) of Theorem 1.7.
With notations as in items f) and g), the Hurwitz genus formula for the extension $`E_n/๐ฝ_q(w)`$ yields
$$\begin{array}{cc}2g(E_n)2\hfill & =2[E_n:๐ฝ_q(w)]+2e_0^{(n)}\mathrm{deg}A^{(n)}\hfill \\ & \\ & +2e_{\mathrm{}}^{(n)}\mathrm{deg}B^{(n)}2(\mathrm{deg}A^{(n)}+\mathrm{deg}B^{(n)})\hfill \\ & \\ & =2[E_n:๐ฝ_q(w)]2(\mathrm{deg}A^{(n)}+\mathrm{deg}B^{(n)}).\hfill \end{array}$$
We have used here that the divisors $`e_{\mathrm{}}^{(n)}B^{(n)}`$ and $`e_0^{(n)}A^{(n)}`$ are the pole divisor and the zero divisor of the function $`w`$ in $`E_n`$, hence their degree is equal to the degree $`[E_n:๐ฝ_q(w)]`$. We have thus proved item h) of Theorem 1.7.
From items e) and h) we see that
$$N(E_n)/g(E_n)\mathrm{}1\text{for all}n1,$$
$`(2.5)`$
Hence $`lim_n\mathrm{}N(E_n)/g(E_n)\mathrm{}1`$. By the Drinfeld-Vladut bound (see Section 1) we also have that $`lim_n\mathrm{}N(E_n)/g(E_n)\mathrm{}1`$, hence equality holds. This proves item j) and finishes the proof of Theorem 1.7.
## 3. Asymptotically Good Transitive Codes
The aim of this section is to prove Theorem 1.5. We use notation as in Theorem 1.5, in particular $`q=\mathrm{}^2`$ is a square. Let $`R,\delta 0`$ with
$$R=1\delta \frac{1}{\mathrm{}1},$$
$`(3.1)`$
and let $`ฯต>0`$. We will construct transitive codes $`C`$ over $`๐ฝ_q`$ of arbitrarily large length such that $`R(C)Rฯต`$ and $`\delta (C)\delta `$; this proves then Theorem 1.5.
Consider the tower $`=(E_0,E_1,E_2,\mathrm{})`$ of function fields over $`๐ฝ_q`$ which was constructed in Theorem 1.7. Choose an integer $`n>0`$ so large that
$$\frac{1}{\mathrm{}^n(\mathrm{}1)}<ฯต.$$
$`(3.2)`$
Let $`N:=[E_n:๐ฝ_q(z)]`$, with the function $`zE_n`$ as in Theorem 1.7, and consider the divisors $`D,G_0`$ of $`E_n`$ which are given by
$$D:=\underset{P|(z=1)}{}P\text{and}G_0:=\underset{Q|(z=\mathrm{})}{}Q.$$
$`(3.3)`$
This means: $`P`$ runs over all places of $`E_n`$ which are zeroes of the function $`(z1)`$, and $`Q`$ runs over all poles of the function $`z`$ in $`E_n`$. By Theorem 1.7 e) all these places $`P`$ are rational, and the degree of $`D`$ is $`\mathrm{deg}D=N`$. With notations as in Theorem 1.7 f), the divisor $`G_0`$ is just the divisor $`G_0=B^{(n)}`$, since the functions $`w`$ and $`z=w^\mathrm{}1`$ have the same poles. The degree of $`G_0`$ satisfies then
$$\mathrm{deg}G_0=\frac{[E_n:๐ฝ_q(w)]}{e_{\mathrm{}}^{(n)}}\frac{[E_n:๐ฝ_q(w)]}{\mathrm{}^n}=\frac{N}{\mathrm{}^n(\mathrm{}1)},$$
by Theorem 1.7 f). Hence we have that
$$(\mathrm{deg}G_0)/N<ฯต,$$
by Inequality (3.2). We choose $`r0`$ such that
$$1\delta r\frac{\mathrm{deg}G_0}{N}>1\delta ฯต$$
$`(3.4)`$
and consider the geometric Goppa code
$$C:=C_{}(D,rG_0)๐ฝ_q^N$$
associated to the divisors $`D`$ and $`rG_0`$. It is defined as follows (cf. \[S1, II.2.1\] or \[T-V\]): If $`(rG_0)E_n`$ denotes the Riemann-Roch space of the divisor $`rG_0`$ and the divisor $`D`$ is defined as $`D=P_1+\mathrm{}+P_N`$, then
$$C_{}(D,rG_0)=\{(f(P_1),\mathrm{},f(P_N))๐ฝ_q^N|f(rG_0)\}.$$
$`(3.5)`$
For the parameters $`k=dimC`$ and $`d=d(C)`$ we have the standard estimates for geometric Goppa codes (see \[S1, II.2.3\]):
$$kr\mathrm{deg}G_0+1g(E_n)\text{and}dNr\mathrm{deg}G_0.$$
Hence the information rate $`R(C)`$ satisfies
$$R(C)=\frac{k}{N}\frac{r\mathrm{deg}G_0}{N}+\frac{1}{N}\frac{g(E_n)}{N}>1\delta ฯต\frac{g(E_n)}{N},$$
by Inequality (3.4). Now observe that
$$\frac{g(E_n)}{N}\frac{1}{\mathrm{}1},$$
by Inequality (2.5), and we obtain using Equality (3.1) the following estimate for $`R(C)`$:
$$R(C)>1\delta ฯต\frac{1}{\mathrm{}1}=Rฯต.$$
For the relative minimum distance $`\delta (C)`$ we get with (3.4):
$$\delta (C)=\frac{d}{N}\frac{Nr\mathrm{deg}G_0}{N}=1\frac{r\mathrm{deg}G_0}{N}\delta .$$
These are the desired inequalities for $`R(C)`$ and $`\delta (C)`$.
It remains to show that the code $`C=C_{}(D,rG_0)`$ that we constructed above is in fact a transitive code. To this end we consider the Galois group of the extension $`E_n/E_0`$,
$$\mathrm{\Gamma }:=\mathrm{Gal}(E_n/E_0).$$
The places $`P_1,\mathrm{},P_N`$ in the support of the divisor $`D`$ are exactly the places of $`E_n`$ lying above the place $`(z=1)`$; hence $`\mathrm{\Gamma }`$ acts transitively on the set $`\{P_1,\mathrm{},P_N\}`$, see \[S1, III.7.1\]. The divisor $`rG_0`$ is obviously invariant under the action of $`\mathrm{\Gamma }`$. Therefore $`\mathrm{\Gamma }`$ acts on the code $`C=C_{}(D,rG_0)`$ as a transitive permutation group in the following way (see \[S1, VII.3.3\]): for $`\sigma G`$ and $`f(rG_0)`$,
$$\sigma (f(P_1),\mathrm{},f(P_N))=(f(\sigma P_1),\mathrm{},f(\sigma P_N)).$$
This completes the proof of Theorem 1.5 โ
###### Remark 3.1.
It is an obvious idea to prove the existence of asymptotically good cyclic codes in a similar manner. One should start with a tower $`=(H_0,H_1,H_2,\mathrm{})`$ of function fields over $`๐ฝ_q`$, where all extensions $`H_n/H_0`$ are cyclic Galois extensions; then one can do the same construction of codes as in the proof of Theorem 1.5 above. However, this method does not work: it is known that the limit $`\lambda ()=lim_n\mathrm{}N(H_n)/g(H_n)`$ of such a โcyclicโ tower $``$ is zero, see \[F-P-S\].
###### Remark 3.2.
The notion of โinformation rateโof a code can be defined also for non-linear codes $`C๐ฝ_q^N`$, by setting $`R(C):=\mathrm{log}_q(|C|/N)`$. Using this definition, one obtains in an obvious manner an analogue of the function $`\alpha _q(\delta )`$ by considering all codes over $`๐ฝ_q`$, not just linear codes. We denote this analoguous function again by $`\alpha _q(\delta )`$. It was shown in \[N-O1\] and \[S-X\] that in a large open subinterval of , the Tsfasman-Vladut-Zink bound
$$\alpha _q(\delta )1\delta A(q)^1$$
$`(3.6)`$
can be improved to
$$\alpha _q(\delta )1\delta A(q)^1+\mathrm{log}_q(1+q^3).$$
$`(3.7)`$
A further slight improvement of Inequality (3.7) was very recently found in \[N-O2\]. However, it seems that the codes which were constructed in \[N-O1,2\] and \[S-X\] in order to prove Inequality (3.6) do not have any algebraic or combinatoric structure. By combining the method of \[S-X\] with our proof of Theorem 1.5 we can now show that the lower bound (3.7) for $`\alpha _q(\delta )`$ is attained by transitive non-linear codes.
###### Theorem 3.3.
Assume that $`q=\mathrm{}^2`$ is a square, and set
$$\delta ^{}:=12/(\mathrm{}1)(4q2)/((q1)(q^3+1)).$$
Then the bound
$$\alpha _q(\delta )1\delta A(q)^1+\mathrm{log}_q(1+q^3)$$
is attained by transitive codes, for all $`\delta `$ in the interval $`(0,\delta ^{})[0,1]`$.
###### Proof.
(Sketch). We recall briefly the code construction given in \[S-X\]. One considers a function field $`F`$ over $`๐ฝ_q`$ of genus $`g`$ and a set $`๐ซ=\{P_1,\mathrm{},P_N\}`$ of $`N`$ distinct rational places of $`F`$. Let $`H0`$ be a divisor of $`F`$ of degree $`\mathrm{deg}H2g1`$ with $`\mathrm{supp}(H)๐ซ=\mathrm{}`$ and consider divisors $`G`$ of the form
$$G=\underset{j=1}{\overset{t}{}}m_{i_j}P_{i_j}\text{with}1i_1<i_2<\mathrm{}<i_tN,m_{i_j}1\text{and}\mathrm{deg}G=s.$$
$`(3.8)`$
Define the set $`M_H(G)`$ as follows:
$$M_H(G):=\{x(H+G)|v_{P_{i_j}}(x)=m_{i_j}\text{for}1jt\}.$$
Choose integers $`s,t`$ with $`1tN`$ and $`st`$, and set
$$S:=S(H,๐ซ,s,t):=\underset{G}{}M_H(G),$$
where $`G`$ runs over all divisors of the form (3.8). It is clear that $`M_H(G_1)M_H(G_2)=\mathrm{}`$ if $`G_1G_2`$. Hence we can define a map $`\phi :S๐ฝ_q^N`$ in the following way: for $`xM_H(G)`$ put $`\phi (x)=(x_1,\mathrm{},x_N)`$ with
$$x_i=\{\begin{array}{cc}x(P_i)& \text{if}P_i\mathrm{supp}(G),\hfill \\ 0& \text{if}P_i\mathrm{supp}(G).\hfill \end{array}$$
Thus we obtain a (non-linear) code $`C=C(H,๐ซ,s,t)`$ by setting
$$C(H,๐ซ,s,t):=\phi (S)๐ฝ_q^N.$$
If the function field $`F`$ runs through a sequence of function fields $`(F_0,F_1,F_2,\mathrm{})`$ over $`๐ฝ_q`$ with $`lim_n\mathrm{}N(F_n)/g(F_n)=\sqrt{q}1`$, one can choose the set $`๐ซ`$, the divisor $`H`$ and the integers $`s,t`$ in such a way that the corresponding codes $`C(H,๐ซ,s,t)`$ reach the bound (3.7), see \[S-X, Prop.3.3 and Thm.3.4.\]
In order to obtain transitive codes with the above construction, we use again the function fields $`E_n`$ of the tower $`=(E_0,E_1,E_2,\mathrm{})`$ from Theorem 1.7. We choose the set $`๐ซ`$ as in the proof of Theorem 1.5; i.e.,
$$๐ซ=\{P|P\text{is a zero of the function}z1\text{in}E_n\},$$
see (3.3). The divisor $`H`$ is chosen as
$$H=m_0G_0,$$
with the divisor $`G_0`$ of $`E_n`$ as in (3.3). Since the set $`๐ซ`$ and the divisor $`G_0`$ are invariant under the action of the group $`\mathrm{\Gamma }=\mathrm{Gal}(E_n/E_0)`$, it follows immediately that the corresponding codes $`C(H,๐ซ,s,t)๐ฝ_q^N`$ are $`\mathrm{\Gamma }`$-invariant; they are therefore transitive codes. โ
## 4. Asymptotically Good Self-Dual and Self-Orthogonal Codes
In this section we shall prove Theorem 1.6. First we recall some definitions and facts.
###### Definition 4.1.
Let $`C๐ฝ_q^N`$ be a linear code, and let $`\underset{ยฏ}{a}=(a_1,\mathrm{},a_N)๐ฝ_q^N`$ with non-zero components $`a_1,\mathrm{},a_N0`$. We set
$$\underset{ยฏ}{a}C:=\{(a_1c_1,\mathrm{},a_Nc_N)๐ฝ_q^N|(c_1,\mathrm{},c_N)C\},$$
and call the codes $`C`$ and $`\underset{ยฏ}{a}C`$ equivalent.
It is clear that equivalent codes have the same parameters (length, dimension, minimum distance, information rate, relative minimum distance). Note however that the automorphism groups $`\mathrm{Aut}(C)`$ and $`\mathrm{Aut}(\underset{ยฏ}{a}C)`$ are in general non-isomorphic.
###### Definition 4.2.
i) A code $`C๐ฝ_q^N`$ is called self-dual if $`C`$ is equal to its dual code $`C^{}`$. The code $`C`$ is called self-orthogonal if $`CC^{}`$.
ii) A code $`C`$ is called iso-dual if $`C`$ is equivalent to its dual code $`C^{}`$, cf. \[P-H\].
iii) A code $`C`$ is called iso-orthogonal if $`C`$ is equivalent to a subcode of $`C^{}`$.
Now let $`F/๐ฝ_q`$ be a function field and let $`P_1,\mathrm{},P_N`$ be distinct rational places of $`F`$. Let $`D=P_1+\mathrm{}+P_N`$ and let $`G`$ be a divisor with $`\mathrm{supp}D\mathrm{supp}G=\mathrm{}`$. As in Section 3, we consider the geometric Goppa code (cf. (3.5))
$$C_{}(D,G):=\{(f(P_1),\mathrm{},f(P_N))๐ฝ_q^N|f(G)\}.$$
$`(4.1)`$
###### Proposition 4.3.
Let $`D`$ and $`G`$ be divisors of the function field $`F/๐ฝ_q`$ as above and consider the code $`C=C_{}(D,G)`$ as defined in (4.1). Suppose that $`\eta `$ is a differential of $`F`$ with the property $`v_{P_i}(\eta )=1`$ for $`i=1,\mathrm{},N`$. Then the dual code $`C^{}=C_{}(D,G)^{}`$ is given by
$$C^{}=\underset{ยฏ}{a}C_{}(D,H),$$
with $`H:=DG+(\eta )`$ and $`\underset{ยฏ}{a}=(\mathrm{res}_{P_1}(\eta ),\mathrm{},\mathrm{res}_{P_N}(\eta ))`$.
###### Proof.
See \[S1, Cor.2.7\]. โ
We want to apply Proposition 4.3 to geometric Goppa codes which are defined by means of the function fields $`E_n`$ in the tower $`=(E_0,E_1,E_2,\mathrm{})`$ of Theorem 1.7. So we must find an appropriate differential $`\eta `$ of $`E_n`$ having the properties as required in Proposition 4.3.
###### Proposition 4.4.
We assume all notations from Theorem 1.7 and consider the differential
$$\eta :=\frac{dw}{1z}$$
of the function field $`E_n`$ (with $`n2`$). Then the following holds:
1. The divisor of $`\eta `$ in $`E_n`$ is given by
$$(\eta )=a_nA^{(n)}+b_nB^{(n)}D^{(n)},$$
where the divisors $`A^{(n)}>0`$ and $`B^{(n)}>0`$ are as in Theorem 1.7 f) , the integers $`a_n>0`$ and $`b_n>0`$ satisfy $`a_nb_n0\mathrm{mod}2`$, and the divisor $`D^{(n)}`$ is the sum over all zeroes of the function $`z1`$ in $`E_n`$; i.e.,
$$D^{(n)}=\underset{P|(z=1)}{}P.$$
2. The residue of the differential $`\eta `$ at a place $`P`$, which is a zero of $`z1`$ in $`E_n`$, is an element of $`๐ฝ_{\mathrm{}}^\times `$; i.e.,
$$\mathrm{res}_P(\eta )=\alpha _P\text{with}\alpha _P^\mathrm{}1=1.$$
###### Proof.
i) By Theorem 1.7 f), the principal divisor of the function $`w`$ in $`E_n`$ is
$$(w)^{E_n}=e_0^{(n)}A^{(n)}e_{\mathrm{}}^{(n)}B^{(n)},$$
and by item g) of Theorem 1.7, the different of $`E_n/๐ฝ_q(w)`$ is
$$\mathrm{Diff}(E_n/๐ฝ_q(w))=2(e_0^{(n)}1)A^{(n)}+2(e_{\mathrm{}}^{(n)}1)B^{(n)}.$$
It follows that the divisor of the differential $`dw`$ in $`E_n`$ is given by (see \[S1, III.4.6\])
$$(dw)=2e_0^{(n)}B^{(n)}+\mathrm{Diff}(E_n/๐ฝ_q(w))=2e_0^{(n)}A^{(n)}2A^{(n)}2B^{(n)}.$$
The divisor of the function $`1z`$ in $`E_n`$ is
$$(1z)^{E_n}=D^{(n)}(\mathrm{}1)e_{\mathrm{}}^{(n)}B^{(n)},$$
and we obtain the divisor of the differential $`\eta =dw/(1z)`$ as follows:
$$\begin{array}{cc}(\eta )\hfill & =2e_0^{(n)}A^{(n)}2A^{(n)}2B^{(n)}D^{(n)}+(\mathrm{}1)e_{\mathrm{}}^{(n)}B^{(n)}\hfill \\ & =a_nA^{(n)}+b_nB^{(n)}D^{(n)},\hfill \end{array}$$
with $`a_n>0,b_n>0`$ and $`a_nb_n0mod2`$.
ii) Let $`P`$ be a place of $`E_n`$ which is a zero of the function $`z1`$. The element $`t:=z1`$ is a $`P`$-prime element. From the equation $`w^\mathrm{}1=z=t+1`$ we obtain
$$dt=(\mathrm{}1)w^\mathrm{}2dw=\frac{w^\mathrm{}1}{w}dw=\frac{1+t}{w}dw,$$
hence
$$\eta =\frac{dw}{1z}=\frac{1}{t}dw=\frac{w}{1+t}\frac{1}{t}dt.$$
Let $`\alpha :=w(P)๐ฝ_q`$ be the residue class of $`w`$ at the place $`P`$; then
$$\frac{w}{1+t}\alpha \mathrm{mod}P\text{and therefore}\mathrm{res}_P(\eta )=\alpha .$$
Since $`\alpha ^\mathrm{}1=w^\mathrm{}1(P)=z(P)=1`$, we conclude that $`\alpha ๐ฝ_{\mathrm{}}\{0\}`$. โ
Now we can construct certain geometric Goppa codes which are associated to the function field $`E_n`$ in the tower $`=(E_0,E_1,E_2,\mathrm{})`$ of Theorem 1.7. For the rest of this section we fix notations as above; in particular we will use without further explanation the divisors $`A^{(n)},B^{(n)}`$ and $`D^{(n)}`$, the differential $`\eta `$ and the integers $`a_n`$ and $`b_n`$ as in Proposition 4.4.
###### Definition 4.5.
For integers $`a,b`$ with $`0aa_n`$ and $`0bb_n`$, we define the code $`C_{a,b}^{(n)}`$ by
$$C_{a,b}^{(n)}:=C_{}(D^{(n)},aA^{(n)}+bB^{(n)}).$$
###### Remarks 4.6.
i) It is clear that the codes $`C_{a,b}^{(n)}`$ are transitive. This follows as in Section 3 from the fact that the Galois group $`\mathrm{\Gamma }=\mathrm{Gal}(E_n/E_0)`$ acts transitively on the places $`P\mathrm{supp}(D^{(n)})`$ and leaves the divisors $`A^{(n)}`$ and $`B^{(n)}`$ invariant.
ii) For $`n\mathrm{}`$, the codes $`C_{a,b}^{(n)}`$ attain the Tsfasman-Vladut-Zink bound $`\alpha _q(\delta )1\delta 1/(\mathrm{}1)`$, for all $`\delta (0,11/(\mathrm{}1))`$. This is proved in the same manner as Theorem 1.5 (see Section 3).
###### Proposition 4.7.
We write $`D^{(n)}=P_1+\mathrm{}+P_N`$, with $`N=[E_n:E_0]`$, and set
$$\underset{ยฏ}{u}:=(\mathrm{res}_{P_1}\eta ,\mathrm{},\mathrm{res}_{P_N}\eta )(๐ฝ_q^\times )^N.$$
Then the dual of the code $`C_{a,b}^{(n)}`$ is given by
$$(C_{a,b}^{(n)})^{}=\underset{ยฏ}{u}C_{a_na,b_nb}^{(n)}.$$
###### Proof.
The differential $`\eta `$ satisfies the condition $`v_{P_i}(\eta )=1`$, for $`i=1,\mathrm{},N`$. Hence we can apply Proposition 4.3 and obtain
$$(C_{a,b}^{(n)})^{}=\underset{ยฏ}{u}C_{}(D^{(n)},H),$$
with
$$\begin{array}{cc}H\hfill & =D^{(n)}(aA^{(n)}+bB^{(n)})+(\eta )\hfill \\ & \\ & =D^{(n)}(aA^{(n)}+bB^{(n)})+(a_nA^{(n)}+b_nB^{(n)}D^{(n)})\hfill \\ & \\ & =(a_na)A^{(n)}+(b_nb)B^{(n)}.\hfill \end{array}$$
We have used here Proposition 4.4 i). โ
The following corollary is an obvious consequence from Proposition 4.7, cf. Definition 4.2.
###### Corollary 4.8.
i) For $`0aa_n/2`$ and $`0bb_n/2`$, the code $`C_{a,b}^{(n)}`$ is transitive and iso-orthogonal.
ii) For $`a=a_n/2`$ and $`b=b_n/2`$, the code $`C_{a,b}^{(n)}`$ is iso-dual.
###### Corollary 4.9.
i) For $`0aa_n/2`$ and $`0bb_n/2`$, the code $`C_{a,b}^{(n)}`$ is equivalent to a self-orthogonal code $`\stackrel{~}{C}_{a,b}^{(n)}`$.
ii) For $`a=a_n/2`$ and $`b=b_n/2`$, the code $`C_{a,b}^{(n)}`$ is equivalent to a self-dual code $`\stackrel{~}{C}_{a,b}^{(n)}`$.
###### Proof.
The components of the vector $`\underset{ยฏ}{u}=(\mathrm{res}_{P_1}\eta ,\mathrm{},\mathrm{res}_{P_N}\eta )`$ in Proposition 4.7 are in $`๐ฝ_{\mathrm{}}^\times `$, by Proposition 4.4 ii). So we can write $`\mathrm{res}_{P_i}\eta =v_i^2`$ with $`v_i๐ฝ_q^\times `$ (note that $`q=\mathrm{}^2`$). We set $`\underset{ยฏ}{v}:=(v_1,\mathrm{},v_N)`$; then the code
$$\stackrel{~}{C}_{a,b}^{(n)}:=\underset{ยฏ}{v}C_{a,b}^{(n)}$$
is self-orthogonal, resp. self-dual. โ
Theorem 1.6 is now an immediate consequence of Corollary 4.9 and Remark 4.6 ii).
###### Remark 4.10.
The existence of asymptotically good sequences $`(C_j)_{j0}`$ of isodual geometric Goppa codes over $`๐ฝ_q`$ (with $`q=\mathrm{}^2`$) was already proved in \[Sch\]. However, the codes which were constructed there attain only the lower bound
$$\underset{j\mathrm{}}{lim}\delta (C_j)\frac{1}{2}\frac{1}{\mathrm{}3}.$$
$`(4.2)`$
The codes $`\stackrel{~}{C}_n:=\stackrel{~}{C}_{a,b}^{(n)}`$ in Corollary 4.9 ii) are not only iso-dual but they are self-dual. They satisfy the bound (see Theorem 1.6 ii))
$$\underset{j\mathrm{}}{lim}\delta (\stackrel{~}{C}_j)\frac{1}{2}\frac{1}{\mathrm{}1},$$
which is better than Inequality (4.2).
## 5. Conclusion
Let $`q=\mathrm{}^2`$ be a square. We have shown in this paper, that the following classes of linear codes over $`๐ฝ_q`$ attain the Tsfasman-Vladut-Zink bound:
1. transitive codes (Theorem 1.5),
2. transitive iso-orthogonal codes (Corollary 4.8),
3. transitive iso-dual codes (Corollary 4.8),
4. self-orthogonal codes (Theorem 1.6),
5. self-dual codes (Theorem 1.6).
In particular, the above classes of codes are better than the Gilbert-Varshamov bound, for all squares $`q49`$. The class of non-linear transitive codes attains an even better bound than the Tsfasman-Vladut-Zink bound (Theorem 3.3).
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# 1 Introduction
## 1 Introduction
Branching fraction and $`CP`$ asymmetry measurements of charmless $`B`$ meson decays provide valuable constraints for the determination of the unitarity triangle constructed from elements of the Cabibbo-Kobayashi-Maskawa quark-mixing matrix . In particular, the angle $`\alpha \mathrm{arg}\left[V_{td}V_{tb}^{}/V_{ud}V_{ub}^{}\right]`$ of the unitarity triangle can be extracted from decays of the $`B`$ meson to $`\rho ^\pm \pi ^{}`$ final states . However, the extraction is complicated by the interference of decay amplitudes with differing weak and strong phases. One strategy to overcome this problem is to perform an SU(2) analysis that uses all $`\rho \pi `$ final states . Assuming isospin symmetry, the angle $`\alpha `$ can be determined free of hadronic uncertainties from a pentagon relation formed in the complex plane by the five decay amplitudes $`B^0\rho ^+\pi ^{}`$, $`B^0\rho ^{}\pi ^+`$, $`B^0\rho ^0\pi ^0`$, $`B^+\rho ^+\pi ^0`$ and $`B^+\rho ^0\pi ^+`$. These amplitudes can be determined from measurements of the corresponding decay rates and $`CP`$ asymmetries. While all these modes have been measured, the current experimental uncertainties need to be reduced substantially for a determination of $`\alpha `$. Here we present an update to a previous measurement of the $`B^\pm \rho ^\pm \pi ^0`$ branching fraction and $`CP`$ asymmetry
$$๐_{CP}=\frac{N(B^{}\rho ^{}\pi ^0)N(B^+\rho ^+\pi ^0)}{N(B^{}\rho ^{}\pi ^0)+N(B^+\rho ^+\pi ^0)}.$$
## 2 Data Set and Candidate Selection
The data used in this analysis were collected with the BABAR detector at the PEP-II asymmetric-energy $`e^+e^{}`$ storage ring at SLAC. Charged-particle trajectories are measured by a five-layer double-sided silicon vertex tracker and a 40-layer drift chamber located within a 1.5-T solenoidal magnetic field. Charged hadrons are identified by combining energy-loss information from tracking with the measurements from a ring-imaging Cherenkov detector. Photons are detected by a CsI(Tl) crystal electromagnetic calorimeter with an energy resolution of $`\sigma _E/E=0.023(E/\mathrm{GeV})^{1/4}0.014`$. The magnetic flux return is instrumented for muon and $`K_L^0`$ identification.
The data sample includes $`232\pm 3`$ million $`B\overline{B}`$ pairs collected at the $`\mathrm{{\rm Y}}(4S)`$ resonance, corresponding to an integrated luminosity of 211 $`\mathrm{fb}^1`$. It is assumed that neutral and charged $`B`$ meson pairs are produced in equal numbers . In addition, 22 $`\mathrm{fb}^1`$ of data collected 40 $`\mathrm{MeV}`$ below the $`\mathrm{{\rm Y}}(4S)`$ resonance mass are used for background studies.
We perform full detector Monte Carlo (MC) simulations equivalent to 460 $`\mathrm{fb}^1`$ of generic $`B\overline{B}`$ decays and 140 $`\mathrm{fb}^1`$ of continuum quark-antiquark production events. In addition, we simulate over 50 exclusive charmless $`B`$ decay modes, including 1.4 million signal $`B^\pm \rho ^\pm \pi ^0`$ decays.
$`B`$ meson candidates are reconstructed from one charged track and two neutral pions, with the following requirements:
Track quality. The charged track used to form the $`B^\pm \rho ^\pm \pi ^0`$ candidate is required to have at least 12 hits in the drift chamber, to have a transverse momentum greater than 0.1 $`\mathrm{GeV}/\mathrm{c}`$, and to be consistent with originating from a $`B`$-meson decay. Its signal in the tracking and Cherenkov detectors is required to be consistent with that of a pion. We remove tracks that pass electron selection criteria based on $`dE/dx`$ and calorimeter information.
$`\pi ^0`$ quality. Neutral pion candidates are formed from two photons, each with a minimum energy of 0.03 $`\mathrm{GeV}`$ and a lateral moment of their shower energy deposition greater than zero and less than 0.6. The angular acceptance of photons is restricted to exclude parts of the calorimeter where showers are not fully contained. We require the photon clusters forming the $`\pi ^0`$ to be separated in space, with a $`\pi ^0`$ energy of at least 0.2 $`\mathrm{GeV}`$ and an invariant mass between 0.10 and 0.16 $`\mathrm{GeV}/\mathrm{c}^2`$.
Kinematic requirements. Two kinematic variables, $`\mathrm{\Delta }E=E_B^{}\sqrt{s}/2`$ and the beam energy substituted mass of the $`B`$-meson $`m_{\mathrm{ES}}=\sqrt{(s/2+๐ฉ_0๐ฉ_B)^2/E_0^2๐ฉ_B^2}`$, are used for the final selection of events. Here $`E_B^{}`$ is the energy of the $`B`$ meson candidate in the center-of-mass frame, $`E_0`$ and $`\sqrt{s}`$ are the total energies of the $`e^+e^{}`$ system in the laboratory and center-of-mass frames, respectively; $`๐ฉ_0`$ and $`๐ฉ_B`$ are the three-momenta of the $`e^+e^{}`$ system and the $`B`$ candidate in the laboratory frame, respectively. For correctly reconstructed $`\rho ^\pm \pi ^0`$ candidates $`\mathrm{\Delta }E`$ peaks at zero, while final states with a charged kaon, such as $`B^\pm K^\pm \pi ^0`$, shift $`\mathrm{\Delta }E`$ by approximately 80 $`\mathrm{MeV}`$ on average. Events are selected with $`5.20<m_{\mathrm{ES}}<5.29\mathrm{GeV}/\mathrm{c}^2`$ and $`|\mathrm{\Delta }E|<0.20\mathrm{GeV}`$. The $`\mathrm{\Delta }E`$ limits help remove background from two- and four-body $`B`$ decays at a small cost to signal efficiency.
Continuum suppression. Continuum quark-antiquark production is the dominant background. To suppress it, we select only those events where the angle $`\theta _{\mathrm{Sph}}^\mathrm{B}`$ in the center-of-mass frame between the direction of the $`B`$ meson candidate and the sphericity axis of the rest of the event satisfies $`|\mathrm{cos}\theta _{\mathrm{Sph}}^\mathrm{B}|<0.9`$. In addition, we construct a non-linear discriminant, implemented as an artificial neural network, that uses three input parameters: the zeroth- and second-order Legendre event shape polynomials $`L_0,L_2`$ calculated from the momenta and polar angles of all charged particle and photon candidates not associated with the $`B`$ meson candidate, and the output of a multivariate, non-linear $`B`$ meson candidate tagging algorithm . The output $`ANN`$ of the artificial neural network is peaked at 0.5 for continuum-like events and at 1.0 for $`B`$ decays. We require $`ANN>0.63`$ for our event selection.
$`\rho `$ mass window. To further improve the signal-to-background ratio we restrict the invariant mass of the $`\rho `$ candidate to $`0.55<m_{\pi \pi }<0.95\mathrm{GeV}/\mathrm{c}^2`$.
Multiple candidates. Neutral pion combinatorics can lead to more than one $`B`$-meson candidate per event. We choose the best candidate based on a $`\chi ^2`$ formed from the measured masses of the two $`\pi ^0`$ candidates within the event compared to the known $`\pi ^0`$ mass . In the case of multiple charged pion candidates the choice is random.
Efficiency. The total $`B^\pm \rho ^\pm \pi ^0`$ selection efficiency is $`15.4\pm 0.1\%`$. In MC studies, the signal candidate is correctly reconstructed 54.9% of the time. The remaining candidates come from self-cross-feed (SCF, 37.5%) and mistag events (7.6%). We define SCF events as those where one or more elements of the $`B`$-candidate reconstruction are incorrect except for its charge. They stem primarily from swapping the low energy $`\pi ^0`$ from the resonance with another from the rest of the event. Signal events reconstructed with the wrong charge are classified as mistag events. Both SCF and mistag events emulate signal events, however the resolution in $`m_{\mathrm{ES}}`$ and $`\mathrm{\Delta }E`$ tends to be worse.
## 3 Background Contributions
MC events are used to study backgrounds from other $`B`$-meson decays. The dominant contribution comes from $`bc`$ transitions; the next most important is from charmless $`B`$-meson decays. The latter tend to be more problematic as the branching fractions are often poorly known, and because they may peak at the same invariant mass as the signal $`B^\pm \rho ^\pm \pi ^0`$ events. Seventeen individual charmless modes show a significant contribution once the event selection has been applied (Table 1). These modes are added into the fit fixed at the yield and asymmetry determined by the simulation, given an assumed branching fraction. Wherever branching fractions are not available, we use half the upper limit. If no charge asymmetry measurement is available, we assume zero asymmetry.
$`\rho ^{}`$ resonances. Although all other states which decay like the $`\rho `$ to $`\pi \pi ^0`$ โ subsequently referred to as $`\rho ^{}`$ โ lie outside our $`\rho (770)`$ mass cut, a contribution to our signal cannot be *a priori* ruled out. The only non-strange vector resonances which can decay to two pions are the $`\rho (1450)`$ and the $`\rho (1700)`$. To account for the possible presence of these modes, a fit to the $`B^\pm \rho ^\pm \pi ^0`$ yield is performed in a sideband of the invariant mass using the three variables $`m_{\mathrm{ES}}`$, $`\mathrm{\Delta }E`$, and $`ANN`$. The mass window is chosen to be as far as possible from the $`\rho (770)`$ mass, centered near the pole of the $`\rho (1700)`$ at $`1.5<m_{\pi \pi }<2.0`$ $`\mathrm{GeV}/c^2`$. The fitted yield for the $`B^\pm \rho ^\pm \pi ^0`$ decay is then extrapolated into the nominal region. Although the choice of mass range is motivated by the $`\rho (1700)`$, any yield seen is attributed entirely to the $`\rho (1450)`$, which is the closer of the two resonances to the signal. From the $`B^\pm \rho ^\pm (1450)\pi ^0`$ MC, the ratio of candidates in the sideband to candidates in the signal mass region is approximately 12.6:1. The fit in the sideband yields $`101\pm 32`$ events, resulting in an estimate of the $`\rho ^{}`$ background of $`8`$ events. We assign a conservative systematic uncertainty of 100% to this number. The $`\rho ^{}`$ then enters into the nominal fit with PDFs constructed from $`B^\pm \rho ^\pm (1450)\pi ^0`$ MC simulation.
Non-resonant decays to $`\pi ^\pm \pi ^0\pi ^0`$. The non-resonant $`B^\pm \pi ^\pm \pi ^0\pi ^0`$ branching fraction has, to date, not been measured. To estimate the significance of its contribution we select a region of the Dalitz plot โ defined by the triangle $`(m_{\pi ^\pm \pi _1^0}^2,m_{\pi ^\pm \pi _2^0}^2)=(6,6),(6,15),(11,11)`$ GeV$`{}_{}{}^{2}/c^4`$ โ that is far from the signal as well as $`\rho (1450)`$ and higher resonances and which has low levels of continuum background. A likelihood fit in this region yields $`5.1\pm 7.6`$ non-resonant events in a data sample of 1100 events. This is consistent with zero. The non-resonant contribution is therefore deemed negligible.
## 4 The Maximum Likelihood Fit
An unbinned maximum likelihood fit to the variables $`m_{\mathrm{ES}}`$, $`\mathrm{\Delta }E`$, $`m_{\pi \pi }`$, and $`ANN`$ is used to extract the total number of signal $`B^\pm \rho ^\pm \pi ^0`$ and continuum background events and their respective charge asymmetries. The likelihood for the selected sample is given by the product of the probability density functions (PDF) for each individual candidate, multiplied by the Poisson factor:
$$=\frac{1}{N!}e^N^{}(N^{})^N\underset{i=1}{\overset{N}{}}๐ซ_i,$$
where $`N`$ and $`N^{}`$ are the number of observed and expected events, respectively. The PDF $`๐ซ_i`$ for a given event $`i`$ is the sum of the signal and background terms:
$`๐ซ_i`$ $`=`$ $`N^{\mathrm{Sig}}\times {\displaystyle \frac{1}{2}}[(1Q_iA^{\mathrm{Sig}})f_{\mathrm{Sig}}๐ซ_i^{\mathrm{Sig}}`$
$`+(1Q_iA^{\mathrm{Sig}})f_{\mathrm{SCF}}๐ซ_{\mathrm{SCF},i}^{\mathrm{Sig}}`$
$`+(1+Q_iA^{\mathrm{Sig}})f_{\mathrm{Mis}}๐ซ_{\mathrm{Mis},i}^{\mathrm{Sig}}]`$
$`+{\displaystyle \underset{j}{}}N_j^{\mathrm{Bkg}}\times {\displaystyle \frac{1}{2}}(1Q_iA_j^{\mathrm{Bkg}})๐ซ_{j,i}^{\mathrm{Bkg}},`$
where $`Q_i`$ is the charge of the pion in the event, $`N^{\mathrm{Sig}}(N_j^{\mathrm{Bkg}})`$ and $`A^{\mathrm{Sig}}(A_j^{\mathrm{Bkg}})`$ are the yield and asymmetry for signal and background component $`j`$, respectively. The fractions of true signal ($`f_{\mathrm{Sig}}`$), SCF signal ($`f_{\mathrm{SCF}}`$), and wrong-charge mistag events ($`f_{\mathrm{Mis}}`$) are fixed to the numbers obtained from MC simulations (Section 2). The $`j`$ individual background terms comprise continuum, $`bc`$ decays, $`\rho ^{}`$, and 17 exclusive charmless $`B`$ decay modes. The PDF for each component, in turn, is the product of the PDFs for each of the fit input variables, $`๐ซ=๐ซ_{\mathrm{m}_{\mathrm{ES}},\mathrm{\Delta }E}๐ซ_{ANN}๐ซ_{m_{\pi \pi }}.`$ Due to correlations between $`\mathrm{\Delta }E`$ and $`m_{\mathrm{ES}}`$, the $`๐ซ_{\mathrm{m}_{\mathrm{ES}},\mathrm{\Delta }E}`$ for signal and all background from $`B`$ decays are described by two-dimensional non-parametric PDFs obtained from MC events. For continuum background, $`๐ซ_{\mathrm{m}_{\mathrm{ES}},\mathrm{\Delta }E}`$ is the product of two orthogonal one-dimensional parametric PDFs; $`m_{\mathrm{ES}}`$ is well described by an empirical phase-space threshold function and $`\mathrm{\Delta }E`$ is parameterized with a second degree polynomial. The parameters of the continuum PDFs are floated in the fit, with $`m_{\mathrm{ES}}`$ constrained to masses below 5.29 $`\mathrm{GeV}/\mathrm{c}^2`$. $`ANN`$ is described by the product of an exponential and a polynomial function for continuum background and by a Crystal Ball function for all other modes. For $`๐ซ_{m_{\pi \pi }}`$, one-dimensional non-parametric PDFs obtained from MC events are used to describe all modes except the signal mode itself, which is described by a non-relativistic Breit-Wigner line-shape. The parameters for this PDF are held fixed to the MC values and varied within errors to estimate systematic uncertainties.
A number of cross checks confirm that the fit is unbiased. In 1000 separate MC pseudo-experiments we generate the expected number of events for the various fit components before using the maximum likelihood fit to extract the yields and asymmetries. The distributions for each component are generated from the componentโs PDF, giving values for the fit variables $`m_{\mathrm{ES}}`$, $`\mathrm{\Delta }E`$, $`ANN`$, and $`m_{\pi \pi }`$. The expected number of events is calculated from the branching fraction and efficiency for each individual mode. The generated number of events for each fit component is determined by fluctuating the expected number according to a Poisson distribution. The test is repeated using samples with differing asymmetry values. We repeat these MC studies using fully simulated signal $`B^\pm \rho ^\pm \pi ^0`$ events instead of generating the signal component from our PDFs. This verifies that the signal component is correctly modeled including correlations between the fit variables. As another cross check we compare the distribution of the helicity angle $`\theta _{\mathrm{Hel}}`$ between the momenta of the charged pion and the $`B`$-meson in the $`\rho `$ rest frame in data with that modeled in MC samples for a variety of cuts. Fig. 1 shows the distribution of $`\mathrm{cos}\theta _{\mathrm{Hel}}`$ for a pseudo-signal-box defined by $`m_{\mathrm{ES}}>5.265`$, $`|\mathrm{\Delta }E|<0.1`$, and $`ANN>0.8`$. We generally find our PDFs in good agreement with the data. Finally, omitting $`m_{\pi \pi }`$ as a fit variable has no significant influence on the signal yield, indicating that our treatment of $`\rho ^{}`$ background is indeed effective.
## 5 Systematic Uncertainties
Individual contributions to the systematic uncertainty are summarized in Table 2.
Absolute uncertainties on yields. We calculate the uncertainty of the continuum background estimation directly from the fit to data. The backgrounds from $`B`$ decays are determined from simulation and fixed according to their efficiencies and branching fractions. The largest individual contribution comes from the $`Ba_1^0\pi ^0`$ channel. For those individual decay modes which have been measured, we vary the number of events in the fit by their measured uncertainty. For all others we vary the amount included in the fit by $`\pm 100\%`$. For the $`bc`$ component, we fix the rate based on the number calculated from MC samples and vary the amount based on the statistical uncertainty of this number. The shifts in the fitted yields are calculated for each mode in turn and then added in quadrature to find the total systematic effect. The largest individual contribution comes from the $`\rho ^{}`$ estimation.
To take into account the variation of the two-dimensional non-parametric PDFs used for $`\mathrm{\Delta }E`$ and $`m_{\mathrm{ES}}`$, we smear the MC-generated distributions from which the PDFs are derived. This is effectively done by varying the kernel bandwidth up to twice its original value. For $`m_{\pi \pi }`$ and $`ANN`$, the parameterizations determined from fits to MC events are varied by one standard deviation. The systematic uncertainties are determined using the altered PDFs and fitting to the final data sample. The overall shifts in the central value are taken as the size of the systematic uncertainty.
We vary the SCF fraction by a conservative estimate of its relative uncertainty ($`\pm 10\%`$) and assign the shift in the fitted number of signal events as the systematic uncertainty of the SCF fraction.
To account for differences in the neutral particle reconstruction between data and MC simulation, the signal PDF distribution in $`\mathrm{\Delta }`$E is offset by $`\pm 5\mathrm{MeV}`$ and the data is then refitted. The larger of the two shifts in the central value of the yield is 2.2 events, which is taken as the systematic uncertainty for this effect.
Relative uncertainties on the branching fraction. Corrections to the $`\pi ^0`$ energy resolution and efficiency, determined using various control samples, add a systematic uncertainty of 7.2%. A relative systematic uncertainty of $`1\%`$ is assumed for the pion identification. A relative systematic uncertainty of 0.8% on the efficiency for a single charged track is applied. Adding all the above contributions in quadrature gives a relative systematic uncertainty on the branching fraction of 7.3%. Another contribution of $`1.1\%`$ comes from the uncertainty on the total number of $`B`$ events.
Uncertainties on the charge asymmetry. To calculate the effects of systematic shifts in the charge asymmetries of background modes, each mode is varied by its measured uncertainty. For contributions with no measurement, we assume zero asymmetry and assign an uncertainty of 20%, motivated by the largest charge asymmetry measured in any mode so far . The individual shifts are then added in quadrature to find the total systematic uncertainty. In addition, the effect of altering the normalizations of the $`B`$ backgrounds affects the fitted asymmetry. The size of the shift on the fitted $`๐_{CP}`$ is taken as the size of the systematic uncertainty.
## 6 Results
The central value of the signal yield from the maximum likelihood fit is $`357\pm 49`$ events, over $`44840\pm 217`$ continuum events and an expected background of $`872\pm 62`$ events from other $`B`$ decays. We find a branching fraction and charge asymmetry of
$$\begin{array}{cc}(B^\pm \rho ^\pm \pi ^0)\hfill & =[10.0\pm 1.4(Stat.)\pm 0.9(Syst.)]\times 10^6\hfill \\ ๐_{CP}(B^\pm \rho ^\pm \pi ^0)\hfill & =0.01\pm 0.13(Stat.)\pm 0.02(Syst.).\hfill \end{array}$$
Compared against the null hypothesis, the statistical significance $`\sqrt{2\mathrm{ln}(_{Null}/_{Max})}`$ of the yield amounts to 8.7 standard deviations.
The results of the fit are illustrated in Fig. 2. The plots are enhanced in signal by selecting only those events which exceed a threshold of 0.1 (0.05 for $`ANN`$) for the likelihood ratio $`R=(N^{\mathrm{Sig}}๐ซ^{\mathrm{Sig}})/(N^{\mathrm{Sig}}๐ซ^{\mathrm{Sig}}+_iN_i^{\mathrm{Bkg}}๐ซ_i^{\mathrm{Bkg}})`$, where $`N`$ are the central values of the yields from the fit and $`๐ซ`$ are the PDFs with the projected variable integrated out. This threshold is optimized by maximizing the ratio $`S=N^{\mathrm{Sig}}ฯต^{\mathrm{Sig}}/\sqrt{N^{\mathrm{Sig}}ฯต^{\mathrm{Sig}}+_iN_i^{\mathrm{Bkg}}ฯต_i^{\mathrm{Bkg}}}`$ where $`ฯต`$ are the efficiencies after the threshold is applied. The PDF components are then scaled by the appropriate $`ฯต`$.
## 7 Conclusions
We have measured the branching fraction and charge asymmetry for the decay $`B^\pm \rho ^\pm \pi ^0`$ using a maximum likelihood fit. We obtain $`(B^\pm \rho ^\pm \pi ^0)=[10.0\pm 1.4\pm 0.9]\times 10^6`$, and $`๐_{CP}=0.01\pm 0.13\pm 0.02`$, respectively, where the first error is statistical and the second error systematic. The statistical significance of the signal is calculated to be 8.7 standard deviations. The results are in good agreement with the previous measurement .
## Acknowledgements
We are grateful for the extraordinary contributions of our PEP-II colleagues in achieving the excellent luminosity and machine conditions that have made this work possible. The success of this project also relies critically on the expertise and dedication of the computing organizations that support BABAR. The collaborating institutions wish to thank SLAC for its support and the kind hospitality extended to them. This work is supported by the US Department of Energy and National Science Foundation, the Natural Sciences and Engineering Research Council (Canada), Institute of High Energy Physics (China), the Commissariat ร lโEnergie Atomique and Institut National de Physique Nuclรฉaire et de Physique des Particules (France), the Bundesministerium fรผr Bildung und Forschung and Deutsche Forschungsgemeinschaft (Germany), the Istituto Nazionale di Fisica Nucleare (Italy), the Foundation for Fundamental Research on Matter (The Netherlands), the Research Council of Norway, the Ministry of Science and Technology of the Russian Federation, and the Particle Physics and Astronomy Research Council (United Kingdom). Individuals have received support from CONACyT (Mexico), the A. P. Sloan Foundation, the Research Corporation, and the Alexander von Humboldt Foundation.
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# Semiclassical transmission across transition states
## 1 Introduction
There has recently been a resurgence of interest in the classical transition state theory of molecular reactions. Results that were historically restricted to two degrees of freedom have been generalised to arbitrary dimensions using the construction of normally hyperbolic invariant manifolds (or NHIMโs) -. It is natural to ask how classical structure such as the NHIM is reflected in the quantum-mechanical problem, which corresponds to scattering from a multidimensional potential barrier.
An answer to this question has been offered in , where a description is given of quantum mechanical transport across a phase-space bottleneck using dynamics linearised around a certain complex periodic orbit. In using linearised dynamics these results are restricted to a classically small region of phase space and energies that are no larger than $`O(\mathrm{})`$ above a transmission threshold. In this paper it is shown how fully nonlinear dynamics may be incorporated in this approach, resulting in a description of transport which is not restricted to a small region of phase space or energy range above threshold. The current approach is based on a quantisation of a classical normal form Hamiltonian, although the final form can be expressed in such a way that explicit calculation of a normal form is not necessary.
To describe the result more concretely, let us consider a waveguide problem with configuration space coordinates $`(x,y)`$ in which $`x`$ represents longitudinal position along the waveguide and $`y`$ represents transverse vibrations. If necessary we can let $`y=(y_1,\mathrm{},y_d)`$ be multidimensional. In chemical applications, $`x`$ might be a reaction coordinate with $`x`$ large and negative corresponding to decoupled reactant molecules and $`x`$ large and positive corresponding to decoupled product molecules (see Figure 1), while $`y`$ describes internal vibrations of the reacting molecules. The quantum mechanics of this problem are described using a scattering matrix, which we write in the form
$$S(E)=\left(\begin{array}{cc}r_{RR}& t_{RP}\\ \\ t_{PR}& r_{PP}\end{array}\right),$$
where, for example, the block $`t_{PR}`$ maps asymptotic incoming states on the reacting side to the corresponding asymptotic outgoing states on the product side. In chemical jargon, $`t_{PR}`$ gives state-selected reaction rates (labelled by the incoming mode number) together with the distribution of product states (labelled by the outgoing mode number). In this paper we will describe a semiclassical approximation for the operator
$$\widehat{}(E)=t_{PR}^{}t_{PR}$$
which gives a probability of transmission for states incoming on the reactant side but which sums over outgoing states and does not give the distribution of product states. Although containing less information than the scattering matrix, this reaction operator has clear experimental relevance and, importantly in this context, admits semiclassical approximations which are considerably simpler. As described in detail in , this is because the orbits used in semiclassical approximation of $`S(E)`$ are singular near the boundary of the reacting subset of phase space whereas those used for $`\widehat{}(E)`$ are not.
The reaction operator $`\widehat{}(E)`$ has a clear relationship with the geometry of the classical transition state. It acts on the Hilbert space $`_R^{\mathrm{in}}`$ of asymptotically propagating incoming states and the classical analogue of this space is a Poincarรฉ section $`\mathrm{\Sigma }_R^{\mathrm{in}}`$ obtained by fixing the reaction coordinate and the total energy $`E`$, for which $`(y,p_y)`$ provide canonical coordinates. Let $`E`$ be greater than threshold so that there is a nonempty reacting subset $`V`$ of $`\mathrm{\Sigma }_R^{\mathrm{in}}`$ โ the boundary of $`V`$ is the intersection with $`\mathrm{\Sigma }_R^{\mathrm{in}}`$ of the stable manifold of the NHIM. Then a phase space representation of $`\widehat{}(E)`$ such as the Weyl symbol $`๐ฒ_\widehat{}(y,p_y)`$ tends to the characteristic function of the reacting region $`V`$
$$๐ฒ_\widehat{}(y,p_y)\chi _V(y,p_y)$$
in the classical limit. Moreover, for finite values of $`\mathrm{}`$, $`๐ฒ_\widehat{}(y,p_y)`$ also incorporates quantum effects such as tunnelling, especially important outside $`V`$ and near its boundary.
An explicit semiclassical approximation was presented for $`\widehat{}(E)`$ in , of the form
$$\widehat{}(E)=\frac{\widehat{๐ฏ}(E)}{1+\widehat{๐ฏ}(E)},$$
(1)
where the operator $`\widehat{๐ฏ}(E)`$ is constructed from the linear stability properties of a complex periodic orbit. This complex periodic orbit has a real initial condition in the interior of $`V`$ and returns to it after encircling the transition state region in a net imaginary time. The formula was derived by using a separable approximation of the Hamiltonian in the transition-state region to match waves propagating in the reactant and product channels. This separable approximation is valid only insofar as the transition state is small on classical scales and the result should therefore work only when the energy $`E`$ is within $`O(\mathrm{})`$ of threshold, where the Liouville volume of $`V`$ is $`O(\mathrm{}^d)`$.
We will now show that, as long as the operator $`\widehat{๐ฏ}(E)`$ is interpreted using fully nonlinear dynamics in a neighbourhood of the periodic orbit, Equation (1) is in fact valid over classical scales. The difference between the approach in and the philosophy applied here has an analogy in the classical treatment of one-dimensional WKB solutions near turning points. The simplest way to treat turning points is to approximate the potential using a truncated Taylor expansion (linear for a single turning point and quadratic for two coalescing turning points) and to use the resulting solutions to match standard WKB approximations on either side. The method of comparison equations , on the other hand, seeks a change of variable which (up to higher-order corrections in $`\mathrm{}`$) maps the potential more globally into a linear or quadratic form as required and this has the advantage of giving uniform results which apply over classical length scales. The approach in is analogous to the method of truncating Taylor series whereas in this publication we pursue a transformation into normal form that is similar in spirit to the method of comparison equations.
The difference is that, for multidimensional problems, a deformation of configuration space variables alone as used in the method of comparison equations does not have sufficient range to put the problem in a solvable form and we must use transformations in phase space . In this way a Hamiltonian which is a quantum analogue of the normal forms in can be used and we arrive at a problem which, while not separable, is simple enough that scattering solutions can be written down. More importantly, the information we need to approximate $`\widehat{}(E)`$ can be formulated in such a way that explicit reference to the normal form can be removed and the end result is a formula (of the same form as Equation (1)) which can be understood simply in terms of complex orbits starting and finishing on $`\mathrm{\Sigma }_R^{\mathrm{in}}`$.
We conclude this section with a brief overview of the paper. The essential features of the classical normal form are described in Section 2 and an overview is given of the corresponding quantum Hamiltonian. The normal form Hamiltonian is not separable but does have scattering solutions that are of separable form and these are described in Section 3. Because these scattering solutions do not fully separate in the eigenvalue equation and because the normal form Hamiltonian is not of kinetic-plus-potential type, using them to describe the transmission properties of a general scattering state is not straightforward. Nevertheless a simple solution to this problem is possible, which is outlined in Section 4. The final step in obtaining a usable result is to present the scattering solution in a basis-independent way, which we do in Section 5 in terms of quantised complex Poincarรฉ mappings. Conclusions are presented in Section 6.
## 2 Classical and quantum normal forms
The basis for the calculation in this paper is a quantisation of the classical normal form for the Hamiltonian around an equilibrium. The classical normal form and its connection with classical transition state theory are described in detail in . In this section we will describe the essential results and adapt some of the notation for our own purposes. We will then describe the important properties of a quantisation of this normal form.
### 2.1 The classical normal form
Let canonical coordinates $`(q_0,p_0,q,p)=(q_0,p_0,q_1,\mathrm{},q_d,p_1,\mathrm{},p_d)`$, be chosen so that the quadratic part of the Hamiltonian is
$$H(q_0,p_0,q,p)=\frac{\lambda }{2}(p_0^2q_0^2)+\underset{i=1}{\overset{d}{}}\omega _i(q_i^2+p_i^2)+\mathrm{h}.\mathrm{o}.\mathrm{t}.$$
We denote by $`f=1+d`$ the number of degrees of freedom and we will refer to $`(q_0,p_0)`$ as the reaction coordinates. In the context of collinear molecular collisions it is useful to let $`q_0=0`$ define a dividing surface between reactants and products with $`q_0<0`$ corresponding to reactants and $`q_0>0`$ corresponding to products. The reaction coordinates $`(q_0,p_0)`$ are useful in interpreting the dynamics in this system, in which reaction amounts to crossing a parabolic potential barrier in that degree of freedom. It will also be useful for computational purposes however to allow alternative coordinates $`(Q,P)`$ defined as a rotation of $`(q_0,p_0)`$ so that
$$I=\frac{1}{2}(q_0^2p_0^2)=QP.$$
We will refer to $`I`$ as the reaction action. It is positive for nonreacting trajectories, which are repelled by the parabolic barrier, and it is negative for the reacting trajectories, which cross over (see Figure 2).
At higher order these canonical coordinates are defined so that the Hamiltonian depends on the reaction coordinates through the reaction action $`I`$ only. To establish notation, we will write
$$H(Q,P,q,p)=H_{\mathrm{rc}}(I)+H_{\mathrm{ts}}(q,p,I)$$
(2)
where
$$H_{\mathrm{ts}}(q,p,I)=H_0(q,p)+IH_1(q,p)+I^2H_2(q,p)+\mathrm{}$$
The leading parts of $`H_{\mathrm{rc}}(I)`$ and $`H_0(q,p)`$ contain the quadratic truncation of the Hamiltonian shown above. That is
$$H_{\mathrm{rc}}(I)=\lambda I+\mathrm{h}.\mathrm{o}.\mathrm{t}.$$
and
$$H_0(q,p)=\underset{i=1}{\overset{d}{}}\omega _i(q_i^2+p_i^2)+\mathrm{h}.\mathrm{o}.\mathrm{t}.$$
We refer to $`H_{\mathrm{rc}}(I)`$ and $`H_{\mathrm{ts}}(q,p,I)`$ respectively as the reaction coordinate part and the transition-state part of the Hamiltonian, whence the subscripts.
We remark that for the following calculation to work, it is not necessary to put the transition-state part of the Hamiltonian in normal form. That is, the functions $`H_m(q,p)`$ need not be written as functions of the transverse actions $`J_i=(q_i^2+p_i^2)/2`$. This allows greater latitude in the normal form construction and might in principle alleviate problems with small denominators. We also remark that since we will be looking later at complex solutions to the equations of motion, there is an implicit assumption throughout this paper that the Hamiltonian is analytic in the transition state region.
### 2.2 The quantum normal form
The transformation to the classical normal form described above is achieved by making a canonical change of coordinates so that the Hamiltonian takes the desired form. The corresponding procedure in quantum mechanics is to use a unitary change of basis to the same effect. While a canonical transformation does not uniquely define a unitary operator in the quantum formalism, it is a well-established feature of the quantum-classical correspondence that such a connection can be achieved within semiclassical approximation. The connection was outlined by Miller in using generating functions to write explicit approximations for corresponding unitary operators up to corrections of relative order $`O(\mathrm{})`$. Explicit and concrete rules describing how unitary transformations may be constructed which achieve normal form to higher order in $`\mathrm{}`$ have recently been published by Cargo et al in and used there to derive compact higher-order Bohr-Sommerfeld rules. In the context of critical transmission, which is the application of interest here, transformation to quantum normal forms have been exploited by Colin de Verdiรจre and Parisse in to provide connection formulas describing the behaviour of wavefunctions near hyperbolic fixed points, and these have been used in to calculate multidimensional quantisation rules which are valid near degenerate tori.
An alternative approach is to work directly with the quantum problem in making the transformation to normal form rather than simply quantising the classical normal form as we do here (see , for example). We do not adopt this viewpoint because we will in any case later need to employ semiclassically constructed nonperturbative unitary transformations to connect the normal form basis to the asymptotic basis used for the scattering matrix and there is no overall advantage in avoiding their use at this stage.
In this work we are interested only in constructing the quantum normal form to leading order semiclassically โ that is, neglecting corrections of relative order $`O(\mathrm{})`$ in wavefunctions or terms of order $`O(\mathrm{}^2)`$ in classical symbols. This is achieved using the โpreliminary transformationโ in the language of and the resulting leading-order quantum normal form can be written straightforwardly as a direct copy of the classical normal form. A detailed discussion of this point would unnecessarily elongate the presentation here and will simply assert a direct equivalence (modulo higher order corrections) between quantum and classical Hamiltonians whenever necessary, referring the reader to for a proper explanation.
We therefore start, in analogy with (2), with a Hamiltonian of the form
$$\widehat{H}=H_{\mathrm{rc}}(\widehat{I})+H_{\mathrm{ts}}(\widehat{q},\widehat{p},\widehat{I})$$
(3)
where $`\widehat{I}`$ and $`H_{\mathrm{ts}}(\widehat{q},\widehat{p},\widehat{I})`$ respectively denote quantisations of the classical symbols $`I`$ and $`H_{\mathrm{ts}}(q,p,\widehat{I})`$. There are ordering issues in this correspondence, of course, but for the purposes of making semiclassical approximation to leading order in $`\mathrm{}`$, it suffices to let $`\widehat{I}`$ and $`H_{\mathrm{ts}}(\widehat{q},\widehat{p},\widehat{I})`$ be Weyl quantisations. The key feature here is that $`\widehat{I}`$ and $`(\widehat{q},\widehat{p})`$ act on different degrees of freedom and therefore commute. For concreteness, it may occasionally help to suppose that the transition-state part can be expanded in the form
$$H_{\mathrm{ts}}(\widehat{q},\widehat{p},\widehat{I})=H_0(\widehat{q},\widehat{p})+\widehat{I}H_1(\widehat{q},\widehat{p})+\widehat{I}^2H_2(\widehat{q},\widehat{p})+\mathrm{}$$
(4)
where we may in particular assume that
$$[\widehat{I},H_m(\widehat{q},\widehat{p})]=0.$$
The central result in this paper will be stated in an invariant way that does not refer explicitly to the normal form construction and the details of how this transformation is performed are not needed to use it. In addition, the essential idea of the calculation is understood simply on the basis of the normal form Hamiltonian itself. We will therefore simply quote the result and refer to Refs. for detailed discussions of various approaches to making this transformation in practice.
## 3 Scattering solutions of the quantum normal form
Because higher-order terms in (4) couple the reaction degree of freedom to the transverse degrees of freedom, the normal form Hamiltonian is not separable in the simple-minded sense of the eigenvalue equation separating into a function of the reaction coordinate plus a function of the transition-state coordinates. Since the Hamiltonian depends on the reaction coordinate only through the reaction action $`I`$, however, it turns out that the eigenvalue equation nevertheless admits solutions which have a separable structure. We will use this property to reduce the transmission problem to one that is effectively one-dimensional and therefore solvable by standard techniques. Technical details of this reduction are given in the present section. In the next, it is shown how the results can be formulated in such a way that they no longer rely on an explicit consideration of the normal form.
### 3.1 The reaction coordinate part
The normal form construction provides us with a transformation to coordinates $`(Q,P)`$ such that $`\widehat{I}`$ takes the form
$$\widehat{I}=\frac{1}{2}\left(\widehat{Q}\widehat{P}+\widehat{P}\widehat{Q}\right),$$
(5)
which is the Weyl quantisation of $`QP`$. For interpretation of the results below in terms of conventional calculations, it may help to suppose that the reaction coordinate part $`\widehat{H}_{\mathrm{rc}}`$ of the total Hamiltonian acts on functions of a coordinate $`x`$ (with conjugate momentum $`p_x`$) so that
$$\widehat{I}=I(\widehat{x},\widehat{p}_x)$$
(6)
and that the problem in the $`x`$-representation is close to a standard barrier-penetration problem. We could, for example, let $`(x,p_x)`$ coincide with the coordinates $`(q_0,p_0)`$ defined in Section 2.1 as a rotation of $`(Q,P)`$ in the reaction-coordinate phase plane. In that case the transformation from (6) to (5) is achieved using a metaplectic rotation which rotates the phase plane clockwise through an angle $`3\pi /4`$ (to give Figure 2).
The operator $`\widehat{I}`$ has continuous spectrum and in $`x`$-representation we write the (improper) eigensolutions in the form
$$\widehat{I}\psi _{}(x)=\psi _{}(x).$$
Then
$$E_{\mathrm{rc}}()=H_{\mathrm{rc}}(I=)$$
is the corresponding reaction-coordinate energy. These eigenfunctions are two-fold degenerate since we can send incoming waves from either the reactant or the product side of the barrier (see A). We will restrict our attention here to states which have an incoming component on the reactant side and outgoing components on both the reactant and product sides, but no incoming component on the product side, in which case there is a unique solution for each $``$.
Either as an inverted parabolic barrier or in the representation implied by (5) , the problem of finding eigensolutions of $`\widehat{I}`$ can be solved exactly and solutions of the scattering problem written in closed form. Details are given in A. For present purposes it is sufficient to note that there is a simple relationship describing the relative fluxes in the incoming and outgoing channels. Let the scattering state $`\psi _{}(x)`$ be normalised so that the incoming flux on the reactant side is normalised to unity (by construction, the incoming flux on the product side is zero). Then the outgoing fluxes on the reactant and product sides are, respectively,
$$T()=\frac{1}{1+\mathrm{}^{2\pi /\mathrm{}}}\text{and}R()=\frac{1}{1+\mathrm{}^{2\pi /\mathrm{}}}.$$
(7)
In the barrier-penetration picture, $`T()`$ and $`R()`$ respectively represent probabilities of transmission and reflection. Note that by writing the transmission probability in the alternative form
$$T()=\frac{\mathrm{}^{2\pi /\mathrm{}}}{1+\mathrm{}^{2\pi /\mathrm{}}}$$
the unitarity condition
$$R()+T()=1$$
becomes self-evident.
We make the following observations concerning this result.
* There is a symmetry between $`R()`$ and $`T()`$ on changing the sign of $``$. This is to be expected on the basis of a phase-space portrait in $`(Q,P)`$ coordinates (see Figure 2) in which changing the sign of $``$ simply exchanges reactants for products in the outgoing channels, but is less obvious in a barrier-penetration picture.
* The transmission and reflection coefficients do not change if we replace $`\widehat{I}`$ by a Hamiltonian $`H(\widehat{I})`$ which is an arbitrary function of $`\widehat{I}`$. This will be obvious after generalised fluxes are defined in the next section and a multidimensional version of this observation will be important in getting a simple formulation of the results in this paper.
* The expressions in (7) give semiclassical approximations to transmission and reflection coefficients for a generic potential barrier and can be derived from the standard representation of the Schrรถdinger equation using the method of comparison equations . In the current calculation they are exact for any Hamiltonian which can be written as a function of the reaction action $`\widehat{I}`$ alone. However, there is in general semiclassical error arising from the transformation to normal form in the first place, during which terms of $`O(\mathrm{}^2)`$ arise in the Hamiltonian which are neglected in the current analysis.
### 3.2 The transition state part
In the transverse degrees of freedom corresponding to $`(\widehat{q},\widehat{p})`$, we suppose a discrete spectrum parametrised by $``$ in the following way
$$H_{\mathrm{ts}}(\widehat{q},\widehat{p},)|\phi _k()=E_{\mathrm{ts}}^k()|\phi _k().$$
Here we suppose that a partial symbol $`H_{\mathrm{ts}}(\widehat{q},\widehat{p},)`$ is defined by replacing $`\widehat{I}`$ by its eigenvalue $``$ in $`H_{\mathrm{ts}}(\widehat{q},\widehat{p},\widehat{I})`$. If $`H_{\mathrm{ts}}(\widehat{q},\widehat{p},\widehat{I})`$ is given as a series of the form (4) then we can write, concretely,
$$H_{\mathrm{ts}}(\widehat{q},\widehat{p},)=H_0(\widehat{q},\widehat{p})+H_1(\widehat{q},\widehat{p})+^2H_2(\widehat{q},\widehat{p})+\mathrm{}.$$
In many interesting chemical applications, problems arise which have a Morse or Van der Waals type potential in the transverse degree of freedom for which the spectrum of $`H_{\mathrm{ts}}(\widehat{q},\widehat{p},)`$ is discrete at the bottom but becomes continuous above a threshold. We will confine ourselves, however, to energies at which asymptotically propagating scattering states correspond to reactant molecules in bound states, and for these cases the continuous spectrum (of $`\widehat{H}_{\mathrm{ts}}`$) does not participate. Rather than adjusting notation here to incorporate the continuous part of the spectrum we simply suppress it notationally and consider states labelled by the discrete index $`k`$ only.
Results like those in can be obtained by ignoring the $``$-dependence of the states $`|\phi _k()`$ and approximating them by $`|\phi _k(0)`$ (or equivalently keeping only the leading part $`H_0(\widehat{q},\widehat{p})`$ in the expansion above), but here we want to investigate the effect of coupling between the reaction and transverse degrees of freedom seen in the full Hamiltonian. For a fixed $``$ we can assume that these states form an orthonormal set but note that we should assume in general that
$$\phi _k^{}(^{})|\phi _k()\delta _{kk^{}}$$
(8)
if $`^{}`$. This is the main point complicating the following analysis and means that we should be wary of assuming โobviousโ results when describing issues of normalisation.
We will now describe how the discrete eigensolutions of the transverse problem combine with the scattering solutions found in the reaction-coordinate degree of freedom.
### 3.3 Eigenfunctions of the total Hamiltonian
With the conventions described in Sections 3.1 and 3.2,
$$|\mathrm{\Psi }_{,k}=\psi _{}(x)|\phi _k()$$
are eigenstates of the full Hamiltonian $`\widehat{H}`$ in mixed position-bra-ket notation. These solutions satisfy
$$\widehat{H}|\mathrm{\Psi }_{,k}=E_k()|\mathrm{\Psi }_{,k},$$
where
$$E_k()=E_{\mathrm{rc}}()+E_{\mathrm{ts}}^k().$$
Note that the states $`|\mathrm{\Psi }_{,k}`$ are separable in form even though the Hamiltonian itself is not strictly speaking separable, as discussed at the beginning of this section. This nonseparability manifests itself through the dependence of the transverse eigenstates $`|\phi _k()`$ on $``$ and in particular through the nonorthogonality condition (8).
Our aim is eventually to express results in such a way that explicit reference to the normal form transformation is unnecessary. For this purpose it is preferable to label states with the total energy instead of the reaction action, since the energy is defined independently of the representation used. For each value $`E`$ of the total energy let, $`_k(E)`$ be defined implicitly as a solution of
$$E=E_k().$$
Although we cannot write explicit expressions for $`_k(E)`$, we can suppose that in an energy range around threshold (where $`E_{\mathrm{rc}}()=\lambda +O(^2)`$), these functions are well defined and single-valued for each $`k`$. We then define scattering states labelled by the total energy $`E`$ and the mode number $`k`$ as follows. Let
$$|\mathrm{\Psi }_{E,k}=\psi _{_k(E)}(x)|\phi _k(E),$$
where for short we write
$$|\phi _k(E)=|\phi _k((_k(E)).$$
Note that in view of (8) we have
$$\phi _k^{}(E)|\phi _k(E)\delta _{kk^{}}$$
(9)
since on changing $`k`$ the action eigenvalue $`=_k(E)`$ changes.
It is possible to normalise these states so that
$$\mathrm{\Psi }_{E^{},k^{}}|\mathrm{\Psi }_{E,k}=\delta (EE^{})\delta _{kk^{}}$$
in the usual way, but this convention turns out not to be particularly useful for our purposes and we will not apply it. Instead we will normalise these states so that they have unit flux in the incoming reaction channel. A discussion of normalisation by flux will also be necessary to appreciate how arbitrary scattering states may be constructed from these separated solutions, so we will defer further consideration of such issues until a method of flux calculation has been outlined in the next section.
## 4 Fluxes and sectional inner products
Since the normal-form Hamiltonian is not of kinetic-plus-potential type, we cannot use the usual definition of current
$$๐=\frac{\mathrm{}}{2im}\left(\mathrm{\Psi }^{}\mathrm{\Psi }\mathrm{\Psi }\mathrm{\Psi }^{}\right)$$
to calculate fluxes. There is a simple generalisation, however, which works for arbitrary Hamiltonians.
Let $`\widehat{\mathrm{\Theta }}`$ be a Hermitian operator which projects to one side of a section $`\mathrm{\Sigma }`$ which has codimension one in phase space. For example, if $`\mathrm{\Sigma }`$ is defined by fixing a configuration space coordinate then in position representation $`\widehat{\mathrm{\Theta }}`$ can represent multiplication by the characteristic function of a region which has boundary $`\mathrm{\Sigma }`$. More generally, $`\widehat{\mathrm{\Theta }}`$ can be an operator for which a classical symbol such as the Weyl symbol rises from zero to unity in a classically small strip around $`\mathrm{\Sigma }`$. Then the flux of a state $`|\mathrm{\Psi }`$ across $`\mathrm{\Sigma }`$ is
$$F=\mathrm{\Psi }//_\mathrm{\Sigma }\mathrm{\Psi },$$
(10)
where we define a sectional overlap by
$$\mathrm{\Phi }//_\mathrm{\Sigma }\mathrm{\Psi }=\frac{1}{i\mathrm{}}\mathrm{\Phi }|[\widehat{\mathrm{\Theta }},\widehat{H}]|\mathrm{\Psi }.$$
(11)
The notation here is adapted from although a factor of $`i\mathrm{}`$ has been introduced which will simplify matters later. Flux defined in this way is an integral part of transition-state calculations in the chemical literature (see for example) and similar ideas are used in . It is easily verified that if $`\widehat{\mathrm{\Theta }}`$ represents multiplication by the characteristic function of a region in configuration space then the flux reduces to the standard case of a surface integral of the current $`๐`$ over the boundary. Although a flux calculation requires only the diagonal case in (11), it is useful to allow the nondiagonal case in the definition of sectional overlap. We will find in particular that the space of scattering solutions with a given total energy $`E`$ can be identified with a quantised surface of section in one of the channels and the sectional overlap then provides a natural inner product for the corresponding Hilbert space.
### 4.1 Fluxes for the separated scattering states
We will now apply this generalised construction to compute fluxes in the normal form representation. Let the section $`\mathrm{\Sigma }`$ be chosen so that $`\widehat{\mathrm{\Theta }}`$ can be constructed in terms of the $`(\widehat{Q},\widehat{P})`$ operators alone and commutes with $`\widehat{q}`$ and $`\widehat{p}`$. In phase space this means that $`\widehat{\mathrm{\Theta }}`$ projects onto a region of phase space defined by a subset of the $`QP`$ plane and independent of the $`(q,p)`$ coordinates. Such a choice is natural if we view the complete system as a pinched waveguide (Figure 1) in which a section obtained by fixing reaction coordinates is used to define fluxes in and out of the reactant and product channels.
With the generalised definition of flux in (10) we are not, however, confined to fluxes across surfaces in configuration space and are free to define sections in phase space which distinguish incoming from outgoing flux in the reactant and product channels. For example, referring to Figure 3, a section $`\mathrm{\Sigma }_R^{\mathrm{in}}`$ defined by a vertical line in the right-half of the $`QP`$-plane measures incoming flux in the reactant channel, whereas a horizontal line $`\mathrm{\Sigma }_R^{\mathrm{out}}`$ in the upper half plane measures outgoing flux in the reactant channel and a horizontal line $`\mathrm{\Sigma }_P^{\mathrm{out}}`$ in the lower half plane measures outgoing flux in the product channel. Fluxes across sections such as these corresponding to horizontal and vertical lines in the $`QP`$ plane are especially easily computed in the $`QP`$ representation. Details are given in A for the one-dimensional scattering solutions described in section 3.1 โ here we simply note that relative fluxes are found to be of the form given in (7). Fluxes in the full system can be understood on the basis of these one-dimensional calculations in the following way.
By an abuse of notation, let us denote a one-dimensional flux for one of these reaction-coordinate sections by
$$\psi _{}/_\mathrm{\Sigma }\psi _{}=\frac{1}{i\mathrm{}}\psi _{}|[\widehat{\mathrm{\Theta }},\widehat{I}]|\psi _{}.$$
(12)
Although we use similar notation to (10), it should be emphasised that this one-dimensional flux differs in having $`\widehat{I}`$ and not a Hamiltonian in the commutator and is not therefore a straightforward physical flux. There is also a minus sign, which compensates for the fact that $`I(Q,P)`$ generates a flow that is opposite in direction to the physical flow generated by the Hamiltonian โ corresponding to $`E_k^{}()<0`$ in the discussion below. Although distinct from the physical flux, this quantity is prominent in the physical answer and the notation is useful for that reason.
We will see that in the full space the separated scattering states $`|\mathrm{\Psi }_{,k}`$ have sectional overlaps of the form
$$\mathrm{\Psi }_{,k^{}}//_\mathrm{\Sigma }\mathrm{\Psi }_{,k}=E_k^{}()\psi _{}/_\mathrm{\Sigma }\psi _{}\delta _{kk^{}}.$$
(13)
Alternatively, using states labelled by the total energy and normalised so that
$$|\stackrel{~}{\mathrm{\Psi }}_{E,k}=\sqrt{_k^{}(E)}|\mathrm{\Psi }_{_k(E),k},$$
(14)
we have
$$\stackrel{~}{\mathrm{\Psi }}_{E,k^{}}//_\mathrm{\Sigma }\stackrel{~}{\mathrm{\Psi }}_{E,k}=\psi __k/_\mathrm{\Sigma }\psi __k\delta _{kk^{}}.$$
(15)
The main point here is that the multidimensional scattering states have fluxes, and therefore probabilities of reflection and transmission, which reduce exactly to the one-dimensional case and we can apply (7) to scattering states of an arbitrary Hamiltonian in normal form. The transmission and reflection probabilities of the states $`|\stackrel{~}{\mathrm{\Psi }}_{E,k}`$ can therefore be written
$$T_k(E)=\frac{1}{1+\mathrm{}^{2\pi _k(E)/\mathrm{}}}\text{and}R_k(E)=\frac{1}{1+\mathrm{}^{2\pi _k(E)/\mathrm{}}}$$
(16)
and we can now use this as a basis with which to treat reactivity of an arbitrary stationary state.
We should emphasise that in view of the nonorthogonality of the transverse modes as expressed in (8) and (9), and in contrast to separable problems, these identities extending the one-dimensional results are not at all obvious. To prove them we consider separately the diagonal and nondiagonal cases, and treat them as follows.
Derivation of (13) in the diagonal case
We treat the diagonal case $`k=k^{}`$ first. Assume that $`\widehat{H}_{\mathrm{ts}}`$ can be expanded in the form (4) and use
$$[\widehat{\mathrm{\Theta }},\widehat{I}^m]=[\widehat{\mathrm{\Theta }},\widehat{I}]\widehat{I}^{m1}+\widehat{I}[\widehat{\mathrm{\Theta }},\widehat{I}]\widehat{I}^{m2}+\mathrm{}+\widehat{I}^{m1}[\widehat{\mathrm{\Theta }},\widehat{I}]$$
to deduce that
$$\psi _{}|[\widehat{\mathrm{\Theta }},\widehat{I}^m]|\psi _{}=m^{m1}\psi _{}|[\widehat{\mathrm{\Theta }},\widehat{I}]|\psi _{}$$
and therefore that
$`\mathrm{\Psi }_{,k}|[\widehat{\mathrm{\Theta }},\widehat{H}_{\mathrm{ts}}]|\mathrm{\Psi }_{,k}`$ $`=`$ $`{\displaystyle \underset{m}{}}m^{m1}\psi _{}|[\widehat{\mathrm{\Theta }},\widehat{I}]|\psi _{}\phi _k()|H_m(\widehat{q},\widehat{p})|\phi _k()`$
$`=`$ $`\psi _{}|[\widehat{\mathrm{\Theta }},\widehat{I}]|\psi _{}\phi _k()|{\displaystyle \frac{\widehat{H}_{\mathrm{ts}}(\widehat{q},\widehat{p},)}{}}|\phi _k()`$
$`=`$ $`\psi _{}|[\widehat{\mathrm{\Theta }},\widehat{I}]|\psi _{}{\displaystyle \frac{E_{\mathrm{ts}}^k()}{}},`$
where in the last line we have invoked the Feynman-Hellman theorem. On adding a similar calculation for $`\widehat{H}_{\mathrm{rc}}`$ we get the claimed result. The important feature is that the factor $`E_k^{}()`$ does not depend the choice of section across which to measure flux and relative fluxes reduce to the one-dimensional case.
Derivation of (13) in the nondiagonal case
To treat the nondiagonal case we use the identity
$`\psi _{^{}}|[\widehat{\mathrm{\Theta }},\widehat{I}^m]|\psi _{}`$ $`=`$ $`\psi _{^{}}|\left([\widehat{\mathrm{\Theta }},\widehat{I}]\widehat{I}^{m1}+\widehat{I}[\widehat{\mathrm{\Theta }},\widehat{I}]\widehat{I}^{m2}+\mathrm{}\widehat{I}^{m1}[\widehat{\mathrm{\Theta }},\widehat{I}]\right)|\psi _{}`$
$`=`$ $`[^{m1}+^{}^{m2}+\mathrm{}^{m1}]\psi _{}^{}|[\widehat{\mathrm{\Theta }},\widehat{I}]|\psi _{}`$
$`=`$ $`{\displaystyle \frac{^m^m}{^{}}}\psi _{}^{}|[\widehat{\mathrm{\Theta }},\widehat{I}]|\psi _{}`$
to deduce that
$$\mathrm{\Psi }_{^{},k^{}}|[\widehat{\mathrm{\Theta }},\widehat{I}^m\widehat{H}_m]|\mathrm{\Psi }_{,k}=\psi _{^{}}|[\widehat{\mathrm{\Theta }},\widehat{I}]|\psi _{}\phi _k^{}(^{})|\left[\frac{^m\widehat{H}_m^m\widehat{H}_m}{^{}}\right]|\phi _k().$$
On summing over $`m`$ and doing a similar calculation for the reaction-coordinate part of the Hamiltonian we find that
$$\mathrm{\Psi }_{^{},k^{}}//_\mathrm{\Sigma }\mathrm{\Psi }_{,k}=\frac{E_k()E_k^{}(^{})}{^{}}\psi _{^{}}/_\mathrm{\Sigma }\psi _{}\phi _k^{}(^{})|\phi _k().$$
In terms of the energy-labelled states $`|\mathrm{\Psi }_{E,k}=|\mathrm{\Psi }_{_k(E),k}`$, we have
$$\mathrm{\Psi }_{E^{},k^{}}//_\mathrm{\Sigma }\mathrm{\Psi }_{E,k}=\frac{EE^{}}{_k(E)_k^{}(E^{})}\psi _{_k^{}(E^{})}/_\mathrm{\Sigma }\psi _{_k(E)}\phi _k^{}(E^{})|\phi _k(E).$$
These sectional overlaps vanish as $`E^{}`$ approaches $`E`$ unless $`k=k^{}`$, in which case we recover the diagonal result. This completes the derivation of (13).
### 4.2 Fluxes for general scattering states and reduced Hilbert space
Let $`_R^{\mathrm{in}}`$ denote the subspace of scattering states of a fixed total energy $`E`$ which are incoming on the reactant side and have no incoming flux on the product side. In the applications we consider the scattering problem will asymptote to a waveguide-type problem in which there are a finite number $`M`$ of propagating channels (corresponding to the energetically accessible bound states of the reacting molecules) and the space $`_R^{\mathrm{in}}`$ will therefore be finite-dimensional. We assert that the space $`_R^{\mathrm{in}}`$ is in fact essentially a quantisation of a classical surface of section (in the incoming reactant channel) in the sense of Bogomolny , with the inner product being given by sectional overlaps of the form in (11).
Let us write a general asymptotically-propagating scattering state in the form
$$|\mathrm{\Phi }=\underset{k=1}{\overset{M}{}}a_k\frac{|\stackrel{~}{\mathrm{\Psi }}_{E,k}}{\sqrt{\psi __k/_{\mathrm{\Sigma }_R^{\mathrm{in}}}\psi __k}}.$$
(17)
Then the total incoming flux for this state is
$$F_R^{\mathrm{in}}=\mathrm{\Phi }//_{\mathrm{\Sigma }_R^{\mathrm{in}}}\mathrm{\Phi }=\underset{k=1}{\overset{M}{}}|a_k|^2$$
while in view of (15) and (16) the outgoing flux in the reactant channel is
$$F_R^{\mathrm{out}}=\mathrm{\Phi }//_{\mathrm{\Sigma }_R^{\mathrm{out}}}\mathrm{\Phi }=\underset{k=1}{\overset{M}{}}\frac{1}{1+\mathrm{}^{2\pi _k(E)/\mathrm{}}}|a_k|^2$$
and the outgoing flux in the product channel is
$$F_P^{\mathrm{out}}=\mathrm{\Phi }//_{\mathrm{\Sigma }_P^{\mathrm{out}}}\mathrm{\Phi }=\underset{k=1}{\overset{M}{}}\frac{1}{1+\mathrm{}^{2\pi _k(E)/\mathrm{}}}|a_k|^2.$$
Furthermore, between any two such scattering states the sectional overlap
$$\mathrm{\Phi }^{}//_{\mathrm{\Sigma }_R^{\mathrm{in}}}\mathrm{\Phi }=\underset{k=1}{\overset{M}{}}a_k^{}a_k$$
provides a natural inner product for $`_R^{\mathrm{in}}`$.
We can regard $`_R^{\mathrm{in}}`$ abstractly as a space spanned by an orthonormal basis $`\{|k\}_{k=1}^M`$ whose elements
$$|k\frac{|\stackrel{~}{\mathrm{\Psi }}_{E,k}}{\sqrt{\psi __k/_{\mathrm{\Sigma }_R}\psi __k}}$$
are in one-to-one correspondence with the scattering states $`|\stackrel{~}{\mathrm{\Psi }}_{E,k}`$, normalised to have unit incoming flux. A general element
$$|\phi =\underset{k=1}{\overset{M}{}}a_k|k$$
of this reduced space can be extended to a scattering state of the form given in (17) for which the total energy is fixed but which is not necessarily an eigenstate of the transverse Hamiltonian $`\widehat{H}_{\mathrm{ts}}`$. The inner product between two reduced states can be defined through sectional overlaps
$$\phi ^{}|\phi =\mathrm{\Phi }^{}//_{\mathrm{\Sigma }_R}\mathrm{\Phi }$$
of the corresponding extended states.
The benefit of this abstraction is that we can compute outgoing fluxes in the product channel using matrix elements
$$F_P^{\mathrm{out}}=\phi |\widehat{}(E)|\phi $$
(18)
of a reaction operator
$$\widehat{}(E)=\underset{k=1}{\overset{M}{}}\frac{|kk|}{1+\mathrm{}^{2\pi _k(E)/\mathrm{}}}=\underset{k=1}{\overset{M}{}}\frac{\mathrm{}^{2\pi _k(E)/\mathrm{}}}{1+\mathrm{}^{2\pi _k(E)/\mathrm{}}}|kk|$$
(19)
which is diagonal in this basis. By writing the outgoing flux as a matrix element of an operator defined on $`_R^{\mathrm{in}}`$ we have in large part achieved the goal of this paper, which is to generalise the construction in so that there is no longer a restriction to states supported in a classically small region of phase space. In fact, the only restriction on incoming states here is that they should be supported in the region of phase space where the normal form in (2) provides an accurate description of dynamics. Although undoubtedly a restricted subset of phase space, this domain has classical dimensions.
The current version is tied to the normal form, however and in order for this operator to be of any practical use, we really need a way of constructing it which does not call on us explicitly to construct the normal form or the basis vectors $`|k`$. As a start in this direction, note that if we denote
$$\widehat{๐ฏ}(E)=\underset{k}{}\mathrm{}^{2\pi _k(E)/\mathrm{}}|kk|,$$
(20)
then $`\widehat{}(E)`$ can be written in the form (1) promised in the introduction. This is precisely the form given in , where $`\widehat{๐ฏ}(E)`$ was a tunnelling operator defined as a quantised surface of section map in the neighbourhood of a complex periodic orbit. We will show in the Section 5 that the same interpretation can be imposed on $`\widehat{๐ฏ}(E)`$ in the present case with the difference that, unlike in , restriction to a neighbourhood of the periodic orbit small enough for linearised dynamics to be used is no longer necessary.
### 4.3 The microcanonical cumulative reaction probability
The trace
$$N(E)=\mathrm{Tr}\widehat{}(E)=\underset{k=1}{\overset{M}{}}\frac{1}{1+\mathrm{}^{2\pi _k(E)/\mathrm{}}}$$
(21)
of the reaction operator is the so-called microcanonical cumulative reaction probability. Results for $`N(E)`$ equivalent to those that would be obtained by using linearised dynamics in $`\widehat{}(E)`$ in the manner of were obtained by Miller in . Semiclassical approximations for $`N(E)`$ have also been given in which include nonlinear effects in the transition state degrees of freedom by using expansions in the transition-state quantum numbers (a related treatment of tunnelling using normal-form coordinates has been given in ). A thermalised version has been given in and see for a semiquantum calculation. A discussion emphasising the fluctuations that occur in $`N(E)`$ as the summands in (21) switch on with increasing $`E`$ can be found in . We also remark that a discussion of $`N(E)`$ has recently been given in which uses the same language of normal forms that we use here.
The major benefit of the current work is that, once a prescription has been given in the next section for computing $`\widehat{}(E)`$ without recourse to the normal form, we will have an explicit prescription for distributing the total reaction probability in $`N(E)`$ among incoming states or, equivalently, using the Wigner-Weyl formalism, for understanding how the reaction probability is distributed in phase space. In principle, this calculation can be implemented simply by computing complex trajectories near the complex periodic orbit used in and does not require a particular deconstruction of the Hamiltonian such as provided by a normal form (although normal forms do help in inverting the operator $`1+\widehat{๐ฏ}(E)`$ in (1) as we will discuss).
## 5 Getting away from the normal form
Equation (1) promises the ability to treat reaction problems without having to deal explicitly with the normal form and separated stationary states. In this section we show how this can be done by interpreting $`\widehat{๐ฏ}(E)`$ as a tunnelling operator in the sense of which can be understood independently of the normal form construction.
### 5.1 The complex return map in normal form coordinates
We first describe how a complex return map is expressed in terms of normal form coordinates. Let $`(I,\theta )`$ be action angle variables in the $`QP`$ plane such that in the positive quadrant
$`Q`$ $`=`$ $`\sqrt{I}\mathrm{}^\theta `$
$`P`$ $`=`$ $`\sqrt{I}\mathrm{}^\theta .`$
We have $`I=QP`$, consistent with previous notation. Although the system is hyperbolic and not normally associated with periodic motion, there is in fact an imaginary period, expressed by the identities
$`Q(\theta +2\pi i)`$ $`=`$ $`Q(\theta )`$
$`P(\theta +2\pi i)`$ $`=`$ $`P(\theta ).`$
In other words, the Hamiltonian flow generated by $`I`$ in the complexified $`QP`$ plane has period $`2\pi i`$. We will now show that an extension of this periodic flow to the full system can be used to define a complex Poincarรฉ return map and that this map has a fixed point corresponding to a complex periodic orbit.
Let us restrict initial conditions to a surface of section $`\mathrm{\Sigma }_R^{\mathrm{in}}`$ defined by fixing $`Q`$ as in Figure 3, along with the total energy $`E`$, and let the condition
$$E=H(q,p,I)$$
(22)
implicitly define the function $`h(q,p,E)`$ on $`\mathrm{\Sigma }_R^{\mathrm{in}}`$ by
$$I=h(q,p,E)e(E).$$
Here we regard $`(q,p)`$ as canonical coordinates for $`\mathrm{\Sigma }_R^{\mathrm{in}}`$ and $`e(E)`$ is defined so that the minimum of $`h(q,p,E)`$ on $`\mathrm{\Sigma }_R^{\mathrm{in}}`$ is zero, as described more explicitly below. Differentiating (22) with respect to the transverse coordinates $`q`$ and $`p`$ while keeping $`E`$ fixed gives
$$0=\frac{H}{I}h+H,$$
which in turn gives
$$X_h=\frac{1}{\dot{\theta }}\stackrel{~}{X}_H=\left(\begin{array}{c}\frac{\mathrm{d}q}{\mathrm{d}\theta }\\ \frac{\mathrm{d}p}{\mathrm{d}\theta }\end{array}\right),$$
where $`\stackrel{~}{X}_H`$ denotes the projection of the full flow vector $`X_H`$ defined by $`H`$ onto the transverse degrees of freedom $`(q,p)`$ and we have used $`\dot{\theta }=H/I`$. The flow defined by $`X_h`$ can therefore be regarded as a restriction to the $`(q,p)`$ degrees of freedom of the full flow, reparametrised so that time $`t`$ is replaced with the angle variable $`\theta `$.
Letting $`\theta `$ evolve from $`0`$ to the final value $`2\pi i`$, trajectories are described in full phase space which start on $`\mathrm{\Sigma }_R^{\mathrm{in}}`$ and return to it โ recall that the $`(Q,P)`$ coordinates which are used to define $`\mathrm{\Sigma }_R^{\mathrm{in}}`$ are periodic under this evolution. Integrating the flow vector $`X_h`$ for a time $`2\pi i`$ then generates a complex symplectic map
$$:\mathrm{\Sigma }_R^{\mathrm{in}}\mathrm{\Sigma }_R^{\mathrm{in}}$$
which we can denote by
$$=\mathrm{exp}[2\pi iX_h]$$
(23)
in Lie-algebraic notation. This is precisely the classical map used to construct the tunnelling operator in .
Before describing explicitly how the quantisation works, it is helpful to see how the complex periodic orbit which provides a fixed point of $``$ arises in normal form coordinates. By construction, the quadratic part of the Hamiltonian $`H(q,p,I)`$ is elliptic in the transverse degrees of freedom. As a result, the sectional Hamiltonian $`h(q,p,E)`$ has a minimum $`(q_e(I),p_e(I))`$, for sufficiently small $`I`$ at least, for which
$$h(q_e(I),p_e(I),E)=0$$
and which coincides with the origin of $`\mathrm{\Sigma }_R^{\mathrm{in}}`$ in the threshold case $`I=0`$. We define $`e(E)`$ above so that $`h(q_e,p_e,E)=0`$ and a Taylor expansion of $`h(q,p,E)`$ about $`(q_e(I),p_e(I))`$ begins with quadratic terms. This minimum is an equilibrium of the flow defined by $`X_h`$ on $`\mathrm{\Sigma }_R^{\mathrm{in}}`$ and is therefore a fixed point of $``$. In full phase space, the trajectory starting with coordinates $`(q_e(I),p_e(I))`$ on $`\mathrm{\Sigma }_R^{\mathrm{in}}`$ evolves so that the coordinates $`I`$ and $`(q,p)`$ are fixed and defines a periodic orbit as $`\theta `$ evolves from $`0`$ to $`2\pi i`$. The time period corresponding to this evolution is $`i\tau (E)`$ where
$$\tau =\frac{2\pi }{\dot{\theta }}=\frac{2\pi }{H(q_e,p_e,I)/I}=2\pi e^{}(E).$$
and its action is an imaginary number $`S(E)=iK_0(E)`$ where
$$K_0=\frac{1}{i}Id\theta =2\pi I=2\pi e(E).$$
In a linearisation of dynamics about this complex periodic orbit was used to approximate $`\widehat{}(E)`$ which here corresponds to truncating a Taylor series of $`h(q,p,E)`$ about $`(q_e(I),p_e(I))`$ at quadratic order. The essential conclusion of this paper is that a complete description of the reaction operator can be achieved simply by replacing this truncation with the full sectional Hamiltonian $`h(q,p,E)`$.
### 5.2 The tunnelling operator
The tunnelling operator is defined to be a quantisation of the classical map $``$ . Using (23) we can write concretely,
$$\widehat{๐ฏ}=\mathrm{}^{2\pi (e\widehat{h})/\mathrm{}},$$
where $`\widehat{h}e`$ is the restriction of the operator $`\widehat{I}`$ to the quantum analogue $`_R^{\mathrm{in}}`$ of $`\mathrm{\Sigma }_R^{\mathrm{in}}`$ defined in the Section 4. In that case
$$(\widehat{h}e)|k=_k(E)|k$$
and
$$\mathrm{}^{2\pi (e\widehat{h})/\mathrm{}}|k=\mathrm{}^{2\pi _k(E)/\mathrm{}}|k$$
and the identification in the previous section of $`\widehat{๐ฏ}(E)`$ in (20) as a tunnelling operator is confirmed.
It should be emphasised that while the sectional Hamiltonian $`h`$ was used in making this identification, it is not necessary to construct it explicitly, or even to refer to it, in order to construct the map $``$ and to approximate its quantum analogue $`\widehat{๐ฏ}(E)`$ semiclassically. The map $``$ can be constructed simply by integrating orbits as described in references . From the dynamical characteristics of these orbits, Van Vleck-type approximations for $`\widehat{๐ฏ}(E)`$ can be written as described in , for example. Alternatively, the Weyl symbol of $`\widehat{๐ฏ}(E)`$ can be obtained as described in . Coherent state representations are also possible ( and references therein). In all of these approaches the dynamical information needed is naturally provided as a result of the orbit computation and $`h(q,p)`$ is not needed explicitly.
We note finally that while it would be quite easy to write a version of (19) that describes the full scattering matrix in the normal form representation, so that we could determine the distribution of product states for each reactant state, it is less obvious how the normal form result could be interpreted in a basis-independent way in that case. Any such reformulation would have to take into account the fact that orbits contributing to the scattering matrix itself depend singularly on initial conditions near the reacting boundary. The simplicity of the normal form representation suggests, however that such a uniformisation might be feasible, although we do not pursue it here.
### 5.3 Asymptotic basis for the reaction operator
Although the reaction operator is diagonal with respect to the stationary states $`|\stackrel{~}{\mathrm{\Psi }}_{E,k}`$ computed in terms of the normal form, it will not in general be diagonal in the basis of asymptotically decoupled stationary states that is used to write the scattering matrix. In fact, since the normal form will in general only provide an accurate description of dynamics in a neighbourhood of the transition state, we need an independent method to describe how these states can be extended to the asymptotic regions of the reactant channel and the operator $`\widehat{}(E)`$ written in the standard asymptotic basis.
To achieve this we note simply that once the transmission problem has been solved locally in the transition state region, the solution can extended to the asymptotic region by applying quantised surface-of-section maps, such as described in , for example. Once outside the transition state region, these mappings can be constructed on the basis of primitive semiclassical approximations and amount in the classical picture simply to conjugating the complex Poincarรฉ map $``$ with standard real ones. As when transforming to the quantum normal form in the first place, the end result of this process is easily stated and the details omitted in the interests of brevity since they follow discussions elsewhere .
To be more specific about this conjugation, let the coordinates $`(x,y,p_x,p_y)`$ be as described in the introduction and let the Hamiltonian decouple asymptotically in the reactant channel according to
$$H(x,y,p_x,p_y)H_{\mathrm{asymp}}(y,p_y)+H_{\mathrm{tr}}(x,p_x).$$
(24)
Here $`H_{\mathrm{asymp}}(y,p_y)`$ describes the internal vibrational motion of the reacting molecules and $`H_{\mathrm{tr}}(x,p_x)`$ is a kinetic energy term for the relative motion of centres of mass. In the simplest atom-diatom collinear case we might have
$$H_{\mathrm{tr}}(x,p_x)=\frac{p_x^2}{2M}$$
where $`M`$ is the atom-diatom relative mass and
$$H_{\mathrm{asymp}}(y,p_y)=\frac{p_y^2}{2m}+V_{\mathrm{AB}}(y)$$
where $`m`$ is the diatomic relative mass and $`V_{\mathrm{AB}}(q)`$ the diatomic interaction potential.
The separated asymptotic solutions are denoted
$$|\mathrm{\Phi }_{n,E}=\chi _{E,n}(x)|\psi _n$$
where $`\chi _{E,n}(x)`$ is a plane-wave eigenfunction of $`\widehat{H}_{\mathrm{tr}}`$ with energy $`EE_n`$ and normalised to have unit incoming flux while $`|\psi _n`$ is an eigenstate of internal dynamics with energy $`E_n`$. There is a phase convention implicit in writing the scattering matrix, which in the present case amounts to specifying the phase of $`\chi _{E,n}(x)`$. In the standard plane-wave case
$$\chi _{E,n}(x)=\frac{\mathrm{}^{ik_nx}}{\sqrt{v_n}},$$
where $`v_n=\mathrm{}k_n`$ ensures unit incoming flux, we can understand the phase of $`\chi _{E,n}(x)`$ as being fixed at $`x=0`$.
The phase convention in the asymptotic regime is then obtained by extending the complete scattering state $`|\mathrm{\Phi }_{n,E}`$ from a section $`\mathrm{\Sigma }_R^0`$ defined by $`x=0`$ to a section $`\mathrm{\Sigma }_R^{\mathrm{asymp}}`$ defined by a large and negative value of $`x`$. This extension is affected in the semiclassical scheme by quantising a surface of section mapping
$$_0:\mathrm{\Sigma }_R^0\mathrm{\Sigma }_R^{\mathrm{asymp}}$$
constructed using the uncoupled dynamics of the Hamiltonian in (24) โ since the plane wave is travelling to the right, this will be achieved in negative time. Scattering of the resulting incoming state then proceeds by first mapping from the section $`\mathrm{\Sigma }_R^{\mathrm{asymp}}`$ to a section $`\mathrm{\Sigma }_R^{\mathrm{ts}}`$ near the transition state using the fully coupled dynamics,
$$_{\mathrm{asymp}}:\mathrm{\Sigma }_R^{\mathrm{asymp}}\mathrm{\Sigma }_R^{\mathrm{ts}}.$$
If $`\mathrm{\Sigma }_R^{\mathrm{ts}}`$ is in the domain of the normal form then its transmission probability is understood using a complex return map
$$_{\mathrm{ts}}:\mathrm{\Sigma }_R^{\mathrm{ts}}\mathrm{\Sigma }_R^{\mathrm{ts}}$$
of the kind described in Section 5.1. Finally, we complete the transformation to the asymptotic basis by mapping back to the asymptotic section $`\mathrm{\Sigma }_R^{\mathrm{asymp}}`$ using the inverse $`_{\mathrm{asymp}}^1`$ of $`_{\mathrm{asymp}}`$ and then mapping to $`\mathrm{\Sigma }_R^0`$ using the inverse $`_0^1`$. The end result is a map
$$=_0^1_{\mathrm{asymp}}^1_{\mathrm{ts}}_{\mathrm{asymp}}_0$$
(25)
which is conjugate to $`_{\mathrm{ts}}`$ but which is adapted to the phase convention used for the scattering matrix. A quantisation of this map gives a tunnelling operator $`\widehat{๐ฏ}(E)`$ appropriate to the asymptotic basis defined by the states $`|\mathrm{\Phi }_{n,E}`$.
Note that by conjugating $`_{\mathrm{ts}}`$ by $`_{\mathrm{asymp}}`$, we can extend the return map far beyond the domain where the normal form applies. The outer conjugation by $`_0`$ is not strictly necessary to calculate reaction probabilities of incoming states $`|\mathrm{\Phi }_{n,E}`$ since the transverse parts $`|\psi _n`$ are eigenfunctions of the quantisation of $`_0`$ and the resulting eigenphases cancel in the transformed version
$$F_P^{\mathrm{out}}=\psi _n|\widehat{}(E)|\psi _n$$
of (18). However phases are important for cross terms if we want to treat general scattering states
$$|\mathrm{\Phi }=\underset{n=1}{\overset{M}{}}c_n|\mathrm{\Phi }_{n,E}$$
and they are also important for representations of $`\widehat{}(E)`$ in phase space. We also remark that $`_{\mathrm{asymp}}`$ will not in general commute with $`_{\mathrm{ts}}`$ and the reaction operator will therefore not be diagonal in the asymptotic basis defined by the states $`|\mathrm{\Phi }_{n,E}`$. On noting that the conjugation in (25) amounts simply to a change of representation so that $`\mathrm{\Sigma }_R^0`$ can be identified with the section $`\mathrm{\Sigma }_R^{\mathrm{in}}`$, however, we see that the tunnelling operator written in this basis is not fundamentally different from the one described in Section 5.2.
Finally, we note that the conjugation in (25) is particular to a waveguide problem in which the plane wave part is written in terms of a coordinate $`x`$ with phases fixed at $`x=0`$. If different conventions are used for the asymptotic coordinates or for the part $`\chi _{E,n}(x)`$ of the scattering state, then the part $`_0`$ of the conjugation must be redefined accordingly.
### 5.4 Limitations of the derivation
No approximations are made in getting to the reaction-operator form of the outgoing flux in (18) and (19) once the quantum normal form of the Hamiltonian in (3) is written down. The sources of error are in transforming to the quantum normal form in the first place and in transforming to the asymptotic scattering basis after the reaction operator has been found in the normal form representation (18) and (19). We now comment on some of the issues affecting this approach.
The classical normal form is a formal series which describes the dynamics locally in a neighbourhood of the transition state region. It is clear that writing a quantum version of it as we do here is a formal step which will ultimately require a more careful justification. Issues of convergence become especially important if we pursue the limit $`\mathrm{}0`$ in its literal sense and we have not addressed such questions here. From a purely practical point of view, however, if oneโs aim is to achieve an approximation that works well for a small but fixed value of $`\mathrm{}`$, then it suffices to describe the dynamics to a corresponding level of accuracy in phase space and the normal form is certainly capable of that in the sorts of parameter regimes that arise in chemical applications .
Results of a numerical investigation are outlined in which indicate that the expression in (1) works very well when $`\widehat{๐ฏ}(E)`$ is computed directly from the complex orbits of the Poincarรฉ return map $``$. We also note that since (1) can be interpreted theoretically without reference to the normal form, it seems natural to expect that it applies independently of the normal form itself. We therefore conjecture that, despite the limitations of the derivation presented here, Equation (1) is in fact โclassically exactโ in the sense that no errors arise from classical dynamics side of the calculation once $`\widehat{๐ฏ}(E)`$ is interpreted as the quantisation of $``$ and the only approximation is the usual semiclassical one which vanishes as $`\mathrm{}0`$.
It is important to add the qualification, however, that even if the conjecture is correct, there are good reasons to expect it to apply only locally, at least in the simple form described in sections 5.1 and 5.2. The return map $``$ describes a unique image for initial conditions in a neighbourhood of the complex periodic orbit and for energies sufficiently close to threshold. In practical terms, this means there is a unique complex solution satisfying the boundary conditions required of orbits by semiclassical approximation of $`\widehat{๐ฏ}(E)`$ in the usual representations . Sufficiently far away, however, bifurcations are likely to occur where this structure breaks down and these are not described by the current formulation. In fact, it has been shown in (and see for related work) that complex orbits contributing to the scattering matrix can be chaotic and are subject to intricate pruning by the Stokesโ phenomenon, while numerical evidence suggests that the same is true of the orbits contributing to $`\widehat{๐ฏ}(E)`$ sufficiently far from the centre of the reacting region. At an even more basic level, the NHIM itself may undergo bifurcation once the energy rises far enough above threshold and in this case the whole bottleneck picture at the basis of our calculation is no longer correct. Global recrossing may occur and resonances arise in the quantum mechanics which are not described by the simple picture of transmission probability we have here. It should be stressed, however, that whatever the limits are on the domain where contributing dynamics are simple, they are independent of $`\mathrm{}`$ and therefore have classical scales.
Finally, we remark that even though the tunnelling operator $`\widehat{๐ฏ}(E)`$ can be routinely approximated using semiclassical approximations, the inverse of $`1+\widehat{๐ฏ}(E)`$ that occurs in (1) is more problematic. Closed form analytic approximations are possible if we know the sectional Hamiltonian $`h(q,p)`$ (see for the harmonic case) but more work is needed to provide an approximation that works directly in terms of the map $``$.
## 6 Conclusion
We have characterised the semiclassical transmission of waves across a phase-space bottleneck using a reaction operator constructed from a complex Poincarรฉ mapping. In contrast to previous work , this construction is not restricted to energies and parts of phase space in a classically small neighbourhood of the transition state at threshold.
A phase-space representation of the reaction operator will be largely supported in the classically reacting subset of phase space, but will also incorporate tunnelling and other quantum effects at the boundary of this region, where trajectories approach the NHIM along its stable manifold. The only dynamical information needed to apply the approximation described in this paper is contained in the complex Poincarรฉ mapping and explicit consideration of normal forms, or other special assumptions regarding the dynamics such separability or adiabatic approximation, are unnecessary. We also note that the particular trajectories used to define this map are well behaved at the reacting boundary. Therefore, despite the fact that the fate of trajectories changes discontinuously across the reacting boundary, the result here uniformly describes the transition from classically allowed transmission inside the reacting region to reaction entirely by tunnelling outside it.
We conclude by noting that in its current form the result here is restricted to collinear problems. In order for the approach to be used in completely realistic models, the marginally stable degrees of freedom associated with rotational symmetry will need to be incorporated. This aspect needs further investigation.
Acknowledgements
This work was supported the European Network MASIE.
## Appendix A Scattering for the one-dimensional normal form
The one-dimensional problem is especially easily solved in a representation in which the reaction action operator takes the form (5) and $`\widehat{P}`$ acts on functions of $`Q`$ according to
$$\widehat{P}\psi (Q)=\frac{\mathrm{}}{i}\psi ^{}(Q).$$
In this case the eigenvalue equation $`\widehat{I}\psi =\psi `$ is a first order differential equation
$$\frac{\mathrm{}}{i}\left(Q\frac{\mathrm{d}}{\mathrm{d}Q}+\frac{1}{2}\right)\psi (Q)=\psi (Q)$$
(26)
and this can be solved in elementary terms without the complication of parabolic cylinder functions that arise in the conventional representation of an inverted oscillator . This approach has in particular been exploited in to treat scattering in one-dimensional networks of tori and we refer to that publication for more detail of the following calculation. It is useful, however, to reiterate some of the main points here and to emphasise flux calculations, which form a basis for the discussion in the main text.
The eigenvalue equation (26) leads to doubly degenerate eigenfunctions $`\psi _{}^\pm (Q)`$ of the forms
$$\psi _{}^+(Q)=\mathrm{\Theta }(Q)Q^{1/2+i/\mathrm{}}$$
and
$$\psi _{}^{}(Q)=\psi _{}^+(Q)$$
respectively. We will concentrate on the solution $`\psi _{}^+(Q)`$, which as we will now show represents bombardment of the equilibrium from the reactant side and is therefore the solution singled out in section 3.1.
The incoming flux is normalised as follows. Let the projection operator $`\widehat{\mathrm{\Theta }}`$ have the $`Q`$-representation
$$\widehat{\mathrm{\Theta }}\psi (Q)=\mathrm{\Theta }(Q_0Q)\psi (Q)$$
so that it measures flux from right to left in the $`QP`$-plane across a section $`\mathrm{\Sigma }_R`$ defined by $`Q=Q_0`$. Then
$$\frac{1}{i\mathrm{}}[\widehat{\mathrm{\Theta }},\widehat{I}]=\frac{1}{2i\mathrm{}}\left(\widehat{Q}[\widehat{\mathrm{\Theta }},\widehat{P}]+[\widehat{\mathrm{\Theta }},\widehat{P}]\widehat{Q}\right)=Q\delta (QQ_0)$$
and sectional overlaps are of the form
$$\psi /_{\mathrm{\Sigma }_R}\psi =\frac{1}{i\mathrm{}}\psi |[\widehat{\mathrm{\Theta }},\widehat{I}]|\psi =Q_0|\psi (Q_0)|^2.$$
The case $`Q_0>0`$ corresponds to incoming flux on the reactant side. We then denote $`\mathrm{\Sigma }_R=\mathrm{\Sigma }_R^{\mathrm{in}}`$ and get
$$\psi _{}^+/_{\mathrm{\Sigma }_R^{\mathrm{in}}}\psi _{}^+=Q_0|Q_0^{1/2+i/\mathrm{}}|^2=1.$$
A section with $`Q_0<0`$ gives a flux in the incoming product channel and this vanishes for the state $`\psi _{}^+(Q)`$, consistent with our interpretation of it as an incoming state in the reactant channel.
Outgoing fluxes are naturally measured in momentum representation
$$\phi (P)=P|\psi $$
using projections of the form
$$\widehat{\mathrm{\Theta }}\phi (P)=\mathrm{\Theta }(PP_0)\phi (P).$$
In this representation we have
$$\frac{1}{i\mathrm{}}[\widehat{\mathrm{\Theta }},\widehat{I}]=\frac{1}{2i\mathrm{}}\left([\widehat{\mathrm{\Theta }},\widehat{Q}]\widehat{P}+\widehat{P}[\widehat{\mathrm{\Theta }},\widehat{Q}]\right)=P\delta (PP_0)$$
and therefore
$$\psi /_{\mathrm{\Sigma }_P}\psi =P_0|\phi (P_0)|^2.$$
To complete the calculation we therefore need to evaluate
$`\phi _{}^+(P)`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2\pi \mathrm{}}}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}\mathrm{}^{iQP/\mathrm{}}\psi _{}^+(Q)dQ`$
$`=`$ $`{\displaystyle \frac{1}{\sqrt{2\pi \mathrm{}}}}\left({\displaystyle \frac{|P|}{\mathrm{}}}\right)^{1/2i/\mathrm{}}{\displaystyle _0^{\mathrm{}}}\mathrm{}^{i\sigma q}q^{1/2+i/\mathrm{}}dq`$
$`=`$ $`{\displaystyle \frac{\mathrm{}^{i/\mathrm{}}}{\sqrt{2\pi }}}\mathrm{}^{i\sigma \pi /4+\sigma \pi /(2\mathrm{})}\mathrm{\Gamma }\left({\displaystyle \frac{1}{2}}+{\displaystyle \frac{i}{\mathrm{}}}\right)|P|^{1/2i/\mathrm{}},`$
where $`\sigma `$ is the sign of $`P`$. Note that $`\phi _{}^+(P)\mathrm{const}\times |P|^{1/2i/\mathrm{}}`$ has the same dependence on its argument as found in the $`Q`$-representation, which is to be expected since there is a symmetry between $`\widehat{Q}`$ and $`\widehat{P}`$ in $`\widehat{I}`$.
In either case, for flux calculations it suffices to know that
$`|\phi _{}^+(P)|^2`$ $`=`$ $`{\displaystyle \frac{\mathrm{}^{\sigma \pi /\mathrm{}}}{\mathrm{}^{\pi /\mathrm{}}+\mathrm{}^{\pi /\mathrm{}}}}{\displaystyle \frac{1}{|P|}},`$
where we have used
$$\left|\mathrm{\Gamma }\left(\frac{1}{2}+\frac{i}{\mathrm{}}\right)\right|^2=\frac{\pi }{\mathrm{cosh}\pi /\mathrm{}}.$$
The $`I`$-flux across a section $`\mathrm{\Sigma }^{\mathrm{out}}`$ defined by $`P=P_0`$ is therefore
$$\psi _{}^+/_{\mathrm{\Sigma }^{\mathrm{out}}}\psi _{}^+=\sigma \frac{\mathrm{}^{\sigma \pi /\mathrm{}}}{\mathrm{}^{\pi /\mathrm{}}+\mathrm{}^{\pi /\mathrm{}}}$$
where here $`\sigma `$ is the sign of $`P_0`$. This is positive in the reactants-out channel ($`P>0`$) and negative in the products-out channel ($`P<0`$) which, when we remember that $`\widehat{\mathrm{\Theta }}`$ is defined so that upward fluxes are positive, is consistent with Figure 2. We have therefore confirmed Equation (7).
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# Fitting Non-Minimally Coupled Scalar Models to Gold SnIa Dataset
## I Introduction
It has been suggested strongly by detailed observations of distant Type Ia Supernovae (SnIa) gold that our universe may be in a phase of accelerating expansion. On the other hand, measurements of cosmic microwave background (CMB) cmb and other surveys lss indicate the spatial flatness of the universe. The simplest explanation for these observations is a dominating cosmological constant, whose equation of state is $`\omega _\mathrm{\Lambda }=p_\mathrm{\Lambda }/\rho _\mathrm{\Lambda }=1`$, with a complementary Cold Dark Matter (LCDM):
$$H^2(z)=H_0^2[\mathrm{\Omega }_\mathrm{\Lambda }+\mathrm{\Omega }_{m0}(1+z)^3]$$
(1)
where $`z`$ is the redshift, $`\mathrm{\Omega }_{m0}`$ is the present value of matters (including the dark matter, whose existence has well been confirmed) and $`\mathrm{\Omega }_\mathrm{\Lambda }`$ that of the cosmological constant. Spatial flatness demands that $`\mathrm{\Omega }_\mathrm{\Lambda }+\mathrm{\Omega }_{m0}=1`$.
LCDM is a simple proposition but requires extreme fine tuning of the cosmological constant $`\rho _\mathrm{\Lambda }`$. More so, its fit with the Gold dataset (157 most updated and reliable set of SnIa) is not very good. Alternative models with evolving dark energy, such as quintessence, phantom, extended gravity, were then put forward de , most of which contain a rolling scalar field. On the other hand, recent analysis of Gold dataset and other observations seem to favor an equation of state $`\omega <1`$ and oscillations of cosmological expansion state ; oscillation .
In normal theories of scalar fields, the equation of state is usually
$$\omega _\varphi =\frac{\frac{1}{2}\dot{\varphi }^2V(\varphi )}{\frac{1}{2}\dot{\varphi }^2+V(\varphi )}.$$
(2)
Obviously, it is not possible to have $`\omega <1`$. One attractive way out of this situation is to consider extended theories in which scalars are non-minimally coupled. These extended theories appear naturally and widely in string theories and theories of extra dimensions in general.
In this paper, we consider simple theories of a non-minimally coupled scalar field $`\varphi `$ which is oscillating around the minimal of the potential $`V(\varphi )`$. In this type of theories, $`\omega <1`$ and oscillations of the cosmological expansion can be obtained readily, as the oscillating scalar field is now strongly coupled to the gravity, and it is well known that non-minimally coupled theories can violate energy conditions.
Fitting to the Gold dataset, these models yield sensible results comparable with other models. The best fit parameters of these models show that the oscillating field $`\varphi `$ may account for both the dark matter and the dark energy in the cosmological expansion. The behavior of $`\varphi `$ depends on the energy density, so $`\varphi `$ could differ significantly within or without galaxies. Properly averaged over space, $`\varphi `$ is scale-dependent. We have concentrated our analysis mainly on large scales, though it is possible that the dark matter appearing on other observations is also due to motion of the same scalar field. Interestingly, our models which have properties of $`\omega <1`$ and oscillating cosmological expansion seem to have better fit, which can be interpreted as circumstantial evidence for oscillations of the cosmological expansion oscillation . For comparison, โ$`H^2(z)z`$โ and โ$`\omega (z)z`$โ relations are plotted for our models, LCDM model and for the best fit model OA Var.(1) in oscillation .
The paper is organized as follows. In section II we present fundamentals of non-minimally coupled theories and outline the procedure of fitting models to the Gold dataset. We then analyze two specific models for illustrations, with results presented in detail in section III. In section IV, we conclude with some discussions.
## II FUNDAMENTALS OF NON-MINIMALLY COUPLED THEORY AND THE FITTING PROCEDURE
Generically, the lagrangian density for a non-minimally coupled scalar field theory assumes the following general form
$$S=d^4x\sqrt{g}\left[\frac{1}{2}F(\varphi )R\frac{1}{2}W(\varphi )^\mu \varphi _\mu \varphi U(\varphi )+L_N\right]$$
(3)
where $`L_N`$ is the total lagrangian density of fields and matters in the universe other than the field $`\varphi `$. In this paper, the theory is parameterized such that $`W(\varphi )=1`$ and $`F(\varphi )`$ of the simple form
$$F(\varphi )=1\xi \varphi ^2$$
(4)
here we have set $`8\pi G=1`$. Assuming a flat Friedmann-Robertson-Walker metric, i.e., we assume a flat prior in all calculations of this paper, the Einstein equations are scalar1 ; scalar2 :
$`H^2`$ $`=`$ $`{\displaystyle \frac{1}{3F}}\left(\rho _N+{\displaystyle \frac{1}{2}}\dot{\varphi }^2+U3H\dot{F}\right)`$ (5)
$`\dot{H}`$ $`=`$ $`{\displaystyle \frac{1}{2F}}\left[(\rho _N+p_N)+\dot{\varphi }^2+\ddot{F}H\dot{F}\right]`$ (6)
where $`\rho _N`$ and $`p_N`$ come from the $`L_N`$, representing, respectively, the energy density and the pressure of matters and fields in the universe other than $`\varphi `$. In what follows, we will approximate these matters and fields as a perfect fluid with $`p_N=0`$.
For simplification, one makes the following redefinitions
$`Q=H^2/H_0^2,V(\varphi )`$ $`=`$ $`U(\varphi )/H_0^2`$ (7)
where $`H_0`$ is the present value of expansion rate $`H`$. From Eqs (5) and (6), one obtains new ; scalar
$`F^{\prime \prime }+\left[{\displaystyle \frac{Q^{}}{2Q}}{\displaystyle \frac{4}{1+z}}\right]F^{}`$ $`+`$ $`\left[{\displaystyle \frac{6}{(1+z)^2}}{\displaystyle \frac{2}{(1+z)}}{\displaystyle \frac{Q^{}}{2Q}}\right]F{\displaystyle \frac{2V}{(1+z)^2Q}}3{\displaystyle \frac{1+z}{Q}}\mathrm{\Omega }_{N0}=0`$ (8)
$`Q(z)`$ $`=`$ $`{\displaystyle \frac{V/3+(1+z)^3\mathrm{\Omega }_{N0}}{FF^{}(1+z)\frac{\varphi ^2}{6}(1+z)^2}}`$ (9)
where $`\mathrm{\Omega }_{N0}=\rho _{N0}/3H_0^2`$. In these two equations, we have replaced derivatives with respect to time by those with respect to the redshift $`z`$, which are in turn denoted by primes. Substitution of Eq (9) into Eq (8) yields a second-order differential equation which includes only field $`\varphi `$.
To be consistent with measurements in the solar system solar , stringent constraints have to be put on these theories. Defining a post-Newtonian parametergr
$$\gamma 1=\frac{(dF/d\varphi )^2}{F+(dF/d\varphi )^2},$$
(10)
solar system tests give the following (present) upper limit solar
$$|\gamma 1|<4\times 10^4,$$
(11)
which is equivalent to
$$\frac{1}{F}\left(\frac{dF}{d\varphi }\right)_0^2<4\times 10^4|\xi \varphi (0)|<10^2$$
(12)
To generate suitable oscillations in the cosmological expansion, $`\xi `$ should be of $`O(1)`$ (see the next section). The constraint is now converted into a constraint on $`\varphi (0)`$. Numerically, $`\varphi (0)<10^3`$ is strict enough.
Tests in the solar system constrain only the behaviors of $`\varphi `$ in our galaxy. According to its equation of motion
$$^2\varphi V_{,\varphi }+\frac{1}{2}F_{,\varphi }R=0$$
(13)
(where $`R`$ is the Ricci scalar) the evolution of $`\varphi `$ depends on the energy density, for $`R`$ is proportional to the energy density. Usually, the potential $`V(\varphi )`$ is tiny, such as the ones appearing in the next section. Within galaxies, where the density of matters is much larger compared with the average density of the universe, the potential term in the above equation can be neglected. Outside of galaxies or averaged over large scales, $`V(\varphi )`$ plays the main role and the frequency of oscillations of $`\varphi `$ would be much smaller. So the magnitude and the evolution of $`\varphi `$ are scale-dependent and sensitive to positions, which produce varied effects at small and large scales. As we are interested in cosmological expansion at large scales, $`\varphi (0)`$ is much less constrained.
For simplicity, we will take $`\varphi (0)=0`$ as one initial condition. Other proper choices of $`\varphi (0)`$ will yield qualitatively similar results. On the other hand, the effective Newton constant
$$G_{\mathrm{eff}}=\frac{1}{F}\frac{2F+4(dF/d\varphi )^2}{2F+3(dF/d\varphi )^2}$$
(14)
changes in time, which is also constrained by present observations
$$\left|\frac{\dot{G}_{\mathrm{eff}}}{G_{\mathrm{eff}}}\right|_{z=0}<6\times 10^{12}\mathrm{yr}^1$$
(15)
By choosing the particular initial condition $`\varphi (0)=0`$, one automatically has
$$\left|\frac{\dot{G}_{\mathrm{eff}}}{G_{\mathrm{eff}}}\right|_{z=0}=0$$
(16)
The second initial condition is obtained by evaluating Eq (9) at $`z=0`$:
$$\varphi ^2(0)=62V(0)6\mathrm{\Omega }_{N0}$$
(17)
From Eqs (8), (9) and these initial conditions, one obtains the current equation of state
$$\omega (0)=\frac{14\xi }{6}\varphi ^2(0)\frac{V(0)}{3}$$
(18)
If $`\xi `$ is large enough, $`\omega (0)`$ can be less than $`1`$ easily.
Given these two initial conditions, the second-order differential equation of $`\varphi `$ can be solved, at least numerically, once the exact form of $`V(\varphi )`$ and the value of $`\mathrm{\Omega }_{N0}`$ are also known. Substituting the solution $`\varphi (z)`$ back into Eq (9), we get $`H(z)`$ for arbitrary redshift $`z`$. Thus, we solve the problem completely. It is then straightforward to compare the theory with observations.
To be definitive, one calculates the goodness of fit of the theory with the observed Gold dataset:
$$\chi ^2(\overline{M},H)=\underset{k=1}{\overset{157}{}}\frac{\left(m^{ob}(z_k)m^{th}(z_k,H,\overline{M})\right)^2}{\sigma _{m^{ob}(z_k)}^2}$$
(19)
where
$$m^{th}(z,H,\overline{M})=\overline{M}+5\mathrm{log}_{10}\left((1+z)_0^z๐z^{}\frac{H_0}{H(z^{})}\right)$$
(20)
is the apparent magnitude and
$$\overline{M}=M+5\mathrm{log}_{10}(\frac{cH_0^1}{Mpc})+25$$
(21)
is the magnitude zero point offset.
To find the best fit parameters of the model, one minimizes $`\chi ^2`$. Since $`\chi ^2`$ is a second-order polynomial of $`\overline{M}`$, the minimization procedure with respect to this parameter is straightforward and simple. The final result of $`\chi _{min}^2`$ can be expressed in terms of quantities independent of $`\overline{M}`$ oscillation :
$$\chi ^2(H)=AB^2/C$$
(22)
where
$`A`$ $`=`$ $`{\displaystyle \underset{k=1}{\overset{157}{}}}{\displaystyle \frac{\left(m^{ob}(z_k)m^{th}(z_k,H,\overline{M}=0)\right)^2}{\sigma _{m^{ob}(z_k)}^2}}`$
$`B`$ $`=`$ $`{\displaystyle \underset{k=1}{\overset{157}{}}}{\displaystyle \frac{\left(m^{ob}(z_k)m^{th}(z_k,H,\overline{M}=0)\right)}{\sigma _{m^{ob}(z_k)}^2}}`$ (23)
$`C`$ $`=`$ $`{\displaystyle \underset{k=1}{\overset{157}{}}}{\displaystyle \frac{1}{\sigma _{m^{ob}(z_k)}^2}}`$
Eq (22) is the starting point of our numerical analysis. By minimizing this quantity, one obtains the best fit values of relevant parameters. As we shall see in the next section, both $`\omega (0)<1`$ and oscillations of cosmological expansion are favored by experiments.
## III NUMERICAL RESULTS
Taking two special potential forms for illustration, we now calculate the fit of non-minimally coupled theory with the Gold dataset and find the best fit parameters in the process. In order to obtain acceleration and oscillations of the cosmological expansion simultaneously, we consider potentials that have non-zero minima $`V_0`$ at the point $`\varphi =0`$. Though we have only analyzed two forms of potentials in this paper, our analysis can easily be extended to other potentials. Other potentials of this type could yield comparable fitting results, as in general one can easily obtain appropriate oscillations and acceleration of the cosmological expansion from this type of models. All results are shown in Table I. The errors are obtained with an increase of $`\chi _{min}^2`$ by 1, while we have assumed that the value of $`\mathrm{\Omega }_{N0}`$ is at lest $`0.045`$, which is the density of all know matters including baryons and CMBR. When estimating the error for each parameter, all other parameters are fixed at their best fit values. For comparison and concreteness, we will plot the evolving of $`\varphi (z)`$, $`H^2(z)`$ and $`\omega (z)`$ for all models in FIG. 5, FIG. 6 and FIG. 7, respectively.
### III.1 $`V(\varphi )=V_0e^{a_1\varphi ^2}`$
This potential appears naturally in supergravity theories (SUGRA) sugr . This form was also reconstructed from SnIa Gold dataset in new , where it is only a good fit of the polynomial parametrization model in reference oscillation and $`\varphi `$ is not oscillating. In total, we have four free parameters: $`\xi ,\mathrm{\Omega }_{N0},V_0`$ and $`a_1`$. Their best fit values are shown in first line of Table I (NMC E(0)).
FIG. 1 shows the dependence of $`\chi _{min}^2`$ on $`\mathrm{\Omega }_{N0}`$, when other parameters are allowed to vary freely. The best fit of $`\mathrm{\Omega }_{N0}`$ is about $`0.066`$. However, Fig. 1 indicates that one gets almost no change in $`\chi ^2`$ if $`\mathrm{\Omega }_{N0}`$ is set to be $`0.045`$. The latter is the value of all known matters including baryons and CMBR baryon . That is to say, the existence of a single non-minimally coupled scalar field $`\varphi `$ may play both the roles which dark matter and dark energy would play in LCDM on the cosmological expansion. Of course we can split this effect into two parts: dark matter (which also generates oscillations of the cosmological expansion) and dark energy. Clearly, this splitting is rather artificial, as there are no obvious distinctions between these two parts in this model. They can just be regarded as a whole as the expansion of the universe is concerned. Now that there is no isolated dark energy, we plot in FIG. 7 the total equation of state of the universe. Dark matters from other observations can also come from this field galaxy , though the behavior of the field should be scale-dependent dd . In this paper we only consider the effect of the field $`\varphi `$ on cosmological expansion at large scales.
Clearly, the bigger $`\mathrm{\Omega }_{N0}`$ is, the worse is the fit to Gold dataset. This may be seen as follows. To obtain enough amplitude of oscillations of field $`\varphi `$, i.e., to obtain enough oscillations of cosmological expansion to fit observations, $`\varphi ^2(0)`$ should be large enough. On the other hand, to have a dominating dark energy, $`V_0`$ should be positive and large enough. According to Eq (17), $`\mathrm{\Omega }_{N0}`$ cannot be too larger. In the rest of this subsection, we will fix $`\mathrm{\Omega }_{N0}=0.045`$ to analyze other parameters. A re-analysis of this model with this value of $`\mathrm{\Omega }_{N0}`$ is given in the second row of Table I.
The parameter $`V_0`$ here plays somewhat the role of cosmological constant, with the effective $`\mathrm{\Omega }_\mathrm{\Lambda }=V_0/3`$. Our best value $`V_02.19`$ is consistent with the best fit value $`\mathrm{\Omega }_\mathrm{\Lambda }=0.69`$ in LCDM model. If $`V_0`$ departs too much from the best fit value, we will have a comparatively large $`\chi ^2`$, which is confirmed from FIG. 2. In FIG. 2 we have plotted the $`\chi ^2V_0`$ relation by fixing the other parameters as: $`\mathrm{\Omega }_{N0}=0.045`$, $`\xi =1.79`$ and $`a_1=25.0`$.
Now we turn to the parameter $`\xi `$, which represents the coupling strength between gravity and $`\varphi `$. As $`\varphi `$ oscillates around the minimum of potential $`V(\varphi )`$, $`\xi `$ affects the amplitude of oscillations of cosmological expansion. If $`\xi `$ is very small, there will be almost no oscillation appearing in the expansion rate $`H^2`$ and it is difficult to cross the $`\omega =1`$ line. Its behaves nearly the same as that of the minimally coupled theory scalar1 . If $`\xi `$ is very larger, the motion of the universe will be drastically changed. So the value of $`\xi `$ should be moderate. In fact, the best fit value of $`\xi `$ is of the order $`O(1)`$ (see Fig. (3)).
Last but not the least, the parameter $`a_1`$ can affect frequency of oscillations, its effects on $`\chi ^2`$ are shown in FIG. 4. As mentioned above, we have plotted $`\varphi (z)`$, $`H^2(z)`$, and $`\omega (z)`$ in terms of $`z`$ for all models, in Figs. 5, 6 and 7. One clearly sees oscillations of $`\varphi `$ and of the cosmological expansion over time in this type of model. The universe may well be just in a period when the $`\omega `$ is in the valley and later on will transit to a period of deceleration slowly.
### III.2 $`V(\varphi )=V_0+a_1\varphi ^n`$
This form of potentials can be regarded as an extension of a constant dark energy, i.e., a cosmological constant $`V_0`$ is added with a field $`\varphi `$ of power potential. The non-minimally coupled scalar field mimics the effects of dark matter and produces oscillations in the expansion rate $`H^2`$. We have considered cases of $`n=2,4,6`$, with their best fit values of parameters shown in Table I. As shown in FIG. 6 and FIG. 7, the expansion patterns of these models are similar to those of LCDM and oscillate around the line of LCDM.
For $`n=2`$, the best fit value $`a_1<0`$ and the potential $`a_1\varphi ^2`$ is negative, but coupling of the scalar field with gravity can drive up the potential. The field $`\varphi `$ is still oscillating around $`\varphi =0`$ if all parameters are within certain range around the best fit values. In this model the constraints on parameters are quite strict, as shown in Table I. In Figs. 5, 6 and 7, we see that oscillations of the cosmological expansion only appear near the $`\varphi =0`$ points, while in other places there is essentially no oscillation.
For $`n=4`$ and more so for $`n=6`$, the cosmological constant $`V_0`$ is quite large. Then there is not much dark matter, which seems to be in conflict with other observations of dark matter. It is possible that we can have enough dark matters at relatively small scales, as the effect of $`\varphi `$ is scale-dependent. For $`n=6`$, the best fit value of $`\xi `$ is $`0`$, but there are still oscillations in the expansion rate $`H^2`$, with a much smaller amplitude. Its best fit value of $`\mathrm{\Omega }_{N0}`$ is negative, much smaller than $`0.045`$, the value of known matters. We set the value $`\mathrm{\Omega }_{N0}=0.045`$ in Table I by hand.
## IV Conclusion
We have numerically analyzed the expansion of the universe in non-minimally coupled scalar field theories with two special forms of potentials: $`V_0e^{a_1\varphi ^2}`$ and $`V_0+a_1\varphi ^n`$. These types of potentials have a no-zero minimum at $`\varphi =0`$, which generates the necessary acceleration of the cosmological expansion. On the other hand, the oscillating $`\varphi `$ produces oscillations of the expansion and yield an equation of state $`\omega <1`$, as favored by the fit with 157 observed SnIa Gold dataset. We have further compared our models with the Gold dataset to find the best fit parameters. The best fit potential is of the form $`V_0e^{a_1\varphi ^2}`$ with $`\chi _{min}^2=170.127`$. The best value of all matter density in universe except $`\varphi `$ is found to be $`\mathrm{\Omega }_{N0}=0.066`$, which is very close to all energy density of the known matters including baryons and CMBR. So the non-minimally coupled scalar field $`\varphi `$ may provide the main source of dark matter as well as dark energy. Of course further investigation will be needed to make these assertions more definitive.
Acknowledgments: This work is supported in part by the National Science Foundation of China (10425525).
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# Polarization Relaxation Induced by Depolarization Field in Ultrathin Ferroelectric BaTiO3 Capacitors
## Abstract
Time-dependent polarization relaxation behaviors induced by a depolarization field $`E_d`$ were investigated on high-quality ultrathin SrRuO<sub>3</sub>/BaTiO<sub>3</sub>/SrRuO<sub>3</sub> capacitors. The $`E_d`$ values were determined experimentally from an applied external field to stop the net polarization relaxation. These values agree with those from the electrostatic calculations, demonstrating that a large $`E_d`$ inside the ultrathin ferroelectric layer could cause severe polarization relaxation. For numerous ferroelectric devices of capacitor configuration, this effect will set a stricter size limit than the critical thickness issue.
preprint: D. J. Kim et al.
With recent breakthroughs in fabricating high-quality oxide films Ahn ; YSKim1 ; HNLee , ultrathin ferroelectric (FE) films have attracted much attention from the scientific as well as application points of view. As the FE film thickness $`d`$ approaches tens of unit cell length, the FE films often show significantly different physical properties from those of bulk FE materials. Some extrinsic effects, especially coming from FE film surfaces and/or interfaces with other materials, could be very important Shaw . For some other cases, intrinsic physical quantities could play vital roles in determining the unique properties of ultrathin films.
Many FE-based electronic devices have the capacitor configuration, where a FE layer is inserted between two conducting electrodes. Then, polarization bound charges will be induced at the surfaces of the FE layer, but compensated by free charge carriers in the conducting electrodes. In real conducting electrodes, however, the compensating charges will be induced with a finite extent, called the screening length $`\lambda `$. This will result in an incomplete compensation of the polarization charges. Such an incomplete charge compensation should induce a depolarization field $`E_d`$ inside the FE layer, with a direction opposite to that of the FE polarization $`P`$ Mehta . Therefore, $`E_d`$ will appear in every FE capacitor, and its effects will becomes larger with the decrease of $`d`$ Mehta . (For a FE film without electrodes, there is no compensation for the polarization bound charge, so the value of $`E_d`$ will become even larger than that of the FE capacitor case.) $`E_d`$ has been known to be important in determining the critical thickness Junquera and domain structure of ultrathin FE films Kornev ; Wu ; Fong , and reliability problems of numerous FE devices Kang1 ; Kang2 .
Recently, using a first principles calculation, Junquera and Ghosez investigated the critical thickness of BaTiO<sub>3</sub> (BTO) layers in SrRuO<sub>3</sub>(SRO)/BTO/SRO capacitor Junquera . For calculations, they assumed that all of the BTO and SRO layers were fully strained with the SrTiO<sub>3</sub> substrate. By taking the real SRO/BTO interfaces into account properly, they showed that $`E_d`$ could make the ferroelectricity vanish for the BTO films thinner than 6 unit cells, i.e. 2.4 nm Junquera . More recently, using pulsed laser deposition with a reflection high energy electron diffraction monitoring system, we fabricated high-quality fully-strained SRO/BTO/SRO capacitors on SrTiO<sub>3</sub> substrates with $`d`$ between 5 and 30 nm YSKim1 ; YSKim2 . With a very low leakage current, we could directly measure their $`P`$-$`E`$ hysteresis loops YSKim1 . In this letter, we report the time-dependent polarization changes of the ultrathin BTO films. We find that the net $`P`$ of the ultrathin BTO films decreases quite rapidly in time. We will show that the $`P`$ relaxation should originate from $`E_d`$. By compensating for $`E_d`$ with an external potential, we can determine the $`E_d`$ values of the SRO/BTO/SRO capacitors experimentally. These measured $`E_d`$ values agree with the values from the electrostatic calculations. Finally, we will discuss the effect of the $`P`$ relaxation on a practical size limitation imposed on actual FE devices.
In our earlier report YSKim1 , we obtained the thickness-dependent remnant polarization $`P_r`$ values from the $`P`$-$`E`$ hysteresis loops, measured at 2 kHz in ultrathin FE films as thin as 5 $``$ 30 nm. With further studies on the frequency dependence of the $`P_r`$ values in $`P`$-$`E`$ hysteresis loops, as shown in Fig. 1(a) for a 15 nm thick BTO capacitor, we found differences in the $`P_r`$ values when the measuring frequency is varied. These results suggest that the FE domain dynamics should play an important role for ultrathin FE films, where the FE domain wall motion is known to be strongly suppressed Tybell . Note that the first principles calculation (FPC) and the Landau-Devonshire calculation (LDC) do not consider the domain dynamics, so their predicted polarization values should be called as spontaneous polarization $`P_s`$.
Since the $`P`$ value significantly affects the subsequent analysis of $`P`$ relaxation, precise determination of $`P_s`$ values is necessary. To determine the precise values of $`P_s`$, we applied pulse trains, which are schematically shown in the inset of Fig. 1(a) Smolenskii . The interval between write and read pulses was set to 1 $`\mu `$s to minimize the effects of the $`P`$ relaxation, and the current responses under the read pulse were measured. The total amount of charge is obtained by integrating the current responses in time. The read pulses with different heights were used to obtain the linear part of the polarization under an external electric field The $`P_s`$ values can be obtained by extrapolating the linear part of the polarization to zero electric field. The triangles (black) and circles (green) in Fig. 1(b) show the $`P_r`$ values measured at 2 and 100 kHz, respectively. Also, the squares (red) show the $`P_s`$ values from the pulse test. The solid (green) and dashed (blue) curves show the theoretical predictions from the FPC Junquera and the LDC Pertsev , respectively, which take account of $`E_d`$. Note that neither of these theories can explain the thickness-dependence of $`P_s`$ quantitatively. However, it is known that the FPC predicts systematically somewhat lower bulk lattice constants compared to real values, so the compressive stress predicted by the FPC could be smaller than that in the fully strained sample, resulting in a smaller $`P_s`$. To avoid this systematic error, we normalized the polarization values to those of a 30 nm thick BTO capacitor. We found that the thickness-dependent scaling of $`P_s`$ also follows the FPC predictions quite well, as shown in Fig. 1(c).
The large difference in $`P`$ values between the 2 and the 100 kHz tests indicates that there should be a strong change in the net $`P`$ between 10 and 500 $`\mu `$s. Time-dependent $`P`$ changes were investigated by applying two kinds of pulse trains, as shown in the inset of Fig. 2(a). For the write and the read pulses with the same (opposite) polarities, the amount of nonswitching (switching) $`P`$ can be determined Kang1 . The difference $`\mathrm{\Delta }P`$, between the switching and the nonswitching $`P`$ should be twice as large as the net $`P`$. As shown in Fig. 2(a), $`\mathrm{\Delta }P`$ decreases quite rapidly for the film with $`d`$ = 15 nm; $`\mathrm{\Delta }P`$ falls to less than $`10\%`$ of the $`P_s`$ value within a relaxation time $`t_{relax}`$ of 1000 s. As shown with the solid squares (black) in Fig. 2(b), $`\mathrm{\Delta }P`$ decay follows a power-law dependence on $`t_{relax}`$. Similar power-law decays of $`\mathrm{\Delta }P`$ were observed for all the BTO films in the thickness range of 5 $``$ 30 nm. Note that such a strong polarization relaxation could pose a serious problem in capacitor-type ultrathin FE devices.
What is the origin of such strong polarization relaxations? We thought that they could be closely related to large $`E_d`$ induced inside the BTO films. To verify this idea, we slowed down the relaxation phenomena by applying an external voltage, as shown in the inset of Fig. 2(a). The values of the applied external electric field $`E_{ext}`$ were obtained by dividing the applied external voltage by the corresponding film thickness. When the external field is applied in the opposite direction of $`E_d`$, the potential gradient inside the FE layer will decrease. Figure 2(b) shows that the slope of the power-law decay becomes smaller, as $`E_{ext}`$ increases. Assuming that the depth of the double-well potential for BTO ferroelectricity can be considered negligible compared to the effect of $`E_d`$, we approximately determined experimental $`E_d`$ values from the applied electric field under which the slope becomes zero. Since $`E_d`$ is proportional to $`P`$, the $`E_d`$ value should increase slightly on application of $`E_{ext}`$. After correcting this minor contribution, we could determine the $`E_d`$ values, which are plotted as solid circles (red) in Fig. 2(c).
From electrostatic calculations on the capacitor geometry, Mehta et al. showed that
$$E_d=\frac{P}{ฯต_0ฯต_F}\left(\frac{2ฯต_F/d}{2ฯต_F/d+ฯต_e/\lambda }\right),$$
(1)
where $`d`$ is the thickness of the FE layer, and $`ฯต_F`$ and $`ฯต_e`$ are the relative dielectric constants of the FE layer and the electrode, respectively Mehta . To obtain theoretical $`E_d`$ values for our SRO/BTO/SRO capacitors, we have to know accurate values of $`ฯต_e`$, $`\lambda `$, and $`ฯต_F`$. Unfortunately, the reported physical parameter values in the literature vary Mehta ; Junquera ; Black ; Dawber . Also, we could not find any definite experimental study on $`ฯต_e`$.
To obtain the value of $`ฯต_e`$ for an SRO electrode, we used optical spectroscopy. We measured the optical reflectivity spectra of epitaxial SrRuO<sub>3</sub> films (thickness: about 0.5 $`\mu `$m) in a wide frequency region between 5 meV and 30 eV and performed a Kramers-Kronig analysis to obtain the frequency-dependent dielectric function, $`ฯต(\omega )`$ \[=$`ฯต^{}(\omega )+iฯต^{\prime \prime }(\omega `$)\]. The details of these measurements and analysis were published elsewhere JSLee1 ; JSLee2 . The open squares in Fig. 3(a) and the inset show experimental values of $`ฯต^{}(\omega )`$ and $`ฯต^{\prime \prime }(\omega )`$, respectively. Note that $`ฯต_e`$ in Eq. (1) represents the dielectric response from the bound charges, namely bound electrons and phonons. Since SRO is metallic, there should be a large contribution from the free Drude carriers, which masks the dielectric response from the bound charges. To obtain $`ฯต_e`$, we decompose $`ฯต(\omega )`$ into a free carrier contribution $`ฯต_{coherent}(\omega )`$ and a bound electron contribution $`ฯต_{bound}(\omega )`$ by fitting the experimental $`ฯต(\omega )`$ with a series of Lorentz oscillators, which are displayed as the dotted (blue) lines in the inset of Fig. 3(a). The dash-dotted (blue) lines indicate the bound electron contribution. From the dc limit of $`ฯต_{bound}(\omega )`$, we could estimate that the bound electron contribution to $`ฯต_e`$ is about 8.17. The phonon contribution to $`ฯต_e`$ was evaluated in a similar way by analyzing the phonon spectra and found to be about 0.28 JSLee1 . Consequently, $`ฯต_e`$ is determined to be about 8.45.
Using the carrier density $`n_01.2\times 10^{22}/`$cm<sup>3</sup> of SRO Shepard , the experimental value of $`ฯต_e`$, and the effective mass of an electron $`m_{eff}7m_e`$, where $`m_e`$ is the mass of a free electron Cao ; Okamoto , we applied the free electron model and obtained $`\lambda =0.8\pm 0.1`$ ร
Mehta ; Kittel . We also measured $`ฯต_F`$ from the capacitance-electric field $`C`$-$`E`$ curves of BTO capacitors. Figure 3(b) shows the $`C`$-$`E`$ curve for the 5 nm BTO capacitor. The $`C`$-$`E`$ curve has the hysteretic behavior typical for a FE capacitor. The BTO capacitors with 5 $``$ 30 nm thickness show almost the same $`ฯต_F`$-$`E`$ curves. The $`ฯต_F`$ values can vary from 70 to 230 depending on the applied $`E`$. Since most of our experiments were performed under a finite applied field, which corresponds to a value between 1 and 2 MV/cm, the $`ฯต_F`$ were estimated to be about 80 interface .
With the measured values of $`ฯต_e`$, $`\lambda `$, and $`ฯต_F`$, we could estimate the theoretical $`E_d`$ values from Eq. (1) with the $`P_s`$ values obtained from the pulse test. The open squares in Fig. 2(c) are the theoretical $`E_d`$ values. The solid (green) line shows the theoretical $`E_d`$ values with the $`P_s`$ values, obtained from the FPC. These theoretical $`E_d`$ values from the electrostatic model agree quite well with the experimental $`E_d`$ values, determined from the polarization relaxation. It should be noted that the $`E_d`$ values are comparable with or even larger than the measured coercive fields (in our samples, 300 $``$ 400 kV/cm). These large $`E_d`$ values can cause $`P`$ reversal and FE domain formation, which will result in a reduction of the net $`P`$ value as time elapses. The fact that two independent determinations provided nearly the same $`E_d`$ values demonstrates that *the polarization relaxation behavior should be dominated by $`E_d`$ inside the FE layer*.
Note that the $`E_d`$-induced $`\mathrm{\Delta }P`$ decay comes intrinsically from the incomplete compensation of the $`P`$ charges (due to the finite screening length of the electrodes) in real conducting electrode, so that it will inevitably pose a fundamental limit for most FE device applications using the capacitor configuration. This limitation should be much more severe than that due to the critical thickness of the FE ultrathin films Junquera . Even if the FE film is thicker than the critical thickness, it is feasible that the $`E_d`$-induced $`\mathrm{\Delta }P`$ decay is large enough to make the net $`P`$ decrease significantly, resulting in retention failures for numerous FE devices. As $`d`$ decreases, $`E_d`$ increases significantly. With the current miniaturization trends in some FE devices, the large value of $`E_d`$ should play a very important role in determining the ultimate size limits of FE devices.
In order to reduce device failure due to the polarization relaxation, we can try to select better electrode and FE materials. Noble metals, such as Pt, have been considered better electrodes because they have high carrier density (resulting in $`\lambda `$ values smaller than that of SRO). However, the $`ฯต_e`$ values of typical noble metals are much smaller than that of SRO, i.e. 8.45 Ehrenreich , so $`E_d`$ in capacitors with noble metal electrodes can be large. For example, $`E_d`$ in the range of 500 $``$ 900 kV/cm is expected for a 15 nm thick BTO film with noble metal electrodes (typically, $`\lambda =0.40.5`$ ร
, $`ฯต_e=24`$). Thus, the $`E_d`$-induced $`P`$ relaxation for the ultrathin BTO capacitors with the noble metal electrodes could be at least equal to or worse than that with SRO electrodes. Proper FE material selection can be another option. Since PbTiO<sub>3</sub> is known to have a much deeper double-well potential than that of BTO Pertsev ; Cohen , the $`P`$ relaxation should occur at a much lower rate even with the same value of $`E_d`$. Optimization of FE materials should be of great importance for the improvement of ultrathin film nanoscale FE device performances.
In summary, we demonstrated that the depolarization field inside the ferroelectric film could cause a severe polarization relaxation. By slowing down the relaxation under an external field, we could determine the depolarization field in a real capacitor of ultrathin SrRuO<sub>3</sub>/BaTiO<sub>3</sub>/SrRuO<sub>3</sub> experimentally, which result is in good agreement with electrostatic calculations. Our investigation demonstrates that the depolarization field originates from intrinsic properties of electrode material such as the finite screening length and that the depolarization field should play an important role in domain dynamics in ultrathin FE films. The polarization relaxation due to the depolarization field could pose a serious size limitation for ultrathin ferroelectric devices.
The authors thank Prof. Sug-Bong Choe in Seoul National University for valuable discussions. This work was financially supported by the Korean Ministry of Science and Technology through the Creative Research Initiative program and by KOSEF through CSCMR.
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# Low-energy dipole excitations towards the proton drip-line: doubly magic 48Ni
## Abstract
The properties of the low-energy dipole response are investigated for the proton-rich doubly magic nucleus <sup>48</sup>Ni, in a comparative study of two microscopic models: fully self-consistent Relativistic Random-Phase Approximation(RRPA) based on the novel density-dependent meson-exchange interactions, and Continuum Random-Phase Approximation(CRPA) using Skyrme-type interactions with the continuum properly included. Both models predict the existence of the low-energy soft mode, i.e. the proton pygmy dipole resonance (PDR), for which the transition densities and RPA amplitudes indicate the dynamics of loosely bound protons vibrating against the rest of the nucleons. The CRPA analysis indicates that the escape width for the proton PDR is rather large, as a result of the coupling to the continuum.
PACS: 21.10.Gv, 21.30.Fe, 21.60.Jz, 24.30.Cz
thanks: Permanent address: Joint Institute of Nuclear Research, Dubna, Russia
One of the major challenges in the region of unstable nuclei is the understanding of soft modes of excitations which involve loosely bound nucleons, nucleon halo, or skin. In particular, in neutron-rich nuclei, nucleons from the neutron skin may give rise to a low-energy dipole mode, known as a pygmy dipole resonance (PDR) . Some experimental evidence about the low-lying dipole excitation phenomena in neutron-rich nuclei is available from (i) electromagnetic excitations in heavy-ion collisions , and (ii) nuclear resonance fluorescence experiments for nuclei with moderate neutron to proton number ratios, e.g. <sup>116,124</sup>Sn , <sup>138</sup>Ba , Pb , and Ca isotopes, and N=82 isotones . On the other side, a variety of theoretical models have in recent years been employed in studies of the low-lying E1 strength: continuum quasiparticle RPA formulated in the coordinate-space Hartree-Fock-Bogoliubov framework , continuum RPA with Woods-Saxon potential for the ground state and Landau-Migdal force in the residual interaction , the time-dependent density-matrix theory , Skyrme Hartree-Fock + quasiparticle RPA with phonon coupling , and the quasiparticle-phonon model . In the relativistic framework, properties of low-energy excitations have been systematically studied in the relativistic RPA (RRPA) , and the relativistic quasiparticle RPA (RQRPA) . Both experimental and theoretical studies qualitatively agree on the global properties of the isovector dipole response, i.e., as the number of neutrons increases along the isotope chain, the transition strength distribution is characterized by the appearance of pronounced low-lying E1 strength. Within the relativistic R(Q)RPA studies, it has been shown in a fully self-consistent way that the low-lying pygmy state represents a genuine structure effect: the neutron skin oscillates against the core, exhibiting a collective nature from medium towards heavier nuclei . A detailed quantitative description of the low-lying E1 strength in neutron-rich nuclei is essential for calculations of radiative neutron capture rates in the r-process, and elemental abundances from nucleosynthesis .
The structure of nuclei on the proton-rich side is equally important for revealing many aspects of the underlying many-body problem and properties of the effective nuclear interactions. At the same time, a quantitative description of proton-rich nuclei is of a particular importance for describing the rapid proton capture process of nucleosynthesis. Proton-rich nuclei are characterized by unique ground-state properties such as $`\beta `$ decays with large $`Q`$ values, and direct emission of charged particles. The two-proton ground state radioactivity has recently been observed as a decay mode of <sup>45</sup>Fe . At the present time, knowledge on dipole excitations in nuclei towards the proton drip-line is rather limited. Only very recently, the first microscopic theoretical prediction of the proton pygmy dipole resonance has been indicated in the framework of RQRPA based on the Relativistic Hartree-Bogoliubov model (RHB) . For nuclei close to the proton drip-line, the model calculations predicted the occurrence of pronounced dipole peaks below 10 MeV in excitation energy, due to collective vibrations of loosely-bound protons against the proton-neutron core.
In this letter, we employ two theoretical approaches in a comparative study of the low-lying E1 strength of the proton-rich nucleus $`{}_{}{}^{N}48`$: the fully self-consistent Relativistic RPA (RRPA) based on novel density-dependent interactions, and the Continuum RPA (CRPA) with Skyrme-type interactions. An essential objective of this study is to ensure that the proton PDR is inherent for different models, based on different assumptions and effective interactions, and to quantify the global properties of this mode, i.e. excitation energies and B(E1) strength, which may also be interesting for the future experimental studies. Furthermore, we employ the CRPA model to investigate the role of the coupling to the continuum for the proton PDR. We choose to analyse a doubly magic nucleus <sup>48</sup>Ni, which is the most proton-rich isotope that has been experimentally discovered . For comparison, we also present results for $`{}_{}{}^{N}56`$, which lies close to the valley of stability. Being doubly magic, these two isotopes can be studied within the CRPA and RRPA models, which do not include pairing correlations. Ground state properties of proton-rich nuclei around <sup>48</sup>Ni have extensively been studied within the shell-model , Hartree-Fock Bogoliubov , and Relativistic Hartree-Bogoliubov (RHB) theory . Proton drip-line nuclei are characterized by a reduction of the spin-orbit term of the effective interaction, outer orbits appear to be very weakly bound, and the Fermi energy level may become positive. Due to the presence of the Coulomb barrier, loosely bound orbits are stabilised, in contrast to the neutron drip-line nuclei where the weakly bound neutron orbits are more spatially extended. The interplay of all these effects in the nuclear ground state will shape the properties of the corresponding dipole excitation response.
In the following, we briefly present the basic theoretical background of the Dirac-Hartree+RRPA and Skyrme-Hartree-Fock+CRPA models and their effective interactions.
In recent years, the relativistic mean-field theory and linear response based on density-dependent interactions, turned out to be very successful in studies of nuclear ground-state properties and excitation phenomena with a minimal set of parameters in a fully microscopic way . Within this framework, the nucleus is described as a system of Dirac nucleons which interact in a relativistic covariant manner by exchange of effective mesons. The model is formulated with the Lagrangian density which explicitely includes the density dependence in $`\sigma `$, $`\omega `$, and $`\rho `$ meson-nucleon vertices. In the present study, we employ the density-dependent effective interactions which are constrained by the properties of finite nuclei and nuclear matter: DD-ME1 , and the new interaction DD-ME2 which provides an improved description of the isovector dipole response . In the small-amplitude limit, the RRPA equations are derived from the equation of motion for the nucleon density . Both in RRPA and CRPA models, we use the same form of the effective isovector dipole operator as in Ref. . The RRPA configuration space is constructed from particle-hole ($`ph`$) pairs composed of the particle states above the Fermi level, and hole states in the Fermi sea. In addition, in the relativistic case, one also needs to include transitions to unoccupied states from the Dirac sea . The resulting RRPA discrete spectra are averaged with the Lorentzian distribution which includes an arbitrary choice for the width, $`\mathrm{\Gamma }_{RRPA}`$=1 MeV. The Dirac-Hartree+RRPA model is fully self-consistent, i.e. both the equations of the ground state, and the residual RPA interaction are derived from the same effective Lagrangian. This is an essential property for an accurate decoupling of the spurious center-of-mass motion without need for including any additional free parameters. For the present study we use the Dirac-Hartree model formulated in the harmonic oscillator basis. Within this approach, the particle continuum is represented by a set of discrete states, which are used to construct the RRPA configuration space. In Ref. it has been verified that for nuclei towards the proton drip-line, an expansion in a large oscillator basis (N=20) provides sufficiently accurate solutions, in complete agreement with the model formulated in the coordinate space.
The second theoretical framework for the present study is the Skyrme-Hartree-Fock (SHF) plus Continuum-RPA (CRPA) model. The HF equations describing the ground state are derived variationally from the Skyrme energy functional. The $`ph`$ residual interaction is derived from the same energy functional. In the present study, the Coulomb interaction, as well as spin-dependent terms are omitted from the residual interaction. The CRPA is formulated in the coordinate space, and the particle continuum is fully taken into account. The transition strength distribution $`R(E)`$ is continuous by construction. A small but finite value of Im$`E\mathrm{\Gamma }/2`$ entering the evaluation of the $`ph`$ Green function ensures that bound transitions acquire a finite width and thereby contribute to the distribution. More details on the CRPA model can be found in Refs. and references therein.
For the purposes of the present study we implement various parameterisations of the Skyrme interaction, corresponding to different nuclear-matter properties. In particular, they have different (isoscalar) effective mass $`m^{}/m`$ and isovector effective mass $`m_v^{}/m`$. These quantities have influence on the evaluated properties of the isovector giant dipole resonance (IVGDR) and the low-energy dipole transitions. In general, interactions with a high effective mass are not able to reproduce the IVGDR properties. However, a high $`m^{}/m`$ may be more appropriate to describe correctly the density of states lying close to the Fermi energy, which are relevant for the present study. Therefore, we use several different parameterisations of the Skyrme interaction, in order to ensure the general validity of our conclusions, at least on a qualitative level. We employ parameterisation MSk7 ($`m^{}/m=m_v^{}/m=1.05`$), based on a Hartree-Fock-BCS model , and two interactions from the recent BSk series, namely BSk8 ($`m^{}/m=0.80`$, $`m_v^{}/m=0.87`$) and BSk2 ($`m^{}/m=1.04`$, $`m_v^{}/m=0.86`$) . Both BSk interactions were parameterised by fitting the values of nuclear masses calculated within the Hartree-Fock-Bogoliubov method to essentially all the measured ones. We also use the traditional parameterisation SkM\* ($`m^{}/m=0.786`$, $`m_v^{}/m=0.875`$) , which was extensively used in previous studies of giant resonances and response of exotic nuclei, e.g. Refs. .
Next, we present the results obtained with the two models. In Fig. 1 we plot the ground-state proton and neutron density distributions for <sup>48</sup>Ni and <sup>56</sup>Ni, calculated with Dirac-Hartree and SHF models, based on the DD-ME1 and BSk8 effective interactions respectively. For <sup>56</sup>Ni ($`Z=28`$,$`N=28`$) the proton and neutron density distributions are similar in the nuclear interior and beyond the surface region the differences completely vanish. However, in the case of <sup>48</sup>Ni ($`Z=28`$,$`N=20`$), due to the excess of loosely-bound protons, the proton density distribution is considerably extended beyond the neutron density distribution. This effect is especially pronounced at radial distances $`r>2`$ fm, and it resembles in structure a proton skin. Due to the presence of the Coulomb barrier, which tends to localise protons within the nuclear interior, the proton skin is not so pronounced an effect as the neutron skin. However, some evidence for increasing of the proton-skin thickness in nuclei towards the proton drip-line is provided both by theoretical and by experimental studies .
In Fig. 2 we present results obtained within the fully self-consistent RRPA model based on Dirac-Hartree ground state with DD-ME1 effective interaction. In the left panel we display the isovector dipole strength distributions for <sup>48</sup>Ni and <sup>56</sup>Ni. In the region of the isovector giant dipole resonance (IVGDR), the difference between the two distributions is very small and mainly consists of small fluctuations of the IVGDR tail. In agreement with the mass dependence of the giant resonance, the distribution of <sup>48</sup>Ni is only slightly pushed to higher energies from the one for <sup>56</sup>Ni. However, in the low-energy region, the E1 strength distributions are rather different: whereas there are no low-lying states for <sup>56</sup>Ni, proton-rich <sup>48</sup>Ni is characterized by the appearance of a pronounced amount of low-lying transition strength. In order to clarify the origin of this strength, in the right panel of Fig. 2 we show the neutron and proton transition densities for two characteristic peaks: the low-lying state at 7.72 MeV, and the giant resonance state at 18.71 MeV. The transition densities of the latter display the dynamics of IVGDR: collective oscillation of neutrons against protons. For the low-energy peak, however, the proton and neutron transition densities are in phase in the nuclear interior, whereas beyond the surface region there are no contributions from the neutrons and proton transition density dominates. This type of nuclear dynamics is characteristic for the proton PDR, where loosely bound protons oscillate against the rest of the nucleons , in analogy to the neutron PDR in neutron-rich nuclei. The RRPA amplitudes for the E1 state at 7.72 MeV reveal the structure of the proton PDR in detail: the main contribution to the strength of the peak consists of transitions from the proton 1f<sub>7/2</sub> state, which is located at 0.11 MeV and weakly bound partly due to the presence of the Coulomb barrier. The contributions from other transitions are at least an order of magnitude smaller. Therefore, the appearance of the low-lying proton PDR strength is directly related to the pronounced proton density distributions from Fig. 1. Protons from the same loosely-bound orbit contribute to the exotic nuclear structure of the ground state, and to the excitation phenomena of the proton PDR. The collectivity of the proton PDR peak considerably increases in open shell-nuclei, due to the increased number of two-quasiparticle configurations composed from many states around the Fermi surface which are, in that case, partially occupied .
The properties of the low-lying dipole transition strength are strongly sensitive on the proton excess. In comparison of the cases of <sup>46</sup>Fe (from Ref. ) and <sup>48</sup>Ni (DD-ME1 interaction is employed in both cases), one can see that the GDR peak energy only weakly changes (0.2 MeV) with addition of two more protons. The peak energies of the proton PDR mode, however, lowers from 9.4 MeV towards 7.7 MeV for <sup>46</sup>Fe and <sup>48</sup>Ni, respectively. Obviously, the properties of the low-lying dipole transition strength are more sensitive to the variations of the nucleon excess than GDR. This type of behaviour indicates the nature of the proton PDR mode: as the number of protons increases, the oscillations of loosely-bound protons acquire lower frequencies due to their weaker binding.
By using the Skyrme Hartree-Fock+CRPA model with BSk8 effective interaction, we repeat the same study of the isovector dipole transition strength for <sup>48</sup>Ni and <sup>56</sup>Ni isotopes. In the left panel of Fig. 3, we notice that for the two Ni isotopes the strength distributions are not very different for E$`>`$10 MeV. In agreement with the RRPA results, the CRPA low-lying transition strength for <sup>48</sup>Ni is strongly enhanced in comparison with the <sup>56</sup>Ni case. In the right panel of Fig. 3 we plot the transition densities for a low energy peak at 9.72 MeV and the IVGDR at 20.28 MeV. The displayed CRPA transition densities are in full agreement with the RRPA calculations, i.e. the high-energy peak corresponds to the collective IVGDR where protons oscillate versus neutrons, while the low-lying transitions reveal the nature of the proton PDR.
Within the CRPA model, we do not obtain only one pronounced low-energy PDR peak for $`{}_{}{}^{N}48`$, but rather a smooth continuum, slightly structured around 9-10 MeV. One of the structures is the 9.72 MeV peak for which the transition density was plotted in Fig. 3. This continuum is not an artifact of the small smearing parameter used, $`\mathrm{\Gamma }=0.05`$ MeV. The particle threshold energy is $`E_{\mathrm{th}}=3.3`$ MeV. We have evaluated the proton- and neutron-transition densities corresponding to various values of excitation energy below $`E=10`$ MeV, and verified that they are characterized by transitions of a similar nature as the 9.72 MeV peak. The above discussion remains valid qualitatively when other Skyrme interactions are used. More numerical results are presented below.
In Table 1 we compare the global properties of the proton PDR for <sup>48</sup>Ni, obtained using the RRPA model with DD-ME1 and DD-ME2 effective interactions, and the CRPA model with Skyrme interactions BSk8, BSk2, MSk7 and SkM\*. We have calculated the summed low-lying strength $`m_0`$, the energy-weighted strength $`m_1`$ and the centroid energy $`m_1/m_0`$ for excitation energies below 10 MeV, where $`m_k`$ corresponds to the $`k`$th moment of the strength distribution. Since the CRPA strength distribution is continuous in this region, the choice of the cut-off value of PDR region at 10 MeV is somewhat arbitrary and the definition of the centroid energy becomes ambiguous. For this reason, the respective results are placed inside brackets. In addition, we list the relative amount of the low-lying strength $`m_1`$ with respect to the classical Thomas-Reiche-Kuhn sum rule TRK$`=14.9(NZ/A)e^2fm^2`$ MeV. The centroid energies of the proton PDR for <sup>48</sup>Ni are obtained around 8 MeV. The two different models are in fair agreement for various interactions, exhausting for the proton PDR from $``$ 0.9$`\%`$ (BSk8) towards $``$ 1.6$`\%`$ (DD-ME1, MSk7) of the classical TRK sum rule.
We found that that BSk8 and SkM\* gave the lowest values of low-lying PDR strength. In general, the Skyrme-type interactions with the lower value of $`m^{}/m`$, result in weaker low-energy transition strength. The amount of the low-lying strength is directly related to the energy $`e_f`$ of the least bound proton single-particle state $`1f_{7/2}`$ and the particle threshold energy $`E_{\mathrm{th}}`$. For <sup>48</sup>Ni, these quantities are as follows: $`e_f=0.88`$ MeV and $`E_{\mathrm{th}}=4.2`$ MeV for SkM\*, $`e_f=0.03`$ MeV and $`E_{\mathrm{th}}=3.3`$ MeV for BSk8, $`e_f=0.26`$ MeV and $`E_{\mathrm{th}}=3.0`$ MeV for BSk2, and $`e_f=0.35`$ MeV and $`E_{\mathrm{th}}=3.0`$ MeV for MSk7. As the energy of $`1f_{7/2}`$ proton state changes to higher values, the corresponding low-lying E1 strength is more enhanced. It is surprising that the Skyrme interaction whose effective-mass properties differ the most from the ones of the relativistic forces, namely MSk7, gives the best agreement with the RRPA results. The opposite holds for BSk8 and SkM\*, whose properties differ the least from the ones of the relativistic forces. We notice, moreover, that the BSk2 corresponds to almost as large an effective mass $`m^{}/m`$ as the MSk7, and practically the same $`E_{\mathrm{th}}`$, yet its predictions for the PDR are closer to the BSk8. Its isovector effective mass, though, is low and almost the same as BSk8. Within the present CRPA model, a Skyrme interaction cannot provide as much low-lying strength as the RRPA model, unless a high value of isovector effective mass is used. It is an interesting trend, but it cannot be established only by the study of a single nucleus. Finally, there does not seem to exist a correlation between the degrees of quantitative agreement of the RRPA and CRPA models on the properties of the PDR on one hand, and the $`e_f`$ value on the other. The absolute difference between the value of $`e_f`$ corresponding to the Skyrme force MSk7 (best agreement with RRPA on PDR strength) and the RRPA result, 0.11 MeV, is almost as large as in the case of the SkM\* (the worst agreement).
In the last row of Table 1, we show an estimate of the PDR width, given by the mean deviation of the strength distribution up to 10 MeV, $`\sigma =\sqrt{(m_2/m_0)(m_1/m_0)^2}`$. Different Skyrme interactions result in similar widths, $`\sigma `$ 1.4 MeV. On the other side, the width evaluated from the RRPA discrete strength distribution with the DD-ME1 interaction with a smaller effective mass, results in $`\sigma =`$ 0.77 MeV. In the RRPA case it provides only a measure of the fragmentation of the low-lying strength, and therefore, it is smaller in comparison with the CRPA width which includes the contributions of the continuum.
In conclusion, we have studied the low-lying dipole response of a representative case of a proton drip-line nucleus, namely the doubly magic <sup>48</sup>Ni. We have employed two different microscopic models, Dirac-Hartree+RRPA, and Skyrme-Hartree-Fock+CRPA, with various effective interactions. Within the latter approach, a proper treatment of the particle continuum is included, enabling us to study the relevance of the continuum for the low-lying strength. The comparison of E1 strength distributions for <sup>48</sup>Ni and <sup>56</sup>Ni, and transition densities, show that the low-lying dipole strength is a fundamental property of the proton-rich nuclei, and it corresponds to the proton pygmy dipole resonance, where loosely bound protons vibrate against the approximately isospin-saturated proton-neutron core. The coupling to the particle continuum results in an enhanced width of the proton PDR mode, estimated around 1.4 MeV. However, it remains unresolved why the best agreement on the proton PDR properties, between the relativistic and nonrelativistic models, is obtained for DD-ME1 and MSk7 interactions which have quite different properties. Whereas DD-ME1 represents an advanced density-dependent interaction appropriate for the studies of both stable and exotic nuclei, MSk7 has large isoscalar and isovector effective mass, and its properties were recently improved in BSk series. A consistent comparison would necessitate a proper inclusion of the continuum in the RRPA model, and on the other side, a fully self-consistent CRPA model without neglecting of terms in the residual interaction. Nevertheless, we have shown that the two different models agree fairly well on the global properties of the proton PDR, and provide a clear theoretical picture for its underlying nature. We hope that the future experimental studies towards the proton drip-line will provide evidence for this exotic mode.
ACKNOWLEDGEMENTS
This work has been supported by the Deutsche Forschungsgemeinschaft (DFG) under contract SFB 634.
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# Shape of the ๐_๐โข(๐๐๐) in ๐ธโข๐ธโ๐
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## I Introduction
Recently, the Belle Collaboration succeeded in observing a clear manifestation of the $`f_0(980)`$ resonance in the reaction $`\gamma \gamma \pi ^+\pi ^{}`$ MBC . This has been made possible owing to the huge statistics and good energy resolution. Evidence for the $`f_0(980)`$ production in $`\gamma \gamma `$ collisions obtained previously by the Mark II E1 , CELLO E2 , ALEPH E3 , Crystal Ball E4 ; E5 , and JADE E6 Collaborations was essentially less conclusive MBC ; BP . The Belle data MBC corresponding to the $`f_0(980)`$ resonance region are shown in Fig. 1. Figure 1(a) shows the distribution of $`e^+e^{}e^+e^{}\pi ^+\pi ^{}`$ and $`e^+e^{}e^+e^{}\mu ^+\mu ^{}`$ events, $`\mathrm{\Delta }N`$, in the invariant mass of the $`\pi ^+\pi ^{}`$ and $`\mu ^+\mu ^{}`$ systems, $`m`$, scanned with a 5-MeV-wide step. A distinct peak due to the $`f_0(980)`$ resonance production in the $`\gamma \gamma \pi ^+\pi ^{}`$ channel can be seen in this plot. The peak position $`m_{f_0}=981.2\pm 0.5`$ MeV and its total width $`\mathrm{\Gamma }=21.7\pm 2.1`$ MeV were determined in Ref. MBC by fitting the $`m`$ dependence of $`\mathrm{\Delta }N`$ in the $`f_0(980)`$ region to the incoherent sum of the resonance and background contributions:
$$\mathrm{\Delta }N=\frac{4.8\pi A\mathrm{\Gamma }}{(m_{f_0}^2m^2)^2+m_{f_0}^2\mathrm{\Gamma }^2}+\mathrm{\Delta }N_{BG},$$
(1)
where $`\mathrm{\Delta }N_{BG}=C_0+C_1m+C_2m^2`$ represents a smooth background and the parameter $`A`$ is the production of the two-photon width $`\mathrm{\Gamma }_{f_0\gamma \gamma }`$, branching ratio $`B(f_0\pi ^+\pi ^{})`$, and known factors connected with the detection efficiency and the setup luminosity MBC . The Belle Collaboration plans to report the information on $`\mathrm{\Gamma }_{f_0\gamma \gamma }`$ after the investigation of the systematic error sources MBC . The Belle data for the $`\gamma \gamma \pi ^+\pi ^{}`$ reaction cross section, $`\sigma (\gamma \gamma \pi ^+\pi ^{})`$, in the region $`|\mathrm{cos}\theta ^{}|<0.6`$, where $`\theta ^{}`$ is the center-of-mass scattering angle of pion, with indication only statistical errors are shown in Fig. 1(b). The comparison of these data with those of the previous Mark II E1 and CELLO E2 experiments is presented in Fig. 1(c).
It should be noted that, according the Belle data, the $`f_0(980)`$ resonance manifests itself in the $`\gamma \gamma \pi ^+\pi ^{}`$ reaction cross section rather as a jump, or a step, with a width of about 15 MeV and a height of about 11 nb, than as a clear peak; see Fig. 1(b). In connection with this โobservationโ, as well as bearing in mind some theoretical reasons (see below), we would like to draw attention, especially of the experimentalists, to the fact that Eq. (1) cannot be used to determine the physical characteristics of the $`f_0(980)`$ resonance from the data on the reaction $`\gamma \gamma \pi ^+\pi ^{}`$. First, due to the proximity of the $`f_0(980)`$ resonance to the $`K\overline{K}`$ thresholds and its strong coupling to the $`K\overline{K}`$ channels, the propagator of the form $`1/(m_{f_0}^2m^2im_{f_0}\mathrm{\Gamma })`$, with the total width independent of $`m`$, cannot be applied in principle to the description of the $`f_0(980)`$ resonance shape. Second, owing to the $`K^+K^{}`$ loop mechanism, the two-photon width of the $`f_0(980)`$ resonance is a sharply varying function of $`m`$ just in the $`f_0(980)`$ peak region. Therefore, it cannot be approximated by a constant. And third, one cannot but take into account that the $`f_0(980)`$ resonance strongly interferes with the considerable $`S`$ wave background contributions in the $`\gamma \gamma \pi ^+\pi ^{}`$ reaction cross section.
In the present paper we analyze in detail the role of basic dynamical mechanisms of the reaction $`\gamma \gamma \pi ^+\pi ^{}`$ in the 1 GeV region and elucidate a possible form of the $`f_0(980)`$ resonance manifestation in this channel. In so doing, we tried to use sufficiently simple, but adequate to the highly not simple physical situation, formulae free of unknown parameters.
The paper is organized as follows. In Sec. II, the $`K^+K^{}`$ loop mechanism of the $`f_0(980)\gamma \gamma `$ decay is discussed. This mechanism not only ensures the appreciable distortion of the
$`f_0(980)`$ resonance shape in the reaction $`\gamma \gamma \pi ^+\pi ^{}`$ but also automatically yields a reasonable estimate for the absolute magnitude of the $`f_0(980)`$ production cross section in this channel with the values of the $`f_0(980)`$ resonance parameters compatible with the data on the other reactions. Thus, we get good reasons to consider the $`K^+K^{}`$ loop mechanism as a major one of the $`f_0(980)`$ production in $`\gamma \gamma `$ collisions. In Sec. III, a simplest dynamical model for the $`S`$ wave amplitude of the reaction $`\gamma \gamma \pi ^+\pi ^{}`$ in the 1 GeV region is examined and the character of the interference between the background and $`f_0(980)`$ resonance contributions, and thereby a possible resulting shape of the $`f_0(980)`$ in the $`\gamma \gamma \pi ^+\pi ^{}`$ channel, is clarified. Most of all the shape obtained resembles a step. This conclusion is supported by the Belle data. A possible manifestation of the $`f_0(980)`$ resonance in the $`\gamma \gamma \pi ^0\pi ^0`$ channel is briefly discussed. The general remarks and conclusions based on the results of our analysis are formulated in Sec. IV.
## II $`๐ฒ^\mathbf{+}๐ฒ^{\mathbf{}}`$ loop mechanism of the $`๐_\mathrm{๐}\mathbf{(}\mathrm{๐๐๐}\mathbf{)}\mathbf{}๐ธ๐ธ`$ decay
Perhaps, none of the known hadronic resonances can โboastโ of such a variety of the forms of its own manifestation that the $`f_0(980)`$ resonance possesses. The $`f_0(980)`$ shape in the two-pion decay channel depends in a crucial way on the reaction and varies from dips to peaks. In many respects this is due to the fact that background contributions, usually accompanying the $`f_0(980)`$ resonance, strongly change in passing from reaction to reaction, which leads in its turn to the change of the interference patterns in the resonance region. But, the even more impressive thing is that there exist reactions in which the $`f_0(980)`$ production amplitude itself sharply changes just in the $`f_0(980)`$ peak region. First of all such a phenomenon takes place in the radiative decays $`\varphi f_0(980)\gamma \pi \pi \gamma `$ AI ; AG1 ; AG2 . As predicted theoretically in Ref. AI and confirmed in the experiments performed at Novosibirsk E7 ; E8 and Frascati E9 , these decays are determined by the $`K^+K^{}`$ loop mechanism of the $`f_0(980)`$ production, $`\varphi K^+K^{}\gamma f_0(980)\gamma \pi \pi \gamma `$, the amplitude of which is large, owing to the strong coupling of the $`f_0(980)`$ to $`K\overline{K}`$, and changes very rapidly as a function of two-pion invariant mass near the $`K^+K^{}`$ threshold. The related decay $`\varphi a_0^0(980)\gamma \eta \pi ^0\gamma `$ is also determined by the $`K^+K^{}`$ loop mechanism AI ; AG1 ; E8 ; E10 ; E11 ; AK1 . It should be also recalled that the important role of this mechanism in the process $`\gamma \gamma a_0^0(980)\eta \pi ^0`$ was shown long ago in Ref. AS1 . The above mentioned manifestations of the $`K^+K^{}`$ loop mechanism present important physical evidences in favor of the four-quark ($`q^2\overline{q}^2`$) nature of the $`f_0(980)`$ and $`a_0^0(980)`$ resonances AI ; AS1 ; A1 ; A11 .
The presentation of high quality data from the Belle Collaboration on the reaction $`\gamma \gamma \pi ^+\pi ^{}`$ provides good reason to discuss in detail the role of the $`K^+K^{}`$ loop mechanism of the $`f_0(980)`$ resonance production in $`\gamma \gamma `$ collisions. As we shall show, it is very important, if not determining at all. Note that the process $`\gamma \gamma K\overline{K}f_0(980)\pi \pi `$ seems to be first mentioned in Ref. ADS1 .
Thus, let us consider the shape of the $`f_0(980)`$ resonance produced in the reaction $`\gamma \gamma \pi ^+\pi ^{}`$ via the $`K^+K^{}`$ loop mechanism. This mechanism corresponds to the following sequence of transitions. At first, there takes place the formation of the $`K^+K^{}`$ pair in $`\gamma \gamma `$ collisions, with the amplitude which near the $`K^+K^{}`$ threshold can be taken in the Born approximation. Then, the $`K^+K^{}`$ system turns into the $`f_0(980)`$ resonance decaying further into $`\pi ^+\pi ^{}`$. According this prescription, the corresponding resonant contribution to the $`\gamma \gamma \pi ^+\pi ^{}`$ reaction cross section can be written as
$$\sigma _{f_0}(\gamma \gamma \pi ^+\pi ^{})=\frac{8\pi }{m^2}\frac{m\mathrm{\Gamma }_{f_0K^+K^{}\gamma \gamma }^{Born}(m)m\mathrm{\Gamma }_{f_0\pi ^+\pi ^{}}(m)}{|D_{f_0}(m)|^2}.$$
(2)
Here
$$\mathrm{\Gamma }_{f_0K^+K^{}\gamma \gamma }^{Born}(m)=\frac{1}{16\pi m}|M_{f_0K^+K^{}\gamma \gamma }^{Born}(m)|^2=\frac{\alpha ^2}{4\pi ^2}|I_{K^+K^{}}(m)|^2\frac{g_{f_0K^+K^{}}^2}{16\pi m}$$
(3)
is the width of the $`f_0(980)\gamma \gamma `$ decay due to the Born $`K^+K^{}`$ loop mechanism, where $`\alpha =e^2/4\pi 1/137`$ and the function $`I_{K^+K^{}}(m)`$ is AS1
$`I_{K^+K^{}}(m)=\{\begin{array}{cc}\frac{m_{K^+}^2}{m^2}\left[\pi +i\mathrm{ln}\frac{1+\rho _{K^+}(m)}{1\rho _{K^+}(m)}\right]^21,\hfill & m2m_{K^+},\hfill \\ \frac{m_{K^+}^2}{m^2}[\pi 2\mathrm{arctan}|\rho _{K^+}(m)|]^21,\hfill & 0m2m_{K^+}.\hfill \end{array}`$ (6)
The propagator of the $`f_0(980)`$ resonance with a mass $`m_{f_0}`$ appearing in Eq. (2) has the form ADS2
$$\frac{1}{D_{f_0}(m)}=\frac{1}{m_{f_0}^2m^2+_{a\overline{a}}[\text{Re}\mathrm{\Pi }_{f_0}^{a\overline{a}}(m_{f_0})\mathrm{\Pi }_{f_0}^{a\overline{a}}(m)]},$$
(7)
where $`\mathrm{\Pi }_{f_0}^{a\overline{a}}(m)`$ is the polarization operator of the $`f_0(980)`$ resonance corresponding to the contribution of the $`a\overline{a}`$ intermediate state ($`a\overline{a}=\pi ^+\pi ^{},\pi ^0\pi ^0,K^+K^{},K^0\overline{K}^0`$). For $`m2m_a`$,
$$\mathrm{\Pi }_{f_0}^{a\overline{a}}(m)=\xi _{a\overline{a}}\frac{g_{f_0a\overline{a}}^2}{16\pi }\rho _a(m)\left[i\frac{1}{\pi }\mathrm{ln}\frac{1+\rho _a(m)}{1\rho _a(m)}\right],$$
(8)
$`\rho _a(m)=(14m_a^2/m^2)^{1/2}`$ \[if $`0m2m_a`$, then $`\rho _a(m)i|\rho _a(m)|`$\], $`\mathrm{\Gamma }_{f_0a\overline{a}}(m)=\text{Im}\mathrm{\Pi }_{f_0}^{a\overline{a}}(m)/m=\xi _{a\overline{a}}g_{f_0a\overline{a}}^2\rho _a(m)/16\pi m`$ is the width of the $`f_0(980)a\overline{a}`$ decay, here $`\xi _{a\overline{a}}=1`$, if $`a\overline{a}`$, and $`\xi _{a\overline{a}}=1/2`$, if $`a=\overline{a}`$, and $`g_{f_0\pi ^+\pi ^{}}^2=g_{f_0\pi ^0\pi ^0}^2=2g_{f_0\pi \pi }^2/3`$, $`g_{f_0K^+K^{}}^2=g_{f_0K^0\overline{K}^0}^2=g_{f_0K\overline{K}}^2/2`$, where $`g_{f_0\pi \pi }`$ and $`g_{f_0K\overline{K}}`$ are the coupling constants of the $`f_0(980)`$ to the $`\pi \pi `$ and $`K\overline{K}`$ channels, respectively. Since we are interested in the $`m`$ region near the $`K\overline{K}`$ thresholds, we take into account the $`K^+`$ and $`K^0`$ meson mass difference.
As for the $`f_0(980)`$ resonance parameters, the available data, together with various model parametrizations, allow wide intervals for their possible values; for example, $`m_{f_0}(0.9650.99)`$ GeV, $`g_{f_0\pi \pi }^2/16\pi (0.0650.3)`$ GeV<sup>2</sup>, and $`g_{f_0K\overline{K}}^2/16\pi (0.31.6)`$ GeV<sup>2</sup>, with the preferred coupling-constant-squared ratio $`R=g_{f_0K\overline{K}}^2/g_{f_0\pi \pi }^246`$, are quite compatible with the data on most reactions of the $`f_0(980)`$ production AI ; AG1 ; E7 ; E8 ; E9 ; Fla ; MOS ; ADS2 ; AS2 ; CF ; Abl ; PDG . For further estimates and illustrations of the role of the $`K^+K^{}`$ loop mechanism, we use, in fact, all the range of possible values of the $`f_0(980)`$ parameters.
Let us now discuss two most important features of the $`K^+K^{}`$ loop mechanism which immediately follow from the above formulae. First, as is seen from Fig. 2(a), the factor $`m\mathrm{\Gamma }_{f_0K^+K^{}\gamma \gamma }^{Born}(m)`$ in Eq. (2) sharply decreases just below the $`K^+K^{}`$ threshold, i.e., directly in the $`f_0(980)`$ resonance region. For instance, it falls relative to the maximum at $`m=2m_{K^+}0.9873`$ GeV by a factor of 1.69, 2.23, 2.75, 3.27, and 6.33 at $`m=`$0.98, 0.97, 0.96, 0.95, and 0.9 GeV, respectively. Such a behavior of $`m\mathrm{\Gamma }_{f_0K^+K^{}\gamma \gamma }^{Born}(m)`$ strongly suppresses the left wing of the $`f_0(980)`$ resonance peak defined by $`1/|D_{f_0}(m)|^2`$ in Eq. (2). Second, from Eqs. (2) and (3) it follows that for the $`K^+K^{}`$ loop mechanism the magnitude of $`\sigma _{f_0}(\gamma \gamma \pi ^+\pi ^{})`$ near the maximum, located between $`m_{f_0}`$ and $`2m_{K^+}`$, is controlled mainly by the parameter $`R=g_{f_0K\overline{K}}^2/g_{f_0\pi \pi }^2`$ and the value of the function $`|I_{K^+K^{}}(m)|^2`$. For example, if $`m_{f_0}<2m_{K^+}`$, then, at $`m=m_{f_0}`$, $`\sigma _{f_0}(\gamma \gamma \pi ^+\pi ^{})=\alpha ^2R|I_{K^+K^{}}(m_{f_0})|^2/[\pi m_{f_0}^2\rho _\pi (m_{f_0})]`$. Furthermore, at fixed $`m_{f_0}`$ and $`R`$, the $`f_0(980)`$ resonance shape in $`\sigma _{f_0}(\gamma \gamma \pi ^+\pi ^{})`$ is very insensitive to the absolute values of the coupling constants $`g_{f_0\pi \pi }^2/16\pi `$ and $`g_{f_0K\overline{K}}^2/16\pi `$. As an illustration we represent in Figs. 2(b) and 2(c) the cross section $`\sigma _{f_0}(\gamma \gamma \pi ^+\pi ^{})`$ for four different sets of the $`f_0(980)`$ resonance parameters: $`m_{f_0}=0.98`$ GeV, $`R=4`$, $`g_{f_0K\overline{K}}^2/16\pi =0.4`$ GeV<sup>2</sup>, and 1.2 GeV<sup>2</sup> (sets A and B), and $`m_{f_0}=0.97`$ GeV, $`R=5.33`$, $`g_{f_0K\overline{K}}^2/16\pi =0.533`$ GeV<sup>2</sup>, and 1.6 GeV<sup>2</sup> (sets C and D). For sets A and D the cross section smoothed with a Gaussian mass distribution with the dispersion of 5 MeV (which we have chosen to be equal to the $`m`$ step in the Belle experiment) is shown in these figures for completeness. Multiplying the resulting cross section values by a factor 0.6 \[in accordance with the fact that the data for $`\sigma (\gamma \gamma \pi ^+\pi ^{})`$ correspond to the region $`|\mathrm{cos}\theta ^{}|<0.6`$\], we obtain that, owing to the $`K^+K^{}`$ loop mechanism, the $`f_0(980)`$ resonance can manifest itself in the measured $`\gamma \gamma \pi ^+\pi ^{}`$ reaction cross section at the level of about $`15.517.5`$ nb at the maximum. As is clear from Fig. 1(b), this estimate for the scale of the enhancement due to the $`f_0(980)`$ resonance contribution to the $`\pi ^+\pi ^{}`$ production cross section is in reasonable (if not excellent) agreement with the Belle data. Thus, we conclude that the $`K^+K^{}`$ loop mechanism, which actually results from the unitarity condition, can be primarily responsible for the $`f_0(980)`$ resonance coupling to photons.
It is clear that in there is no sense in speaking about a two-photon width at the resonance point if the two-photon decay width of the resonance varies rapidly within its hadronic width, see Fig. 2(a). For the $`K^+K^{}`$ loop mechanism, it is of interest to evaluate the $`f_0(980)\gamma \gamma `$ width averaged by the resonance mass distribution in the $`\pi \pi `$ channel, $`\mathrm{\Gamma }_{f_0K^+K^{}\gamma \gamma }^{Born}_{\pi \pi }`$ AS1 . By definition,
$$\mathrm{\Gamma }_{f_0K^+K^{}\gamma \gamma }^{Born}_{\pi \pi }=\underset{m_1}{\overset{m_2}{}}\mathrm{\Gamma }_{f_0K^+K^{}\gamma \gamma }^{Born}(m)\frac{3}{2}\left[\frac{m\mathrm{\Gamma }_{f_0\pi ^+\pi ^{}}(m)}{\pi |D_{f_0}(m)|^2}\right]2m๐m=\frac{3}{2}\underset{m_1}{\overset{m_2}{}}\frac{m^2}{4\pi ^2}\sigma _{f_0}(\gamma \gamma \pi ^+\pi ^{})๐m$$
(9)
\[see also Eq.(2)\]. This averaged width can serve as an adequate, working characteristic of the $`f_0(980)`$ coupling to $`\gamma \gamma `$. Substituting $`\sigma _{f_0}(\gamma \gamma \pi ^+\pi ^{})`$, shown in Figs. 2(b) and 2(c) in Eq. (7) and integrating, for example, over two $`m`$ regions 0.93 GeV $`m`$ 1.03 GeV and $`2m_\pi m<\mathrm{}`$, we obtain $`\mathrm{\Gamma }_{f_0K^+K^{}\gamma \gamma }^{Born}_{\pi \pi }0.114`$ keV and 0.191 keV, respectively, for set A, 0.132 keV and 0.351 keV for set B, 0.129 keV and 0.211 keV for set C, and 0.152 keV and 0.377 keV for set D. Defining also $`\mathrm{\Gamma }_{f_0K^+K^{}\gamma \gamma }^{Born}_{K\overline{K}}`$ in a similar way, we find that the total averaged width of the $`f_0\gamma \gamma `$ decay $`\mathrm{\Gamma }_{f_0K^+K^{}\gamma \gamma }^{Born}=\mathrm{\Gamma }_{f_0K^+K^{}\gamma \gamma }^{Born}_{\pi \pi }+\mathrm{\Gamma }_{f_0K^+K^{}\gamma \gamma }^{Born}_{K\overline{K}}0.14`$ keV and 0.359 keV for the two above mentioned integration regions, respectively, for set A, 0.164 keV and 0.884 keV for set B, 0.158 keV and 0.439 keV for set C, and 0.189 keV and 1.094 keV for set D. It is worth to point out for comparison that $`\mathrm{\Gamma }_{f_0K^+K^{}\gamma \gamma }^{Born}(m)`$ at the maximum, i.e., $`\mathrm{\Gamma }_{f_0K^+K^{}\gamma \gamma }^{Born}(2m_{K^+})`$, is approximately equal to 0.589, 1.766, 0.785, and 2.355 keV for $`g_{f_0K\overline{K}}^2/16\pi `$ from sets A, B, C, and D, respectively.
Certainly, the real situation in the $`\gamma \gamma \pi ^+\pi ^{}`$ channel is more complicated because the $`f_0(980)`$ resonance in this channel is by no means a solitary one. It is accompanied by the considerable coherent background, and therefore the interference effects are of great significance. Their role will be analyzed in detail below in Sec. III.
We wish to conclude this section with a general remark concerning the $`f_0(980)`$ resonance propagator, see Eq. (5), which we utilize throughout here. As is shown in Ref. AK2 , this propagator rigorously satisfies the Kรคllen-Lehmann representation, i.e., it possesses the analytic properties required in field theory. Thus, the resonance mass distributions calculated with the use of this propagator are automatically normalized to the corresponding branching ratios of the $`f_0(980)a\overline{a}`$ decays, the sum of which is exactly equal to unit, i.e.,
$$\underset{2m_a}{\overset{\mathrm{}}{}}\left[\frac{m\mathrm{\Gamma }_{f_0a\overline{a}}(m)}{\pi |D_{f_0}(m)|^2}\right]2m๐m=B(f_0(980)a\overline{a}),\underset{a\overline{a}}{}B(f_0(980)a\overline{a})=1.$$
## III $`๐บ`$ wave in the reaction $`๐ธ๐ธ\mathbf{}๐
^\mathbf{+}๐
^{\mathbf{}}`$ near 1 GeV
Let us consider a simplest dynamical model for the $`S`$ wave amplitude of the reaction $`\gamma \gamma \pi ^+\pi ^{}`$ in the 1 GeV region. There are no arbitrary, free parameters in this model (that is the parameters which would be unknown from other reactions), and within its framework the character of the interference between the background and $`f_0(980)`$ resonance contributions, and thus a possible resulting $`f_0(980)`$ shape in the $`\gamma \gamma \pi ^+\pi ^{}`$ channel, will be fully elucidated. The results obtained in this way will be useful, in particular, as to estimate the potentialities and โpriceโ of the more complicated model constructions.
Using the conventional normalization, we write the $`S`$ wave cross section of the reaction $`\gamma \gamma \pi ^+\pi ^{}`$, together with the corresponding amplitude $`A_S(m)`$, in the form:
$$\sigma _S(\gamma \gamma \pi ^+\pi ^{})=\frac{\rho _\pi (m)}{32\pi m^2}|A_S(m)|^2,$$
(10)
$`A_S(m)=M_{\gamma \gamma \pi ^+\pi ^{}}^{Born}(m)+8\alpha I_{\pi ^+\pi ^{}}(m)T_{\pi ^+\pi ^{}\pi ^+\pi ^{}}(m)+8\alpha I_{K^+K^{}}(m)T_{K^+K^{}\pi ^+\pi ^{}}(m).`$ (11)
Here
$$M_{\gamma \gamma \pi ^+\pi ^{}}^{Born}(m)=\frac{16\pi \alpha m_\pi ^2}{m^2\rho _\pi (m)}\mathrm{ln}\frac{1+\rho _\pi (m)}{1\rho _\pi (m)}=\frac{8\alpha }{\rho _\pi (m)}\text{Im}I_{\pi ^+\pi ^{}}(m)$$
(12)
is the $`S`$ wave Born amplitude of the process $`\gamma \gamma \pi ^+\pi ^{}`$, the function $`I_{\pi ^+\pi ^{}}(m)`$ results from Eq. (4) by replacing $`m_{K^+}`$ and $`\rho _{K^+}(m)`$ by $`m_\pi `$ and $`\rho _\pi (m)`$, respectively, and $`T_{\pi ^+\pi ^{}\pi ^+\pi ^{}}(m)`$ and $`T_{K^+K^{}\pi ^+\pi ^{}}(m)`$ are the $`S`$ wave amplitudes of hadronic reactions indicated in their subscripts. Hence it is obvious that the second and third terms on the right-hand side of Eq. (9) correspond to the contributions from the $`\gamma \gamma \pi ^+\pi ^{}`$ and $`\gamma \gamma K^+K^{}`$ Born amplitudes modified by the final state interactions. Such a structure of the amplitude $`A_S(m)`$ can be easily obtained within the framework of the field-theoretical model in which the electromagnetic Born amplitudes are the only primary sources of the $`\pi ^+\pi ^{}`$ and $`K^+K^{}`$ pairs, and the strong amplitudes, used for unitarization of the Born contributions, are constructed by summing up all the $`s`$ channel bubble diagrams. In so doing, the strong amplitudes can involve, in principle, any number of resonances plus background contributions to describe the relevant data on the phase shifts and inelasticities. The resulting strong and electromagnetic amplitudes in such a model are unitary. This model has a very old history ZGMZ ; Tir and up to now was successfully used, together with its dispersive modifications, as the effective tool in analyzing dynamics of electromagnetic and strong interaction processes, see for example AG2 ; BG ; Men ; Joh ; GRR ; AS3 ; AS4 .
The amplitude $`T_{\pi ^+\pi ^{}\pi ^+\pi ^{}}(m)`$ is related to the phase shifts $`\delta _0^I(m)`$ and inelasticities $`\eta _0^I(m)`$ of the $`S`$ wave $`\pi \pi `$ scattering amplitudes with definite isospin $`I=0,2`$ in the conventional way: $`T_{\pi ^+\pi ^{}\pi ^+\pi ^{}}(m)=\frac{2}{3}T_0^0(m)+\frac{1}{3}T_0^2(m)`$, where $`T_0^I(m)=\{\eta _0^I(m)\mathrm{exp}[2i\delta _0^I(m)]1\}/[2i\rho _\pi (m)]`$. As is well known, the only, strongly coupled $`S`$ wave channels in the 1 GeV region are the $`\pi \pi `$ and $`K\overline{K}`$ channels with $`I=0`$. Therefore we set $`\eta _0^2(m)=1`$ for all $`m`$ of interest and $`\eta _0^0(m)=1`$ for $`m<2m_{K^+}`$. Then, for $`m<2m_{K^+}`$, the amplitude $`A_S(m)`$ can be rewritten as, see Eqs. (9) and (10),
$`A_S(m)=e^{i\delta _0^0(m)}\{A_{S,0}(m)+A_{S,2}(m)\mathrm{cos}[\delta _0^2(m)\delta _0^0(m)]+iA_{S,2}(m)\mathrm{sin}[\delta _0^2(m)\delta _0^0(m)]\},`$ (13)
where the amplitudes $`A_{S,I}(m)`$ with $`I=0`$ and 2 have the form:
$`A_{S,0}(m)={\displaystyle \frac{2}{3}}M_{\gamma \gamma \pi ^+\pi ^{}}^{Born}(m)\mathrm{cos}\delta _0^0(m)`$
$`+[8\alpha /\rho _\pi (m)]\text{Re}[I_{\pi ^+\pi ^{}}(m)]{\displaystyle \frac{2}{3}}\mathrm{sin}\delta _0^0(m)+8\alpha I_{K^+K^{}}(m)T_{K^+K^{}\pi ^+\pi ^{}}(m)e^{i\delta _0^0(m)},`$ (14)
$`A_{S,2}(m)={\displaystyle \frac{1}{3}}\{M_{\gamma \gamma \pi ^+\pi ^{}}^{Born}(m)\mathrm{cos}\delta _0^2(m)+[8\alpha /\rho _\pi (m)]\text{Re}[I_{\pi ^+\pi ^{}}(m)]\mathrm{sin}\delta _0^2(m)\}.`$ (15)
Because for $`m<2m_{K^+}`$ the imaginary part of the function $`I_{K^+K^{}}(m)`$ vanishes, see Eq. (4), and the phase of the amplitude $`T_{K^+K^{}\pi ^+\pi ^{}}(m)`$ reduces to $`\delta _0^0(m)+n\pi `$ (where $`n=0`$ or 1) in accordance with unitarity, it is easy to see that all the terms in the amplitudes $`A_{S,0}(m)`$ and $`A_{S,2}(m)`$ are real.
Moreover, all of these terms have well definite signs. Begin with the amplitude $`T_{K^+K^{}\pi ^+\pi ^{}}(m)=T_{\pi ^+\pi ^{}K^+K^{}}(m)`$ in Eq. (12). Its sign, $`(1)^n`$, is known experimentally and it is positive MOS ; Mor ; Wet ; Est ; Coh . In terms of the $`f_0(980)`$ coupling constants this means that if we parametrize the amplitude $`T_{K^+K^{}\pi ^+\pi ^{}}(m)`$ in the 1 GeV region as MOS ; ADS2 ; AS2
$$T_{K^+K^{}\pi ^+\pi ^{}}(m)=\frac{g_{f_0\pi ^+\pi ^{}}g_{f_0K^+K^{}}}{16\pi D_{f_0}(m)}e^{i\delta _B(m)},$$
(16)
where $`\delta _B(m)`$ is a smooth and large phase (of about 90 for $`m1`$ GeV) of the elastic background in the $`I=0`$ $`S`$ wave $`\pi \pi `$ channel, then the production $`g_{f_0\pi ^+\pi ^{}}g_{f_0K^+K^{}}`$ is positive MOS ; Mor ; Est . Recall that with such a parametrization the $`\pi \pi `$ scattering amplitude $`T_0^0(m)`$ has the form Fla ; MOS ; ADS2 ; AS2 :
$$T_0^0(m)=\frac{\eta _0^0(m)e^{2i\delta _0^0(m)}1}{2i\rho _\pi (m)}=\frac{1}{\rho _\pi (m)}\left[\frac{e^{2i\delta _B(m)}1}{2i}+e^{2i\delta _B(m)}\frac{m\mathrm{\Gamma }_{f_0\pi \pi }(m)}{D_{f_0}(m)}\right],$$
(17)
and that the $`f_0(980)`$ resonance appears as a dip in $`|T_0^0(m)|`$ fn1 . Equations (14) and (15) will be used in the following. Thus, the last term on the right-hand side of Eq. (12) is positive because, according Eq. (4), $`I_{K^+K^{}}(m)>0`$ for $`0<m2m_{K^+}`$. Now we take into account the following circumstances. For $`0.85`$ GeV$`<m<2m_{K^+}`$, the phase shift $`\delta _0^0(m)`$ increases with $`m`$ from 90 to about 200 sharply flying up near the $`K^+K^{}`$ threshold; see, for example, Ref. Hya . In the same region of $`m`$, the phase shift $`\delta _0^2(m)`$ is of about $`\text{(19โ24)}^{}`$; see, for example, Ref. Hoo . Moreover, $`\text{Re}[I_{\pi ^+\pi ^{}}(m)]<0`$ for $`m>0.376`$ GeV. So, for $`0.85`$ GeV$`<m<2m_{K^+}`$, the first term on the right-hand side of Eq. (12) is negative, the second term is also negative at least up to 0.98 GeV, and it is small in magnitude for $`0.98`$ GeV$`<m<2m_{K^+}`$. Finally, it is easy to check that the amplitude $`A_{S,2}(m)\mathrm{cos}[\delta _0^2(m)\delta _0^0(m)]`$, see Eqs. (11) and (13), is also negative for $`0.85`$ GeV$`<m<2m_{K^+}`$.
Thus, one can conclude that, for $`m<2m_{K^+}`$, the sharply increasing with $`m`$, $`f_0(980)`$ production amplitude due to the $`K^+K^{}`$ loop mechanism has to interfere destructively with the accompanying background contributions in $`\sigma _S(\gamma \gamma \pi ^+\pi ^{})`$. Such an interference is able to suppress the left wing of the $`f_0(980)`$ resonance practically in full. A detailed illustration of the described general picture is presented in Fig. 3. In constructing the curves shown in this figure, we used set A for the values of the $`f_0(980)`$ resonance parameters and approximated the smooth phase shifts $`\delta _B(m)`$ and $`\delta _0^2(m)`$ by the following expressions: $`\delta _B(m)=\rho _\pi (m)_{n=0}^3q_\pi ^{2n}(m)a_{2n}=\rho _\pi (m)[0.1243+q_\pi ^2(m)16.32q_\pi ^4(m)73.50+q_\pi ^6(m)118.3]`$ and $`\delta _0^2(m)=q_\pi (m)b_0/[1+_{n=1}^3q_\pi ^{2n}(m)b_{2n}]=q_\pi (m)0.9098/[1+q_\pi ^2(m)2.629q_\pi ^4(m)13.19+q_\pi ^6(m)18.83]`$, where $`\delta _B(m)`$ and $`\delta _0^2(m)`$ in radians and $`q_\pi (m)=m\rho _\pi (m)/2`$ in units of GeV. Note that in this way we obtain the excellent description of the $`S`$ wave $`\pi \pi `$ scattering data Hya ; Hoo ; Ros at least in the $`m`$ region from $`2m_\pi `$ up to 1.2 GeV (for example, according to our fit, the $`S`$ wave $`\pi \pi `$ scattering length $`a_0^00.229/m_\pi `$). In Fig. 3(a) the solid curve shows that the real part of the amplitude $`\sqrt{\rho _\pi (m)/(32\pi m^2)}e^{i\delta _0^0(m)}A_S(m)`$, see Eqs. (8), (9), and (11) โ (13), vanishes at $`m0.967`$ GeV as a result of the compensation of the resonance and background contributions. As is seen from Fig. 3(b), this leads to a minimum in the cross section at the place of the left wing of the $`f_0(980)`$ resonance. As a whole, the resulting cross section $`\sigma _S(\gamma \gamma \pi ^+\pi ^{})`$ near 1 GeV in Fig. 3(b) resembles a step. Furthermore, we verified that sets B, C, and D for the $`f_0(980)`$ resonance parameters yield very similar results for $`A_S(m)`$ and $`\sigma _S(\gamma \gamma \pi ^+\pi ^{})`$.
To compare the model with the data pertaining to the partial solid angle, one must yet take into account the interference of the amplitude $`A_S(m)`$ with the higher partial waves. Usually, the measurements of the reaction $`\gamma \gamma \pi ^+\pi ^{}`$ are performed in the angular region $`|\mathrm{cos}\theta ^{}|<Z_0<1`$. The $`\gamma \gamma \pi ^+\pi ^{}`$ cross section is presented as the sum of the cross sections $`\sigma _{\lambda =0}(\gamma \gamma \pi ^+\pi ^{},|\mathrm{cos}\theta ^{}|<Z_0)`$ and $`\sigma _{|\lambda |=2}(\gamma \gamma \pi ^+\pi ^{},|\mathrm{cos}\theta ^{}|<Z_0)`$, where $`\lambda `$ is a photon helicity difference. In the $`Z_0<1`$ case, all the partial waves interfere between themselves in both cross sections. The cross section with $`|\lambda |=2`$ is dominated by the $`D`$ wave Born contribution and the well known $`f_2(1270)`$ resonance MBC ; E1 ; E2 ; Joh . The $`f_2(1270)`$ coupling to the $`\gamma \gamma `$ system in the $`\lambda =0`$ state is small E1 ; Joh . Therefore, we assume for estimate that in the 1 GeV region all the higher partial waves with $`\lambda =0`$ are defined simply by the corresponding $`\gamma \gamma \pi ^+\pi ^{}`$ Born amplitude. Then, $`\sigma _{\lambda =0}(\gamma \gamma \pi ^+\pi ^{},|\mathrm{cos}\theta ^{}|<Z_0)`$ can be written in the form:
$`\sigma _{\lambda =0}(\gamma \gamma \pi ^+\pi ^{},|\mathrm{cos}\theta ^{}|<Z_0)={\displaystyle \frac{\rho _\pi (m)}{32\pi m^2}}\{Z_0|\stackrel{~}{A}_S(m)|^2+C\text{Re}[\stackrel{~}{A}_S(m)].`$
$`.\times {\displaystyle \frac{1}{\rho _\pi (m)}}\mathrm{ln}{\displaystyle \frac{1+Z_0\rho _\pi (m)}{1Z_0\rho _\pi (m)}}+C^2[{\displaystyle \frac{Z_0/2}{1Z_0^2\rho _\pi ^2(m)}}+{\displaystyle \frac{1}{4\rho _\pi (m)}}\mathrm{ln}{\displaystyle \frac{1+Z_0\rho _\pi (m)}{1Z_0\rho _\pi (m)}}]\},`$ (18)
where the amplitude $`\stackrel{~}{A}_S(m)=A_S(m)M_{\gamma \gamma \pi ^+\pi ^{}}^{Born}(m)`$, see Eq. (9), and $`C=32\pi \alpha m_\pi ^2/m^2`$. With the use of Eq. (16), one can easily verify that for the typical value of $`Z_0=0.6`$ the higher partial wave influence, certainly, exists, but it is not too large.
Figures 4(a) and 4(b) illustrate the comparison of the model predictions for $`\sigma (\gamma \gamma \pi ^+\pi ^{},|\mathrm{cos}\theta ^{}|<0.6)=\sigma _{\lambda =0}(\gamma \gamma \pi ^+\pi ^{},|\mathrm{cos}\theta ^{}|<0.6)+\sigma _{|\lambda |=2}(\gamma \gamma \pi ^+\pi ^{},|\mathrm{cos}\theta ^{}|<0.6)`$ with the Belle data in the $`f_0(980)`$ resonance region. To obtain the curves in Fig. 4(a), we performed the simultaneous fit to the Belle data MBC and the well known $`S`$ wave $`\pi \pi `$ scattering data from Refs. Hya ; Ros . In so doing, we used Eqs. (9), (10), and (14) โ (16), the above mentioned expression for $`T_0^2(m)`$, and the approximation of the cross section with $`|\lambda |=2`$ by a linear function of $`m`$, $`C_1+C_2m`$ \[of course, this is a reasonable approximation only in the considered, narrow region of $`m`$ around the $`f_0(980)`$ resonance\]. The parameters obtained (set E) are $`m_{f_0}=0.9676`$ GeV, $`g_{f_0\pi \pi }^2/16\pi =0.07017`$ GeV<sup>2</sup>, $`g_{f_0K\overline{K}}^2/16\pi =0.3442`$ GeV<sup>2</sup> (R=4.9), $`C_1=57.69`$ nb, $`C_2=23.45`$ nb/GeV, and $`a_{2n=0,2,4,6}=0.1404`$, 17.17, $`80.17`$, 127.4, respectively. In order to illustrate that the Belle data tolerate, in fact, the wide range for the $`f_0(980)`$ coupling constant values, we fixed $`g_{f_0K\overline{K}}^2/16\pi =1.6`$ GeV<sup>2</sup> and performed once again the fit to the above mentioned data. For this case, the parameters obtained (set F) are $`m_{f_0}=0.968`$ GeV, $`g_{f_0\pi \pi }^2/16\pi =0.2438`$ GeV<sup>2</sup> (R=6.56), $`C_1=57.05`$ nb, $`C_2=24.62`$ nb/GeV, and $`a_{2n=0,2,4,6}=0.01903`$, 18.13, $`96.71`$, 173.2, respectively, and the resulting picture is shown in Fig. 4(b). As a whole, we obtain the quite satisfactory, qualitative agreement with the data in both the magnitude and shape of the $`f_0(980)`$ resonance manifestation. The strong difference of the $`f_0(980)`$ resonance shape in the $`\gamma \gamma \pi ^+\pi ^{}`$ reaction cross section from the shape of the solitary Breit-Wigner resonance is a result of fine interference effects between the different contributions. As we have made sure, the considered dynamical model provides a fairly good basis for understanding these effects. The model unambiguously points to the destructive interference pattern between the resonance and background contributions in the $`m`$ region below the $`K^+K^{}`$ threshold.
Now we discuss, in brief, a possible manifestation of the $`f_0(980)`$ resonance in the $`S`$ wave $`\gamma \gamma \pi ^0\pi ^0`$ reaction cross section. In the considered model we have:
$$\sigma _S(\gamma \gamma \pi ^0\pi ^0)=\frac{\rho _\pi (m)}{64\pi m^2}|B_S(m)|^2,$$
(19)
$$B_S(m)=8\alpha I_{\pi ^+\pi ^{}}(m)T_{\pi ^+\pi ^{}\pi ^0\pi ^0}(m)+8\alpha I_{K^+K^{}}(m)T_{K^+K^{}\pi ^0\pi ^0}(m),$$
(20)
where $`T_{\pi ^+\pi ^{}\pi ^0\pi ^0}(m)=\frac{2}{3}T_0^0(m)\frac{2}{3}T_0^2(m)`$ and $`T_{K^+K^{}\pi ^0\pi ^0}(m)=T_{K^+K^{}\pi ^+\pi ^{}}(m)`$. In comparison with the amplitude $`A_S(m)`$, see Eq. (9), the amplitude $`B_S(m)`$ does not contain the Born term and the $`T_0^2(m)`$ amplitude contribution is doubled and has the opposite sign. These differences are essential. As is seen from Fig. 5, the $`f_0(980)`$ resonance in the $`\gamma \gamma \pi ^0\pi ^0`$ channel has to manifest itself as a distinct peak. In this respect, the reaction $`\gamma \gamma \pi ^0\pi ^0`$, generally speaking, is more preferred than the reaction $`\gamma \gamma \pi ^+\pi ^{}`$. Unfortunately, in the Crystal Ball E4 ; E5 and JADE E6 experiments, the $`\gamma \gamma \pi ^0\pi ^0`$ cross section was scanned with a 50-MeV and 30-MeV-wide step, respectively. Such a mass resolution is still lacking to discover the $`f_0(980)`$ peak. Figures 5(a) and 5(b) show, in particular, that a Gaussian smearing with the dispersion of 30 MeV leaves nothing from the specific features of the $`f_0(980)`$ peak in $`\sigma _S(\gamma \gamma \pi ^0\pi ^0)`$. Notice, that there are no contradictions between the presented estimate for the smoothed $`\sigma _S(\gamma \gamma \pi ^0\pi ^0)`$ and the normalized Crystal Ball data E4 ; E5 for $`\sigma (\gamma \gamma \pi ^0\pi ^0,|\mathrm{cos}\theta ^{}|<0.8,0.7)`$.
Of course, the considered model allows us to predict the $`S`$ wave $`\gamma \gamma \pi \pi `$ reaction cross sections for a more wide region of $`m`$ than the neighborhood of the $`f_0(980)`$ resonance. The corresponding cross sections $`\sigma _S(\gamma \gamma \pi ^+\pi ^{})`$ and $`\sigma _S(\gamma \gamma \pi ^0\pi ^0,|\mathrm{cos}\theta ^{}|<0.8)`$ in the $`m`$ region from $`2m_\pi `$ to 1.2 GeV are shown in Fig. 6 for the model parameters corresponding to sets E and F. Unfortunately, in such a wide $`m`$ interval we cannot directly compare the predictions for $`\sigma _S(\gamma \gamma \pi \pi )`$ with experiment, because this requires the accurate $`S`$ wave data obtained by separating highest partial waves with the use of a partial wave analysis of the reaction events. For example, in the reaction $`\gamma \gamma \pi ^+\pi ^{}`$, the $`D`$ wave contribution with $`|\lambda |=2`$ can constitute from 75% to 90% of the total cross section for $`m>0.5`$ GeV. In the $`\gamma \gamma \pi ^0\pi ^0`$ cannel, the $`D`$ wave contribution, caused in the main by the $`f_2(1270)`$ resonance, is also very important for $`m>0.85`$ GeV, as is clear from Fig. 6(c). Hence, the thorough separation of the large $`D`$ wave background is of crucial importance for the extraction of the $`S`$ wave in both reactions.
Finally, we wish to say a few words about ambiguities which, in fact, inevitably occur in theoretical models for the amplitudes of electromagnetic interactions of hadrons. Concretely, we bear in mind rather evident possibilities of the incorporation of some unknown, free parameters into the aforesaid model. One of these parameters is the so-called direct $`f_0(980)\gamma \gamma `$ coupling constant, $`g_{f_0\gamma \gamma }^0`$. Taking account of this constant, the corresponding total amplitude of the $`f_0\gamma \gamma `$ decay can be written as $`M_{f_0\gamma \gamma }(m)=M_{f_0K^+K^{}\gamma \gamma }^{Born}(m)+g_{f_0\gamma \gamma }^0`$, where $`M_{f_0K^+K^{}\gamma \gamma }^{Born}(m)`$ is the amplitude due to the Born $`K^+K^{}`$ loop mechanism from Eq. (3). Of course, for any mechanism, the two-photon decay amplitude of any scalar meson must be proportional to $`m^2`$ for $`m0`$, as, for example, the Born amplitude $`M_{f_0K^+K^{}\gamma \gamma }^{Born}(m)`$. However, if we are interested in only the narrow $`m`$ region around 1 GeV, the adding of the constant $`g_{f_0\gamma \gamma }^0`$ to $`M_{f_0K^+K^{}\gamma \gamma }^{Born}(m)`$ is a quite reasonable approximation. About the coupling constant $`g_{f_0\gamma \gamma }^0`$ one can say as follows. It can have neither the value comparable in magnitude and coincident in sign with the value of $`M_{f_0K^+K^{}\gamma \gamma }^{Born}(m)`$ at the maximum, i.e., with $`M_{f_0K^+K^{}\gamma \gamma }^{Born}(2m_{K^+})=\alpha (\pi ^2/41)g_{f_0K^+K^{}}/2\pi `$, nor the value comparable in magnitude but opposite in sign with $`M_{f_0K^+K^{}\gamma \gamma }^{Born}(2m_{K^+})`$, since otherwise the $`\gamma \gamma \pi ^+\pi ^{}`$ reaction cross section in the $`f_0(980)`$ region would be in sharp contradiction with the data, in both magnitude and shape. Moreover, there are no evidence for the presence of the pointlike $`f_0(980)\varphi \gamma `$ interaction from the data on the $`\varphi \pi \pi \gamma `$ decays E7 ; E8 ; E9 ; A11 . Actually, experiment tells us that the direct $`f_0\gamma \gamma `$ coupling seems to be small. Any reliable theoretical estimates for this coupling have not existed yet. Serious experimental and theoretical search for its signs together with those of the direct coupling of the $`f_0(600)/\sigma `$ to $`\gamma \gamma `$ are still a matter of the future.
## IV Conclusion
The present analysis was stimulated by the Belle data MBC . The main results consist in the following.
(i) It has been shown that the $`K^+K^{}`$ loop mechanism provides the absolutely natural and reasonable scale of the $`f_0(980)`$ resonance manifestation in the $`\gamma \gamma \pi ^+\pi ^{}`$ and $`\gamma \gamma \pi ^0\pi ^0`$ reaction cross sections.
(ii) It has been shown that the shape of the $`f_0(980)`$ resonance in the reaction $`\gamma \gamma \pi ^+\pi ^{}`$ has nothing to do with the shape of a solitary Breit-Wigner resonance. This result is supported by the Belle data. In so doing, the observed pattern of the $`f_0(980)`$ peak distortion can be easily explained with use of the simple dynamical model.
Certainly, for the more full understanding of the situation, the information based on a partial wave analysis of the $`\gamma \gamma \pi ^+\pi ^{}`$ reaction events in the $`f_0(980)`$ resonance region would be extremely useful. The huge statistics collected in the Belle experiment MBC , in principle, allows one to hope for the successful performance of such an analysis.
It is clear from the preceding discussion that high quality data on the reaction $`\gamma \gamma \pi ^0\pi ^0`$ would be also highly desirable, because the relative role of the background contributions in the $`f_0(980)`$ region in this channel is considerably smaller than in the charged one.
The new stage of high statistics measurements of the processes $`\gamma \gamma \pi ^+\pi ^{}`$, $`\gamma \gamma \pi ^0\pi ^0`$, $`\gamma \gamma \eta \pi ^0`$, $`\gamma \gamma K^+K^{}`$, and $`\gamma \gamma K^0\overline{K}^0`$, begun by the Belle Collaboration, undoubtedly, will serve the further progress of physics of light scalar mesons.
ACKNOWLEDGMENTS
This work was supported in part by the Presidential Grant No. 2339.2003.2 for the support of Leading Scientific Schools.
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# Transmission matrix of a uniaxial optically active crystal platelet
## 1 Introduction
Optical activity in crystals is a phenomenon which has been known for a long time (see Ref. for an historical review). Two methods are currently used to describe the effect of the crystal optical activity on plane wave propagation. The first one is phenomenological and is valid under normal incidence . The second one is a first principle method and consists in solving the Maxwell equations with modified constitutive equations (i.e. the relationship between the electric and magnetic fields vectors and the electric displacement and magnetic induction field vectors). The choice of the constitutive equations is of prior importance and was itself of phenomenological nature until the link between the optical activity as well as the magnetic-dipole and electric-octupole responses of the crystal medium was consistently formalised. This theory leads to constitutive equations different from those previously used. To the authorโs knowledge, the constitutive equations of Ref. have not yet been used to calculate the transmission matrix of an optical active uniaxial platelet. They were only applied to describe the reflection of a plane wave by an interface between an isotropic achiral and a uniaxial chiral media , but for an optical axis parallel or perpendicular to the plane of incidence. It is the purpose of this article to provide this transmission matrix in a simple and usable form. We further restrict ourselves to two geometric configurations often encountered experimentally, namely the optical axis parallel and perpendicular to the plane of interface.
Recent experimental results have shown that the contribution of the optical activity to the transmission matrix of a quarter wave plate can reach the percent level . As an application of our transmission matrix formula, the origin of this effect is investigated and the case of half wave plates is also considered. Our formula being valid under oblique incidence, numerical calculations of the variations of the platelet transmission matrix as a function of the angle of incidence are presented.
In section 2 we describe the method. In section 3, we derive the explicit expressions of the platelet transmission matrix for oblique and normal incidences. The contribution of the optical activity to the transmission matrix is studied numerically for various quartz platelet thicknesses in section 4.
## 2 Formalism
A monochromatic plane wave of wavelength $`\lambda `$ impinging a uniaxial optically active crystal slab of thickness $`\mathrm{}`$ is considered. The crystal is assumed to be non-absorbing, non-magnetic and surrounded by an isotropic achiral ambient medium.
A direct axis system $`x,y,z`$ is defined such that the $`z`$ axis is perpendicular to the plane of interface. The origin of the $`z`$ axis (i.e. $`z=0`$) is fixed on the first plane of interface so that $`z=\mathrm{}`$ corresponds to the second plane of interface (see Fig. 1). The unit vector basis attached to the system axes is denoted $`\{\widehat{๐ฑ},\widehat{๐ฒ},\widehat{๐ณ}\}`$.Without loss of generality, the plane of incidence is taken to be $`yz`$. To determine the transmission and reflection matrices of the interfaces between the ambient medium and the crystal faces, we follow the method of Ref. .
Near the first interface $`z=0`$, the electric and magnetic field vectors read:
$`\text{Incident: }๐=`$ $`\left[A_s\widehat{๐ฌ}+A_p\widehat{๐ฉ}\right]\mathrm{exp}(i[\omega t๐ค\text{ }๐ซ])`$ (1)
$`๐=`$ $`{\displaystyle \frac{1}{\mu _0\omega }}๐ค\times \left[A_s\widehat{๐ฌ}+A_p\widehat{๐ฉ}\right]\mathrm{exp}(i[\omega t๐ค\text{ }๐ซ])`$ (2)
$`\text{Reflected: }๐=`$ $`\left[B_s\widehat{๐ฌ}+B_p\widehat{๐ฉ}^{}\right]\mathrm{exp}(i[\omega t๐ค^{}\text{ }๐ซ])`$ (3)
$`๐=`$ $`{\displaystyle \frac{1}{\mu _0\omega }}๐ค^{}\times \left[B_s\widehat{๐ฌ}+B_p\widehat{๐ฉ}^{}\right]\mathrm{exp}(i[\omega t๐ค^{}\text{ }๐ซ])`$ (4)
$`\text{Refracted: }๐=`$ $`\left[C_o๐จ_+\mathrm{exp}(i๐ค_{๐จ}^{}{}_{+}{}^{}\text{ }๐ซ)+C_e๐_+\mathrm{exp}(i๐ค_{๐}^{}{}_{+}{}^{}\text{ }๐ซ)\right]\mathrm{exp}(i\omega t)`$ (5)
$`๐=`$ $`{\displaystyle \frac{1}{\mu _0\omega }}[C_o(๐ค_{๐จ}^{}{}_{+}{}^{}\times ๐จ_+i\mu _0\omega ๐ฏ๐จ_+)\mathrm{exp}(i๐ค_{๐จ}^{}{}_{+}{}^{}\text{ }๐ซ)+`$
$`C_e(๐ค_{๐}^{}{}_{+}{}^{}\times ๐_+i\mu _0\omega ๐ฏ๐_+)\mathrm{exp}(i๐ค_{๐}^{}{}_{+}{}^{}\text{ }๐ซ)]\mathrm{exp}(i\omega t)`$ (6)
and near the second interface $`z=\mathrm{}`$:
$`\text{Reflected: }๐=`$ $`\left[C_o^{}๐จ_{}\mathrm{exp}(i๐ค_{๐จ}^{}{}_{}{}^{}\text{ }๐ซ)+C_e^{}๐_{}\mathrm{exp}(i๐ค_{๐}^{}{}_{}{}^{}\text{ }๐ซ)\right]\mathrm{exp}(i\omega t)`$ (7)
$`๐=`$ $`{\displaystyle \frac{1}{\mu _0\omega }}[C_o^{}(๐ค_{๐จ}^{}{}_{}{}^{}\times ๐จ_{}i\mu _0\omega ๐ฏ๐จ_{})\mathrm{exp}(i๐ค_{๐จ}^{}{}_{}{}^{}\text{ }๐ซ)+`$
$`C_e^{}(๐ค_{๐}^{}{}_{}{}^{}\times ๐_{}i\mu _0\omega ๐ฏ๐_{})\mathrm{exp}(i๐ค_{๐}^{}{}_{}{}^{}\text{ }๐ซ)]\mathrm{exp}(i\omega t)`$ (8)
$`\text{Refracted: }๐=`$ $`\left[A_s^{}\widehat{๐ฌ}+A_p^{}\widehat{๐ฉ}\right]\mathrm{exp}(i[\omega t๐ค\text{ }๐ซ])`$ (9)
$`๐=`$ $`{\displaystyle \frac{1}{\mu _0\omega }}๐ค\times \left[A_s^{}\widehat{๐ฌ}+A_p^{}\widehat{๐ฉ}\right]\mathrm{exp}(i[\omega t๐ค\text{ }๐ซ])`$ (10)
with the incident field vectors given by Eqs. (5,6).
In these expressions, $`๐ค`$ and $`๐ค^{}`$ are the wave vectors in the ambient medium; the three vectors $`\widehat{๐ฌ}`$, $`\widehat{๐ฉ}`$ and $`๐ค`$ form a direct basis, they are given by $`\widehat{๐ฌ}=\widehat{๐ฉ}\times ๐ค/|\widehat{๐ฉ}\times ๐ค|`$ and $`\widehat{๐ฌ}=\widehat{๐ฉ}^{}\times ๐ค^{}/|\widehat{๐ฉ}^{}\times ๐ค^{}|`$ with $`\widehat{๐ฌ}`$ perpendicular to the plane of incidence (i.e. $`\widehat{๐ฌ}=\widehat{๐ฑ}`$); $`๐ฏ`$ is a second-rank tensor coming from the constitutive equations , its components are related to those of the gyration tensor coefficients $`g_{ij}`$ ; the electric vectors $`๐จ_\pm `$ and $`๐_\pm `$ are the solutions of the wave equation inside the crystal; the wave vectors inside the medium, $`๐ค_{๐จ}^{}{}_{\pm }{}^{}`$ and $`๐ค_{๐}^{}{}_{\pm }{}^{}`$, are given by the Snellโs law and by the condition of existence of solutions of the wave equation.
The wave equation inside the crystal is given by
$$\underset{\beta }{}\left(k_\alpha k_\beta \delta _{\alpha \beta }k^2+\mu _0\omega ^2ฯต_{\alpha \beta }+i\mu _0\omega \underset{\gamma }{}k_\gamma ๐_{\alpha \beta \gamma }\right)E_\beta =0$$
(11)
where $`\alpha ,\beta `$ and $`\gamma `$ stand for $`x`$ or $`y`$ or $`z`$ and where the magnetic-dipole and electric-quadrupole response of the medium is embodied in the third-rank tensor $`๐`$ whose components are related to those of the gyration tensor. In Eq. (11), $`k_\alpha `$ stands for the components of one of the four wave vectors $`๐ค_{๐จ}^{}{}_{\pm }{}^{}`$ and $`๐ค_{๐}^{}{}_{\pm }{}^{}`$; $`E_\beta `$ stands for one of the components of the four electric vectors $`๐จ_\pm `$ and $`๐_\pm `$; $`ฯต_{\alpha \beta }`$ are the elements of the dielectric tensor.
The wave amplitudes are related by matrix relations. The refraction by the first (i.e. $`z=0`$) and second (i.e. $`z=\mathrm{}`$) interfaces are described by :
$$\left(\begin{array}{c}C_o\\ C_e\end{array}\right)=T\left(\begin{array}{c}A_s\\ A_p\end{array}\right)\text{ and }\left(\begin{array}{c}A_s^{}\\ A_p^{}\end{array}\right)=T^{}P_+\left(\begin{array}{c}C_o\\ C_e\end{array}\right)$$
(12)
respectively. As for the reflections by the first and second crystal-ambient medium interfaces (inside the crystal), one gets
$$\left(\begin{array}{c}C_o\\ C_e\end{array}\right)=R\left(\begin{array}{c}C_o^{}\\ C_e^{}\end{array}\right)\text{ and }\left(\begin{array}{c}C_o^{}\\ C_e^{}\end{array}\right)=P_{}^1RP_+\left(\begin{array}{c}C_o\\ C_e\end{array}\right)$$
(13)
respectively. In the previous equations, the following $`2\times 2`$ phase matrices were introduced:
$$P_\pm =\left(\begin{array}{cc}\mathrm{exp}(i\phi _{o_\pm })& 0\\ 0& \mathrm{exp}(i\phi _{e_\pm })\end{array}\right),$$
(14)
with $`\phi _{o_\pm }=\mathrm{}(๐ค_{๐จ}^{}{}_{\pm }{}^{}\widehat{๐ณ})`$ and $`\phi _{e_\pm }=\mathrm{}(๐ค_{๐}^{}{}_{\pm }{}^{}\widehat{๐ณ})`$.
The $`2\times 2`$ matrices $`T`$, $`T^{}`$ and $`R`$ are obtained from the continuity conditions of the projections of the electric and magnetic field vectors in the planes $`z=0`$ and $`z=\mathrm{}`$. Once known, these matrices are used to compute the platelet transmission matrix $`M`$ such that $`๐_{\mathrm{๐จ๐ฎ๐ญ}}=M๐_{\mathrm{๐ข๐ง}}`$ where $`๐_{\mathrm{๐ข๐ง}}`$ and $`๐_{\mathrm{๐จ๐ฎ๐ญ}}`$ are the electric field vectors before and after the crystal slab respectively. Note that $`M`$, $`๐_{\mathrm{๐ข๐ง}}`$ and $`๐_{\mathrm{๐จ๐ฎ๐ญ}}`$ are thus represented in the $`\{\widehat{๐ฌ},\widehat{๐ฉ},\widehat{๐ค}\}`$ basis, so that $`M`$ reduces to a $`2\times 2`$ matrix and $`๐_{\mathrm{๐ข๐ง}}`$ and $`๐_{\mathrm{๐จ๐ฎ๐ญ}}`$ to two-dimension vectors when plane waves are considered.
Taking into account the multiple reflections inside the crystal, one gets:
$$M=T^{}P_+T+T^{}P_+RP_{}^1RP_+T+\mathrm{}=T^{}P_+\left[\underset{ยฏ}{1}RP_{}^1RP_+\right]^1T$$
(15)
where we have used $`_{i=0}^{\mathrm{}}X^i=[\underset{ยฏ}{1}X]^1`$ and where $`\underset{ยฏ}{1}`$ is the $`2\times 2`$ identity matrix.
We shall now choose a crystal symmetry class group in order to specify the two tensors $`๐`$ and $`๐ฏ`$.
## 3 Transmission matrix of a platelet with the optical axis parallel and perpendicular to the plane of interface
Quartz crystal is widely used in crystallography to test the reliability of experimental setups and to manufacture retardation plates. We thus restrict ourselves to the non-centrosymmetric crystals belonging to the symmetry classes 32, 422 and 622 for which the tensors $`๐`$ and $`๐ฏ`$ have similar expressions . In the crystallographic reference frame, i.e for the crystal optical axis aligned along the $`z`$ axis, the non-vanishing components of these tensors are :
$`\stackrel{~}{๐ฏ}_{xx}=\stackrel{~}{๐ฏ}_{yy}=g_{33}/2,\stackrel{~}{๐ฏ}_{zz}=g_{11}g_{33}/2`$ (16)
$`\stackrel{~}{๐}_{xyz}=\stackrel{~}{๐}_{yxz}=g_{33},\stackrel{~}{๐}_{yzx}=\stackrel{~}{๐}_{zyx}=g_{11},\stackrel{~}{๐}_{zxy}=\stackrel{~}{๐}_{xzy}=g_{11}`$ (17)
with $`\stackrel{~}{๐ฏ}_{lm}=\mu _0c๐ฏ_{lm}`$ and $`\stackrel{~}{๐}_{lmn}=\mu _0c๐_{lmn}`$ and where $`g_{11}`$ and $`g_{33}`$ are the two independent coefficients of the gyration tensor. Note that for another orientation of the optical axis, standard tensor transformations are used to determine the representations of $`\stackrel{~}{๐ฏ}`$ and $`\stackrel{~}{๐}`$.
We shall now consider two geometrical configurations often encountered experimentally: a platelet with its optical axis parallel and perpendicular to the plane of interface. The two cases are treated separately because different approximations are made during the calculations.
### 3.1 Optical axis parallel to the plane of interface
The wave equation (11) can be written in a matrix form
$$(+\delta )๐=0$$
(18)
where the contribution of the optical activity is contained in the second term $`\delta `$ and where $`๐=(E_x,E_y,E_z)^T`$. Writing the wave vector $`๐ค=(\omega /c)(0,\stackrel{~}{\beta },\stackrel{~}{k}_z)^T`$, with $`\stackrel{~}{\beta }=n_a\mathrm{sin}\theta `$ and where $`\theta `$ is the angle of incidence and $`n_a`$ the optical index of the ambient medium, we obtain:
$$=\left(\begin{array}{ccc}n_e^2\mathrm{cos}^2\varphi +n_o^2\mathrm{sin}^2\varphi \stackrel{~}{k}_{z}^{}{}_{}{}^{2}\stackrel{~}{\beta }^2& \mathrm{\Delta }\mathrm{sin}\varphi \mathrm{cos}\varphi & 0\\ \mathrm{\Delta }\mathrm{sin}\varphi \mathrm{cos}\varphi & n_o^2\mathrm{cos}^2\varphi +n_e^2\mathrm{sin}^2\varphi \stackrel{~}{k}_{z}^{}{}_{}{}^{2}& \stackrel{~}{\beta }\stackrel{~}{k}_z\\ 0& \stackrel{~}{\beta }\stackrel{~}{k}_z& n_o^2\stackrel{~}{\beta }^2\end{array}\right)$$
and
$$\delta =\left(\begin{array}{ccc}0& ig_{11}\stackrel{~}{k}_z& i\stackrel{~}{\beta }(g_{11}\mathrm{cos}^2\varphi +g_{33}\mathrm{sin}^2\varphi )\\ ig_{11}\stackrel{~}{k}_z& 0& i\stackrel{~}{\beta }\mathrm{sin}\varphi \mathrm{cos}\varphi (g_{33}g_{11})\\ i\stackrel{~}{\beta }(g_{11}\mathrm{cos}^2\varphi +g_{33}\mathrm{sin}^2\varphi )& i\stackrel{~}{\beta }\mathrm{sin}\varphi \mathrm{cos}\varphi (g_{33}g_{11})& 0\end{array}\right)$$
with $`\mathrm{\Delta }=n_e^2n_o^2`$ and where $`n_o=\sqrt{ฯต_o/ฯต_0}`$ and $`n_e=\sqrt{ฯต_e/ฯต_0}`$ are the ordinary and extraordinary optical indices respectively. In the previous expressions, $`\varphi `$ is the azimuth angle describing the orientation of the optical axis (note that $`\varphi =0`$ corresponds to the optical axis perpendicular to the plane of incidence).
Since optical activity is induced by the electric-quadrupole and magnetic-dipole responses of the crystal medium, we assume that $`|g_{ij}||n_en_o|`$. The condition of existence of a solution for Eq. (18) reads $`det(+\delta )=0`$. This condition gives us the four possible values for $`\stackrel{~}{k}_z`$: $`\stackrel{~}{k}_{oz_\pm }=\pm (\stackrel{~}{k}_{oz}^{(0)}+\delta \stackrel{~}{k}_{oz})`$ and $`\stackrel{~}{k}_{ez_\pm }=\pm (\stackrel{~}{k}_{ez}^{(0)}+\delta \stackrel{~}{k}_{ez})`$ which correspond to the ordinary and extraordinary waves respectively (the sign $`\pm `$ refers to the direction of propagation). The zero order terms $`\stackrel{~}{k}_{oz}^{(0)}`$ and $`\stackrel{~}{k}_{ez}^{(0)}`$ are given by the solution of $`det()=0`$:
$`\stackrel{~}{k}_{oz}^{(0)}`$ $`=(n_o^2\stackrel{~}{\beta }^2)^{1/2},`$ (19)
$`\stackrel{~}{k}_{ez}^{(0)}`$ $`=\left(n_e^2\stackrel{~}{\beta }^2(\mathrm{cos}^2\varphi +{\displaystyle \frac{n_e^2}{n_o^2}}\mathrm{sin}^2\varphi )\right)^{1/2},`$ (20)
and the first non-vanishing contributions in $`g_{ij}`$ read as
$`\delta \stackrel{~}{k}_{oz}`$ $`={\displaystyle \frac{1}{4n_o^2\stackrel{~}{k}_{oz}^{(0)}}}(g_{11}^2n_o^2\stackrel{~}{\beta }^2\mathrm{sin}^2\varphi (g_{33}g_{11})^2\delta k)`$ (21)
$`\delta \stackrel{~}{k}_{ez}`$ $`={\displaystyle \frac{1}{4n_o^2\stackrel{~}{k}_{ez}^{(0)}}}(g_{11}^2n_o^2\stackrel{~}{\beta }^2\mathrm{sin}^2\varphi (g_{33}g_{11})^2+\delta k)`$ (22)
with
$`\delta k={\displaystyle \frac{1}{(n_o^2\stackrel{~}{\beta }^2\mathrm{sin}^2\varphi )\mathrm{\Delta }}}[n_o^4g_{11}^2(n_e^2+n_o^2)`$
$`+\stackrel{~}{\beta }^2n_o^2\mathrm{sin}^2\varphi [2g_{11}(g_{33}g_{11})(n_e^2+n_o^2)+g_{33}^2\mathrm{\Delta }]`$
$`+\stackrel{~}{\beta }^4\mathrm{sin}^4\varphi (g_{33}g_{11})^2(n_e^2+n_o^2)].`$ (23)
As expected , Eqs. (21) and (22) are of second order in $`g_{ij}`$.
Accordingly, the ordinary and extraordinary electric field vectors are decomposed as follows: $`๐จ_\pm =๐จ_\pm ^{(0)}+\delta ๐จ_\pm `$ and $`๐_\pm =๐_\pm ^{(0)}+\delta ๐_\pm `$. The zero order terms $`๐จ_\pm ^{(0)}`$ and $`๐_\pm ^{(0)}`$ are solutions of $`๐จ_\pm ^{(0)}=0`$ and $`๐_\pm ^{(0)}=0`$:
$`๐จ_\pm ^{(0)}`$ $`=N_o(\stackrel{~}{k}_{oz}^{(0)}\mathrm{sin}\varphi ,\stackrel{~}{k}_{oz}^{(0)}\mathrm{cos}\varphi ,\pm \stackrel{~}{\beta }\mathrm{cos}\varphi )^T,`$ (24)
$`๐_\pm ^{(0)}`$ $`=N_e(n_o^2\mathrm{cos}\varphi ,(\stackrel{~}{k}_{oz}^{(0)})^2\mathrm{sin}\varphi ,\stackrel{~}{k}_{ez}^{(0)}\mathrm{sin}\varphi )^T.`$ (25)
where $`N_o`$ and $`N_e`$ are the normalisation factors such that $`|๐จ_\pm ^{(0)}|=1`$ and $`|๐_\pm ^{(0)}|=1`$. $`\delta ๐จ_\pm `$ and $`\delta ๐_\pm `$, the first order contributions in $`g_{ij}/\mathrm{\Delta }`$, are the solutions of $`\delta ๐จ_\pm =\delta ๐จ_\pm ^{(0)}`$ and $`\delta ๐_\pm =\delta ๐_\pm ^{(0)}`$. Although $`\stackrel{~}{k}_{oz\pm }`$ and $`\stackrel{~}{k}_{ez\pm }`$ are of second order in $`g_{ij}`$, implying that $`det()=0`$ at first order, it turns out that the conditions of existence of solutions for these equations are fulfilled. We obtain
$`\delta ๐จ_\pm `$ $`={\displaystyle \frac{\pm iN_o}{\mathrm{\Delta }}}\left(\begin{array}{c}\mathrm{cos}\varphi [g_{11}n_o^2+\stackrel{~}{\beta }^2\mathrm{sin}^2\varphi (g_{33}g_{11})]\\ \mathrm{sin}\varphi [g_{11}n_o^2+\stackrel{~}{\beta }^2\mathrm{sin}^2\varphi (g_{33}g_{11})]\\ \stackrel{~}{\beta }\mathrm{sin}\varphi [g_{11}n_o^2+g_{33}\mathrm{\Delta }+\stackrel{~}{\beta }^2\mathrm{sin}^2\varphi (g_{33}g_{11})]/(n_o^2\stackrel{~}{\beta }^2)\end{array}\right)`$ (26)
$`\delta ๐_\pm `$ $`={\displaystyle \frac{\pm iN_e[g_{11}n_o^2+\stackrel{~}{\beta }^2\mathrm{sin}^2\varphi (g_{33}g_{11})]}{\mathrm{\Delta }}}\left(\begin{array}{c}\stackrel{~}{k}_{ez}^{(0)}\mathrm{sin}\varphi n_o^2/(n_o^2\stackrel{~}{\beta }^2)\\ \stackrel{~}{k}_{ez}^{(0)}\mathrm{cos}\varphi \\ \stackrel{~}{\beta }\mathrm{cos}\varphi [(\stackrel{~}{k}_{ez}^{(0)})^2\mathrm{\Delta }]/(n_o^2\stackrel{~}{\beta }^2)^2\end{array}\right)`$ (27)
Once the ordinary and extraordinary wave and electric vectors are known, the matrices $``$, $`๐ฏ`$ and $`๐ฏ^{}`$ are determined by supplying the boundary conditions in the planes $`z=0`$ and $`z=\mathrm{}`$. For an oblique incidence, these calculations are more simply performed numerically. Compact analytical expressions are obtained for the special case of normal incidence and will be given bellow.
To define a platelet transmission matrix at first order in $`g_{ij}/\mathrm{\Delta }`$, we introduce the three matrices $`R_0`$, $`T_0`$ and $`T_0^{}`$ obtained by setting $`g_{11}=g_{33}=0`$ to zero in the expressions of $`R`$, $`T`$ and $`T`$ respectively. We then define the matrices $`\delta R`$, and $`\delta T`$ and $`\delta T^{}`$ such that $`R=R_0+\delta R`$, $`T=T_0+\delta T`$ and $`T^{}=T_0^{}+\delta T`$. Having thus isolated the zero order from the first order terms in $`g_{ij}/\mathrm{\Delta }`$, we define the platelet transmission matrix $`M=M_0+\delta M`$ and thereby
$$๐_{\mathrm{๐จ๐ฎ๐ญ}}=๐_{\mathrm{๐จ๐ฎ๐ญ}}^{(\mathrm{๐})}+\delta ๐_{\mathrm{๐จ๐ฎ๐ญ}}\text{ such }๐_{\mathrm{๐จ๐ฎ๐ญ}}^{(\mathrm{๐})}=M_0๐_{\mathrm{๐ข๐ง}},\delta ๐_{\mathrm{๐จ๐ฎ๐ญ}}=\delta M๐_{\mathrm{๐ข๐ง}}.$$
(28)
Proceeding as in Eq. (15), we obtain:
$`\delta M=`$ $`\delta T^{}(T_0^{})^1M_0+M_0T_0^1\delta T+`$
$`T_0^{}\left[\underset{ยฏ}{1}P_+R_0P_{}^1R_0\right]^1P_+R_0P_{}^1\delta R\left[\underset{ยฏ}{1}P_+R_0P_{}^1R_0\right]^1P_+T_0+`$
$`T_0^{}\left[\underset{ยฏ}{1}P_+R_0P_{}^1R_0\right]^1P_+\delta R\left[\underset{ยฏ}{1}P_{}^1R_0P_+R_0\right]^1P_{}^1R_0P_+T_0`$ (29)
and
$$M_0=T_0^{}P_+\left[\underset{ยฏ}{1}R_0P_{}^1R_0P_+\right]^1T_0.$$
(30)
Note that $`M_0`$ only depends on the optical activity through the phase matrices $`P_+`$ and $`P_{}^1`$.
Under normal incidence, one gets $`\stackrel{~}{k}_{oz_\pm }=\pm n_o[1g_{11}^2/(2\mathrm{\Delta })]`$, $`\stackrel{~}{k}_{ez_\pm }=\pm n_e[1+g_{11}^2/(2\mathrm{\Delta })]`$ to the first order in $`g_{11}/\sqrt{\mathrm{\Delta }}`$ and the corresponding electric vectors
$`๐จ_\pm `$ $`=(\mathrm{sin}\varphi \pm {\displaystyle \frac{in_og_{11}\mathrm{cos}\varphi }{\mathrm{\Delta }}},\mathrm{cos}\varphi \pm {\displaystyle \frac{in_og_{11}\mathrm{sin}\varphi }{\mathrm{\Delta }}},0)^T,`$ (31)
$`๐_\pm `$ $`=(\mathrm{cos}\varphi \pm {\displaystyle \frac{in_eg_{11}\mathrm{sin}\varphi }{\mathrm{\Delta }}},\mathrm{sin}\varphi {\displaystyle \frac{in_eg_{11}\mathrm{cos}\varphi }{\mathrm{\Delta }}},0)^T.`$ (32)
The reflection and transmission matrices are:
$`R`$ $`=\left(\begin{array}{cc}\frac{n_on_a}{n_o+n_a}& in_e\frac{2g_{11}(n_o^2n_a^2)+g_{33}\mathrm{\Delta }}{(n_a+n_o)(n_a+n_e)\mathrm{\Delta }}\\ in_o\frac{2g_{11}(n_o^2n_a^2)+g_{33}\mathrm{\Delta }}{(n_a+n_o)(n_a+n_e)\mathrm{\Delta }}& \frac{n_en_a}{n_e+n_a}\end{array}\right)`$ (33)
$`T`$ $`=\left(\begin{array}{cc}in_a\frac{2g_{11}(n_en_a+n_o^2)+g_{33}\mathrm{\Delta }}{(n_a+n_o)(n_a+n_e)\mathrm{\Delta }}& \frac{2n_a}{n_o+n_a}\\ \frac{2n_a}{n_e+n_a}& in_a\frac{2g_{11}n_o(n_a+n_o)+g_{33}\mathrm{\Delta }}{(n_a+n_o)(n_a+n_e)\mathrm{\Delta }}\end{array}\right)(\varphi )`$ (34)
$`T^{}`$ $`=(\varphi )\left(\begin{array}{cc}in_o\frac{2g_{11}(n_en_a+n_o^2)+g_{33}\mathrm{\Delta }}{(n_a+n_o)(n_a+n_e)\mathrm{\Delta }}& \frac{2n_e}{n_e+n_a}\\ \frac{2n_o}{n_o+n_a}& in_e\frac{2g_{11}n_o(n_a+n_o)+g_{33}\mathrm{\Delta }}{(n_a+n_o)(n_a+n_e)\mathrm{\Delta }}\end{array}\right)`$ (35)
where $`(\varphi )`$ is the $`2\times 2`$ matrix representing a rotation of an angle $`\varphi `$ around the $`z`$ axis.
In these expressions, terms in $`g_{11}`$ are enhanced, unlike those in $`g_{33}`$, by a factor $`1/(n_en_o)1`$. This is a major difference with the case of an optical axis perpendicular to the plane of interface (see next section). From Eqs. (26-27) one also sees that the contribution of $`g_{33}`$ to the transmission matrix becomes noticeable only under oblique incidence.
Let us mention that Eq. (34) agrees with the result of Ref. for an isotropic chirality, i.e. $`g_{11}=g_{33}`$. It is however in disagreement with the expression of Ref. obtained using phenomenological constitutive equations. The expression of our platelet transmission matrix therefore disagrees with the one of Refs. .
The zero and first order platelet transmission matrices are given by Eqs. (30) and (3.1) respectively. For Eq. (30), we obtain the following simple expression:
$$M_0=(\varphi )\left(\begin{array}{cc}\frac{4n_on_a\mathrm{exp}(i\phi _o)}{(n_a+n_o)^2(n_an_o)^2\mathrm{exp}(2i\phi _o)}& 0\\ 0& \frac{4n_en_a\mathrm{exp}(i\phi _e)}{(n_a+n_e)^2(n_an_e)^2\mathrm{exp}(2i\phi _e)}\end{array}\right)(\varphi )$$
(36)
with $`\phi _o=(\omega /c)\mathrm{}\stackrel{~}{k}_{oz_+}`$ and $`\phi _e=(\omega /c)\mathrm{}\stackrel{~}{k}_{ez_+}`$.
### 3.2 Optical axis perpendicular to the plane of interface
When the optical axis is parallel to the $`z`$ axis, one can write Eq. (11) as follows:
$$\left(\begin{array}{ccc}n_o^2\stackrel{~}{k}_z^2\stackrel{~}{\beta }^2& ig_{33}\stackrel{~}{k}_z& ig_{11}\stackrel{~}{\beta }\\ ig_{33}\stackrel{~}{k}_z& n_o^2\stackrel{~}{k}_z^2& \stackrel{~}{\beta }\stackrel{~}{k}_z\\ ig_{11}\stackrel{~}{\beta }& \stackrel{~}{\beta }\stackrel{~}{k}_z& n_e^2\stackrel{~}{\beta }^2\end{array}\right)\left(\begin{array}{c}E_x\\ E_y\\ E_z\end{array}\right)=0.$$
To first order in $`g_{ij}`$, the four solutions for $`k_z`$ are given by $`\stackrel{~}{k}_{oz_\pm }=\pm (ab)^{1/2}`$ and $`\stackrel{~}{k}_{ez_\pm }=\pm (a+b)^{1/2}`$ with
$`a`$ $`=n_o^2{\displaystyle \frac{\stackrel{~}{\beta }^2}{2n_e^2}}(n_o^2+n_e^2)`$ (37)
$`b`$ $`=\left(n_o^2g_{33}^2+{\displaystyle \frac{n_o^2g_{33}\stackrel{~}{\beta }^2}{2n_e^2}}(4g_{11}3g_{33})+{\displaystyle \frac{\stackrel{~}{\beta }^4}{4n_e^4}}[2(n_o^2+n_e^2)(g_{33}g_{11})^2+(n_e^2n_o^2)^2]\right)^{1/2}`$ (38)
and the corresponding electric vectors read as
$$\widehat{๐ฏ}=N\left(\begin{array}{c}n_e^2(n_o^2\stackrel{~}{k}_z^2)n_o^2\stackrel{~}{\beta }^2\\ i\stackrel{~}{k}_zg_{33}n_e^2+i\stackrel{~}{\beta }^2\stackrel{~}{k}_z(g_{33}g_{11})\\ i\stackrel{~}{\beta }[\stackrel{~}{k}_z^2(g_{33}g_{11})+g_{11}n_o^2]\end{array}\right)$$
where $`\widehat{๐ฏ}๐จ_\pm ,๐_\pm `$ for $`\stackrel{~}{k}_z\stackrel{~}{k}_{oz_\pm },\stackrel{~}{k}_{ez_\pm }`$ respectively and where $`N`$ is the normalisation factor such that $`|\widehat{๐ฏ}|=1`$.
As in the previous case, the transmission matrix is more simply computed numerically. Under normal incidence, one gets $`\stackrel{~}{k}_{oz_\pm }=\pm (n_og_{33}/2)`$ and $`\stackrel{~}{k}_{ez_\pm }=\pm (n_o+g_{33}/2)`$ to the first order in $`g_{33}`$. The corresponding electric vectors are the two circular polarisation states that we write $`๐จ_\pm ^T=(i,1,0)/\sqrt{2}`$ and $`๐_\pm ^T=(\pm i,1,0)/\sqrt{2}`$. For the transmission and reflection matrices, one gets:
$$T=\frac{\sqrt{2}n_a}{n_a+n_o}\left(\begin{array}{cc}i& 1\\ i& 1\end{array}\right),T^{}=\frac{n_o\sqrt{2}}{n_a+n_o}\left(\begin{array}{cc}i& i\\ 1& 1\end{array}\right),R=\frac{n_on_a}{n_o+n_a}\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right).$$
Using Eq. (15), the following platelet transmission matrix is obtained:
$$M=\frac{2n_on_a[\mathrm{exp}(i\phi _o)+\mathrm{exp}(i\phi _e)]}{(n_a+n_o)^2(n_an_o)^2\mathrm{exp}(i[\phi _o+\phi _e])}\left(\begin{array}{cc}1& \mathrm{tan}((\phi _e\phi _o)/2)\\ \mathrm{tan}((\phi _e\phi _o)/2)& 1\end{array}\right)$$
(39)
with $`\phi _o=\mathrm{}k_{oz_+}`$ and $`\phi _e=\mathrm{}k_{ez_+}`$. This expression agrees with the result obtained using a $`4\times 4`$ matrix method .
Note that the interface matrices of eq. (39) do not depend explicitly of the coefficients $`g_{11}`$ and $`g_{33}`$. Thereby, the variations of $`M`$ with these coefficients come only from the phases $`\phi _o`$ and $`\phi _e`$. Note as well that the coefficient $`g_{11}`$ contributes only under oblique incidence.
## 4 Numerical results
To estimate the effects of the gyration coefficients on the platelet transmission, an experience closed to the HAUP is considered: a linearly polarised He-Ne laser beam crosses a quartz platelet and then a supposed perfect analyser. The incident electric vector reads $`๐_{\mathrm{๐ข๐ง}}=(\mathrm{cos}\psi ,\mathrm{sin}\psi )^T`$ in the $`\{\widehat{๐ฌ},\widehat{๐ฉ}\}`$ basis (i.e. in the basis $`\{\widehat{๐ฑ},\widehat{๐ฒ}\}`$ under normal incidence, see Fig. 1). After the analyser, whose eigen axis corresponds to, say, the $`x`$ axis, the beam intensity is measured as a function of $`\psi `$ and $`\varphi `$.
For our wavelength, $`\lambda =0.6328\mu `$m, the optical indices are $`n_o=1.542637`$ and $`n_e=1.551646`$ and the gyration coefficients $`g_{11}=5.910^5`$ and $`g_{33}=10.110^5`$.
To first order in $`g_{ij}/\mathrm{\Delta }`$, the beam intensity after the analyser is given by:
$$I_x=|๐_{\mathrm{๐จ๐ฎ๐ญ}}^{(\mathrm{๐})}\widehat{๐ฑ}|^2+[(๐_{\mathrm{๐จ๐ฎ๐ญ}}^{(\mathrm{๐})}\widehat{๐ฑ})(\delta ๐_{\mathrm{๐จ๐ฎ๐ญ}}\widehat{๐ฑ})^{}+(๐_{\mathrm{๐จ๐ฎ๐ญ}}^{(\mathrm{๐})}\widehat{๐ฑ})^{}(\delta ๐_{\mathrm{๐จ๐ฎ๐ญ}}\widehat{๐ฑ})]$$
(40)
where the field vectors are defined in Eq. (28).
We first consider the normal incidence and a crystal optical axis parallel to the plane of interface. As for the plate thickness, two values are chosen: $`\mathrm{}_1=87.80\mu `$m and $`\mathrm{}_2=1000.93\mu `$m (i.e. a first and a fourteen order quarter wave plates).
The intensity $`I_x`$ is shown as a function of $`\psi `$ and $`\varphi `$ in Fig. 2 (it is similar for the two plate thicknesses $`\mathrm{}_1`$ and $`\mathrm{}_2`$). Defining $`I_{0x}=|\widehat{๐ฑ}๐_{\mathrm{๐จ๐ฎ๐ญ}}^{(\mathrm{๐})}|^2=|\widehat{๐ฑ}M_0๐_{\mathrm{๐ข๐ง}}|^2`$, the difference $`\mathrm{\Delta }_x=I_xI_{0x}`$ is shown as a function of $`\varphi `$ and $`\psi `$ for the plate thickness $`\mathrm{}_1`$ in Fig. 3. The effect due to the medium gyrotropy is surprisingly high for $`\psi \pi /4`$ and the same behaviour is indeed observed for the second plate thickness $`\mathrm{}_2`$. In the expression of $`I_{0x}`$, the contribution of the optical activity comes only from the phases $`\phi _o`$ and $`\phi _e`$.
To distinguish the chiral contribution coming from the phases $`\phi _o`$ and $`\phi _e`$ from the one coming from the interface matrices, we define $`\stackrel{~}{I}_{0x}=|\widehat{๐ฑ}\stackrel{~}{M}_0๐_{\mathrm{๐ข๐ง}}|^2`$ where $`\stackrel{~}{M}_0`$ is given by Eq. 36, but fixing $`\phi _o=2\pi n_o/\lambda `$ and $`\phi _e=2\pi n_e/\lambda `$ (i.e. by neglecting the optical activity). The difference $`\mathrm{\Delta }_{0x}=I_{0x}\stackrel{~}{I}_{0x}`$ is shown in Fig. 4 as a function of $`\psi `$ and $`\varphi `$ for the plate thicknesses $`\mathrm{}_1`$. These figures indicate that the large variations observed in Fig. 3 come from the chiral dependence of the interface matrices. For the plate thickness $`\mathrm{}_2`$, the shape of $`\mathrm{\Delta }_{0x}`$ is identical to Fig. 4 but scaled by a factor of ten. As expected, the chiral dependence induced by the phase shift alone becomes only sizeable for thick plates.
The previous discussion holds for quarter wave plates. Let us consider another plate thickness $`\mathrm{}_3=1018.49\mu `$m, i.e. a fourteen order half wave plate. $`I_x`$, $`\mathrm{\Delta }_x`$ and $`\mathrm{\Delta }_{0x}`$ are shown as a function of $`\varphi `$ and $`\psi `$ in Figs. 5, 6 and 7 respectively. The effect induced by the optical activity is here reduced by two orders of magnitude with respect to a quarter wave plate and is dominated by the variation of the phases $`\phi _o`$ and $`\phi _e`$ with $`g_{11}`$. For a first order half wave plate ($`\mathrm{}_4=105.36\mu `$m), the effect on $`\mathrm{\Delta }_{0x}`$ is reduced by an order of magnitude whereas $`\mathrm{\Delta }_x`$ is essentially unchanged.
As mentioned in section 3.1, under normal incidence and for the optical axis parallel to the plane of interface, the contribution of the optical activity to the platelet transmission matrix comes essentially from the coefficient $`g_{11}`$. The contribution of the coefficient $`g_{33}`$ appears under oblique incidence and can, a priori, be distinguished from the other one by considering various angles of incidence. We now consider the ratio $`r=I_s/I_p`$ where $`I_s`$ and $`I_p`$ are obtained by substituting $`\widehat{๐ฌ}`$ and $`\widehat{๐ฉ}`$ to $`\widehat{๐ฑ}`$ in Eq. (40) respectively. To reproduce the experimental results of Ref. , we also fix $`\psi =\pi /4`$, i.e. $`๐_{\mathrm{๐ข๐ง}}=(1,1)^T/\sqrt{2}`$. The ratio $`r`$ is shown as a function of $`\varphi `$ and $`\theta `$ for the plate thickness $`\mathrm{}_1`$ in Fig. 8. To distinguish the contributions from $`g_{11}`$ and $`g_{33}`$, we show in Figs. 9 and 10 the relative differences $`\delta r=(rr_0)/r`$ and $`\delta ^{}r=(rr_1)/r`$ where $`r_0`$ and $`r_1`$ are obtained by setting $`g_{11}=g_{33}=0`$ and $`g_{11}=0`$ in the expression of $`r`$ respectively. From these figures, one sees that, as reported in Ref. , the contribution of the optical activity reach the percent level for certain values of $`\varphi `$.
In addition, we observe here that the contribution of $`g_{33}`$ is negligible under normal incidence but reaches the few $`10^3`$ level for reasonable values of the angle of incidence. The shape of the contribution of $`g_{11}`$ and $`g_{33}`$ being different, it seems possible to determine both coefficients by measuring $`r`$ as a function of $`\varphi `$ for various angles of incidence.
We finally consider the case of an optical axis perpendicular to the plane of interface. We choose $`\mathrm{}_5=1000\mu `$m for the plate thickness and, as previously, $`๐_{\mathrm{๐ข๐ง}}=(1,1)^T/\sqrt{2}`$ for the incident electric vector. In addition to $`\delta r`$, we introduce the phase angle $`\alpha `$ defined by the argument of the complex number $`(๐_{\mathrm{๐จ๐ฎ๐ญ}}\widehat{๐ฌ})/(๐_{\mathrm{๐จ๐ฎ๐ญ}}\widehat{๐ฉ})`$ and the difference $`\mathrm{\Delta }\alpha =\alpha \alpha _0`$ where $`\alpha _0`$ is obtained by setting $`g_{11}=0`$ in the expression of $`๐_{\mathrm{๐จ๐ฎ๐ญ}}`$. Figs. 11 and 12 show $`\delta r`$ and $`\mathrm{\Delta }\alpha `$ as a function of the angle of incidence $`\theta `$.
As in the case of an optical axis parallel to the plane of interface, the measurements of $`r`$ and $`\mathrm{\Delta }\alpha `$ as function of the angle of incidence lead to the determination of both $`g_{11}`$ and $`g_{33}`$. The sensitivity to these parameters is even higher when the optical axis is perpendicular to the plane of interface.
## 5 Conclusion
The transmission matrix for a uniaxial and chiral platelet has been calculated for an optical axis parallel and perpendicular to the plane of interface. The optical activity has been taken into account by considering the electric- quadrupole and magnetic-dipole responses of the crystal medium. Simple and usable expressions have been derived for normal incidence and for the optical axis parallel and perpendicular to the plane of interface. The case of an oblique incidence has been treated partly analytically and partly numerically.
As an application, numerical studies of the transmission of quartz platelets of various thicknesses have been provided. Under normal incidence, it was shown that the contribution of the optical activity can reach the percent level for quarter wave plates, independently of the plateโs order. Such a high contribution was found to be due to the transmission and reflection interface matrices. The case of half wave plates was also investigated and, unlike quarter wave plates, the main contribution of the optical activity came from the phase shift induced by the optical path inside the crystal. The variation of the intensity after a quartz platelet was studied as a function of the angle of incidence and of the orientation of the optical axis. A sensitivity to both gyration coefficients $`g_{11}`$ and $`g_{33}`$ was observed under oblique incidence. This suggests that varying the angle of incidence one could determine experimentally these two coefficients with a unique crystal cut. With respect to the results of Ref. , oblique incidence also offers the possibility to determine very accurately the platelet thickness and thereby the gyration coefficients.
## Acknowledgement
I would like to thank H. Feumi-Jentou and P.H. Burgat-Charvillon for their support. I would also like to thank F. Marechal for careful reading.
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# Post-Newtonian Parameters from Alternative Theories of Gravity
## I Introduction
The most striking and recent experimental discovery regarding Cosmology and the structure of the universe is related with the evidence of the acceleration of the universe, which is supported by experimental data deriving from different tests: i.e., from Type-Ia Supernovae, from CMWB and from the large scale structure of the universe Perlmu . Standard General Relativity is not able to provide a theoretical explanation to these experimental results unless some exotic and invisible matter is admitted to exist in the universe (Dark Energy). Proposals to explain the cosmic acceleration also arise from higher dimensional theories of Gravity dvali . Alternative theories to explain the acceleration of the universe have been recently proposed in the framework of higher order theories of Gravity carrol2 , already introduced in the framework of cosmological models to explain the early time inflation staro . Different models have been then studied both in the standard metric formalism metricfR and in the first order Palatini formalism palatinifR . Higher order theories of Gravity have been studied also in a quantum framework and a quantization of $`L(R)`$ theories has been performed in zerbo .
To test the theoretical consistence of these theories with observational data is however necessary to examine and to fit the standard tests for General Relativity: in particular solar system experiments and the tests of gravity at galactic scales. General Relativity reproduces with an excellent precision the experimental results obtained at the solar system scale solar . This naturally implies that each theory which pretends to be consistent with experimental results should surely reproduce General Relativity in this limit.
The aim of this paper is to provide a general theoretical framework to test the reliability of alternative theories of Gravity with solar system experiments. Such a problem was already studied from a different viewpoint in the standard metric formalism in solarmet and in the Palatini formalism in palatinifR , solarpal . Some debate is still open on the accordance of experimental results with solar system experiments and some authors erroneously claim that only theories which do not differ too much from General Relativity do the job (see palatinifR ); however, as we shall see, this is not true and, moreover, it is known that the Palatini formalism can naturally provide accordance with solar system experiments (see e.g. solarpal , Barraco , Barraco2 and references quoted therein). Some interesting results are also present in literature regarding the accordance of alternative theories of Gravity with rotational curves of galaxies capogal . In this paper we shall study the problem of the reliability of alternative theories of Gravity with solar system experiments and give also some hints regarding the galactic scale tests of Gravity (which will be considered in a forthcoming paper MGnew ). We shall do this from a purely theoretical viewpoint, trying to understand which Newtonian or Post-Newtonian modifications to standard General Relativity arise from specific modifications of the Hilbert-Einstein Lagrangian. In particular we consider $`L(R)`$ theories where the Lagrangian depends on an arbitrary analytic function $`L`$ of the scalar curvature $`R`$. Starting from the results already obtained in buchdahl and FFV we find an exact solution to field equations in vacuum. In that case field equations are controlled by a scalar-valued equation called the structural equation. It is relevant that modifications to the standard general relativistic gravitational field arise, and they turn out to be directly related to solutions of the structural equation and, consequently, to the particular form of the Lagrangian chosen (the choice of $`L(R)`$). We shall show how these modifications can be suitably interpreted as Post-newtonian parameters related to the non linearity of the theory.
We consider furthermore field equations in the case of matter universes (i.e. when the stress energy tensor is non vanishing). Considering a linear approximation of the metric, either with respect to a Minkowski flat space-time, or with respect to a de Sitter or an anti de Sitter space-times, the non-linear structure of the theory influences the gravitational field. We stress however that, in the first order approximation of the Palatini formalism, the presence of non-linear terms in the Lagrangian only influences the definition of $`R`$, while field equations remain unchanged. We finally derive the gravitational field for the particular Lagrangians $`R+\alpha f(R)`$ where $`\alpha `$ is an adimentional parameter. The corresponding gravitational potential contains then a term which is directly proportional to $`\alpha `$, such that General Relativity is reproduced in the limit $`\alpha =0`$, as it should be expected. This implies that the parameter $`\alpha `$ behaves as a sort of scale parameter which becomes relevant at large scales and it can be interpreted as a Post-Newtonian parameter ensuing from the non linearity of the Lagrangian.
Our approach, of course does not completely solve the problem of the generic reliability of alternative theories of gravity at solar system and galactic scale. However by introducing some Post-Newtonian parameters, it shows that General Relativity is certainly reproduced at small scales (as it is expected) for large families of Lagrangians. Further comparisons with other classical tests of General Relativity and applications to more general cases, as well as tests of Gravity at large (galactic) scale will be presented in the forthcoming paper MGnew .
## II The theoretical framework of $`L(R)`$ gravity
We deal with a $`4`$-dimensional gravitational theory on a Lorentzian manifold ($`M,g`$) with signature $`(,+,+,+)`$<sup>1</sup><sup>1</sup>1If not otherwise stated, we use units such that $`G=c=1`$.. The action is chosen to be:
$$A=A_{\mathrm{grav}}+A_{\mathrm{mat}}=[\sqrt{g}L(R)+2\kappa L_{\mathrm{mat}}(\psi ,\psi )]d^4x$$
(1)
where $`RR(g,\mathrm{\Gamma })=g^{\alpha \beta }R_{\alpha \beta }(\mathrm{\Gamma })`$, $`R_{\mu \nu }(\mathrm{\Gamma })`$ being the Ricci tensor of any torsionless connection $`\mathrm{\Gamma }`$ independent on a metric $`g`$ is assumed to be the physical metric. The gravitational part of the Lagrangian is represented by any real analytic function $`L(R)`$ of one real variable, which is assumed to be the scalar curvature $`R`$. The total Lagrangian contains also a first order matter part $`L_{\mathrm{mat}}`$ functionally depending on yet unspecified matter fields $`\mathrm{\Psi }`$ together with their first derivatives, equipped with a gravitational coupling constant $`\kappa =\frac{8\pi G}{c^4}`$ (see e.g. buchdahl ).
Equations of motion ensuing from the first order รก la Palatini formalism are (see Barraco ; palatinifR ; FFV )
$`L^{}(R)R_{(\mu \nu )}(\mathrm{\Gamma }){\displaystyle \frac{1}{2}}L(R)g_{\mu \nu }`$ $`=`$ $`\kappa T_{\mu \nu }^{mat}`$ (2)
$`_\alpha ^\mathrm{\Gamma }[\sqrt{g}L^{}(R)g^{\mu \nu })`$ $`=`$ $`0`$ (3)
where $`T_{mat}^{\mu \nu }=\frac{2}{\sqrt{g}}\frac{\delta L_{\mathrm{mat}}}{\delta g_{\mu \nu }}`$ denotes the matter source stress-energy tensor and $`^\mathrm{\Gamma }`$ means covariant derivative with respect to the connection $`\mathrm{\Gamma }`$, which we recall to be independent on the metric $`g`$. In this paper the metric $`g`$ and its inverse are used for lowering and raising indices.
We denote by $`R_{(\mu \nu )}`$ the symmetric part of $`R_{\mu \nu }`$, i.e. we set $`R_{(\mu \nu )}\frac{1}{2}(R_{\mu \nu }+R_{\nu \mu })`$. From (3) it follows that $`\sqrt{g}L^{}(R)g^{\mu \nu }`$ is a symmetric twice contravariant tensor density of weight $`1`$, so that it can be used (if non degenerate) to define a new metric $`h_{\mu \nu }`$ by the prescription:
$$\sqrt{g}L^{}(R)g^{\mu \nu }=\sqrt{h}h^{\mu \nu }$$
(4)
which is generically invertible. This means that the two metrics $`h`$ and $`g`$ are conformally equivalent so that space-time $`M`$ can be a posteriori endowed with a bi-metric structure $`(M,g,h)`$ FFV equivalent to the original metric-affine structure $`(M,g,\mathrm{\Gamma })`$. The corresponding conformal factor can be easily found to be $`L^{}(R)`$, since (4) gives:
$$h_{\mu \nu }=L^{}(R)g_{\mu \nu }$$
(5)
Therefore, as it is well known, equation (3) implies that $`\mathrm{\Gamma }=\mathrm{\Gamma }_{LC}(h)`$, i.e. the dynamical connection turns out a posteriori to be the Levi-Civita connection of the newly defined metric $`h`$, so that $`R_{(\mu \nu )}(\mathrm{\Gamma }_{LC}(h))=R_{\mu \nu }(h)R_{\mu \nu }`$ is now the metric Ricci tensor of the new metric $`h`$.
Equation (2) can be supplemented by the scalar-valued equation obtained by taking the $`g`$-trace of (2), where we set $`\tau =\mathrm{tr}T=g^{\mu \nu }T_{\mu \nu }^{mat}`$:
$$L^{}(R)R2L(R)=\kappa \tau $$
(6)
Equation (6) is called the structural equation and it controls the solutions of equation (2). For any real solution $`R=F(\tau )`$ of (6) we have in fact that both $`L(R)=L(F(\tau ))`$ and $`L^{}(R)=L^{}(F(\tau ))`$ can be seen as functions of $`\tau `$. For notational convenience we shall use the abuse of notation $`L(\tau )=L(F(\tau ))`$ and $`L^{}(\tau )=L^{}(F(\tau ))`$.
Now we are in position to introduce the generalized Einstein equations under the form<sup>2</sup><sup>2</sup>2Provided that $`L^{}(\tau )0`$: see below.
$$R_{\mu \nu }\left(h\right)=\frac{L(\tau )}{2L^{}(\tau )}g_{\mu \nu }+\frac{\kappa }{L^{}(\tau )}T_{\mu \nu }$$
(7)
with $`h_{\mu \nu }`$ defined by (5) for a given $`g_{\mu \nu }`$ and $`T_{\mu \nu }^{mat}`$ (see also Barraco ; FFV ; palatinifR ).
## III Some exact Solution of the Field Equations in vacuum
In this Section we look for a spherically symmetrical solution of the generalized Einstein equations in vacuum, starting from the results obtained in buchdahl and FFV . To this end first notice that eqs. (2-3), in vacuum, can be written under the form
$`[L^{}(R)]R_{(\mu \nu )}(\mathrm{\Gamma }){\displaystyle \frac{1}{2}}[L(R)]g_{\mu \nu }`$ $`=`$ $`0`$ (8)
$`_\alpha ^\mathrm{\Gamma }(\sqrt{g}[L^{}(R)]g^{\mu \nu })`$ $`=`$ $`0`$ (9)
Furthermore, the structural equation (6) becomes
$$L^{}(R)R2L(R)=0$$
(10)
In order to solve (8)-(9), we follow the discussion outlined in FFV . Let us suppose that the structural equation (10) is not identically satisfied and has a countable set of (real) solutions ($`i=1,2\mathrm{}.`$):
$$R=c_i$$
(11)
Then, we have two possibilities, depending on the value of the first derivative $`L^{}(R)`$ evaluated at the point $`R=c_i`$:
1. $`L^{}(c_i)=0`$
2. $`L^{}(c_i)0`$
In the first case, eqs. (10) implies that also $`L(c_i)=0`$, and, hence, the equations of motion (8)-(9) are identically satisfied. The only relation between $`g`$ and $`\mathrm{\Gamma }`$ is the following
$$R(g,\mathrm{\Gamma })=c_i$$
(12)
Indeed, this equation is not sufficient in this case to determine an explicit relation between the metric and the connection. Hence, in what follows, we shall suppose that $`L^{}(c_i)0`$.
We remark that if the Lagrangian is in the form $`L(R)=R^n`$, with $`n2`$, $`n`$, $`R=0`$ is solution of eq. (10), and, moreover one has $`L^{}(R=0)=0`$. Consequently we exclude such Lagrangians.
If $`L^{}(c_i)0`$ then the solution of the equations of motion (8)-(9) is given by the Levi-Civita connection of the metric $`h`$, which, in turn turns out to be equivalent to Levi-Civita connection of the physical metric $`g`$ (owing to the relation $`h=L^{}(c_i)g`$). Accordingly, the metric $`g`$ is the solution of the generalized Einstein equations
$$R_{\mu \nu }\left(g\right)=\mu g_{\mu \nu }$$
(13)
where
$$\mu =c_i/4$$
(14)
We look for a static solution of the field equations (13) describing the field outside a spherically symmetric mass distribution. Hence we may write the metric in the form
$$ds^2=e^{\mathrm{\Phi }(r)}dt^2+e^{\mathrm{\Lambda }(r)}dr^2+r^2d\vartheta ^2+r^2\mathrm{sin}^2\vartheta d\phi ^2$$
(15)
It is easy to check that the field equations (13) are satisfied if we set
$$e^{\mathrm{\Phi }(r)}=g_{tt}=1+\frac{C}{r}\frac{\mu r^2}{3}$$
(16)
and
$$e^{\mathrm{\Lambda }(r)}=g_{rr}=\left(1\frac{C}{r}+\frac{\mu r^2}{3}\right)^1$$
(17)
where $`C`$ is an arbitrary constant; in particular, the metric defined by (16)-(17) corresponds to the so called Schwarzschild-de Sitter space-time (see he ,ruffi ). The physical meaning of the constant $`C`$ becomes clear when considering the limit of weak gravitational field. We know that in General Relativity in this limit we have
$$g_{tt}\left(1+2\varphi \right)$$
(18)
where
$$\varphi =\frac{M}{r}$$
(19)
is the Newtonian potential, $`M`$ being the mass of the spherically symmetric source of the gravitational field. Consequently, in order to obtain the Newtonian limit we must set $`C=2M`$. Moreover, from (16) it is evident that a further contribution to the standard Newtonian potential is present in higher order theories of gravity. In particular, this contribution is proportional to the values of the Ricci scalar, owing to the proportionality between $`\mu `$ and $`c_i`$ (see (11) and (14)). This implies that the higher order contribution to the gravitational potential should be small enough not to contradict the known tests of gravity. In the case of small values of $`R`$ (which surely occur at solar system scale) the Einsteinian limit (i.e. the Schwarzschild solution) and the Newtonian limit are recovered, as it is evident from (16). In this context, $`\mu `$ can be naturally thought of as a Post-Newtonian parameter, ensuing from the non linearity of the theory ($`\mu =0`$ for the Hilbert-Einstein Lagrangian).
On the other hand this Post-Newtonian correction could play some role at larger scales and it could be interesting to test higher-order theories at galactic scales, as already done in the metric formalism in capogal .
## IV Field Equations in linear approximation
We aim at writing the field equations for Lagrangians $`L(R)=R+\alpha f(R)`$ in linear approximation: that is, we are going to solve the field equation at first order approximation with respect to a given background. In other words, we suppose to know a background solution of field equations (2), (3) determined by the affine connection $`{}_{}{}^{(0)}\mathrm{\Gamma }`$ and the metric $`{}_{}{}^{(0)}g`$ <sup>3</sup><sup>3</sup>3Here and henceforth, the superscripts <sup>(0)</sup> and <sup>(1)</sup> refer to the background and perturbed quantities, respectively.. We now perturb this solution by writing
$`\mathrm{\Gamma }_{\mu \nu }^\alpha `$ $`=`$ $`{}_{}{}^{(0)}\mathrm{\Gamma }_{\mu \nu }^{\alpha }+^{(1)}\mathrm{\Gamma }_{\mu \nu }^\alpha `$ (20)
$`g_{\mu \nu }`$ $`=`$ $`{}_{}{}^{(0)}g_{\mu \nu }^{}+^{(1)}g_{\mu \nu }`$ (21)
Furthermore, the matter source stress-energy tensor is written with respect to this perturbation in the form:
$$T_{\mu \nu }^{mat}=^{(0)}T_{\mu \nu }^{mat}+^{(1)}T_{\mu \nu }^{mat}$$
(22)
As a consequence, the equation
$$L^{}(R)R_{(\mu \nu )}(\mathrm{\Gamma })\frac{1}{2}L(R)g_{\mu \nu }=\kappa T_{\mu \nu }^{mat}$$
(23)
can be written under the form<sup>4</sup><sup>4</sup>4We have taken into account the fact that, on the background,
$$L^{}(^{(0)}R)^{(0)}R_{(\mu \nu )}(^{(0)}\mathrm{\Gamma })\frac{1}{2}L(^{(0)}R)^{(0)}g_{\mu \nu }=\kappa ^{(0)}T_{\mu \nu }^{mat}$$
(24)
$$L^{}(^{(0)}R)^{(1)}R_{\mu \nu }+L^{}(^{(1)}R)^{(0)}R_{\mu \nu }\frac{1}{2}^{(1)}g_{\mu \nu }L(^{(0)}R)\frac{1}{2}^{(0)}g_{\mu \nu }L(^{(1)}R)=\kappa ^{(1)}T_{\mu \nu }^{mat}$$
(25)
The Ricci curvature $`{}_{}{}^{(0)}R_{\mu \nu }^{}`$ (and the corresponding Ricci scalar $`{}_{}{}^{(0)}R`$) refer to the background solution; in terms of the perturbation of this solution we may write
$$R_{\mu \nu }=^{(0)}R_{\mu \nu }+^{(1)}R_{\mu \nu }$$
(26)
and
$$R=^{(0)}R+^{(1)}R+^{(1)}g^{\mu \nu }{}_{}{}^{(0)}R_{\mu \nu }^{}$$
(27)
So, in order to explicitly write the perturbed field equations (25) we have to evaluate the perturbed Ricci curvature and scalar in terms of the fields $`g`$ and $`\mathrm{\Gamma }`$. In general, we have Barraco :
$$R_{\mu \nu }(\mathrm{\Gamma })R_{\mu \nu }(g)=_{(\mu }Q_{\nu )\alpha }^\alpha _\alpha Q_{\mu \nu }^\alpha +Q_{\beta (\mu }^\alpha Q_{\nu )\alpha }^\beta Q_{\mu \nu }^\alpha Q_{\alpha \beta }^\beta $$
(28)
where
$$Q_{\mu \nu }^\alpha \{{}_{\mu \nu }{}^{\alpha }\}\mathrm{\Gamma }_{\mu \nu }^\alpha =\frac{1}{2}g^{\alpha \beta }(_\mu g_{\nu \beta }+_\nu g_{\mu \beta }_\beta g_{\mu \nu })$$
(29)
in terms of the Christoffel symbols
$$\{{}_{\mu \nu }{}^{\alpha }\}=\frac{1}{2}g^{\alpha \beta }(g_{\nu \beta ,\mu }+g_{\mu \beta ,\nu }g_{\mu \nu ,\beta })$$
(30)
Notice that $`_\mu _\mu ^\mathrm{\Gamma }`$ here and henceforth and we denote moreover with $`g_{\mu \beta ,\nu }`$ the partial derivative $`_\nu g_{\mu \beta }`$. The second set of field equations
$$_\alpha (\sqrt{g}[L^{}(R)]g^{\mu \nu })=0$$
(31)
can be now written in the form
$$_\alpha g_{\mu \nu }=b_\alpha g_{\mu \nu }$$
(32)
We have here defined:
$$b_\alpha _\alpha \left[\mathrm{ln}L^{}(R)\right]$$
(33)
ยฟFrom the structural equation (6) and from (33), we obtain then
$$b_\alpha \kappa \frac{L^{\prime \prime }(R)}{L^{}(R)}\frac{\tau _{,\alpha }}{L^{\prime \prime }(R)RL^{}(R)}$$
(34)
As a consequence, from equation (29) we may write
$$Q_{\mu \nu }^\alpha =\frac{1}{2}g^{\alpha \beta }\left(b_\mu g_{\nu \beta }+b_\nu g_{\mu \beta }b_\beta g_{\mu \nu }\right)$$
(35)
The expression of the Ricci tensor of the affine connection reads then
$$R_{\mu \nu }(\mathrm{\Gamma })=R_{\mu \nu }(g)+_{(\mu }b_{\nu )}\frac{1}{2}b_\mu b_\nu +g_{\mu \nu }b_\alpha b^\alpha +\frac{1}{2}_\alpha b^\alpha g_{\mu \nu }=R_{\mu \nu }(g)+B_{\mu \nu }$$
(36)
by introducing the tensor
$$B_{\mu \nu }_{(\mu }b_{\nu )}\frac{1}{2}b_\mu b_\nu +g_{\mu \nu }b_\alpha b^\alpha +\frac{1}{2}_\alpha b^\alpha g_{\mu \nu }$$
(37)
This expression (36) holds in the exact theory, so the task of writing its linear approximation is fulfilled by separately approximating the metric Ricci tensor $`R_{\mu \nu }(g)`$ and the $`B_{\mu \nu }`$ tensor; the latter, in particular, depends on the analytic expression of $`L(R)`$.
The perturbation of the metric Ricci tensor is given by (see Barraco ; straumann04 ):
$${}_{}{}^{(1)}R_{\mu \nu }^{}(g)=\frac{1}{2}^{(0)}g^{\alpha \beta }\left({}_{}{}^{(1)}g_{\beta \mu |\nu \alpha }^{}^{(1)}g_{\alpha \beta |\mu \nu }+^{(1)}g_{\beta \nu |\mu \alpha }^{(1)}g_{\mu \nu |\beta \alpha }\right)$$
(38)
where <sub>|</sub> stands for the (metric) covariant derivative with respect to the background<sup>5</sup><sup>5</sup>5The covariant derivative defined by<sub>|</sub> is such that $`{}_{}{}^{(0)}g_{\mu \nu |\alpha }^{}=0`$.. By perturbing the $`B_{\mu \nu }`$ tensor, we obtain
$$B_{\mu \nu }=^{(0)}B_{\mu \nu }+^{(1)}B_{\mu \nu }$$
(39)
where
$${}_{}{}^{(0)}B_{\mu \nu }^{}=^{(0)}b_{(\mu ;\nu )}\frac{1}{2}^{(0)}b_\mu ^{(0)}b_\nu +^{(0)}g_{\mu \nu }^{(0)}b_\alpha ^{(0)}b^\alpha +\frac{1}{2}^{(0)}b_{;\alpha }^\alpha g_{\mu \nu }$$
(40)
and
$`{}_{}{}^{(1)}B_{\mu \nu }^{}`$ $`=`$ $`{}_{}{}^{(1)}b_{\mu ,\nu }^{}^{(0)}\mathrm{\Gamma }_{\nu \mu }^\alpha {}_{}{}^{(1)}b_{\alpha }^{}^{(1)}\mathrm{\Gamma }_{\nu \mu }^\alpha {}_{}{}^{(0)}b_{\alpha }^{}{\displaystyle \frac{1}{2}}^{(0)}b_\mu ^{(1)}b_\nu {\displaystyle \frac{1}{2}}^{(1)}b_\mu ^{(0)}b_\nu +`$ (41)
$`+h_{\mu \nu }^{(0)}b_\alpha ^{(0)}b^\alpha +^{(0)}g_{\mu \nu }^{(0)}b_\alpha ^{(1)}b^\alpha +^{(0)}g_{\mu \nu }^{(1)}b_\alpha ^{(0)}b^\alpha +{\displaystyle \frac{1}{2}}^{(0)}g_{\mu \nu }^{(1)}b_{,\alpha }^\alpha +`$
$`+{\displaystyle \frac{1}{2}}^{(0)}g_{\mu \nu }^{(0)}\mathrm{\Gamma }_{\alpha \gamma }^\alpha {}_{}{}^{(1)}b_{}^{\gamma }+{\displaystyle \frac{1}{2}}^{(0)}g_{\mu \nu }^{(1)}\mathrm{\Gamma }_{\alpha \gamma }^\alpha {}_{}{}^{(0)}b_{}^{\gamma }+{\displaystyle \frac{1}{2}}h_{\mu \nu }^{(0)}b_{;\alpha }^\alpha `$
Notice that <sub>;</sub> stands here for the covariant derivative with respect to the unperturbed connection $`{}_{}{}^{(0)}\mathrm{\Gamma }`$.
### IV.1 Perturbation of flat space-time
We recall here that field equations are
$`L^{}(R)R_{(\mu \nu )}(\mathrm{\Gamma })`$ $``$ $`{\displaystyle \frac{1}{2}}L(R)g_{\mu \nu }=\kappa T_{\mu \nu }^{mat}`$ (42)
$`_\alpha g_{\mu \nu }`$ $`=`$ $`b_\alpha g_{\mu \nu }`$ (43)
where $`b_\alpha `$ is defined by formula (34). It is easy to check that the pair $`g_{\mu \nu }=\eta _{\mu \nu }`$, $`\mathrm{\Gamma }=0`$, i.e. the Minkowski flat space-time is a solution of the field equations (42), (43) iff
$$L(R=0)=0$$
(44)
In fact, in vacuum $`T0`$, hence $`b_\alpha =0`$<sup>6</sup><sup>6</sup>6Notice that in order to have a well posed definition of $`b_\alpha `$, we must have $`L^{}(R)0`$, and $`L^{\prime \prime }(R)RL^{}(R)0`$..
Now, we have to solve the field equations in terms of a perturbation of the Minkowski flat solution. In particular, we look for solutions in the form
$`\mathrm{\Gamma }_{\mu \nu }^\alpha `$ $`=`$ $`{}_{}{}^{(0)}\mathrm{\Gamma }_{\mu \nu }^{\alpha }+^{(1)}\mathrm{\Gamma }_{\mu \nu }^\alpha =^{(1)}\mathrm{\Gamma }_{\mu \nu }^\alpha `$ (45)
$`g_{\mu \nu }`$ $`=`$ $`{}_{}{}^{(0)}g_{\mu \nu }^{}+^{(1)}g_{\mu \nu }=\eta _{\mu \nu }+^{(1)}g_{\mu \nu }`$ (46)
In what follows we use Cartesian coordinates adapted to the background metric $`{}_{}{}^{(0)}g_{\mu \nu }^{}=\eta _{\mu \nu }`$; furthermore, the latter is used to raise and lower indices. The matter source stress-energy tensor is written in the form:
$$T_{\mu \nu }^{mat}=^{(0)}T_{\mu \nu }^{mat}+^{(1)}T_{\mu \nu }^{mat}=^{(1)}T_{\mu \nu }^{mat}$$
(47)
The Ricci curvature is written in the form
$$R_{\mu \nu }=^{(0)}R_{\mu \nu }+^{(1)}R_{\mu \nu }=^{(1)}R_{\mu \nu }$$
(48)
and the corresponding Ricci Scalar (owing to $`{}_{}{}^{(0)}R_{\mu \nu }^{}(\eta )=0`$):
$$R=^{(0)}R+^{(1)}R+^{(1)}g^{\mu \nu }{}_{}{}^{(0)}R_{\mu \nu }^{}=^{(1)}R$$
(49)
Notice that both the Ricci curvature and the Ricci scalar, when it is not explicitly stated (like in the above equations), refer to the connection $`\mathrm{\Gamma }`$. As a consequence, equation (24) can be written in the form
$$L^{}(0)^{(1)}R_{\mu \nu }(\mathrm{\Gamma })\frac{1}{2}\eta _{\mu \nu }L(^{(1)}R)=\kappa ^{(1)}T_{\mu \nu }^{mat}$$
(50)
ยฟFrom eq. (36), the perturbed Ricci tensor is made of two contributions:
$${}_{}{}^{(1)}R_{\mu \nu }^{}(\mathrm{\Gamma })=^{(1)}R_{\mu \nu }(g)+^{(1)}B_{\mu \nu }$$
(51)
The perturbation of the metric part of the Ricci tensor is obtained by replacing the covariant derivative <sub>|</sub> with the ordinary derivative in (38), since our background is Minkwowski flat space-time:
$${}_{}{}^{(1)}R_{\mu \nu }^{}(g)=\frac{1}{2}^{(0)}g^{\alpha \beta }\left({}_{}{}^{(1)}g_{\beta \mu ,\nu \alpha }^{}^{(1)}g_{\alpha \beta ,\mu \nu }+^{(1)}g_{\beta \nu ,\mu \alpha }^{(1)}g_{\mu \nu ,\beta \alpha }\right)$$
(52)
On the other hand, since on the background one has $`{}_{}{}^{(0)}b_{\alpha }^{}0`$, the perturbation of the $`B_{\mu \nu }`$ tensor reads now as:
$${}_{}{}^{(1)}B_{\mu \nu }^{}=^{(1)}b_{(\mu ,\nu )}+\frac{1}{2}\eta _{\mu \nu }^{(1)}b_{,\alpha }^\alpha $$
(53)
Hence, the perturbed Ricci tensor turns out to be
$${}_{}{}^{(1)}R_{\mu \nu }^{}(\mathrm{\Gamma })=\frac{1}{2}\eta ^{\alpha \beta }\left({}_{}{}^{(1)}g_{\beta \mu ,\nu \alpha }^{}^{(1)}g_{\alpha \beta ,\mu \nu }+^{(1)}g_{\beta \nu ,\mu \alpha }^{(1)}g_{\mu \nu ,\beta \alpha }\right)+^{(1)}b_{(\mu ,\nu )}+\frac{1}{2}\eta _{\mu \nu }^{(1)}b_{,\alpha }^\alpha $$
(54)
By exploiting gauge freedom, we may arbitrarily impose the following gauge condition
$$g^{\mu \nu }\mathrm{\Gamma }_{\mu \nu }^\alpha =0$$
(55)
which, in linear approximation and by forgetting vanishing terms, becomes:
$$\eta ^{\mu \nu }{}_{}{}^{(1)}\mathrm{\Gamma }_{\mu \nu }^{\alpha }=0$$
(56)
By taking the linear approximations of eqs. (29), (30), (35), condition (55) simply becomes
$${}_{}{}^{(1)}g_{\mu \alpha ,}^{\alpha }\frac{1}{2}^{(1)}g_{\alpha ,\mu }^\alpha +^{(1)}b_\mu =0$$
(57)
The gauge condition (57) allows us to write the perturbed Ricci tensor and the corresponding scalar curvature under the form
$${}_{}{}^{(1)}R_{\mu \nu }^{}(\mathrm{\Gamma })=\frac{1}{2}^{(1)}g_{\mu \nu ,\alpha }^\alpha +\frac{1}{2}\eta _{\mu \nu }^{(1)}b_{,\alpha }^\alpha $$
(58)
$${}_{}{}^{(1)}R(\mathrm{\Gamma })=\eta ^{\mu \nu }{}_{}{}^{(1)}R_{\mu \nu }^{}(\mathrm{\Gamma })=\frac{1}{2}^{(1)}g_{\mu ,\alpha }^{\mu \alpha }+2^{(1)}b_{,\alpha }^\alpha $$
(59)
Now we are in position to explicitly write the field equation (50). By taking into account (44) ,we may now suppose that the function $`L(R)`$ has the explicit form
$$L(R)=R+\alpha f(R)$$
(60)
where $`\alpha `$ is a constant parameter, and $`f(R)`$ is some function that for simplicity we may think to be as a polynomial of degree higher than one. A similar analysis holds for thr non-polynomial but still real analytic functions $`f(R)`$. Consequently, up to linear order we have
$$L^{}(0)=1L(^{(1)}R)^{(1)}R$$
(61)
and field equations become:
$${}_{}{}^{(1)}R_{\mu \nu }^{}\frac{1}{2}\eta _{\mu \nu }^{(1)}R=\kappa ^{(1)}T_{\mu \nu }^{mat}$$
(62)
By substituting the expression of the Ricci tensor and the scalar curvature (58), (59), we obtain
$$\frac{1}{2}^{(1)}g_{\mu \nu ,\alpha }^\alpha +\frac{1}{4}\eta _{\mu \nu }^{(1)}g_{\mu ,\alpha }^{\mu \alpha }=\kappa ^{(1)}T_{\mu \nu }^{mat}+\frac{1}{2}\eta _{\mu \nu }^{(1)}b_{,\alpha }^\alpha $$
(63)
Now, from (34), up to first order we may write
$${}_{}{}^{(1)}b_{,\alpha }^{\alpha }\kappa \frac{L^{\prime \prime }(0)}{\left(L^{}(0)\right)^2}^{(1)}T_{,\alpha }^{mat\alpha }$$
(64)
Furthermore we may introduce the tensor $`\overline{h}_{\mu \nu }`$ defined by
$$\overline{h}_{\mu \nu }^{(1)}g_{\mu \nu }\frac{1}{2}\eta _{\mu \nu }^{(1)}g_\alpha ^\alpha $$
(65)
Then by means of (64) and (65) the field equations (63) simplify to
$$\mathrm{}\left[\overline{h}_{\mu \nu }\eta _{\mu \nu }\kappa L^{\prime \prime }(0)^{(1)}\tau \right]=2\kappa ^{(1)}T_{\mu \nu }^{mat}$$
(66)
If we set
$$L(R)=R+\alpha f(R)=R+\alpha R^2+P(R)$$
(67)
where $`P(R)`$ is a polynomial of degree higher than 2, the field equations (66) can consequently be written in the form
$$\mathrm{}\left[\overline{h}_{\mu \nu }2\eta _{\mu \nu }\kappa \alpha ^{(1)}T\right]=2\kappa ^{(1)}T_{\mu \nu }^{mat}$$
(68)
By setting
$$_{\mu \nu }\overline{h}_{\mu \nu }2\eta _{\mu \nu }\kappa \alpha \tau ^{mat}$$
(69)
the field equations take the simple expression
$$\mathrm{}_{\mu \nu }=2\kappa ^{(1)}T_{\mu \nu }^{mat}$$
(70)
The solution of (70) can now be written in terms of retarded potentials:
$$_{\mu \nu }=4\frac{G}{c^4}\frac{T_{\mu \nu }^{mat}(t|๐ฑ๐ฑ^{}|/c,๐ฑ^{})}{|๐ฑ๐ฑ^{}|}d^3x^{}$$
(71)
where we have explicitly written $`\kappa =8\pi G/c^4`$. Hence
$$\overline{h}_{\mu \nu }=4\frac{G}{c^4}\frac{T_{\mu \nu }^{mat}(t|๐ฑ๐ฑ^{}|/c,๐ฑ^{})}{|๐ฑ๐ฑ^{}|}d^3x^{}+\frac{16\pi G}{c^4}\eta _{\mu \nu }\alpha \tau ^{mat}$$
(72)
ยฟFrom the above calculations, which have been performed step by step to exactly clarify what happens, it turns out that the first order perturbation does not influence the form of field equations (50), while it enters into the definition of the perturbed Ricci tensor (51) and consequently of the scalar curvature.
Specifying to the case of the Lagrangian (67) we see that the solution of the perturbed field equations, written in terms of the retarded potentials, contains two terms: i) the first one, in the weak field approximation, reduces to the standard Newtonian potential, ii) the second one is related to the Lagrangian chosen for the alternative theory of gravity. In particular it vanishes in the limit $`\alpha 0`$, i.e. exactly reproducing the weak field limit of standard General Relativity. It is thus clear that $`\alpha `$ can be identified with a scale parameter, vanishing at small (solar system) scales and consequently reproducing General Relativity. The same reasoning can be done by supposing that all the term $`\alpha f(R)`$ becomes in fact irrelevant at solar system scales.
### IV.2 Perturbation of the de Sitter space-time
The calculation performed in the previous sub-section can be generalized to the case of space-times which do not admit a Minkowski background solution. However, as we have seen in the previous section, having a Minkowski solution heavily constrains the available Lagrangians, since it implies that $`L(R=0)=0`$: in particular, these Lagrangians are not interesting for cosmological applications (see carrol2 and Barraco2 ). This case was already studied in Barraco2 , where it is shown that theories with singular $`L(R)`$ and $`\frac{d^2L}{dR^2}(^{(0)}R)=0`$ provide the correct Newtonian limit and they are good candidates to explain the cosmic acceleration.
We want hereafter to comment this case and to apply it to the particular Lagrangian $`L(R)=R+\alpha f(R)`$. We skip calculations as they can be reproduced step by step following the headlines of the previous chapter and moreover they have been already performed in Barraco2 . We consider as a background metric the (anti) de Sitter metric:
$${}_{}{}^{(0)}g=dt^2+e^{2t\sqrt{\frac{\mathrm{\Lambda }}{3}}}(dr^2+r^2d\vartheta ^2+r^2\mathrm{sin}^2\vartheta d\phi ^2)$$
(73)
which satisfies the field equations:
$${}_{}{}^{(0)}R_{\mu \nu }^{}=\mathrm{\Lambda }{}_{}{}^{(0)}g_{\mu \nu }^{}$$
(74)
Considering the Lagrangian $`L(R)=R+\alpha f(R)`$ we obtain that the resulting Newtonian potential is (see Barraco2 for details on calculations):
$$V(๐ฑ)=e^{2t\sqrt{\frac{\mathrm{\Lambda }}{3}}}C\frac{\rho (๐ฑ^{})\mathrm{exp}(|๐ฑ๐ฑ^{}|e^{2t\sqrt{\frac{\mathrm{\Lambda }}{3}}})}{๐ฑ๐ฑ^{}}d^3x^{}+A\rho (๐ฑ^{})$$
(75)
where $`\rho (๐ฑ)`$ represents the energy density while:
$$\{\begin{array}{cc}A=\alpha \frac{\kappa {}_{}{}^{(0)}f_{}^{\prime \prime }}{2{}_{}{}^{(0)}L_{}^{}(4\mathrm{\Lambda }\alpha {}_{}{}^{(0)}f_{}^{\prime \prime }+{}_{}{}^{(0)}L_{}^{})}\hfill & \\ & \\ C=\frac{8\mathrm{\Lambda }\alpha {}_{}{}^{(0)}f_{}^{\prime \prime }{}_{}{}^{(0)}L_{}^{}+({}_{}{}^{(0)}L_{}^{})^2}{({}_{}{}^{(0)}L_{}^{})^2(4\mathrm{\Lambda }\alpha {}_{}{}^{(0)}f_{}^{\prime \prime }+{}_{}{}^{(0)}L_{}^{})}\hfill & \end{array}$$
Also in this case it is evident that in the limit $`\alpha 0`$, for any current experiment and observation, the first term in (75) reduces to the standard Newtonian potential Barraco2 , and once again we obtain the weak field limit of standard General Relativity. Moreover $`\alpha `$ naturally behaves as a Post-Newtonian parameter in the potential, which is supposed to vanish at small scales. Once more $`\alpha `$ behaves like a scale parameter and the accordance with experimental results is supported.
## V Conclusions
In this paper we have shown that, both in the case of vacuum universes and in the case of matter universes, solar system experiments can be theoretically explained and reproduced in the framework of alternative theories of Gravity for specific classes of Lagrangians. The gravitational potential of alternative theories of Gravity reduces, under suitable hypotheses, to the standard Newtonian potential at the solar system scale. This has been proven both in the case of vacuum and matter universes (with a flat or an (anti) de-Sitter background). Gravitational effects due to the (alternative) form of the Lagrangian generate Post-Newtonian parameters appearing in the gravitational potential, which vanish when the corrections to the standard Hilbert Lagrangian are cancelled. Moreover we stress that these corrections are negligible when we consider values of the scale parameter $`\alpha `$ which is necessary to explain cosmic acceleration (see e.g. carrol2 ). These contributions become however relevant when considering larger scales (cosmology palatinifR and, hopefully, galactic scales). This implies that higher order corrections to the standard Hilbert-Einstein theory could behave as a scale effect, ruled by a scale parameter which vanishes at solar system scales. General Relativity is consequently reproduced at the solar system scale, as it has to be surely expected.
The results obtained here slightly differ from some results already presented in literature palatinifR and in particular from recent results obtained by G.J. Olmo, (see again palatinifR ). It was there argued that only small corrections to the Hilbert-Einstein Lagrangian can pass the solar system experiments. However, calculations were there performed by means of a conformal transformation on a flat Minkowski background spacetime. We have here shown that in the particular case of a flat spacetime, the theory $`\frac{1}{R}`$ is not viable, owing to the conditions (44). This implies that the results obtained by G.J. Olmo does not exclude the reliability of $`\frac{1}{R}`$-like theories, which should however be examined in the (anti)de-Sitter background framework. In fact $`R=0`$ is singular for all Lagrangians which contain inverse powers or logarithms. Moreover, it is not at all evident that our universe should be asymptotically flat. We have here proven the accordance of such theories with solar system experiments, at least when the scale (post-Newtonian) parameter becomes small enough.
## VI Acknowledgements
We are very grateful to Prof. S.D. Odintsov and Dr. Andrzej Borowiec for helpful comments and suggestions.
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# Dephasing representation of quantum fidelity for general pure and mixed states
## I Introduction
Time evolution in classical mechanics is very sensitive to perturbations of both initial conditions of a trajectory and the Hamiltonian. Because of the unitarity of quantum evolution, on the other hand, the overlap of two different quantum states remains constant in time. However, we can still define sensitivity of quantum evolution to perturbations of the Hamiltonian. This is usually done using the notion of *quantum fidelity* (sometimes called *Loschmidt echo*), defined for pure states as Peres (1984)
$$M\left(t\right)=\left|\psi \left|e^{+iH^ฯตt/\mathrm{}}e^{iH^0t/\mathrm{}}\right|\psi \right|^2.$$
(1)
Here $`|\psi `$ is the initial state, $`H^0`$ and $`H^ฯต=H^0+ฯตV`$ are the unperturbed and perturbed Hamiltonians, respectively. In words, fidelity is the overlap at time $`t`$ of two identical initial states evolved with two slightly different Hamiltonians. Because of its relevance in theories of decoherence and in experimental realizations of quantum computation Nielsen and Chuang (2000), quantum Jalabert and Pastawski (2001); Jacquod et al. (2001); Cerruti and Tomsovic (2002); Prosen (2002); Prosen and ลฝnidariฤ (2002); Jacquod et al. (2003); Silvestrov et al. (2003); Prosen and ลฝnidariฤ (2003, 2005); Bevilaqua and Heller (2004); Vanรญฤek and Heller (2003); Vanรญฤek (2004a, b); Cucchietti et al. (2002a, b); Wisniacki et al. (2002); Wisniacki and Cohen (2002); Wang and Li (2002); Prosen and Seligman (2002); Jacquod et al. (2002); Emerson et al. (2002); Wisniacki (2003); Cerruti and Tomsovic (2003); Cucchietti et al. (2003); Adamov et al. (2003); Kottos and Cohen (2003); ลฝnidariฤ and Prosen (2003); Hiller et al. (2004); Gorin et al. (2004); Jacquod (2004); Iomin (2004); Wang et al. (2004, 2005); Combescure (2005); Weinstein and Hellberg (2005) as well as classical Benenti et al. (2003a); Eckhardt (2003); Benenti et al. (2003b); Garcia-Mata et al. (2003); Veble and Prosen (2004) fidelity has been extensively studied in the last few years. Many universal regimes of fidelity decay have been found in different limiting cases Jalabert and Pastawski (2001); Jacquod et al. (2001); Cerruti and Tomsovic (2002); Prosen (2002); Prosen and ลฝnidariฤ (2002); Jacquod et al. (2003); Silvestrov et al. (2003); Prosen and ลฝnidariฤ (2003, 2005); Bevilaqua and Heller (2004). Many of these works used a semiclassical approach, but before Ref. Vanรญฤek and Heller (2003) only as a starting point for various approximations, because of difficulties in treating an exponentially growing number of terms in the general semiclassical expression for fidelity, especially in chaotic systems. This problem was solved in Ref. Vanรญฤek and Heller (2003) by a uniform expression for fidelity which implicitly summed over all these contributions using an integral over initial conditions, similar in spirit to Millerโs initial value representation Miller (1970, 2001). This surprisingly simple and accurate expression, although limited to wave packets localized in position, has been successfully applied as a starting point to derive fidelity decay in the deep Lyapunov regime Wang et al. (2005) and the plateau of fidelity in neutron scattering Bevilaqua and Heller (2004). Five other known regimes of fidelity can also be simply described by this method Vanรญฤek (2004b).
In a recent Rapid Communication Vanรญฤek (2004a), the uniform expression for fidelity was justified by the shadowing theorem of classical mechanics Hammel et al. (1987); Grebogi et al. (1990) and a more general and, in fact, *always* more accurate expression, valid for arbitrary pure states, was stated. One purpose of the present article is to provide (in Sec. II) a detailed derivation of this general semiclassical expression (18) for fidelity of arbitrary pure, i. e., also nonlocal states. Fidelity is expressed as an interference sum of dephasing trajectories weighed by the Wigner function of the initial state. The general derivation provides an alternative and more explicit justification of the validity of this *dephasing representation* (DR). Interestingly, in Sec. III it is shown that the same dephasing representation is valid also for general mixed states. Section IV shows how the general expression reduces to the original form Vanรญฤek and Heller (2003) and other specialized forms for position and momentum states or Gaussian wave packets localized in position or momentum. In Sec. V, the general dephasing representation is tested on a non-local stateโa coherent superposition of two separated wave packetsโand on two two types of mixed stateโan incoherent superposition of two wave packets and a completely random state. It is also shown that the general expression is superior to the original form Vanรญฤek and Heller (2003) even for a single Gaussian wave packet. All numerical calculations are done for a system with a finite Hilbert basis. In such systems, quantum phase space can be rigorously defined if the original Wigner function Wigner (1932) is replaced by the discrete Wigner transform Wooters (1987); Leonhardt (1995); Hannay and Berry (1980); Rivas and de Almeida (1999); Bouzouina and Bievre (1996). Since this discrete transform can be defined in a general abstract Hilbert space with finite basis, the present approach should be applicable to problems of quantum computation if phase space approach is used Miquel et al. (2002). In Sec. VI, DR is compared to other โWignerโ methods. The main conclusions of the paper are summarized in Sec. VII.
## II Dephasing representation for a general pure state
Fidelity amplitude for a general pure state $`|\psi `$ can be written as
$$O\left(t\right)=\psi \left|e^{+iH^ฯตt/\mathrm{}}e^{iH^0t/\mathrm{}}\right|\psi .$$
(2)
In order to derive the general dephasing representation of fidelity, we could start by replacing the two quantum propagators in Eq. (2) by the corresponding semiclassical Van Vleck propagators Vleck (1928), as in Refs. Vanรญฤek and Heller (2003); Vanรญฤek (2004a). However, we will save some effort if we start directly from the semiclassical initial value representation (IVR) Miller (1970, 2001) for the two Van Vleck propagators,
$`e^{iH^0t/\mathrm{}}`$ $`\left(2\pi i\mathrm{}\right)^{d/2}{\displaystyle ๐๐ซ_0^{}๐๐ฉ_0^{}\left|๐ซ_t^{}(๐ซ_0^{},๐ฉ_0^{})/๐ฉ_0^{}\right|^{1/2}}`$
$`\times e^{iS^0(๐ซ_0^{},๐ฉ_0^{};t)/\mathrm{}}|๐ซ_t^{}๐ซ_0^{}|,`$ (3)
$`e^{+iH^ฯตt/\mathrm{}}`$ $`\left(2\pi i\mathrm{}\right)^{d/2}{\displaystyle ๐๐ซ_0^{\prime \prime }๐๐ฉ_0^{\prime \prime }\left|๐ซ_t^{\prime \prime }(๐ซ_0^{\prime \prime },๐ฉ_0^{\prime \prime })/๐ฉ_0^{\prime \prime }\right|^{1/2}}`$
$`\times e^{iS^ฯต(๐ซ_0^{\prime \prime },๐ฉ_0^{\prime \prime };t)/\mathrm{}}|๐ซ_0^{\prime \prime }๐ซ_t^{\prime \prime }|.`$
Here $`๐ซ_0^{},๐ฉ_0^{}`$ and $`๐ซ_0^{\prime \prime },๐ฉ_0^{\prime \prime }`$ are the initial conditions of trajectories of $`H^0`$ and of $`H^ฯต`$, respectively, and $`๐ซ_t^{},๐ฉ_t^{}`$ and $`๐ซ_t^{\prime \prime },๐ฉ_t^{\prime \prime }`$ are the corresponding coordinates and momenta at time $`t`$. Action $`S^0`$ of a trajectory of the unperturbed Hamiltonian $`H^0`$, is given by
$$S^0(๐ซ_0^{},๐ฉ_0^{};t)=_0^t๐\tau \left[๐ฉ_\tau ^{}\dot{๐ซ}_\tau ^{}H^0(๐ซ_\tau ^{},๐ฉ_\tau ^{};\tau )\right].$$
(4)
Similar expression holds for the action $`S^{ฯต\prime \prime }(๐ซ_0^{\prime \prime },๐ฉ_0^{\prime \prime };t)`$ of a trajectory of the perturbed Hamiltonian $`H^ฯต`$. In the simplified notation above, the square roots of the determinants in Eq. (3) also include the appropriate Maslov indices Gutzwiller (1990). Using the IVR expressions (3), fidelity amplitude (2) becomes
$`O_{\text{IVR}}\left(t\right)`$ $`=`$ $`\left(2\pi \mathrm{}\right)^d{\displaystyle ๐๐ซ_0^{}๐๐ฉ_0^{}๐๐ซ_0^{\prime \prime }๐๐ฉ_0^{\prime \prime }\left|\frac{๐ซ_t^{}}{๐ฉ_0^{}}\right|^{1/2}}`$ (5)
$`\times `$ $`\left|{\displaystyle \frac{๐ซ_t^{\prime \prime }}{๐ฉ_0^{\prime \prime }}}\right|^{1/2}\psi |๐ซ_0^{\prime \prime }๐ซ_t^{\prime \prime }|๐ซ_t^{}๐ซ_0^{}|\psi e^{i\left(S^0S^{ฯต\prime \prime }\right)/\mathrm{}}.`$
### II.1 Uniform semiclassical expression for fidelity
If we further expand the $`\delta `$ function in integral (5) as an integral over a dummy momentum $`๐ช`$,
$$๐ซ_t^{\prime \prime }|๐ซ_t^{}=\delta \left(\mathrm{\Delta }๐ซ_t\right)=\left(2\pi \mathrm{}\right)^d๐๐ชe^{i๐ช๐ซ๐ซ_t/\mathrm{}},$$
we obtain a โfullโ uniform semiclassical expression for fidelity,
$`O_{\text{unif}}\left(t\right)`$ $`=`$ $`\left(2\pi \mathrm{}\right)^{2d}{\displaystyle ๐๐ซ_0^{}๐๐ฉ_0^{}๐๐ซ_0^{\prime \prime }๐๐ฉ_0^{\prime \prime }๐๐ช}`$ (6)
$`\times `$ $`\left|{\displaystyle \frac{๐ซ_t^{}}{๐ฉ_0^{}}}\right|^{1/2}\left|{\displaystyle \frac{๐ซ_t^{\prime \prime }}{๐ฉ_0^{\prime \prime }}}\right|^{1/2}\psi ^{}\left(๐ซ_0^{\prime \prime }\right)\psi \left(๐ซ_0^{}\right)`$
$`\times `$ $`\mathrm{exp}\left\{{\displaystyle \frac{i}{\mathrm{}}}\left[S^0S^{ฯต\prime \prime }+๐ช๐ซ๐ซ_t\right]\right\}.`$
This integral is, formally, semiclassically โexact.โ In particular, it is free of caustics, unlike, e. g., the Van Vleck propagator. Because it is expressed only in terms of initial conditions (and dummy momentum $`๐ช`$), it appears to be ready for numerical evaluations. Unfortunately, this integral is highly oscillatory, and very difficult to compute, especially in many-dimensional or chaotic systems. Therefore we will take an alternative route, using a further approximation, but obtain an integral much easier to tackle numerically.
### II.2 Dephasing representation
First, let us make a change of variables $`\{๐ซ^{},๐ซ^{\prime \prime },๐ฉ^{},๐ฉ^{\prime \prime }\}\{๐ซ,\mathrm{\Delta }๐ซ,๐ฉ,\mathrm{\Delta }๐ฉ\}`$ in integral (5). It should be emphasized that we do not assume $`\mathrm{\Delta }๐ซ`$ or $`\mathrm{\Delta }๐ฉ`$ to be small. New variables (averages and differences) are defined for all times from $`0`$ to $`t`$ as
$`๐ซ`$ $`={\displaystyle \frac{1}{2}}\left(๐ซ^{}+๐ซ^{\prime \prime }\right),`$ (7)
$`\mathrm{\Delta }๐ซ`$ $`=๐ซ^{\prime \prime }๐ซ^{},`$
$`๐ฉ`$ $`={\displaystyle \frac{1}{2}}\left(๐ฉ^{}+๐ฉ^{\prime \prime }\right),`$
$`\mathrm{\Delta }๐ฉ`$ $`=๐ฉ^{\prime \prime }๐ฉ^{},`$
The Jacobian of this transformation is unity. If we intend to perform integrals over $`\mathrm{\Delta }๐ซ`$ and $`\mathrm{\Delta }๐ฉ`$ first, we can consider $`๐ซ_0`$ and $`๐ฉ_0`$ as fixed for the moment, and write
$$\left|\frac{๐ซ_t^{}}{๐ฉ_0^{}}\right|^{1/2}\left|\frac{๐ซ_t^{\prime \prime }}{๐ฉ_0^{\prime \prime }}\right|^{1/2}=\left|\frac{\left(\mathrm{\Delta }๐ซ_t\right)}{\left(\mathrm{\Delta }๐ฉ_0\right)}\right|^{1/2}\left|\frac{\mathrm{\Delta }๐ซ_t}{\mathrm{\Delta }๐ฉ_0}\right|^{1/2}=\left|\frac{\mathrm{\Delta }๐ซ_t}{\mathrm{\Delta }๐ฉ_0}\right|,$$
(8)
$`O\left(t\right)`$ $`=`$ $`\left(2\pi \mathrm{}\right)^d{\displaystyle ๐๐ซ_0๐๐ฉ_0๐\mathrm{\Delta }๐ซ_0๐\mathrm{\Delta }๐ฉ_0\left|\frac{\mathrm{\Delta }๐ซ_t}{\mathrm{\Delta }๐ฉ_0}\right|}`$ (9)
$`\times `$ $`\psi ^{}\left(๐ซ_0^{\prime \prime }\right)\delta \left(\mathrm{\Delta }๐ซ_t\right)\psi \left(๐ซ_0^{}\right)\mathrm{exp}\left[{\displaystyle \frac{i}{\mathrm{}}}\left(S^0S^{ฯต\prime \prime }\right)\right].`$
Next we change variables from $`\mathrm{\Delta }๐ฉ_0`$ to $`\mathrm{\Delta }๐ซ_t`$ and eliminate the $`\delta `$ function,
$`O\left(t\right)`$ $`=`$ $`\left(2\pi \mathrm{}\right)^d{\displaystyle ๐๐ซ_0๐๐ฉ_0๐\mathrm{\Delta }๐ซ_0\psi ^{}\left(๐ซ_0^{\prime \prime }\right)\psi \left(๐ซ_0^{}\right)}`$ (10)
$`\times `$ $`\mathrm{exp}\left[{\displaystyle \frac{i}{\mathrm{}}}\left(S^0S^{ฯต\prime \prime }\right)\right]|_{\mathrm{\Delta }๐ซ_t=0}.`$
The present form is equivalent to Eq. (6). On one hand, the present form appears much simpler (a $`3d`$\- vs. $`5d`$-dimensional integral), on the other hand it is not an integral over independent variables because it contains a constraint on the final positions ($`\mathrm{\Delta }๐ซ_t=0`$).
While we do not intend to evaluate this integral by the stationary phase (SP) approximation, it is instructive to check where the action difference $`S^0S^{ฯต\prime \prime }`$ is stationary because those regions give the main contributions to the integral. Variation of action $`S^0`$ gives
$$\delta S^0=๐ฉ_0^{}\delta ๐ซ_0^{}+๐ฉ_t^{}\delta ๐ซ_t^{}$$
and a similar expression holds for $`\delta S^{ฯต\prime \prime }`$. Due to the $`\mathrm{\Delta }๐ซ_t=0`$ constraint, we have a constraint $`\delta ๐ซ_t^{}=\delta ๐ซ_t^{\prime \prime }`$ on the variation of endpoints, and therefore
$$\delta \left(S^0S^{ฯต\prime \prime }\right)=๐ฉ_0^{}\delta ๐ซ_0^{}+๐ฉ_0^{\prime \prime }\delta ๐ซ_0^{\prime \prime }\mathrm{\Delta }๐ฉ_t\delta ๐ซ_t^{}.$$
Expanding variation $`\delta ๐ซ_t^{}`$ in terms of variations $`\delta ๐ซ_0^{}`$ and $`\delta ๐ฉ_0^{}`$, we find
$`\delta \left(S^0S^{ฯต\prime \prime }\right)`$ $`=`$ $`\left(\mathrm{\Delta }๐ฉ_0\mathrm{\Delta }๐ฉ_t{\displaystyle \frac{๐ซ_t^{}}{๐ซ_0^{}}}\right)\delta ๐ซ_0^{}`$ (11)
$``$ $`\mathrm{\Delta }๐ฉ_t{\displaystyle \frac{๐ซ_t^{}}{๐ซ_0^{}}}\delta ๐ฉ_0^{}+๐ฉ_0^{\prime \prime }\delta \mathrm{\Delta }๐ซ_0.`$
Note again that so far we have not assumed anything about closeness of the two trajectories. Since we can easily shift integration variables $`๐ซ_0`$ and $`๐ฉ_0`$ to $`๐ซ_0^{}`$ and $`๐ฉ_0^{}`$ in Eq. (10), variation (11) indeed tells us where the action difference would be stationary. There are three stationary phase conditions,
$`\mathrm{\Delta }๐ฉ_0\mathrm{\Delta }๐ฉ_t{\displaystyle \frac{๐ซ_t^{}}{๐ซ_0^{}}}`$ $`=0,`$ (12)
$`\mathrm{\Delta }๐ฉ_t{\displaystyle \frac{๐ซ_t^{}}{๐ซ_0^{}}}`$ $`=0,`$ (13)
$`๐ฉ_0^{\prime \prime }\delta \mathrm{\Delta }๐ซ_0`$ $`=0.`$ (14)
The third SP condition was intentionally written in the full form. In general, all three conditions would be satisfied only for a discrete set of trajectories ($`3d`$ equations for $`3d`$ unknowns). However, if the perturbation were $`ฯต=0`$, one could immediately guess that there is one continuous set of solutions satisfying $`\mathrm{\Delta }๐ฉ_0=\mathrm{\Delta }๐ฉ_t=\mathrm{\Delta }๐ซ_0`$. The first two conditions are satisfied exactly, the third one approximately for small variations $`\delta \mathrm{\Delta }๐ซ_0`$. Even though the third condition is satisfied only approximately, we obtain the correct resultโidentical trajectories $`\mathrm{\Delta }๐ซ_\tau =0`$ for all times $`\tau `$, $`0<\tau <t`$โand as we shall see below, also the final result for fidelity will become exact in this limit ($`ฯต=0`$). If we add the perturbation, these precise solutions break down, due to the exponential sensitivity of classical dynamics. However, as was shown in Ref. Vanรญฤek (2004a), if the shadowing theorem Hammel et al. (1987); Grebogi et al. (1990) is applicable in the given system (for a given perturbation $`ฯต`$ and up to time $`t`$), there will be a very near solution with $`\mathrm{\Delta }๐ซ_\tau 0`$ for all times $`\tau `$, $`0<\tau <t`$. Putting off a discussion of the shadowing theorem until later, suffice it to say that this theorem, completely counterintuitively, guarantees that we can compensate one exponential sensitivity (to perturbations of $`H^0`$) by another exponential sensitivity (to initial conditions) and get a trajectory which remains very close to the unperturbed trajectory up to time $`t`$. In fact, these approximate (โdiagonalโ) solutions with $`\mathrm{\Delta }๐ซ_\tau 0`$ will be by far the most dominant ones because for short times no other solutions exist and for long times the diagonal solutions dephase much slower than the remaining (โoff-diagonalโ) solutions with different trajectories. Again this will be justified later in this section. Assuming the validity of shadowing, the โdiagonalโ solutions dephase as
$`S^0S^{ฯต\prime \prime }`$ $`ฯต{\displaystyle _0^t}๐\tau V(๐ซ_\tau ,\tau )\mathrm{\Delta }๐ซ_t๐ฉ_t+\mathrm{\Delta }๐ซ_0๐ฉ_0`$ (15)
$`=\mathrm{\Delta }S_t\mathrm{\Delta }๐ซ_t๐ฉ_t+\mathrm{\Delta }๐ซ_0๐ฉ_0.`$ (16)
The first term is due to the perturbing potential $`ฯตV`$ along the unperturbed trajectory, the other two terms are due to the small difference of trajectories at time $`t`$ and at time $`0`$. Substituting this action difference into integral (10), we obtain the dephasing representation
$`O_{\text{DR}}\left(t\right)`$ $`=`$ $`\left(2\pi \mathrm{}\right)^d{\displaystyle ๐๐ซ_0๐๐ฉ_0๐\mathrm{\Delta }๐ซ_0}`$ (17)
$`\times `$ $`\psi ^{}\left(๐ซ_0+{\displaystyle \frac{1}{2}}\mathrm{\Delta }๐ซ_0\right)\psi \left(๐ซ_0{\displaystyle \frac{1}{2}}\mathrm{\Delta }๐ซ_0\right)`$
$`\times `$ $`\mathrm{exp}\left[{\displaystyle \frac{i}{\mathrm{}}}\left(\mathrm{\Delta }S_t+\mathrm{\Delta }๐ซ_0๐ฉ_0\right)\right].`$
The final result is more succinctly written as
$$O_{\text{DR}}\left(t\right)=๐๐ซ_0๐๐ฉ_0\rho _W(๐ซ_0,๐ฉ_0)\mathrm{exp}\left(i\mathrm{\Delta }S_t/\mathrm{}\right),$$
(18)
using the Wigner function of the initial state $`|\psi ,`$
$`\rho _W(๐ซ,๐ฉ)`$ $`=`$ $`\left(2\pi \mathrm{}\right)^d{\displaystyle ๐\mathrm{\Delta }๐ซ\psi ^{}\left(๐ซ+\frac{1}{2}\mathrm{\Delta }๐ซ\right)}`$ (19)
$`\times `$ $`\psi \left(๐ซ{\displaystyle \frac{1}{2}}\mathrm{\Delta }๐ซ\right)\mathrm{exp}\left(i\mathrm{\Delta }๐ซ๐ฉ/\mathrm{}\right).`$
The general expression (18) expresses fidelity as an interference integral over initial positions $`๐ซ_0`$ and momenta $`๐ฉ_0`$. Because of this property it was called dephasing representation in Ref.Vanรญฤek (2004a). The amplitude of each term is given by the Wigner function $`\rho _W(๐ซ_0,๐ฉ_0)`$ and the phase by the integral of the perturbing potential along the unperturbed trajectory, $`\mathrm{\Delta }S_t(๐ซ_0,๐ฉ_0)`$. This is a very intuitive and simple picture that differs from the simplest โsemiclassicalโ picture only in using the Wigner function instead of the classical phase space distribution $`\rho _{\text{class}}(๐ซ_0,๐ฉ_0)`$.
For zero perturbations, $`ฯต=0`$, expression (18) correctly reduces to the obvious exact result,
$$O_{\text{DR}}^{ฯต=0}\left(t\right)=๐๐ซ_0๐๐ฉ_0\rho _W(๐ซ_0,๐ฉ_0)=1$$
(20)
for all times $`t`$, where the basic property of the Wigner function was used.
Although we started our derivation for a pure state, we ended up with a dephasing representation in terms of the Wigner function. Since this function can also be defined for mixed states, it appears that expression (18) should remain valid for mixed states, with appropriate generalization of the notion of fidelity. In Section III, it will be shown that this is indeed the case.
### II.3 Shadowing theorem and its double use
#### II.3.1 Trajectories of $`H^0`$ and $`H^ฯต`$
Shadowing theorems in general state that (under certain detailed conditions) for small enough $`ฯต`$ there is a time $`t`$ such that for a trajectory of $`H^0`$ with initial condition $`๐ซ_0^{}`$, $`๐ฉ_0^{}`$ there exists a trajectory of $`H^ฯต`$ with initial condition $`๐ซ_0^{\prime \prime }`$, $`๐ฉ_0^{\prime \prime }`$ remaining within a certain small distance from the first trajectory up to time $`t`$. In uniformly hyperbolic systems this shadowing time $`t`$ is infinite Anosov (1967); Bowen (1987), in more general systems at least finite Grebogi et al. (1990). Since it is very difficult to find the maximum shadowing time $`t`$ and the corresponding bound on the closeness of trajectories for a specific system, the derivation of dephasing representation of fidelity assumed that shadowing was applicable for a given perturbation and time: the numerical results will provide the final verification.
#### II.3.2 Numerical evaluation
In order to use DR in numerical applications, one only needs to generate initial conditions $`๐ซ_0`$, $`๐ฉ_0`$ from a distribution given by the Wigner function $`\rho _W`$, run trajectories with the unperturbed Hamiltonian $`H^0`$ and compute the action difference $`\mathrm{\Delta }S_t=ฯต_0^t๐\tau V(๐ซ_\tau ,\tau )`$ along this trajectory. There is no need to compute Van Vleck determinants or Maslov indices as in many other semiclassical applications. Because the Wigner function, unlike classical probability, can be negative, some care must be taken to sample from its distribution. The simplest possible recipe would be to sample according to the probability $`\left|\rho _W\right|`$ and attach a sign afterward together with the dephasing factor. As we will see from the analysis of special cases in Sec. IV, Wigner function is particularly simple for position and momentum eigenstates (just a delta function), for Gaussian wave packets (a Gaussian in both position and momentum), or for a random mixed state (a constant over the whole phase space). These distributions can be easily sampled using standard methods. For general pure or mixed states one can resort to a Monte-Carlo procedure, e. g., using the Metropolis algorithm, which is frequently done for the IVR approximation Miller (2001).
One might object that numerical computation of trajectories, due to the exponential sensitivity of classical evolution, will destroy the validity of the DR (18). However, here the shadowing theorem helps againโin fact in its original form Hammel et al. (1987); Grebogi et al. (1990) where the perturbation was indeed due to errors of numerical propagation. The shadowing idea, as stated in Refs. Hammel et al. (1987); Grebogi et al. (1990) guarantees that for each numerical (noisy) trajectory there will be a nearby exact trajectory of $`H^0`$.
### II.4 Comparison of diagonal and off-diagonal terms
Let us attempt to quantify the validity of the DR by comparing the importance of diagonal and off-diagonal terms in fidelity amplitude. These should be distinguished from the โdiagonalโ and โoff-diagonalโ terms in the fidelity itself (i.e., the amplitude squared), which have been frequently discussed in the literature, see, e. g. Refs. Jalabert and Pastawski (2001); Jacquod et al. (2003) where the off-diagonal terms in the fidelity amplitude are already neglected.
For short enough times $`t`$, it is clear why DR is accurateโthere will be no off-diagonal contributions because there will be no off-diagonal SP solutions of Eqs. (12). For long times, the number of off-diagonal solutions increases, but in the semiclassical limit (small $`\mathrm{}`$) and for small perturbations $`ฯต`$, their contribution is again negligible, due to their much faster dephasing. Let us see in detail how this happens.
If the unperturbed potential is denoted $`W`$, then the off-diagonal solutions dephase as
$$\mathrm{\Delta }S_{\text{off-diag}}=S^0S^{0\prime \prime }=_0^t๐\tau \left[W\left(๐ซ_\tau ^{}\right)W\left(๐ซ_\tau ^{\prime \prime }\right)\right]$$
(21)
because for small enough $`ฯต`$, the perturbation $`V`$ is really unimportant in dephasing of off-diagonal terms. Assuming for simplicity that the off-diagonal terms have the same weight and that their action differences are Gaussian distributed, their average will be
$$e^{i\mathrm{\Delta }S_{\text{off-diag}}/\mathrm{}}\mathrm{exp}\left[\left(\mathrm{\Delta }S_{\text{off-diag}}\right)^2/2\mathrm{}^2\right].$$
(22)
In chaotic systems, $`\left(\mathrm{\Delta }S_{\text{off-diag}}\right)^2=2K_Wt`$, where the diffusion coefficient is $`K_W=_0^{\mathrm{}}๐tC_W\left(t\right)`$ and $`C_W`$ is the potential correlator, $`C_W\left(t\right)=W\left[๐ซ\left(t\right)\right]W\left[๐ซ\left(0\right)\right]`$.
Similar analysis can be done for the diagonal terms Jalabert and Pastawski (2001); Cerruti and Tomsovic (2002); Vanรญฤek and Heller (2003). Their average is then given by a formula analogous to Eq. (22), except with $`\mathrm{\Delta }S_{\text{diag}}`$ given by expression (15). The variance is now given by $`\left(\mathrm{\Delta }S_{\text{diag}}\right)^2=2K_Vฯต^2t`$. Because the diagonal contributions are weighed by the Wigner function, their total contribution is roughly equal to the average. The number of discrete off-diagonal semiclassical contributions should for long times grow as $`e^{\gamma t}`$ where $`\gamma `$ is the topological entropy. Then in the worst possible scenario, where each off-diagonal term contributes by its full weight (as if the Wigner functionโin the case of diagonal termsโwere unity everywhere), the ratio of the sum of the off-diagonal contributions to the total contribution of the diagonal terms should be
$$\frac{\text{off-diag.}}{\text{diag.}}\mathrm{exp}\left\{\left[\left(K_WK_Vฯต^2\right)/\mathrm{}^2+\gamma \right]t\right\}.$$
For small enough $`ฯต`$ and small enough $`\mathrm{}`$, the off-diagonal terms will become negligible. Namely, the diagonal terms will give a smaller contribution if both $`\mathrm{}^2<K_W/\gamma `$ and $`ฯต^2<\left(K_W\mathrm{}^2\gamma \right)/K_V`$.
Similar analysis is possible for integrable systems Prosen and ลฝnidariฤ (2002); Vanรญฤek (2004b). There the number of off-diagonal contributions grows only algebraically, $`t^\alpha `$ and variance of their action difference $`\left(\mathrm{\Delta }S_{\text{off-diag}}\right)^2=C_W^{\mathrm{}}t^2`$ where $`C_W^{\mathrm{}}=lim_t\mathrm{}t^1_0^t๐\tau C_W\left(\tau \right)`$. Similarly, for diagonal terms, $`\left(\mathrm{\Delta }S_{\text{diag}}\right)^2=C_V^{\mathrm{}}t^2ฯต^2`$ Prosen and ลฝnidariฤ (2002). In this case, the ratio of the two types of contributions is
$$\frac{\text{off-diag.}}{\text{diag.}}t^\alpha \mathrm{exp}\left[\left(C_W^{\mathrm{}}C_V^{\mathrm{}}ฯต^2\right)t/\left(2\mathrm{}^2\right)\right]$$
and the condition for negligibility of the off-diagonal terms in the limit $`t\mathrm{}`$ is $`ฯต^2<C_W^{\mathrm{}}/C_V^{\mathrm{}}`$.
## III Dephasing representation for a general mixed state
There are several ways to generalize the pure-state definition (2) of fidelity to mixed states. The simplest generalization is
$$O\left(t\right)=tr\left(e^{iH^0t/\mathrm{}}\rho e^{+iH^ฯตt/\mathrm{}}\right)$$
(23)
where $`\rho `$ is the density matrix of the mixed state, normalized such that $`tr\rho =1`$ Prosen and ลฝnidariฤ (2002). For pure states $`\rho =|\psi \psi |`$, this general definition reduces to the pure-state definition (2). One interpretation of the general expression (23) is that the ket vectors evolve with the unperturbed Hamiltonian $`H^0`$ and the bra vectors with the perturbed Hamiltonian $`H^ฯต`$. Another interpretation is that expression (23) is simply an average of fidelity amplitudes of pure-state components of the given mixed state. This should be distinguished from the often studied averaged fidelity.
The second possible generalization of the notion of fidelity to mixed states replaces the expression for fidelity (1), rather than fidelity amplitude (2) by an expression
$$M\left(t\right)=tr\left[\rho ^0\left(t\right)\rho ^ฯต\left(t\right)\right]=tr\left[\rho (0)\rho (t)\right],$$
(24)
where $`\rho ^0\left(t\right)`$, $`\rho ^ฯต\left(t\right)`$ are the evolved density operators,
$$\rho ^ฯต\left(t\right)=e^{iH^ฯตt/\mathrm{}}\rho e^{+iH^ฯตt/\mathrm{}},$$
or, alternatively, $`\rho (t)`$ is the evolved operator
$$\rho \left(t\right)=e^{+iH^ฯตt/\mathrm{}}e^{iH^0t/\mathrm{}}\rho e^{+iH^0t/\mathrm{}}e^{iH^ฯตt/\mathrm{}}.$$
Again for pure states $`\rho =|\psi \psi |`$, definition (24) reduces to the pure-state definition (1).
Finally there is another, more intuitive but also more complicated generalization, which uses the notion of โpurity fidelityโโthe trace of the squared reduced density matrix Prosen and Seligman (2002),
$$P_F\left(t\right)=\underset{S}{tr}\left[\underset{E}{tr}\rho \left(t\right)\right]^2,$$
(25)
where subscripts $`E`$ or $`S`$ denote that the trace operation is performed on the environment or system degrees of freedom, respectively. For details see Ref. Prosen and Seligman (2002). Purity fidelity (25) does not, of course, reduce to the definition of fidelity for pure states (1).
While dephasing representation expressions are possible for the last two generalizations, in what follows the simplest generalization (23) is assumed. With the mixed-state definition (23), the semiclassical derivation in Eqs. (3)-(18), can be followed closely for mixed states, if we replace the product $`\psi |๐ซ_0^{\prime \prime }๐ซ_0^{}|\psi =\psi ^{}\left(๐ซ_0^{\prime \prime }\right)\psi \left(๐ซ_0^{}\right)`$ in Eqs. (5), (9), and (17) by the matrix element $`๐ซ_0^{}|\rho |๐ซ_0^{\prime \prime }`$ of the density operator. For instance, Eq. (9) will become
$`O\left(t\right)`$ $`=`$ $`\left(2\pi \mathrm{}\right)^d{\displaystyle ๐๐ซ_0๐๐ฉ_0๐\mathrm{\Delta }๐ซ_0๐\mathrm{\Delta }๐ฉ_0\left|\frac{\mathrm{\Delta }๐ซ_t}{\mathrm{\Delta }๐ฉ_0}\right|}`$
$`\times \delta \left(\mathrm{\Delta }๐ซ_t\right)๐ซ_0^{}|\rho |๐ซ_0^{\prime \prime }\mathrm{exp}\left[{\displaystyle \frac{i}{\mathrm{}}}\left(S^0S^ฯต\right)\right].`$
At the end, we obtain the same final result (18), only the Wigner function of a pure state (19) must be replaced by the Wigner-Weyl transform of the density operator,
$`\rho _W(๐ซ,๐ฉ)`$ $`=`$ $`\left(2\pi \mathrm{}\right)^d{\displaystyle ๐\mathrm{\Delta }๐ซ๐ซ+\frac{1}{2}\mathrm{\Delta }๐ซ\left|\rho \right|๐ซ\frac{1}{2}\mathrm{\Delta }๐ซ}`$ (26)
$`\times \mathrm{exp}\left(i\mathrm{\Delta }๐ซ๐ฉ/\mathrm{}\right).`$
## IV Special cases
For a *position state* $`|๐`$, $`\psi \left(๐ซ\right)=\delta \left(๐ซ๐\right)`$, the Wigner function (19) is
$`\rho _W^{\text{pos.st.}}(๐ซ,๐ฉ)`$ $`=\left(2\pi \mathrm{}\right)^d{\displaystyle ๐๐ฑ\delta \left(๐ซ+\frac{1}{2}๐ฑ๐\right)}`$
$`\times \delta \left(๐ซ{\displaystyle \frac{1}{2}}๐ฑ๐\right)\mathrm{exp}\left(i๐ฑ๐ฉ/\mathrm{}\right)`$
$`=\left(2\pi \mathrm{}\right)^d\delta \left(๐ซ๐\right).`$ (27)
Substituting Eq. (19) into the general dephasing representation (18), we find
$$O_{\text{DR}}^{\text{pos.st.}}\left(t\right)=\left(2\pi \mathrm{}\right)^d๐๐ฉ_0\mathrm{exp}\left[i\mathrm{\Delta }S_t(๐,๐ฉ_0)/\mathrm{}\right],$$
in agreement with Eq. (1) from Ref. Vanรญฤek (2004a) and with Vanรญฤek and Heller (2003); Wang et al. (2004).
For a *momentum state* $`|๐`$, $`\psi \left(๐ซ\right)=\left(2\pi \mathrm{}\right)^{d/2}\mathrm{exp}\left(i๐๐ซ/\mathrm{}\right)`$, the Wigner function (19) becomes
$`\rho _W^{\text{mom.st.}}(๐ซ,๐ฉ)`$ $`=`$ $`\left(2\pi \mathrm{}\right)^{2d}{\displaystyle ๐๐ฑ\mathrm{exp}\left[i\left(๐ฉ๐\right)๐ฑ/\mathrm{}\right]}`$ (28)
$`=`$ $`\left(2\pi \mathrm{}\right)^d\delta \left(๐ฉ๐\right)\text{,}`$
and the general DR of fidelity (18) reduces to
$$O_{\text{DR}}^{\text{mom.st.}}\left(t\right)=\left(2\pi \mathrm{}\right)^d๐๐ซ_0\mathrm{exp}\left[i\mathrm{\Delta }S_t(๐ซ_0,๐)/\mathrm{}\right].$$
A *general Gaussian wave packet* with average position $`๐`$, average momentum $`๐`$, and position spread $`\sigma `$,
$$\psi \left(๐ซ\right)=\left(\pi \sigma ^2\right)^{d/4}\mathrm{exp}\left[i๐\left(๐ซ๐\right)/\mathrm{}\left(๐ซ๐\right)^2/2\sigma ^2\right]$$
has Wigner function
$`\rho _W^{\text{gen.G.w.p.}}(๐ซ,๐ฉ)=\left(\pi \sigma ^2\right)^{d/2}\left(2\pi \mathrm{}\right)^{2d}{\displaystyle ๐๐ฑ}`$
$`\times \mathrm{exp}\left\{{\displaystyle \frac{i}{\mathrm{}}}\left(๐ฉ๐\right)๐ฑ\left[\left(๐ซ๐\right)^2+\left(๐ฑ/2\right)^2\right]/\sigma ^2\right\}`$
$`=\left(\pi \mathrm{}\right)^d\mathrm{exp}\left[\left(๐ซ๐\right)^2/\sigma ^2\left(๐ฉ๐\right)^2\sigma ^2/\mathrm{}^2\right].`$ (29)
In general the dephasing representation of a Gaussian wave packet is (18) with the Wigner function (29) where we must include dephasing trajectories with varying *both* positions *and* momenta. Only in special cases, such as when the wave packet is strongly *localized in position* (i. e., when $`\sigma \mathrm{}^{1/2}`$), can we make a further simplification by replacement of $`\mathrm{\Delta }S_t(๐ซ_0,๐ฉ_0)`$ by $`\mathrm{\Delta }S_t(๐,๐ฉ_0)`$ in Eq. (18). Then we can compute the $`๐ซ_0`$ integral in Eq. (18) analytically and obtain
$`O_{\text{DR}}^{\text{pos.G.w.p.}}\left(t\right)=\left(\sigma ^2/\pi \mathrm{}^2\right)^{d/2}{\displaystyle ๐๐ฉ_0}`$
$`\times \mathrm{exp}\left[i\mathrm{\Delta }S_t(๐,๐ฉ_0)/\mathrm{}\left(๐ฉ๐\right)^2\sigma ^2/\mathrm{}^2\right],`$ (30)
in agreement with Eq. (8) in Ref.Vanรญฤek and Heller (2003). There the same result was obtained by linearizing the Van Vleck semiclassical propagator about the central trajectory. In Section V it will be shown that the symmetric expression (29) based on the general DR (18) is superior to the specialized form (30). Similarly, if the initial Gaussian wave packet is *localized in momentum* (i. e., when $`\sigma \mathrm{}^{1/2}`$), we can replace $`\mathrm{\Delta }S_t(๐ซ_0,๐ฉ_0)`$ by $`\mathrm{\Delta }S_t(๐ซ_0,๐)`$ and obtain
$`O_{\text{DR}}^{\text{mom.G.w.p.}}\left(t\right)=\left(\pi \sigma ^2\right)^{d/2}{\displaystyle ๐๐ซ_0}`$
$`\times \mathrm{exp}\left[i\mathrm{\Delta }S_t(๐ซ_0,๐)/\mathrm{}\left(๐ซ๐\right)^2/\sigma ^2\right].`$ (31)
For general (non-Gaussian) wave packets, which are nevertheless localized either in position (about $`๐`$) or momentum (about $`๐`$), we can use the general property of the Wigner function
$`{\displaystyle ๐๐ซ\rho _W(๐ซ,๐ฉ)}`$ $`=\left|\psi \left(๐ฉ\right)\right|^2,`$
$`{\displaystyle ๐๐ฉ\rho _W(๐ซ,๐ฉ)}`$ $`=\left|\psi \left(๐ซ\right)\right|^2,`$
and obtain, upon substitution into the general DR (18),
$`O_{\text{DR}}^{\text{pos.w.p.}}\left(t\right)`$ $`={\displaystyle ๐๐ฉ_0\mathrm{exp}\left[i\mathrm{\Delta }S_t(๐,๐ฉ_0)/\mathrm{}\right]\left|\psi \left(๐ฉ_0\right)\right|^2},`$ (32)
$`O_{\text{DR}}^{\text{mom.w.p.}}\left(t\right)`$ $`={\displaystyle ๐๐ซ_0\mathrm{exp}\left[i\mathrm{\Delta }S_t(๐ซ_0,๐)/\mathrm{}\right]\left|\psi \left(๐ซ_0\right)\right|^2}.`$ (33)
Finally, for a completely *random state*, i. e., an incoherent superposition of all pure basis states, the density operator as well as its Wigner function (26) is just a constant (independent of position or momenta), and for a system with a finite phase space volume $`\mathrm{\Omega }`$, the DR becomes
$$O_{\text{DR}}^{\text{random st.}}\left(t\right)=\frac{1}{\mathrm{\Omega }}๐๐ซ_0๐๐ฉ_0\mathrm{exp}\left(i\mathrm{\Delta }S_t/\mathrm{}\right).$$
(34)
It should be pointed out that while names like โpositionโ or โmomentumโ states have been used to describe the special cases, they do not necessarily need to be eigenstates of the usual position or momentum operator. In the case of abstract Hilbert space with a finite basis, โpositionโ states are simply the basis states (called *computational* states in the setting of quantum information, could be, e. g. spin eigenstates), and โmomentumโ states are simply the states defined by the discrete Fourier transform of the original basis states Miquel et al. (2002). In Ref. Miquel et al. (2002), this generalized phase-space representation is used to show that for quite a few interesting operations on computational states, the Wigner function evolves classically. In all these cases, the dephasing representation described in Secs. II-V should be applicable if discrete Wigner function Miquel et al. (2002) is used and simple other modifications are made to account for the finite-size of phase space. In fact, this is done in the numerical examples in the following section.
## V Numerical tests
Now let us apply the theoretical analysis from previous sections to a specific system, the Chirikov standard map. Its advantage is that it is discrete, coordinate space is only one-dimensional, but at the same time standard map already contains generic complexities of classical dynamics. Specifically, the phase space is mixed and so various simplifications applicable in quasi-integrable or strongly chaotic systems are in general not applicable. Standard map is a symplectic map defined on a compact two-dimensional phase spaceโtorus, as follows,
$`q_{j+1}`$ $`=q_j+p_j\text{ (mod }2\pi \text{)}`$
$`p_{j+1}`$ $`=p_jW^{}\left(q_{j+1}\right)ฯตV^{}\left(q_{j+1}\right)\text{ (mod }2\pi \text{)},`$
where $`q`$ and $`p`$ are position and momentum on the torus, potential $`W\left(q\right)=k\mathrm{cos}q`$, and the perturbation is $`V\left(q\right)=\mathrm{cos}2q`$. Using an $`n`$-dimensional Hilbert space for the quantized map fixes the effective Planck constant to be $`\mathrm{}=\left(2\pi n\right)^1`$. (We are using letter $`q`$ for the coordinate to distinguish this special system from the general considerations. Similarly, we will use letter $`Q`$ to denote the position of a position state or center of a wave packet.) Parameter $`ฯต`$ controls the strength of perturbation. For $`ฯต1`$, the map is close to being integrable, for $`ฯต1`$, the map is strongly chaotic. The goal of this section is not to use the dephasing representation to explore various universal regime that occur in these two limits and have been carefully studied in the literature. This was already done in Refs. Vanรญฤek and Heller (2003); Vanรญฤek (2004b). The goal of this section is rather to explore the detailed features of fidelity in non-universal regimes. The optimal region of parameter space is in the vicinity of $`ฯต=1`$, since there phase space has a significant amount of chaotic as well as integrable regions. Mixed phase space is in general the hardest to treat and therefore this setting is chosen here because it provides the most challenging test for any approximation.
### V.1 Gaussian wave packets
One might think that the general dephasing representation (18) is only useful for highly non-local states and that the original expression (27) from Ref. Vanรญฤek and Heller (2003) is good enough at least for Gaussian wave packets. This subsection demonstrates that even for Gaussian wave packets, the general dephasing representation (18) is superior to the original expression (27) from Ref. Vanรญฤek and Heller (2003).
Figure 1 compares three approximations to compute fidelity of Gaussian wave packets with the exact result: the expression (27) from Ref. Vanรญฤek and Heller (2003) for wave packets localized in position (red dashed line), corresponding expression (31) for wave packets localized in momentum (blue dotted line), and the general DR (18), with the Wigner function (29), symmetrically treating position and momentum (black solid line).
The exact fidelity, computed by exact quantum evolution using the Fast Fourier Transform algorithm, is represented by solid dots. The parameters are $`n=1000`$, $`k=0.95`$, $`ฯต=0.015`$, and the wave packet is localized at $`Q=0.7\pi `$ and $`P=0.4\pi `$. The number of classical trajectories used in the calculations is $`1000`$. Wave packets used in parts a), b), and c) of Fig. 1 have position spread $`\sigma `$ equal to 0.004$`\pi `$, 0.16$`\pi `$, and 0.04$`\pi `$, respectively. For a wave packet localized in position in Fig. 1a), the original expression (27) from Ref. Vanรญฤek and Heller (2003) works very well and is almost indistinguishable from the general DR (18), as expected, whereas Eq. (31) for momentum wave packets fails. For a wave packet localized in momentum in Fig. 1b), the momentum-wave-packet expression (31) works well and it is almost indistinguishable from the general DR (18), but the original position-wave-packet expression (27) from Ref. Vanรญฤek and Heller (2003) fails completely. The general DR works very well in both cases. It might seem that either the momentum or position versions could cover the whole range of Gaussian wave packets, because one might think that the intermediate case, i. e., a fairly symmetric wave packet, is localized enough in both position and momentum. That this is not so is provided by the final test in Fig. 1c): both specialized expressions (30) and (31) give a significant error in comparison with exact fidelity, but the general DR (18) gives very accurate results, as expected because of its โfairโ treatment of position and momentum. To conclude, expression (18), is accurate for the whole range of Gaussian wave packets, from position-like to symmetric to momentum-like, even in the presence of mixed dynamics.
### V.2 Nonlocal states
For nonlocal states, there is even less hope that the position-wave-packet expression for fidelity (30) from Ref. Vanรญฤek and Heller (2003) would work. One might think that for a superposition of localized wave packets it is enough to simply add the terms (30) for the fidelity amplitude. This is not the case which can be seen by considering a wave packet $`\psi `$ that is a superposition of two Gaussian wave packets $`\psi _1`$ and $`\psi _2`$, centered at phase space points $`(๐_1,๐_1)`$ and $`(๐_2,๐_2)`$. The resulting wave packet has a Wigner function that is not just a simple sum of the Wigner functions of the two Gaussian wave packets. The correct Wigner function has in addition an *interference term* localized in the vicinity of the phase-space point $`((๐_1+๐_2)/2,(๐_1+๐_2)/2)`$. We will demonstrate now the importance of this interference term and show that if it is taken into account, the general DR (18) will still give excellent results, even for nonlocal states.
Being motivated by the quantum computation applications, let us consider a superposition of computational states (i. e., position states in the abstract phase space), instead of Gaussian wave packets. Our initial state is a *coherent* superposition
$$|\psi =\frac{1}{\sqrt{2}}\left(|๐_1+|๐_2\right),$$
(35)
with a Wigner distribution,
$`\rho _W^{\text{coh}}(๐ซ,๐ฉ)={\displaystyle \frac{1}{2}}\left(2\pi \mathrm{}\right)^d\{\delta (๐ซ๐_1)+\delta (๐ซ๐_2)`$
$`+2\delta [๐ซ(๐_1+๐_2)/2]\mathrm{cos}[(๐_1๐_2)๐ฉ/\mathrm{}]\}.`$ (36)
If the interference term is neglected, we obtain a Wigner function of the *incoherent* superposition (38),
$$\rho _W^{\text{incoh}}(๐ซ,๐ฉ)=\frac{1}{2}\left(2\pi \mathrm{}\right)^d\left[\delta \left(๐ซ๐_1\right)+\delta \left(๐ซ๐_2\right)\right]$$
(37)
Figure 2 compares two approximate ways to compute fidelity with the exact quantum result: both approximations use the general DR (18), but whereas one uses the correct full Wigner function (36) (black solid line), the other uses the incorrect Wigner function (37), neglecting the interference term (purple dashed-dotted line). Again, the exact result is represented by solid dots. The parameters used in Fig. 2 are $`n=200`$, $`k=0.7`$, $`ฯต=0.02`$, $`Q_1=0.4\pi `$, and $`400`$ classical trajectories were used. The position of the other component state varies in the two parts.
If the positions $`๐_1`$ and $`๐_2`$ are largely separated, the oscillations in the interference term have a high frequency. Because nearby initial conditions follow similar trajectories and have similar actions, the phase factor in the DR (18) varies slowly. Therefore the fast oscillations in the weight factor given by the interference term in the Wigner function can completely cancel out the contribution of the interference part to the DR integral. (Incidentally, this situation is in a way opposite to the usual semiclassical considerations where the weight is a slowly varying function and the phase factor is the fast oscillating factor.) Figure 1a) shows an example of situation where this cancellation occurs: $`Q_1=0.4\pi `$ and $`Q_2=1.2\pi `$. Because the interference term is negligible, both approximations give the same and very accurate result.
If the initial states are closer, as in Fig. 2b), where $`Q_1=0.4\pi `$ and $`Q_2=0.42\pi `$, the interference term is important, and only the correct Wigner function (36) agrees well with the exact result. This shows that for coherent nonlocal states, the general DR (18) must be used instead of some approximate versions which neglect quantum coherence of the initial state.
### V.3 Mixed states
Wigner function (37) was wrong for the coherent state (35), but it does correctly describe a certain mixed state, namely the incoherent superposition of computational states $`|๐_1`$ and $`|๐_2`$,
$$\rho ^{\text{incoh}}=\frac{1}{2}\left(|๐_1๐_1|+|๐_2๐_2|\right).$$
(38)
In Sec. III it was shown that if the generalized definition (23) of fidelity for mixed states is used, dephasing representation (18) remains valid, as long as the Wigner transform of the density operator (26) is used. For the incoherent mixture with density operator (38), Wigner distribution is precisely that given by Eq. (37). Figure 3 compares DR (18) with the Wigner function (37) with the exact fidelity for the state (38). Fidelity computed by the DR is drawn with a black solid line, exact fidelity with solid dots. The parameters are the same as in Fig. 2b), in particular $`Q_1=0.4\pi `$ and $`Q_2=0.42\pi `$. Although now only $`200`$ classical trajectories were used, the agreement is again excellent.
Last but not least we consider the completely random mixed state. It is an incoherent superposition of all computational states and in a finite-dimensional Hilbert space, its density operator is
$$\rho ^{\text{random}}=\frac{1}{n}\underset{i=1}{\overset{n}{}}|Q_iQ_i|=\frac{1}{n}\widehat{1}.$$
Figure 4 compares the random-state version (34) of DR (black solid line) with the exact result (solid dots). Parameters in this calculation are $`n=100`$, $`k=2`$, $`ฯต=0.03`$ and $`1000`$ classical trajectories were used. Again, it is reassuring that even in the case that the whole phase space is important, with just $`1000`$ trajectories, dephasing representation still works so wellโdespite the fact that it was derived solely from semiclassical arguments and requires only classical information.
## VI Relation to other โWignerโ methods
It should be noted that the Wigner distribution has been used in various other approximate methods, especially in chemical physics. For instance, it was used to compute photodissociation cross-sections Heller (1976); Brown and Heller (1981), to treat inelastic scattering Lee and Scully (1980), or to compute thermal correlation functions using the linearized semiclassical IVR method Miller (1974); Wang et al. (1998); Miller (2001). In all these applications, there was just one Hamiltonian, but the two states (or more generally, density or other operators) were different. The quantity of interest was a general correlation function of the type
$$C_{AB}\left(t\right)=tr\left(AU^{}BU\right)$$
(39)
where $`A`$ and $`B`$ are general operators and $`U=๐ฏe^{i{\scriptscriptstyle H๐\tau }}`$ is the time evolution operator. Using various approximations, all authors Heller (1976); Brown and Heller (1981); Lee and Scully (1980); Miller (1974); Wang et al. (1998); Miller (2001) obtain the same final result, expressed as an overlap of two Wigner distributions, one at time 0, the other evolved classically to time $`t`$,
$$C_{AB}^{\text{Wigner}}\left(t\right)=(2\pi \mathrm{})^d๐๐ซ_0๐๐ฉ_0A_W(๐ซ_0,๐ฉ_0)B_W(๐ซ_t,๐ฉ_t)$$
(40)
Here $`A_W`$ and $`A_W`$ are the Wigner transforms (26) of operators $`A`$ and $`B`$.
Because there is only one Hamiltonian, there is no dephasing factor $`e^{i\mathrm{\Delta }S/\mathrm{}}`$, as in the DR. In fact we could apply one of these older approaches to the second generalized definition (24) of fidelity for mixed states because that definition is in the form of Eq. (39) with $`A=B=\rho `$ and the time evolution operator $`U=e^{+iH^ฯตt/\mathrm{}}e^{iH^0t/\mathrm{}}`$. Then we would obtain a very different result from the DR,
$$M^{\text{Wigner}}\left(t\right)=(2\pi \mathrm{})^d๐๐ซ_0๐๐ฉ_0\rho _W(๐ซ_0,๐ฉ_0)\rho _W(๐ซ_t,๐ฉ_t).$$
(41)
Although appearing as elegant as the dephasing representation, there is a problem with this expression. First, it will be much more sensitive to numerical errors. We can see that already by considering zero perturbation. Correctly, for each initial condition $`๐ซ_0,๐ฉ_0`$, we should have $`๐ซ_0=๐ซ_t`$ and $`๐ฉ_t=๐ฉ_0`$. In systems with nonlinear dynamics, particularly chaotic systems, numerical errors in forward and backward propagation will yield exponentially growing errors. If the initial state is a localized wave packet, expression (41) would give a numerically decaying overlap even for zero perturbations when exact fidelity is constant $`M(t)=1`$. Indeed, numerical test not presented here showed that instead of staying at unity, $`M^{\text{Wigner}}`$ quickly decays to a plateau and remains there for some time, and finally decays exponentially again. (This is the same behavior as observed in literature for physical perturbations Prosen and ลฝnidariฤ (2005); Bevilaqua and Heller (2004).
Even if numerical errors did not exist, equation (41) would have problems. It can describe some decay due to dephasing, but only that in the fast oscillating parts of the initial state. For simple Gaussian wave packets, the fidelity decay in Eq. (41) is completely due to the decay of classical overlaps, i. e., classical fidelity. To conclude, the โWignerโ form (41) is apparently not as good as the dephasing representation, but it does deserve further study, especially because it might shed further light on the question of importance of various contributions to fidelity. Preliminary studies show that $`M^{\text{Wigner}}`$ correctly describes exact fidelity in both chaotic and quasi-integrable systems for large perturbations (i. e., in Lyapunov and algebraic regimes, respectively), when dephasing is not important Vanรญฤek (2004b). It gives wrong results in both chaotic and quasi-integrable systems for small perturbations (in the FGR and Gaussian regimes), when dephasing is important Vanรญฤek (2004b).
## VII Conclusion
This paper has presented a derivation of a general semiclassical expression for fidelity of pure and mixed states. This dephasing representation expresses fidelity as an interference integral, with weight of each term given by the Wigner function and the phase by the integrated perturbation along an unperturbed trajectory. In particular, no analog of the Van Vleck determinant is needed. As the original specialized expression (30) from Ref. Vanรญฤek and Heller (2003), dephasing representation avoids searching for the exponentially growing number of terms in the standard semiclassical expressions Jalabert and Pastawski (2001). It also avoids the ubiquitous divergences in Van Vleck determinants present in the usual semiclassical expressions.
The advantage of dephasing representation lies in that it does not require the original state to be localized. Its form suggests that it should be applicable to general pure and mixed states. This claim was supported by the following numerical evidence: First, it was shown, on the example of Gaussian wave packets, that position and momentum must be treated symmetrically. This was the flaw of the expression from Ref. Vanรญฤek and Heller (2003) and is apparently corrected in the DR. Second, on the example of coherent superpositions of states, it was shown that oscillatory patterns in the Wigner function are important: therefore classical phase space distribution, resulting from incoherent superposition of component Wigner distributions (for states for which these are the same as classical distributions). This may shed some further light on the controversial issue of importance of sub-Plank structures on decoherence Zurek (2001); Jacquod et al. (2002). Finally, it was shown that DR is also accurate for mixed states: incoherent superpositions and completely random states. All tests were performed on a system with mixed phase space: with both integrable and chaotic regions.
While the numerical tests were quite successful, a further study is needed to determine precisely all situations where the dephasing representation breaks down. The analysis provided in Sec. II of this paper should simplify that task. Also, a more rigorous formulation of the precise conditions of validity of the dephasing representation is needed.
###### Acknowledgements.
The author wishes to thank the Department of Chemistry at the University of California, Berkeley for support.
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# Cosmic-ray strangelets in the Earthโs atmosphere
## I Introduction
A dense ball of up, down, and strange quarks, called a โstrangeletโ, may be the true ground state of nuclear matter. A strangelet is stabilized by its filled quark energy levels; the extra flavor degree of freedom provided by the strange quark lowers the Fermi energy of the up and down quark seas. The strangelet may be absolutely stable if this decrease in Fermi energy compensates for the extra mass of the strange quark; this may be true above some threshhold baryon number $`A_t`$. (Throughout this paper, we will speak as if $`A_t`$ is โsmallโ, in the range 50โ4000, and therefore that โsmallโ strangelets with are stable. Of course, this is the unknown we wish to study.)
Small stable strangelets have approximately equal numbers of u, d, and s quarks. A slight deficiency of s quarks gives the strangelet a net positive charge of approximately $`Z=0.3A^{2/3}`$ for the color-flavor locked (CFL) model or $`Z=0.1A`$, turning over to $`8A^{(1/3)}`$ for large A, for the MIT bag model. (We discuss both models with a strange quark mass of 150 MeV.) Ordinary nuclei cannot decay into strangelets in the the lifetime of the Universe. However, a neutron star may quickly convert into a โstrange starโ, essentially a large strangelet with $`A10^{50}`$. Strange star collisions may release smaller strangelets with a range of sizes, including low-mass strangelets with charges $`Z`$ where $`Z_{min}Z<100`$ and baryon number $`A`$ where $`A_tA<10^5`$. These strangelets are accelerated by astrophysical shocks and propagate throughout the Galaxy. Therefore, if small strangelets are stable at all, we expect to find them among energetic cosmic rays and in cosmic-ray targets like the Earth and Moon. (Although this same cosmic-ray flux would also contribute strangelets to Earthโs progenitor materialโthe presolar nebula, etc.โwe neglect this contribution, and aim for a conservative estimate based on the flux over the Earthโs history.)
At rest, a strangelet of charge $`Z`$ behaves like a heavy isotope of atomic number $`Z`$; it has the same number of electrons, chemical and thermodynamic behavior, etc.. Thus, we give a โPeriodic Table of Strangeletsโ in Figure 2. Its chemistry is modified by isotope mass effects which are discussed in section IV.2.2 and considered to be small; we will neglect these effects for all discussions of geology and atmospheric chemistry.
Strangelets can be studied by a mass-spectroscopic search for rare, ultra-heavy isotopes of ordinary elements. In order to search for strangelets with high sensitivity, we must find samples where the strangelet concentration is as high as possible. Because we can purify a sample chemically before doing mass spectroscopy, strangelets of charge Z will be mixedโor diluted, or contaminatedโonly with nuclei of charge Z. The chemical abundances of strangelets are uncorrelated with the abundances of ordinary elements on Earth; therefore, it is helpful to search for strangelets among rare elements. To illustrate, calcium (Z=20) is abundant and scandium (Z=21) is rare, on Earth, due to details of nuclear structure and nucleosynthesis. We have no reason to expect that โstrange scandiumโ (<sup>s</sup>Sc) is any less abundant than <sup>s</sup>Ca. Broadly speaking, then, a search for <sup>s</sup>Sc in a Sc sample is more likely to find high concentrations than a search for <sup>s</sup>Ca in a Ca sample.
Several searches have been done along these lines, most commonly in meteorites, Earth rock, and moon rock. However, rock of any sort is a very dirty environment: arriving strangelets are immediately mixed with large numbers of ordinary nuclei of all charges. Before strangelets reach the Earthโs crust, they stop in a very clean environment: the atmosphere. Studying the chemical and circulation properties of the atmosphere suggests new places to search for small stable strangelets. Surprisingly high concentrations can be found in the atmosphere itself. In this paper, we suggest a series of atmospheric strangelet searches which together cover a wide range of allowable strangelet charges and masses, in some cases making it possible to probe astrophysical models using existing mass spectroscopy techniques.
### I.1 Expected strangelet flux
We derive abundance estimates from the detailed cosmic-ray propagation model of . This work assumes a strangelet production rate of $`10^{10}`$ solar masses (M) per year in our Galaxy, with $`10^5`$M of strange matter released per collision, and a collision every 30,000 y. It includes the effects of acceleration, interstellar propagation, and solar modulation using several phenomenological models. If we assume that all of the released strange-star material is converted into strangelets of charge Z and and mass A (GeV/c<sup>2</sup>), we derive the maximum flux F
$$F=2\times 10^5\mathrm{m}^2\mathrm{y}^1\mathrm{ster}\times A^{0.467}Z^{1.2}R_{\mathrm{cutoff}}^{1.2}$$
(1)
where $`F`$ is the flux in particles m<sup>-2</sup> ster<sup>-1</sup> y<sup>-1</sup>, and $`R_{\mathrm{cutoff}}`$ is the either the geomagnetic cutoff or the solar-modulation cutoff, whichever is greater. Notice that, although a factor of $`A^1`$ is expected when a finite pool of strange matter is partitioned into strangelets, heavier strangelets propagate more efficiently, giving the shallow A dependence. Equation 1 gives an upper limit on the flux because it assumes that *all* of the produced strangelets have the same charge. More realistically, we expect strange star collisions to produce a wide range of strangelet charges, thereby lowering the number of strangelets with any particular charge.
We compute the average flux at Earth by integrating Equation 1 over the Stormer dipole approximation for the geomagnetic cutoff, and approximating the solar modulation cutoff as $`R=(A/Z)^{1/2}(\mathrm{\Phi }/500MV)^{1/2}`$, where the solar modulation parameter $`\mathrm{\Phi }`$ is taken to be 500 MV. These fluxes are given in Table 1 and Figure 1.
## II Strangelets on the Earth, moon, and meteorites
As a basis for comparison, we review the expected abundances of strangelets in the Earthโs crust, on the moon, and in meteorite material.
The moon is the simplest case. Most of the moonโs surface, called the โhighlandsโ, has been essentially stationary for about $``$4 billion years. Its surface has been turned over and reburied to a depth of $``$10 meters by the action of meteor impacts. It has no substantial magnetic field. If conservatively assume that the surface turnover/reburial rate is constant over time, then the expected strangelet per-atom concentration in the lunar highlands is $`c_{\mathrm{moon}}=F\times 3\times 10^{20}`$, where $`F`$ is the Earth-averaged annual flux per m<sup>2</sup> y from equation 1. For example, for <sup>s</sup>O (Z=8, A=130) the concentration is $`2\times 10^{16}`$. The turnover rate has been slower in the modern epoch; the top 10cm of soil may have been exposed, at many sites, for 500 My, leading to a proportionally higher concentration.
Meteorites are more complex as strangelet targets; cosmic-ray exposure measurements of most meteorite samples show quite short ages $`T_{CRE}`$, often of order 10 My, suggesting that the sampled material spent most of its life โburiedโโand shielded from cosmic ray exposureโinside a larger object. Furthermore, much of this exposure is due to cosmic-ray protons, which penetrate much more deeply than high-Z strangelets. If we ignore these and other complications<sup>1</sup><sup>1</sup>1For example, a meteoroid in the asteroid belt would encounter a lower solar modulation cutoff., then the strangelet abundance is roughly proportional to the cosmic-ray exposure age, $`c_{\mathrm{meteor}}=F\times 10^{23}\times \frac{T_{CRE}}{1My}`$.
The Earthโs crust and oceans are a particularly complex environment, some details of which we will address in a forthcoming paper. Over the Earthโs 4 billion year history, strangelets have come out of the atmosphere and been deposited on land and in the ocean. Initially, strangelets are probably incorporated into sedimentary rock. This rock can be buried, subducted, metamorphosed, uplifted, and eroded by various processes, while new (possibly strangelet-free) rock wells up from the mantle below. For a rough estimate, we may neglect the details and suppose that the Earthโs strangelet inventory is evenly mixed through all sedimentary rocks. The average depth of sedimentary material is about 2 km over the Earthโs surface, and this layer has re-cycled with the mantle about 5 times. This suggests that Earthโs integrated strangelet flux is evenly mixed with a 10km-deep column of rock, where the strangelet abundance is approximately $`c_{\mathrm{crust}}=F\times 1.5\times 10^{23}`$.
## III Strangelets in the atmosphere
A cosmic ray strangelet has less penetrating power than a nucleus of the same $`Z`$ and energy<sup>2</sup><sup>2</sup>2Stopping distances can be calculated with SRIM, ยฉ2000 by James Zeigler, http://www.srim.org due to a lower velocity and higher dE/dx. Therefore, most strangelets will lose energy by ionization and stop the mesosphere and upper stratosphere, 50โ80 km above the surface. After stopping, their fate is determined by atmospheric chemistry and dynamics. We will consider three broad classes of chemical species the atmosphere: noble gases, volatile strangelets, and metallic strangelets. These categories are defined in Figure 2.
### III.1 Effects of atmospheric circulation
Several details of atmospheric physics are particularly relevant to the behavior of metallic strangelets.
The mesosphere is mixed vertically by convection, with a timescale of order $``$1 y, as well as north-south in a โsloshingโ motionโtowards the poles in winter, towards the equator in summer. Since most cosmic rays stop in the mesosphere and must participate in this circulation, we expect that strangelets do *not* descend through the atmosphere primarily at the geomagnetic poles. Before crossing from the mesosphere to the stratosphere, they redistribute somewhat evenly over all latitudes. The extent of the redistribution is unclear.
Below the mesosphere is the stratosphere ($`10`$$`50`$ km). The stratosphere does not convect, since the air temperature increases with altitude; metal atoms and clusters injected at the top of the stratosphere may take $`5`$ y to descend to the base of the stratosphere, called the โtropopauseโ. Other than meteoritic and cosmic-ray metals, the stratosphere is extremely clean; water vapor<sup>3</sup><sup>3</sup>3The lower stratosphere has typically $`10^6`$ H<sub>2</sub>O., dust, and soot are nearly absent. This is due to the lack of air circulation across most of the tropopause. Ground-level air, carrying terrestrial dust, enters the stratosphere only in the tropics; tropical rainfall removes most of these contaminants as the air rises. This air flows gradually towards the poles, in particular towards the winter pole. Due to this pattern, we expect condensible material in the stratosphere, including metallic strangelets, cosmic ray metals, and meteorite dust, to accumulate in the polar regions towards springtime; we discuss this point further below. In the spring, โtropopause foldsโ, primarily at midlatitudes, mix lower-stratospheric material back into the troposphere.
Air circulation in the troposphere is rapid and complex. Complete vertical turnover of the troposphere occurs on a scale of weeks; we note that this precludes any gravitational sorting of gaseous strangelets. Water vapor contents are high ($`0.01`$$`10^4`$ by volume), and dust and soot levels are high and higly variable. Non-gaseous materials in the troposphere are mostly scavenged by rain and snow.
### III.2 Strange noble gases
A strangelet with charge 2, 10, 18, 36, 54, or 86 will have the chemical properties of helium, neon, argon, krypton, xenon, or radon respectively. These strange noble gas atoms will stay in the atmosphere forever unless a) lost into space or b) dissolved in water or rock and subducted. Both removal effects are minor, so we expect that the atmosphere has accumulated most of the strangelets from 4.5 billion years of CR flux. Thus, the atom abundance of a noble-gas-like strangelet in the bulk atmosphere is $`c_{\mathrm{gas}}=F\times 3\times 10^{20}`$. This concentration factor is as high as that of the moon.
As discussed in Section I, the concentration of strangelets in a pure element sample is inversely proportional to the normal elementโs abundance. In particular, He and Xe are quite scarce in the atmosphere, so <sup>s</sup>He and <sup>s</sup>Xe nuclei may have high concentrations<sup>4</sup><sup>4</sup>4In the Earthโs atmosphere, helium-like strangelets are gravitationally bound while ordinary <sup>4</sup>He is not. This phenomenon means that the Earthโs <sup>s</sup>He/<sup>4</sup>He ratio is greatly enhanced relative to the *primordial* <sup>s</sup>He/<sup>4</sup>He ratio. However, in terms of cosmic-ray accumulation, the Earthโs <sup>s</sup>He/<sup>4</sup>He ratio is, as for all other noble gases, simply the integrated flux of CR <sup>s</sup>He divided by the total modern amount of background <sup>4</sup>He. The fact that the modern <sup>4</sup>He reservoir turns over rapidly is immaterial, as long as we recognize that CR <sup>s</sup>He, once collected, does not turn over.. A <sup>s</sup>Rn atom in the atmosphere would be extremely unique, because all ordinary Rn is short-lived. The atmospheric abundances and strangelet concentrations for each noble gas are summarized in Table 1.
### III.3 Strange volatile elements
Several elements, which we label โvolatilesโ, can cycle back and forth between the atmosphere, ocean, and crust: H, C, N, O, S, and the halogens. A volatile strangelet resides in the atmosphere for a time longer than the atmospheric mixing time, but shorter than the age of the Earth; volatile strangelets on the ground may be re-volatilized and reenter the atmosphere. In the case of H, C, S, we do not expect strangelets to accumulate in the atmosphere; turnover is fairly rapid and there is much dilution by ordinary matter, particularly H<sub>2</sub>O, CO<sub>2</sub>, and SO<sub>X</sub>. The halogens case for halogens is quite similar; all halogens are present in the atmosphere, both naturally and as pollutants: fluorine as manmade fluorocarbons; sea-spray HCl in the troposphere, natural CH<sub>3</sub>Cl and manmade CFCs in the stratosphere; bromine as BrO over oceans, and CH<sub>3</sub>Br, CHB<sub>3</sub>, and halons throughout; iodine as various iodocarbons. The fully-halogenated species have atmospheric lifetimes of order 100 y; all other species have short lifetimes, usually $`<`$ 1 y, so we do not expect halogen-like strangelets to be at all abundant in atmospheric samples.
Nitrogen is more interesting. Most of the Earthโs entire nitrogen inventory ($`4\times 10^{18}`$ kg) is in the atmosphere as N<sub>2</sub>. Although nitrogen cycles back and forth into solid forms, mainly due to biological fixation, at a rate of $`10^{11}`$ kg/y, it is very slow ($`10^{10}`$ kg/y) to be sequestered in the crust<sup>5</sup><sup>5</sup>5Nitrogen sequestration is mainly via deposition (diatom shells, etc.) on the ocean floor. Interestingly, industrial-scale agriculture has nearly doubled the global fixation budget over the past few decades.. (Much of this, too, eventually reenters the atmosphere via uplift and erosion.) Although it is difficult to project the nitrogen cycle into the distant geological past, the modern nitrogen cycle suggests that the lifetime of a <sup>s</sup>N atom in the atmosphere is of order $`5\times 10^8`$ years. The effective lifetime may be much longer, depending on how much sequestered nitrate is eventually re-exposed rather than subducted. Conservatively, taking the net <sup>s</sup>N flux of the past $`5\times 10^8`$ y to remain in the atmosphere, we expect a <sup>s</sup>N concentration of $`5\times 10^{18}`$.
Oxygen is one of the most abundant elements on Earth, making up much of the mass of the biosphere, oceans, crust, and mantle in the form of water, metal oxides, and organic molecules. However, the most likely chemical fate of an an arriving <sup>s</sup>O atom is to be incorporated into O<sub>2</sub>. O<sub>2</sub> has an atmospheric residence time of $`3\times 10^6`$ y, suggesting a <sup>s</sup>O abundance of $`10^{19}`$ in atmospheric O<sub>2</sub><sup>6</sup><sup>6</sup>6Aerobic respiration turns atmospheric O<sub>2</sub> into H<sub>2</sub>Oโthe oxygen in other biological compounds, like glucose, is derived from H<sub>2</sub>O and CO<sub>2</sub>โso high <sup>s</sup>O concentrations are not expected in the biosphere.
### III.4 Metallic strangelets in the stratosphere
Most possible strangelet charges correspond to metallic elements. A metallic strangelet stopped in the mesosphere will mix immediately with the metals left behind by vaporizing micrometeors. These metals agglomerate into nanometer-scale โsmokeโ clusters within a few months of arrival. This smoke has been observed to drift down into the lower stratosphere and collect in sulfuric acid aerosol droplets; 1/2 of such droplets found at midlatitudes carry 0.5โ1% by mass of dissolved metal, with approximately chondritic composition. From this, and the known flux of sulfur to the stratosphere (100โ160 $`\times 10^6`$ kg/y), the flux of vaporized meteoritic material to the stratosphere is inferred to be 4โ19$`\times 10^6`$ kg/y.<sup>7</sup><sup>7</sup>7Note that this measurement specifically probes the amount of meteoritic material which is left in the stratosphere, not the total amount which reaches the Earth.
For further discussion, let us consider $`10^7`$ kg/y (roughly the middle of the plausible range) to be the meteorite flux. Of this amount, approximately 44% is oxygen, and the remainder amounts to $`1.7\times 10^8`$ moles of metals entering the stratosphere per year<sup>8</sup><sup>8</sup>8Three other sources of metal are available to the stratosphere. Major volcanic eruptions inject ash into the lower stratosphere, where it remains for $`1`$ year. Recent eruptions of this power were Mt. St. Helens (1981), El Chichon (1982), Nevado del Ruiz (1985), Mt. Augustine (1986), and Mt. Pintubo (1991). Large aboveground nuclear tests may also affect the stratosphere, but we hope these are no longer an issue. High-altitude aircraft, including the Concorde, and rockets also contribute.. Thus, the atom abundance of strangelets mixed in with these metals is $`c_{\mathrm{strat}}=F\times 1.5\times 10^{17}`$ where F is the flux from equation 1. This is a factor of 1000 higher than the concentrations on the moon, and $`10^6`$ higher than concentrations in the Earthโs crust. Due to this high concentration, stratospheric aerosols may be a target for a strangelet search, discussed below.
We also note that the stratospheric concentration factor provides an upper limit, albeit a weak one, for cosmic-ray strangelet concentrations in rainfall or in geological samples; no Earth environment is exposed to metallic strangelets without also being exposed to micrometeorite smoke.
### III.5 Thermodynamic and chemical effects of strangelet mass
We might be concerned that the unusually large isotope shifts affecting strangelets could change their chemistry or thermodynamic behavior in such a way as to invalidate our atmospheric residence times, or else to defeat some sample collection techniques. These effects can be estimated, albeit with considerable uncertainty. From such a calculation, we conclude that the errors in our analysis due to isotope effects are less than a factor of 2, and thus much smaller than the other (e.g., astrophysical) uncertainties relevant to a strangelet search.
The effect of isotopic mass on vapor pressure is discussed in section IV.2.2; the vapor pressure shift is very similar in principle to shifts in chemical equilibria, solubility, etc.. We show in Table 2 that strange gases should have only slightly depressed vapor pressures. Although these numbers are approximate, we may conclude that strangelets are not removed from noble gases by commercial distillation plants<sup>9</sup><sup>9</sup>9Strangelets may vary in concentration *within* a particular noble gas fraction in a column; however, the normal operation of a commercial distillation plantโin which the โwasteโ gas from a given stage is not discarded, but rather re-injected at an earlier stageโshould not maintain this separation.. We can also conclude that strangelets do not behave significantly differently than ordinary gases in nature. Ordinary isotope effects are sometimes observed to cause 1%-2% variations in natural isotope abundances; since strangelet vapor pressure shifts are only a factor of O(10) larger, and natural processes do not typically cascade or amplify these effects, the variation in strangelet abundances between natural samples should be of order 10โ50% at most.
## IV Aspects of strangelet searches
Now that we have identified these environments where strangelets may be present at high concentrations, we can design experiments to search for them. The three steps in any search are: sample collection, purification/preconcentration, and mass spectroscopy.
We would like to emphasize that we do *not* know, if strangelets are stable: what is the minimum mass $`A_t`$? What is the exact charge/mass relation? What is the distribution of masses produced by strange star collisions? What is the mass distribution for a sample with a given Z? In order to comprehensively address the question of the existence of strangelets, we should perform experiments over a wide range of charges, and each search should be sensitive to a wide range of possible masses. Although some particular searches stand out for their high sensitivity, we wish to investigate more of the periodic table in order to cover the search space completely.
### IV.1 Past atmospheric strangelet searches
Several searches for strangelets in meteorite, Earth rock and in cosmic rays have set upper abundance limits. Stellar structure places strong constraints on the strangelet abundance in the Sun, in particular probing <sup>s</sup>H and <sup>s</sup>He. The current state of experiments is reviewed by Klingenberg. We describe the results most relevant to atmospheric strangelets.
Mueller et. al. searched for <sup>s</sup>He in the Earthโs atmosphere. This search, using absorbtion spectroscopy, limits the isotopic abundance of <sup>s</sup>He to be $`<10^8`$ for strangelets of mass $`>20`$ amu. This is not far off from the predictions of the Madsen model. Vandegriff et. al. searched for <sup>s</sup>He in the mass range 42โ82 and found abundances $`<2\times 10^5`$. Both authors focus on primordial strangelets, for which the <sup>s</sup>He/He ratio on Earth is enhanced by a factor of $`10^7`$ over the average galactic <sup>s</sup>He/He ratio by gravitational trapping; this enhancement factor is not relevant for cosmic-ray strangelets. Interpreted as a limit on strangelets in CR, Mueller result constrains the strange star production rate of <sup>s</sup>He to be $`<5\times 10^9\text{M}\text{}`$ y<sup>-1</sup>.
Hemmick et. al. searched for <sup>s</sup>O using a commercial <sup>18</sup>O sample, as discussed in Section IV.2. They find strangelet abundances no greater than $`4\times 10^{17}`$ in bulk oxygen, or $`1\times 10^{15}`$ in their enriched sample<sup>10</sup><sup>10</sup>10We presume that their enriched sample of <sup>18</sup>O, like ours, comes from distillation of CO containing atmospheric O. using a tandem accelerator with an all-electrostatic beam line and a segmented gas chamber. This constrains the strange star production rate of <sup>s</sup>O to be $`<3\times 10^7\text{M}\text{}`$ y<sup>-1</sup>.
Holt et. al. searched for โcollapsedโ Rn nuclei using an atmospheric Xe sample. They collected Rn from the equivalent of $`10^4`$ l of atmospheric Xe, using the reaction (IF$`{}_{6}{}^{})^+`$(SbF$`{}_{6}{}^{})^{}`$ \+ Rn<sub>(g)</sub> $``$ (RnF)<sup>+</sup>(SbF<sub>6</sub>)$`{}_{(s)}{}^{}{}_{}{}^{}`$ \+ IF<sub>5</sub>. This sample was illuminated with a high flux of thermal neutrons, and monitored for the emission of 30โ250 MeV gamma rays. They determined the number of โcollapsedโ Rn atoms to be $`<3\times 10^{10}`$. If this result applied to <sup>s</sup>Rn, it would limit the strange star production rate of <sup>s</sup>Rn to be $`<2\times 10^{17}\text{M}\text{}`$ y<sup>-1</sup>. However, we do *not* expect such high-energy gamma emission from strangelet neutron capture, so this limit does not apply. Nevertheless, this search illustrates the potential for a <sup>s</sup>Rn search to limit strangelet production rates.
### IV.2 Obtaining samples
The standard technique for obtaining noble gases is fractional distillation of air. It appears that most strange gases will distill out along with the normal gases: <sup>s</sup>Ne with Ne, <sup>s</sup>Kr with Kr, etc.<sup>11</sup><sup>11</sup>11It should be noted that the details of commercial distillation plants are proprietary; one can imagine design details, albeit odd ones, that would cause strangelets to be discarded.. <sup>s</sup>Rn comes out in the Xe fraction. To search for strange noble gases, it should be sufficient start with any commercial gas sample.
The exception is helium; commercial helium supplies are extracted from natural gas, where it accumulates largely due to alpha decay rather than due to any primordial or atmospheric source. Therefore, there is no good reason for these samples to contain cosmic-ray strangelets. A strangelet search must be performed on a specially-collected atmospheric He sample.
We also note that <sup>15</sup>N (natural abundance 0.0037) is extracted from atmospheric N<sub>2</sub>, by distillation or chemical exchange, on a commercial scale. Strangelets will follow the heavy fraction in this distillation, so the <sup>s</sup>N abundance will be 270 times greater in a <sup>15</sup>N sample than in a natural N sample. 99% pure <sup>18</sup>O (natural abundance 0.002) is separated commercially via distillation of CO; this concentrate will have 500 times the strangelet concentration of natural O. We include O, N and concentrated <sup>18</sup>O and <sup>15</sup>N in Table 1.
In addition, <sup>s</sup>N should be present at atmospheric concentrations in most solid nitrogen samples, which come from the modern atmosphere via the Haber process or via the biosphere. (The same is not true of oxygen.)
#### IV.2.1 Stratospheric metals
The idea of searching for strangelets in stratospheric metals is somewhat speculative; the primary difficulty is obtaining a large enough material sample. The sample should be gathered from the low polar stratosphere (10-20 km altitude) in late winter. An appropriate collecting device could be launched on a weather balloon or a high-altitude research aircraft.
A large, lightweight electrostatic precipitator could in principle collect aerosol droplets in large quantities from stratospheric air. Although power, weight, and efficiencies need to be evaluated, if one could construct a precipitator which collected all of the aerosol materials in a 1 m<sup>2</sup> area, while making a single pass through a 20 km aerosol layer in the stratosphere, would collect about 1 mg of aerosols, carrying 10<sup>15</sup>โ10<sup>16</sup> atoms of metals and of order $``$ 10โ100 strangelets. Clearly many such flights, or a few more complex missions, would be needed to collect an analyzable sample. However, we emphasize that the stranglet concentration in this sample, if not the sample size, is within the reach of some modern mass spectrometers. Furthermore, we note that such a sample may contain any of a large number of strangelet charges and masses, in constrast to gas-based searches which require some strangelets to have charge 2, 7, 8, 10, 18, 36, etc.
#### IV.2.2 Distillation
Distillation separates isotopes by taking advantage of the small differences in vapor pressure between isotopic species. This difference can be calculated, for monatomic gases, from the elementโs Debye temperature. For a gas with a light isotope of mass $`M^{}`$, vapor pressure $`P^{}`$, and Debye temperature $`T_d`$, and a heavy isotope with $`M`$ and $`P`$, the vapor pressure at temperature $`T`$ obeys
$$\mathrm{ln}(\frac{P^{}}{P})=(\frac{T_d}{T})^2\frac{M^{}}{24}(\frac{1}{M^{}}\frac{1}{M})$$
(2)
Compared to monatomic gases, the calculation of vapor-pressure isotope effects in polyatomic molecules is difficult; vibrational and rotational excitations make large contributions to the isotope shift. For example, the pressure change at 77K, relative to <sup>12</sup>C<sup>16</sup>O, differs by a factor of three between <sup>12</sup>C<sup>18</sup>O, <sup>13</sup>C<sup>17</sup>O, and <sup>14</sup>C<sup>16</sup>O, all of which have the same mass. We make no attempt to calculate them, and use Eq. 2 to give a ballpark figure.
For ideal mixtures, distillation relies on the elementary separation factor $`q_0=p_1/p_2`$, where $`p_x`$ is the vapor pressure of isotope $`x`$ in the pure state. In general, the ultimate separation power of a column is a complicated function of this separation factor. In the special case of strangelets, where the heavy species always has a very small abundance, the behavior simplifies. In a distillation column (or any repeated batch process) with $`N`$ equilibration steps, the maximum obtainable concentration factor is $`c_{max}=q_0^N`$.
For the purposes of this study, rather than designing a distillation column from first principles and optimizing it for strangelet separation, we will consider several columns used for isotope separation in the past. Consider the column used by Johns and London for the production of <sup>13</sup>C from CO. A column 32 feet long used 600 equilibration stages to enrich <sup>13</sup>C from a natural abundance of 0.011% to approximately 0.6%, a factor of 400. The feed rate was 625g/day, and the enriched product was extracted at 0.4g/day. The single-stage separation factor is 1.01. Running the same column to separate <sup>s</sup>Xe from Xe (single-stage separation factor 1.006 at the boiling point) would give an enrichment of a factor of 40. Running on <sup>s</sup>Kr and Kr (single-stage separation factor 1.01) would give an enrichment of $`10^2`$$`10^3`$<sup>12</sup><sup>12</sup>12A small uncertainty on the single-stage separation leads to a very large uncertaintly in the column performance. For lighter gases, the single-stage enrichments are larger; operating the column at the original feed rate would give a maximum enrichment factor of $`1.5\times 10^3`$ (a number which depends on the amount of material in the reboiler to which the strangelets are transported.) Reducing the enriched product withdrawal rate, or increasing the feed rate, would increase the enrichment factor approximately linearly.
A smaller column, like that used by Clusius and Mayer for the separation of <sup>36</sup>Ar, would be entirely adequate for work with <sup>s</sup>Ne, <sup>s</sup>Ar, <sup>s</sup>N, or <sup>s</sup>O. This column produced an enrichment factor of 2 using 160 stages with a single-stage separation of 1.005; running this column on <sup>s</sup>Ar would give a strangelet enrichment of $`10^3`$. The column produced 0.7 moles of enriched product per day.
For <sup>s</sup>Ne, due to the very low temperatures involved, an even smaller column would be desired; an apparatus with 20 stages would be sufficient to concentrate strangelets by a factor 1000. Even the fairly small apparatus of Keesom et. al. , with 60 separation stages, is overdesigned for <sup>s</sup>Ne separation; as operated on 22Ne, this column would enrich <sup>s</sup>Ne by a factor $`500`$ in a single pass ( 4d), and would increase that enrichment factor approximately linearly with additional running time.
#### IV.2.3 Thermal diffusion separation
Thermal diffusion separation is very powerful and efficient at separating heavy gases from light mixtures. A typical column might consist of a hollow tube 25mm in diameter and 3m long, water-cooled on the outside, with a coaxial wire electrically heated to 600K. Light isotopes preferentially migrate towards the hotter gas and are swept upwards by convection; heavy isotopes concentrate at the bottom. It is straightforward to operate several short columns in series so that they behave as a single long column.
Consider a column of length $`L`$, filled with a gas of molecular mass $`m`$, and a trace concentration $`c_0`$ of strangelets of mass $`M`$. At equilibrium, the strangelet concentration $`c_L`$ at the bottom of the column is increased by
$$q_e\frac{c_L}{c_0}=ke^{L/\lambda }$$
(3)
where the characteristic length $`\lambda `$ scales approximately as $`\lambda (Mm)/(M+m)`$, and $`k`$ depends in detail on diffusion parameters, temperature, viscosity, etc. For the heavier gases, this equation predicts extremely large values for $`q_e`$ even for fairly short columns. In practice, very large separation constants suggest that strangelets are transported down the column unidirectionally. There is unavoidably some dead space at the bottom of the column; the columnโs whole inventory of strangelets will accumulate in the dead space in the equilibration time $`\tau `$. This allows us to compute the strangelet production rate $`P`$, the number of strangelets transported down the column per day.
The final strangelet concentration at the bottom of the column will depend on $`q_e`$, $`P`$, the run time, and the amount of dead space from which the final sample is drawn. The dead space consideration may be removed by mixing the initial sample with an intermediate-mass buffer gas. For example, consider a long column containing <sup>s</sup>O<sup>16</sup>O and <sup>18</sup>O<sub>2</sub>. A sample drawn from the bottom of the column may contain 60 times as much <sup>s</sup>O<sup>16</sup>O as the source gas, but the bulk of the sample will still be <sup>18</sup>O<sub>2</sub>. Compare this to a column containing <sup>s</sup>O<sup>16</sup>O , <sup>18</sup>O<sub>2</sub> and Kr as a buffer gas. In this case, a bottom sample will contain <sup>s</sup>O<sup>16</sup>O mixed with bulk Kr, to the near-complete exclusion of <sup>18</sup>O<sub>2</sub>. This sample is subject to additional chemical purification. This allows the preparation of extremely concentrated strangelet samples from almost any gas, with appropriate choice of buffer.
#### IV.2.4 Chemical separation of Rn
We can collect very large atmospheric Rn samples, including <sup>s</sup>Rn, by extraction from Xe over (IF<sub>6</sub>)<sup>+</sup>(SbF<sub>6</sub>)<sup>-</sup>, as described in section IV.1. An even simpler technique would be chromatography on activated carbon. A stream of mixed Xe and Rn is passed over activated carbon at room temperature. Rn is strongly retained. Later, the Rn is released by heating under a stream of He or N<sub>2</sub>. In either case, 1 mole of Xe (22.4 lSTP) with the Table 1 strangelet concentrations would yield $`3\times 10^{14}`$ <sup>s</sup>Rn atoms.
#### IV.2.5 Isotope separation by chromatography
If it is necessary to pre-concentrate strangelets from a stratospheric metal sample, this must be done with very high efficiency for the strangelets; very small initial samples cannot be wasted. However, since the strangelet concentrations in these samples are high to begin with, we may be content with a low-throughput technique like chromatography. Ion-exchange chromatography is known to show small isotope separation effects. We show some data in Table 5; more can be found in and . Assuming these effects to be due solely to nuclear mass effects<sup>13</sup><sup>13</sup>13As always, this assumption may be distorted by nuclear size, spin, and quantum effects., and following the Bigeleisen-Mayer prescription, we extrapolate these small effects to strangelets by assuming that the isotope effect $`ฯต`$ (basically, the fractional difference in the solubility product between isotopes of mass $`M`$ and $`m`$) depends on the masses as $`ฯต=ฯต_0\frac{Mm}{M\times m}`$. The separation power of a chromatography column scales of length $`L`$ scales as $`Lฯต^2`$. Large separation factors, like those predicted for <sup>s</sup>Li and <sup>s</sup>Mg, would allow near-complete strangelet purification in a single column. Smaller separation factors, as for <sup>s</sup>Vโ<sup>s</sup>U, permit concentration by factors of 10-100 in simple setup. In principle, such a process can purify and recover nearly 100% of a strangelet sample, although this requires testing.
##### Electromagnetic separation
It is always possible to separate gas-phase ions by mass using electromagnetic fields. Traditional โcalutronsโ separate isotopes using a spark ion source, an accelerating potential of order 10โ20 kV, and large-area ($`>`$1m<sup>2</sup>) magnetic field. Ions with different m/z have different trajectories in the magnetic field; beams of various masses are made to implant in collector foils at different positions, from which the atoms are later extracted. Consider a calutron which produces a beam of 10 $`\mu `$A of ions from a strangelet-containing oxygen sample. Supplying this source with straight commercial <sup>18</sup>O, with the nominal strangelet concentration, produces $`10^5`$ separated strangelets per month. The collection system is necessarily complicated, in a search experiment, by the fact that the strangelet mass is unknown. The purity of the resulting sample will be limited by, among other things, the nonzero beam size and halo, which will send some fraction of the non-strange ions into the strangelet collector area. This suggests that electromagnetic separation is not promising compared to thermal diffusion separation.
Because typical ion sources have very low source ionization/transmission efficiencies ($`10^5`$), electromagnetic separation is unsuitable for the already-small stratospheric metal samples. It should be kept in mind for future terrestrial and lunar work.
## V Mass spectroscopy
After collecting and pre-concentrating a strangelet-bearing sample, one must perform mass spectroscopy to try to identify strangelets. A detailed survey of the spectroscopy techniques available is beyond the scope of this paper. However, we observe that many extremely-sensitive techniques are available off the shelf. There are three major factors influencing oneโs choice of mass spectrometer: sensitivity, sample consumption, and mass range. For example, accelerator mass spectrometry (AMS) is often capable of detecting elements or isotopes at concentrations of 10<sup>-14</sup>, but consumes comparatively large amounts of sample, and is only sensitive to $``$ 1 amu at a time. Fourier transform ion cyclotron resonance mass spectroscopy (FTICR-MS) can process very tiny samples, and observes a wide range of candidate masses, but cannot deal with isotopic abundances below about $`10^6`$. Table 6 gives a brief overview of some mass spectroscopy techniques which could be used for strangelet searches.
## VI Summary and conclusions
The Earthโs atmosphere serves as a beam dump for cosmic ray strangelets. In the upper atmosphere, recently-fallen strangelets accumulate and concentrate in aerosol droplets, free from contamination by normal nuclei. In the atmosphere as a whole, strange isotopes of noble gases, N, and O should have accumulated to quite high concentrations. These concentrations are within the reach of modern mass spectroscopy, even if strangelet fluxes are quite small.
We sum up strangelet accumulation in a given environment by a โconcentration factorโ, shown in Table 7. This is the per-atom abundance of strangelets in this environment, divided by the flux in m<sup>-2</sup>y<sup>-1</sup>. This factor shows that stratospheric aerosols yield the highest strangelet concentrations known; the Earthโs atmosphere and the moon are equally good places to obtain large quantities of less-concentrated strangelets. The Earthโs crust has a very low concentration factor, so effective strangelet searches in rock will require additional effort.
Taken together, these suggestions constitute a fairly comprehensive strangelet search program which could be executed mainly with off-the-shelf equipment. An optical search for <sup>s</sup>Rn offers an excellent opportunity to discover high-mass strangelets. A deep search for low-mass strangelets can be carried out on atmospheric <sup>s</sup>He, <sup>s</sup>N, <sup>s</sup>O, and <sup>s</sup>Ne, by carrying out thermal-diffusion separation and accelerator mass spectroscopy. Intermediate-mass strangelets might be found in strange transition metals in the stratosphere, searched for by resonance ionization mass spectroscopy, or in <sup>s</sup>Kr, with an ATTA search. An overview of the expected concentrations is shown in Figure 4. This illustrates that a very large strangelet search space can be addressed using the atmospheric sampling and preconcentration techniques described here.
The author would like to thank Sarah Bagby, Edward Boyle, Alan Davison, Peter Fisher, Ben Lane, David Mohrig, Alan Plumb, David Pritchard, Joseph L. Smith, and Forest White of MIT, Dick Majka and Jack Sandweiss of Yale, Jes Madsen of Aarhus, Alan Marshall of Florida State University, and Wallis Calaway and Zheng-Tian Lu of Argonne National Lab for useful discussions on many aspects of this research.
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# The Recent and Continuing Assembly of Field Ellipticals by Red Mergers
## 1 Introduction
Although greatly outnumbered by spiral galaxies, elliptical galaxies contain $`20`$ % of the stellar mass in the present-day Universe (Fukugita, Hogan, & Peebles 1998) and dominate the high mass end of the galaxy mass function (e.g., Nakamura et al. 2003). Their status at the top of the food chain is a strong motivation for their study as it provides insight in the processes governing the build-up of galaxies over cosmic time, the nature of the most luminous galaxies at $`z3`$ and beyond, the formation of the first stars, and the origins of supermassive black holes.
Despite vigorous research efforts it is still not known how and when elliptical galaxies were formed. Separating the time of the formation of their stars from the time of their assembly, a partial answer is that most of the stars in the most massive ellipticals were formed at high redshift. This result is supported by a large number of studies of nearby and distant galaxies, and applies to ellipticals in the general field as well as those in rich clusters (e.g., Bower, Lucey, & Ellis 1992; Ellis et al. 1997; Bernardi et al. 1998, 2003; van Dokkum et al. 1998a, 1998b; Treu et al. 1999, 2002; van Dokkum et al. 2001a; Holden et al. 2005). Even when accounting for possible selection effects due to morphological evolution (โprogenitor biasโ; van Dokkum & Franx 2001), it is very difficult to fit the observed evolution of massive ellipticals with formation redshifts substantially lower than $`z2`$. Although there is general agreement that the mean age is high, there are strong indications that not all stars in all elliptical galaxies formed at high redshift: Trager et al. (2000) and others (e.g., Yi et al. 2005) find evidence for a โsprinklingโ of recent star formation in otherwise old ellipticals, and recent studies of the evolution of the Fundamental Plane relation (Djorgovski & Davis 1987) find evidence that low mass early-type galaxies have lower luminosity-weighted ages than high mass ones (Treu et al. 2005; van der Wel et al. 2005).
Much less is known about the assembly of elliptical galaxies, that is, when and how the galaxies were โput togetherโ. In the literature, a distinction is often made between models in which elliptical galaxies form from the collapse of primordial gas clouds (โmonolithic collapseโ; e.g., Eggen, Lynden-Bell, & Sandage 1962; Jimenez et al. 1999), or in mergers of smaller galaxies (โhierarchicallyโ; e.g. Toomre & Toomre 1972; White & Frenk 1991). Within the framework of a merger model the question of assembly can be specified as, first, the typical age of the last major merger; and second, the nature of the progenitor galaxies. In many individual cases we can readily answer both questions: abundant evidence leaves little doubt that some ellipticals formed recently through the merger of spiral galaxies, or will do so in the near future (e.g., Toomre & Toomre 1972; Schweizer 1982; Hibbard & van Gorkom 1996). However, it is not clear whether such mergers are exceptional or responsible for the formation of the majority of ellipticals.
Among the observational arguments for recent mergers are the detection of โfine-structureโ (e.g., shells and ripples) in a large number of nearby ellipticals (Malin & Carter 1983; Schweizer et al. 1990); the presence of kinematically-decoupled cores (Franx & Illingworth 1988; Bender 1988); small-scale dust seen in the majority of ellipticals (e.g., van Dokkum & Franx 1995); (some) studies of the redshift evolution of galaxy morphologies, and of close pairs (e.g., Le Fรจvre et al. 2000; Patton et al. 2002; Conselice et al. 2003); and strong evolution of the mass density of red galaxies to $`z1`$, inferred from the COMBO-17 survey (Bell et al. 2004). Furthermore, and perhaps most importantly, the hierarchical assembly of galaxies is a central aspect of current models for galaxy formation in $`\mathrm{\Lambda }`$CDM cosmologies (e.g., Kauffmann, White, & Guiderdoni 1993; Cole et al. 2000; Meza et al. 2003).
Among arguments against the merger hypothesis are the low probability of mergers in virialized clusters, which are dominated by ellipticals (Ostriker 1980; Makino & Hut 1997); the fact that the central densities of ellipticals are too high to be the result of dissipationless mergers of pure disks (Carlberg 1986; Hernquist 1992); the low scatter in the color-magnitude relation and Fundamental Plane, which is inconsistent with the diversity in stellar populations expected from spiral galaxy mergers (e.g., Bower, Kodama, & Terlevich 1998, Peebles 2002); the much higher specific frequency of globular clusters for ellipticals than for spirals (although some globulars may be formed during mergers; e.g., Schweizer et al. 1996); the observation that known remnants of spiral galaxy mergers have underluminous X-ray halos (Sansom, Hibbard, & Schweizer 2000); the existence of the $`M_{}\sigma `$ relation (Gebhardt et al. 2000; Ferrarese & Merritt 2000); and the presence of massive galaxies at high redshift (e.g., Franx et al. 2003; Daddi et al. 2004; Glazebrook et al. 2004)
These apparently contradictory lines of evidence may be largely reconciled by postulating that most elliptical galaxies formed through (nearly) dissipationless (or โdryโ) mergers of red, bulge-dominated galaxies rather than mergers of spiral disks. Motivated by the comparatively red colors of ellipticals exhibiting fine-structure Schweizer & Seitzer (1992) discuss this possibility, but argue that mergers between early-type galaxies are statistically unlikely as the โmedianโ field galaxy is an Sb spiral (see also Silva & Bothun 1998). However, most mergers likely occur in groups, where the early-type galaxy fraction is much higher than in the general field (e.g., Zabludoff & Mulchaey 1998). Furthermore, semianalyical models of galaxy formation have predicted that the most recent mergers of bright ellipticals were between gas-poor, bulge-dominated galaxies (Kauffmann & Haehnelt 2000; Khochfar & Burkert 2003), and several studies have shown that such merging is consistent with the observed Fundamental Plane relation (e.g., Gonzรกlez-Garcรญa & van Albada 2003; Boylan-Kolchin, Ma, & Quataert 2005). Finally, mergers between red galaxies may be common in young, unvirialized galaxy clusters at $`z1`$ (van Dokkum et al. 1999, 2001b, Tran et al. 2005).
Although some early-type/early-type mergers are known to exist in the local universe (e.g., Davoust & Prugniel 1988; Combes et al. 1995), little is known about their frequency. Here we investigate the relevance of dry merging by analyzing the frequency and nature of tidal distortions associated with a well-defined sample of bright red galaxies. This study was motivated by the results of Kauffmann & Haehnelt (2000), Khochfar & Burkert (2003), and Bell et al. (2004), the observations of red mergers in $`z1`$ clusters, and the advent of very deep photometric surveys over wide areas. We assume $`\mathrm{\Omega }_m=0.3`$, $`\mathrm{\Omega }_\mathrm{\Lambda }=0.7`$, and $`H_0=70`$ km s<sup>-1</sup> Mpc<sup>-1</sup> (Spergel et al. 2003).
## 2 Morphological signatures
Mergers between gas-poor, bulge-dominated galaxies are generally more difficult to recognize than mergers between gas-rich disks. Elliptical-elliptical mergers do not develop prominent tidal tails dotted with star forming regions, but are instead characterized by the ejection of broad โfansโ of stars, and in certain cases an asymmetric deformation of the inner isophotes (e.g., Rix & White 1989; Balcells & Quinn 1990; Combes et al. 1995). Tails may develop if one of the progenitors rotates or has a disk component (Combes et al. 1995), but as there is no cold, young component these are expected to be more diffuse than those seen in encounters between late-type galaxies.
These effects are illustrated in Fig. 1, which shows the evolution of an off-axis collision between two hot stellar systems. The simulation was performed with an implementation of the Barnes & Hut (1986) hierarchical tree code, using 65,536 particles. The galaxies have a 3:1 mass ratio and are represented by Plummer models. The simulation illustrates the rapid merging of the two bodies, the lack of tidal tails, and the development of a large fan of material. The central parts of the merger remnant quickly settle in a new equilibrium configuration, whereas the material in the outer parts โ where dynamical timescales are long โ gradually becomes more diffuse. For a more general analysis of the development of elliptical-elliptical mergers, including the effects of rotation and different mass ratios, we refer to Combes et al. (1995) and subsequent papers.
The surface brightness of tidal features associated with this type of mergers is comparatively low: the features consist of old, red stellar populations, and have higher $`M/L`$ ratios than the blue tails associated with mergers of late-type galaxies. Furthermore, they are typically more diffuse, due to the lack of a cold disk in the progenitors. After the merger, the evidence of the past encounter becomes increasingly difficult to detect as dynamical evolution reduces the surface brightness further (e.g., Mihos 1995). No detailed simulations have been done of the expected surface brightness evolution of tidal features resulting from early-type mergers, but a very rough estimate can be obtained from the simulation shown in Fig. 1. Calibrating the output such that the average surface brightness within the effective radius of the remnant $`\mu _R=20`$ mag arcsec<sup>-2</sup> (Jรธrgensen, Franx, & Kjรฆrgaard 1995), we find that the โfanโ of material to the upper left has $`\mu _R25`$ at $`t=10`$, i.e., immediately after the merger. The following timesteps show a continued rapid decrease of the brightness, reaching $`\mu _R27`$ at $`t=15`$, the final time step of the simulation.
It is difficult to reach such low surface brightness levels observationally: extreme enhancements of photographic plates enabled Malin & Carter (1983) to detect features to $`\mu 26.5`$, similar to the limits reached by Schweizer & Seitzer (1992) with their deep CCD exposures on the Kitt Peak 0.9 m telescope. Detecting features with surface brightness $`2728`$ requires several hours of integration time on 4 m class telescopes, and to image a sizeable sample of nearby galaxies would require hundreds of hours. Furthermore, flat fielding would have to be accurate to $`0.1`$ % over the entire $`10^{}30^{}`$ field.
These practical difficulties can be solved by exploiting the fact that the $`(1+z)^4`$ cosmological surface brightness dimming is a very slow function of redshift for $`z1`$. Whereas moving a galaxy from $`z=0.01`$ to $`z=0.1`$ reduces its total brightness by 5 magnitudes, its surface brightness changes by only 0.35 mag. In practice it is easier to detect very faint surface brightness features around a galaxy at $`z=0.1`$ than at $`z=0.01`$ as it occupies a $`70\times `$ smaller area of the detector. Furthermore, and most importantly, a study of tidal features at $`z0.1`$ does not require a dedicated survey but can be done with existing imaging data.
## 3 Data
We use data from two deep extra-galactic surveys: the Multiwavelength Survey by Yale-Chile (MUSYC; Gawiser et al. 2005), and the NOAO Deep Wide-Field Survey (NDWFS; Januzzi & Dey 1999, Januzzi et al. 2005, Dey et al. 2005). The data products from the optical imaging components of the two survey are quite similar, as they both use the $`8192\times 8192`$ pixel MOSAIC cameras on the NOAO 4m telescopes. The full dataset covers 10.5 square degrees of sky, reaching $`1\sigma `$ surface brightness levels of $`\mu 29`$ mag arcsec<sup>-2</sup>.
### 3.1 MUSYC
MUSYC is a deep optical/near-infrared imaging and spectroscopic survey of two Southern and two equatorial $`30^{}\times 30^{}`$ fields, chosen for their low background and the availability of data at other wavelengths. The fields are the Extended Chandra Deep Field South (E-CDFS), the Extended Hubble Deep Field South (E-HDFS), the field centered on the $`z=6.3`$ Sloan Digital Sky Survey quasar SDSS 1030+05 (Becker et al. 2001), and the blank low background field CW 1255+01. The primary goal of the survey is the study of normal and active galaxies at $`z>2`$.
The survey design and data analysis procedures are described in Gawiser et al. (2005). Briefly, the optical imaging data consist of deep $`UBVRIz`$ exposures obtained with MOSAIC-I on the Mayall 4m telescope on Kitt Peak and MOSAIC-II on the Blanco 4m telescope on Cerro Tololo. The images have a scale of $`0\stackrel{}{\mathrm{.}}267`$ pix<sup>-1</sup>, and typical seeing of $`1\stackrel{}{\mathrm{.}}1`$. The data reach depths of $`26`$ in $`U`$, $`26.5`$ in $`B`$, $`V`$, and $`R`$, $`25`$ in $`I`$, and $`24`$ in $`z`$ (AB magnitudes, $`5\sigma `$ point source detections), with some variation between the fields. The full exposure time is typically realized over an area of $`33^{}\times 33^{}`$, and the total optical survey area is $`1.2`$ deg<sup>2</sup>.
For detection of low surface brightness features we used the co-added โBVRโ images (see Gawiser et al. 2005), which are combinations of the $`B`$, $`V`$, and $`R`$ frames. The total effective integration time of these frames is typically $`8`$ hours. The final depth reaches magnitude $`27`$ (AB; $`5\sigma `$ point source detection), with slight variations from field to field.
### 3.2 NOAO Deep Wide-Field Survey
The NDWFS is a public optical and near-infrared imaging survey over two 9.3 deg<sup>2</sup> fields: a $`3\mathrm{ยฐ}\times 3\mathrm{ยฐ}`$ field in Boรถtes, centered at 14<sup>h</sup>32<sup>m</sup>, $`+34\mathrm{ยฐ}17\mathrm{}`$, and a $`2.3\mathrm{ยฐ}\times 4\mathrm{ยฐ}`$ field in Cetus, centered at $`2^h`$10<sup>m</sup>, $`4\mathrm{ยฐ}30\mathrm{}`$ (see Januzzi & Dey 1999; Januzzi et al. 2005). The primary goal of the survey is to study large scale structure at $`z>1`$. We use the optical images available in the third data release (DR3; October 22 2004). This release comprises the entire $`9.3`$ deg<sup>2</sup> Boรถtes field.
Compared to MUSYC, the NDWFS sacrifices multi-band coverage for area. Only three optical filters are used: $`B_w`$, $`R`$, and $`I`$. The $`B_w`$ filter was designed specifically for the NDWFS: it is a very broad blue filter, effectively a combination of the standard $`U`$ and $`B`$ filters. Exposure times vary across the field; the median values are 8400 s in $`B_w`$, 6000 s in $`R`$, and $`11400`$ s in $`I`$. Details of DR3 can be found on the NDWFS webpages<sup>1</sup><sup>1</sup>1http://www.noao.edu/noao/noaodeep/DR3/.
The data were not pasted together in a mosaic by NOAO but released as 27 partially overlapping $`35^{}\times 35^{}`$ MOSAIC pointings. We obtained all 27 pointings in $`B_w`$, $`R`$, and $`I`$ from the NOAO archive. In order to increase the signal-to-noise (S/N) ratio we created co-added โBRIโ images, in the following way. First, all images were normalized such that 1 count corresponds to AB$`=30`$ mag. Next, the summed BRI images were created from the $`B`$, $`R`$, and $`I`$ frames according to BRI$`=B+R+0.5\times I`$. This procedure optimizes the S/N ratio for objects with the approximate SED of nearby elliptical galaxies.
### 3.3 Depth of Summed Images
The detection of diffuse tidal features depends on the limiting surface brightness level of the data, not on the point source detection limit. To assess the (AB) surface brightness limits we placed random apertures with an area of $`1\mathrm{}`$ in empty regions, and determined the rms fluctuations. The co-added images were used, i.e., the โBVRโ images for MUSYC and the โBRIโ images for NDWFS.
The MUSYC images have a typical $`1\sigma `$ surface brightness limit of $`\mu 29.5`$ mag arcsec<sup>-2</sup> whereas this limit is $`29`$ mag arcsec<sup>-2</sup> for the NDWFS images. These numbers are approximate: the exact limit depends on the field, the local flat fielding accuracy, and on the assumed spectral energy distribution of the source. The greater depth of MUSYC was expected, as more exposure time was devoted to the efficient $`B`$, $`V`$ and $`R`$ filters.
The ability to detect low surface brightness emission is influenced by the extent of the features, confusion with neighboring objects, the contrast with the smooth emission from the galaxy, and other effects; in practice, we find we can confidently identify features down to $`\mu 28`$ mag arcsec<sup>-2</sup>.
## 4 Sample Selection
The initial sample selection is based on color and magnitude only. No morphological criteria are applied, for three reasons: first, a simple color cut is straightforward to reproduce; second, as the galaxies are typically at $`z0.1`$ it can be difficult to determine morphologies from ground-based data (see ยง 5.1); and third, the classification process itself might bias the sample toward or against galaxies with tidal features.
### 4.1 Catalogs
The combination of the two surveys comprises 31 multi-band tiles, each covering $`0.3`$ deg<sup>2</sup>. Catalogs were created using the SExtractor software (Bertin & Arnouts 1996). SExtractor catalogs were available for both surveys but these are optimized for faint point sources, not for large, nearby galaxies. SExtractor was run with its default settings, with the following changes. The detection threshold was set at $`20\sigma `$, with the added requirement that 10 adjacent pixels are above the threshold, and the mesh size for background subtraction was set at 400 pixels rather than the default 64 to avoid oversubtraction of the background for large galaxies. Matched photometry in each band was obtained by running SExtractor in dual image mode, always using the $`R`$-band images for detection.
SExtractorโs AUTO magnitudes are used as best estimates for total magnitudes. Colors were initially determined in fixed $`5\mathrm{}`$ diameter apertures. However, we find that colors measured in fixed apertures are not well suited for these large objects as they introduce distance-dependent selection biases. Specifically, when selecting on aperture color the sample invariably contains large, presumably very nearby spiral galaxies with blue disks and red bulges. In the following, $`BR`$ colors refer to $`B_{\mathrm{AUTO}}R_{\mathrm{AUTO}}`$, measured in dual image mode.
Finally, the measured $`B_wR`$ colors in the NDWFS are converted to the standard Johnson system. The $`B_w`$ filter was created for the NDWFS and has a bluer central wavelength than the $`B`$ filter. Using templates from Coleman, Wu, & Weedman (1980) we derive
$$(BR)=0.27+0.62(B_wR)+0.07(B_wR)^2.$$
(1)
This transformation holds to a few percent for $`1<(BR)<2.5`$ and redshifts $`z<0.2`$, and was applied to the NDWFS photometry prior to the sample selection. Total magnitudes and colors are given on the Vega system to facilitate comparisons to previous studies.
### 4.2 Color Selection
Galaxies are selected on the basis of their total $`R`$ magnitude and $`BR`$ color. Figure 2 shows color-magnitude diagrams for the four MUSYC fields (left) and for the 27 NDWFS fields (right). Only objects with a star/galaxy classification $`<0.06`$ are shown (see Bertin & Arnouts 1996). This classification is fairly robust at the magnitudes of interest, and effectively removes both non-saturated and saturated stars from the sample. Visual inspection revealed that the software assigned a star/galaxy classification $`0.06`$ to four galaxies in our sample; for these objects the value was manually set to zero. We note that the cores of some bright galaxies can be saturated, which influences the magnitudes and (particularly) the colors in unpredictable ways. As this effect is difficult to quantify and only influences a handful of objects we did not attempt a correction. The distribution of points is similar in the two surveys. The number of objects in the NDWFS is an order of magnitude larger than in MUSYC, as expected from the difference in area.
We selected galaxies with the colors and magnitudes of $`L>L_{}`$ early-type galaxies at $`0.05<z<0.2`$. The two solid lines show the colors of the Coleman et al. (1980) E/S0 template redshifted from $`z=0.04`$ to $`z=0.20`$, normalized to $`L_{}`$ and $`3L_{}`$ (Blanton et al. 2001). Dashed lines delineate the adopted selection region: $`R<17`$, $`1.6(BR)2.2`$, and $`(BR)>1.6+0.12\times (R15)`$. The galaxies in the dashed region are the brightest and reddest in $`10.5`$ square degrees of sky.
### 4.3 Additional Steps
After visual inspection of the initial sample of 155 red objects 32 were removed for various reasons. The discarded objects fall in a wide range of categories: mis-classified stars (usually blends with other stars or galaxies); objects on the edge of a field; spurious objects (always near very bright stars); severely saturated galaxies (although interesting in their own right, as they are likely very bright active nuclei); and galaxies that have more than one catalog listing. The latter occurs because the 27 NDWFS tiles have some overlap; in these cases the data from the overlapping pointings were added, unless the S/N is much higher in one of the two. We also removed 12 galaxies that are likely members of known galaxy clusters. Eight are located within $`8^{}`$ of NSC J142841+323859, a cluster at $`z=0.127`$, and four are located within $`4^{}`$ of the poorer, more compact cluster NSC J142701+341214 at $`z0.2`$ (Gal et al. 2003).
Finally, we inspected all galaxies with $`17R17.75`$ to identify merging pairs (see ยง 5.2) that are split into two objects by SExtractor, and whose combined luminosities and colors would place them in our selection region. Three such pairs were identified, and the brightest galaxy of each pair was added to the sample: 4-1190, 4-1975, and 22-2252. The final sample thus consists of 126 red field galaxies in MUSYC and the third data release of the NDWFS.
### 4.4 Median Redshift
The median $`R`$ magnitude of galaxies in our sample is 16.4, and the median $`BR`$ color is 1.92. As can be seen in Fig. 2 these values suggest that the median galaxy is an $`L1.2L_{}`$ early-type galaxy at redshift $`z0.11`$. Currently we have redshifts for only nine of the 123 galaxies: four in the CDF-S from COMBO-17 (Wolf et al. 2005), two in the SDSS 1030+05 field from the SDSS (York et al. 2000), and three in the CW 1255+01 field from the SDSS. The mean redshift of these nine objects $`z=0.13`$ with a spread of 0.02, consistent with the expectations from our color selection.
As a further test on the selection we obtained photometric and spectroscopic data in ten randomly selected $`3^{}\times 3^{}`$ fields from the SDSS. Galaxies were selected in the $`gr`$ vs. $`r`$ plane in the same manner as in our study, using the Fukugita et al. (1995) transformations to convert our $`R`$ and $`BR`$ limits to SDSS $`r`$ and $`gr`$ limits. The median redshift is 0.098, with 90 % of the galaxies at $`z>0.05`$. The rms field-to-field variation of the median is only 0.008.
In the following we will adopt $`z=0.1`$ as the median redshift; our conclusions change very little if we were to use, e.g., $`z=0.08`$ or $`z=0.13`$ instead. For $`z0.1`$ our observed $`R`$-band limit corresponds to $`M_R21`$, and the median galaxy has $`M_R22`$.
## 5 Analysis
### 5.1 Morphologies
All galaxies were assigned a morphological type by visually inspecting their summed images (i.e., the โBVRโ images for MUSYC galaxies and the โBRIโ images for NDWFS objects). The images span a large range in surface brightness levels, going from the nearly saturated central regions to very low surface brightness features up to $`>50`$ kpc away from the center. As the relevant dynamic brightness range is $`10^4`$ each galaxy was displayed at four different contrast levels simultaneously.
The morphological types are necessarily broad. Although the S/N ratio is very high the spatial resolution is quite poor: the typical seeing is $`1\stackrel{}{\mathrm{.}}1`$, which corresponds to 3 kpc at $`z=0.1`$. Therefore, we only have $`23`$ resolution elements within the half-light radii of many galaxies (see, e.g., Jรธrgensen, Franx, & Kjรฆrgaard 1995). The assigned types are spiral (S), indicating the presence of spiral arms and/or star forming regions in a disk; S0, indicating an early-type galaxy with an unambiguous disk component; and E/S0, indicating a bulge-dominated early-type galaxy. We cannot securely separate elliptical galaxies from bulge-dominated or face-on S0 galaxies. Contrary to the usual definition the E/S0 class therefore encompasses ellipticals, bulge-dominated S0s, E/S0s, and S0/Es.
Morphological classifications are listed in Table 1, and images of all 126 galaxies are shown in the Appendix. As expected, the sample is dominated by early-type galaxies: of 126 objects, 10 (8 %) are classified as spirals, 30 (24 %) as S0s (i.e., disk-dominated early-type galaxies), and 86 (68 %) as E/S0s (i.e., bulge-dominated early-type galaxies). Most of the spiral galaxies have large red bulges and faint blue arms. Based on their morphologies at surface brightness levels $`\mu 25`$ mag arcsec<sup>2</sup> we infer that virtually all galaxies in our sample are red because of their evolved stellar populations, and not because of dust (as is well known from many previous studies of bright red galaxies in the local Universe; see, e.g., Sandage & Visvanathan 1978, Strateva et al. 2001).
### 5.2 Tidal Features
Tidal features were first identified by visual inspection of the full sample of 126 galaxies; a quantitative analysis of disturbances in the restricted sample of 86 bulge-dominated early-type galaxies follows in ยง 5.3. The flag describing tidal features can have one of four values: 0 for no tidal features, 1 for weak features, 2 for strong features, and 3 for an ongoing interaction with another galaxy. Galaxies in the โ2โ class are generally highly deformed merger remnants, whereas the โ1โ class indicates more subtle features. The difference is obviously quite subjective, and in the subsequent analysis these two classes will generally be combined into one. A โ3โ classification implies that both the primary and secondary galaxy show clear distortions or tidal tails.
The key result of our analysis is the ubiquity of tidal features, particularly among bulge-dominated early-type galaxies. In the full sample of 126 galaxies, 44 (35 %) show clear signs of past interactions and in an additional 23 cases (18 %) the interaction is still in progress. Only 59 galaxies (47 %) appear undisturbed, showing no unambiguous tidal features at the surface brightness limit of the survey. The fraction of disturbed objects is lowest among galaxies with a clear disk component. Among 40 galaxies classified as S or S0, only six (15 %) show evidence for past or present interactions. In contrast, among the 86 galaxies classified as E/S0 undisturbed objects are the exception, as 61 (71 %) show tidal features.
The nature and extent of the disturbances span a wide range. Some galaxies have clearly defined tidal tails while others show broad fans of stars similar to the final frames of the simulation shown in Fig. 1. In most cases the disturbances are subtle and only visible at large radii and very faint surface brightness levels, although some objects are strongly disturbed throughout. Often we see a mixture, e.g., a well defined tail in addition to a broad, smooth disturbance at very faint levels. This large range of properties is not surprising as it reflects the variation in the age of the interaction, the viewing angles, and the properties of the progenitors.
The tidal features are almost always red. In many cases the features are only visible in the $`R`$ and $`I`$ frames despite the substantial depth of the blue exposures. They also appear smooth, showing no or very little evidence for clumps and condensations. In these respects the features are very different from the blue, clumpy tidal tails seen in spiralโspiral interactions (e.g., Mirabel, Dottori, & Lutz 1992; Hunsberger, Charlton, & Zaritsky 1996), and from the sharp shells and ripples detected in unsharp-masked images of ellipticals (e.g., Schweizer & Seitzer 1992, Colbert et al. 2001). We note here that our data are not well suited to identify narrow features, as the FWHM resolution of our data is about 2 kpc at $`z=0.1`$.
The remarkable nature of these red mergers and their remnants is illustrated in Fig. 3. The figure shows two examples of ongoing mergers, an example of a strongly disturbed merger remnant, and an example of a galaxy with more subtle distortions at faint surface brightness levels, arranged in a plausible red merger sequence. In all cases the tidal features are smooth and red, and quite different from the highly structured blue tails associated with known mergers between gas rich disk galaxies. Figure 4 shows summed images of the same objects, highlighting the faintest features and providing an indication of the surface brightness levels reached by the observations. We stress that these objects are fairly typical examples: as can be inferred from the smaller images shown in the Appendix they are by no means unique within our sample.
### 5.3 Quantitative Identification of Tidal Features
Quantitative characterization of the tidal features is important for testing the robustness of the visual classifications, measuring the flux associated with the features, assessing the effects of changing the S/N ratio, and examining correlations between distortions and other properties of the galaxies. Quantitative criteria are also useful as a tool for future studies of larger samples. Standard measures of asymmetry (e.g., Abraham et al. 1996; Conselice et al. 2003) are not applicable to these galaxies as the features comprise only a small percentage of the total luminosity of the galaxies. Instead a method was developed which determines distortions with respect to a model light distribution (see, e.g., Colbert et al. 2001). The method is only meaningful for bulge-dominated early-type galaxies, which make up the bulk of the sample.
First, the galaxies are fitted by an elliptical galaxy model using the โellipseโ task in IRAF, in three iterations. After each iteration a mask file is updated using the residuals from the previous fit. The center position, ellipticity, and position angle are allowed to vary with radius. In cases where a second galaxy overlaps the primary object the two objects are fitted iteratively. The final model image is denoted $`M`$. Next, a โcleanโ galaxy image $`G`$ is produced in the following way. Objects bluer or redder than the primary galaxy are masked, by dividing the $`R`$ and $`B`$ band images, median filtering, and identifying pixels deviating more than a factor two from the median color of the galaxy. To remove small foreground and background objects and retain the smooth galaxy light a โreverseโ unsharp masking technique is used: the galaxy images are compared to Gaussian-smoothed versions of themselves, and pixels deviating more than a factor two are masked (excluding the galaxy centers). Pixels in the vicinity of masked pixels are also masked. Finally, a fractional distortion image $`F`$ is created by dividing $`G`$ by $`M`$. The distortion image is convolved with a $`5\times 5`$ median filter to reduce pixel-to-pixel variations. The parameter $`t`$ describing the level of distortion is defined as
$$t=\overline{\left|F_{x,y}\overline{F_{x,y}}\right|},$$
(2)
The tidal parameter thus measures the median absolute deviation of the (fractional) residuals from the model fit.
The procedure is illustrated in Fig. 5, for a galaxy pair with no visible distortions, a galaxy with a weak tidal feature, and a strongly disturbed object. The visually identified tidal features are isolated and emphasized in the distortion images $`F`$ shown at right. The corresponding values of $`t`$ vary from 8 % for the undistorted object to 24 % for the strongly disturbed object.
In Fig. 5.3(a) the values of $`t`$ are compared to the visual tidal classifications for all 86 bulge-dominated early-type galaxies. There is a strong correlation: the median value of $`t`$ is 0.08 for galaxies classified as undisturbed, 0.13 for weakly disturbed galaxies, and 0.19 for strongly disturbed galaxies and ongoing mergers. Tidal features are usually visually identified if $`t>0.1`$: of 23 galaxies with $`t<0.1`$ only 5 (22 %) were visually classified as tidally distorted, compared to 56 out of 63 galaxies with $`t0.1`$ (89 %). We conclude that the visually identified distortions are robust and imply median absolute deviations from an ellipse fit of $`10`$ %.
The distortion maps $`F`$ can be used to estimate the amount of star light associated with the tidal features. Visual inspection of the distortion maps suggests that pixels deviating more than 15 % from the model are usually associated with visible features. A mask $`F^{}`$ was created by setting pixels with values $`0.15`$ in $`F`$ to 1 and all other pixels to zero. The relative flux in the tidal features was then estimated as follows:
$$f_t=\frac{{\displaystyle F_{x,y}^{}\times (G_{x,y}M_{x,y})}}{{\displaystyle M_{x,y}}},$$
(3)
with the summations over all pixels $`x,y`$. As expected, the median value of $`f_t`$ is low for the 25 E/S0 galaxies with $`t<0.1`$: 0.01, about 1 % of the galaxy light. The median is 0.04 for the 63 galaxies with $`t0.1`$ and 0.07 for the 14 most disturbed galaxies with $`t0.2`$. From varying the cutoff in $`F`$ the systematic uncertainty in these values is estimated at $`30`$ %. We infer that โ despite their large extent โ the tidal features typically contain only about $`5`$ % of the total light of the galaxies.
### 5.4 Ongoing Mergers
There are 23 galaxies in the sample that show a clear tidal connection to a secondary object. In four cases the secondary object is also in the sample of 126, leaving nineteen unique systems. The fact that tidal features have developed implies that dynamical friction is already at work. Assuming that the galaxies were initially on nearly parabolic orbits, which is a reasonable assumption for field galaxies, the implication is that virtually all these galaxies will eventually merge (J. Barnes, private communication). We note that our selection of merger pairs by tidal features is different from โstandardโ photometric and spectroscopic selection of close pairs, where the number of mergers is always smaller than the number of close pairs (e.g., Patton et al. 2002): as expected, there are several close pairs in our sample which do not show tidal features and are not classified as ongoing mergers (see, e.g., the top panels of Fig. 5).
The 19 merger systems are shown in Fig. 6; color images of these objects are in the Appendix, along with all other galaxies. In most cases the primary galaxy is in the sample of 126 and the secondary galaxy is not (because it has $`R>17`$). In cases where both galaxies are in the sample, and in cases where the primary galaxy has a double nucleus, we designated the brightest object or nucleus as the primary object. The primary galaxies in the merging pairs are equally bright as other galaxies in the sample: their median $`R`$ magnitude is 16.3, compared to 16.4 for the full sample of red galaxies.
The merger sample comprises objects connected by a tidal โbridgeโ (e.g., 1-2874, 5-2345, 18-485/522); double nuclei in a highly disturbed common envelope (cdfs-374, cdfs-1100, 4-1190, 11-962); and objects of similar brightness with disturbed isophotes and tidal tails or fans (e.g., 2-3070/3102, 4-1975, 7-4247, 17-596/681, 19-2206/2242). Visual inspection of the merging systems indicates that the mergers are usually not between blue, disk-dominated systems but between red, bulge-dominated systems. Furthermore, in many cases the secondary galaxy appears to be of similar brightness as the primary galaxy. Eight of the 126 galaxies are merging with each other: the four bright, red pairs are 2-3070/3102, 17-596/681, 18-485/522, and 19-2206/2242.
We quantified these effects in the following way. The luminosity ratio is defined as
$$\frac{L_2}{L_1}=10^{(R_1R_2)/2.5},$$
(4)
and the color difference as
$$\mathrm{\Delta }(BR)=(B_2R_2)(B_1R_1).$$
(5)
In the 16 cases where the primary and the secondary galaxy are both in the SExtractor catalog we calculated the luminosity ratio and the color difference directly (see ยง 4.1). In three cases the pair is too close to be separated into two objects by SExtractor (cdfs-374, cdfs-1100, and 11-962). These objects are perhaps better described as single galaxies with double nuclei. Photometry for the nuclei was obtained from aperture photometry, with the radius of the aperture equal to half the distance between the nuclei. We note that the double nucleus of 11-962 is elongated but unresolved in the NDWFS images: the enlarged image shown in Fig. 6 was obtained with the OPTIC camera (Tonry et al. 2005) on the WIYN telescope, and has a resolution of $`0\stackrel{}{\mathrm{.}}45`$ FWHM.
As a check on the robustness of the results we also obtained aperture photometry in fixed $`5\mathrm{}`$ apertures for the 16 pairs that are well separated. This method underestimates the luminosity differences, as the secondary galaxies are usually more compact than the primary galaxies. Nevertheless, the results are very similar, giving a median luminosity ratio that is only $`20`$ % higher than derived from the SExtractor AUTO magnitudes.
The pair photometry is given in Table 1. Luminosity ratios range from 0.04 to 0.80, i.e., from relatively minor 1:25 accretion events to nearly equal mass mergers. The median luminosity ratio is 0.31, or a 1:3 merger. We conclude that approximately half the ongoing interactions are major mergers. The median color difference is only $`0.07`$, i.e., companions are typically 0.07 magnitudes bluer than the primary galaxy. As the companions are (by definition) fainter than the primary galaxies, this small difference may be an effect of the existence of the color-magnitude relation. Figure 7(a) shows the relation between color difference and luminosity ratio. There is a correlation, with the faintest companions being the bluest. The dashed line shows the expected correlation for elliptical galaxies on the (galaxy cluster) color-magnitude relation (Lรณpez-Cruz, Barkhouse, & Yee 2004). In Fig. 7(b) we show the distribution of residual color differences after subtracting this relation. The median of the distribution is $`0.02`$ and the rms scatter is only $`0.13`$ magnitudes.
These results show that the ongoing mergers in our sample occur โwithinโ the red sequence with very small scatter. This result cannot be attributed to our selection criteria: as 16 of the 19 pairs are separated by SExtractor there is no obvious bias against detecting blue companions to the primary galaxies. There is a hint that the observed relation in Fig. 7(a) is steeper than expected from the color-magnitude relation, but this is largely due to the two pairs with $`L_2/L_10.05`$ and the spiral/S0 merger 18-485/522. Among pairs with $`L_2/L_1>0.1`$ whose primary galaxy is classified as early-type the residual scatter in $`\mathrm{\Delta }(BR)`$ is only 0.08 magnitudes.
TABLE 1
Photometry of Merger Pairs
| Object | $`L_2/L_1`$ | $`\mathrm{\Delta }\left(BR\right)`$ |
| --- | --- | --- |
| cdfs-374 | 0.39 | $`0.08`$ |
| cdfs-1100 | 0.35 | $`0.12`$ |
| cdfs-6976 | 0.23 | $`0.14`$ |
| 1-2874 | 0.05 | $`0.50`$ |
| 2-3070 | 0.63 | $`0.05`$ |
| 4-1190 | 0.39 | $`0.02`$ |
| 4-1975 | 0.51 | $`0.16`$ |
| 4-2713 | 0.13 | $`0.29`$ |
| 5-2345 | 0.04 | $`0.39`$ |
| 7-4247 | 0.21 | $`0.03`$ |
| 11-962 | 0.69 | $`0.00`$ |
| 11-1278 | 0.15 | $`0.07`$ |
| 11-1732 | 0.20 | $`0.29`$ |
| 14-1401 | 0.24 | $`0.03`$ |
| 17-596 | 0.65 | $`0.02`$ |
| 18-485 | 0.80 | $`0.22`$ |
| 19-2206 | 0.41 | $`0.08`$ |
| 22-2252 | 0.31 | $`0.07`$ |
| 26-2558 | 0.20 | $`0.10`$ |
## 6 Discussion
### 6.1 Why Were They Missed?
The key results of our analysis are the large number of ongoing red mergers and the ubiquity of tidal features associated with (in particular) bulge-dominated early-type galaxies (galaxies classified as โE/S0โ). As discussed in ยง 5.2 the large scale, low surface brightness features that we see are different from the ripples and shells that have been reported previously in unsharp-masked images of ellipticals at $`z0`$. Our sample is essentially local, in the sense that the merger rate is not expected to evolve significantly from $`z0.1`$ to $`z0`$. Nearby galaxies have been studied in great detail over the past decades, and before discussing the consequences of our findings we address the question why this preponderance of smooth red tidal features has not been seen before.
The subjective nature of the visual classifications can be ruled out as a cause given the consistency of visual and quantitative classifications: 73 % of E/S0 galaxies have $`t>0.1`$. Small number statistics also play a minor role given the large number of objects in the sample. Given the fact that more than 90 % of the sample comes from the NDWFS we cannot rule out that the fraction of interacting galaxies is unusually high in that region of the sky. Random $`3^{}\times 3^{}`$ SDSS fields (see ยง 4.4) show substantial peaks in the redshift distribution within each field, reflecting the fact that most red galaxies live in groups and filaments. We note, however, that the red galaxies are distributed rather uniformly over the NDWFS area, and that some of the most spectacular mergers are in the unrelated MUSYC fields.
The most likely reason why the ubiquity of red mergers in the local universe was missed so far is the depth and uniformity of the available imaging. Classic imaging studies of nearby elliptical galaxies (e.g., Franx, Illingworth, & Heckman 1989; Peletier et al. 1990) used exposure times in the red of only 180 s โ 600 s on 1 m โ 2 m class telescopes, and such short exposures have been the norm ever since (e.g., Zepf, Whitmore, & Levison 1991; Pildis, Bregman, & Schombert 1995; Jansen et al. 2000; Colbert et al. 2001). Similarly, the effective exposure time of the Sloan Digital Sky Survey is about 51 s per filter (York et al. 2000).
Among the deepest imaging surveys are those of Malin & Carter (1983), who performed extreme enhancements of photographic plates to bring out low surface brightness features, and Schweizer & Seitzer (1992), who used 2400 s โ 3600 s integrations with the 0.9 m telescope on Kitt Peak. These studies reach AB surface brightness limits of $`\mu 26.5`$ mag arcsec<sup>-2</sup>. Although both studies find sharp distortions in the form of ripples and shells in a large fraction of ellipticals, their depth and field-of-view are not sufficient to find the large low surface brightness features that we report here. The exposure time that went into each of the images discussed here is $`27,000`$ s on 4 m class telescopes equipped with modern CCDs (equivalent to 120 hours on a 1 m class telescope), and to our knowledge such long exposures have never been obtained of a significant sample of nearby red galaxies.
Figure 6.1 shows what Fig. 4 would look like if we had exposed for 600 s on a 1 m class telescope. Very little remains of the dramatic tidal features evident in Fig. 4, illustrating the extreme depth of the NDWFS and MUSYC surveys. We quantified the effect of the S/N ratio on the detectability of tidal features by artificially increasing the noise in our images by a factor of ten, corresponding to a decrease in exposure time of a factor of 100. The degraded images are of similar depth as those of Malin & Carter (1983) and Schweizer & Seitzer (1992), and are still substantially deeper than virtually all other imaging studies of nearby elliptical galaxies. As shown in Fig. 5.3(b) the distribution of $`t`$ changes substantially when the S/N is decreased. Only 29 % of E/S0s in the degraded images have $`t>0.1`$, compared to 73 % in the original images.
We conclude that the high incidence of tidal features in our sample is a direct consequence of the faint surface brightness levels reached by the data. We also note that flat fielding uncertainties may prohibit the detection of broad tidal features around local galaxies irrespective of exposure time. The faint debris in Fig. 4 has a very large extent in comparison to the high surface brightness regions visible in Fig. 6.1, and moving the galaxies to $`z0.01`$ would change the surface brightness very little but increase the sizes of the debris fields to $`10^{}`$ or more. Finally, galaxy surveys in blue filters (e.g., Arp 1966) would classify the majority of objects in our sample as undisturbed even if they met the surface brightness and flat fielding requirements.
### 6.2 Recent Merger History of Bulge-Dominated Galaxies
We first determine the implications of the observed interactions only, without applying corrections for the short duration of the mergers. We define a sample which contains both current and future bulge-dominated galaxies, consisting of the 86 bulge-dominated galaxies (classified as E/S0), minus half of the six E/S0s which are interacting with each other, plus the disk-dominated galaxies (classified as S or S0) which are involved in a major merger. Among disk-dominated galaxies the only major merger is the spiral/S0 pair 18-485/522, whose constituent galaxies have roughly equal luminosity. The total sample of current and future bulge-dominated galaxies is therefore $`863+1=84`$. Among this sample of 84 galaxies there are eighteen current mergers and 41 remnants, and we infer that 70 % of current and future bulge-dominated galaxies experienced a merger or accretion event in the recent past.
The red ongoing mergers and the red tidal features associated with many of the E/S0 galaxies very likely sample the same physical process at different times. On average, galaxies with strong tidal features are probably observed shortly after the merger and galaxies with weak features are observed at later times. Assuming that the ongoing mergers are representative for the progenitors of all remnants we can directly infer the mass ratios of the progenitors of the full sample of 59 current and future bulge-dominated early-type galaxies. The median luminosity ratio of the ongoing mergers is 0.31, and the median color difference is negligible after correcting for the slope of the color-magnitude relation. Assuming that $`M/LM^{0.2}`$ (e.g., Jรธrgensen et al. 1996) a luminosity ratio of 0.31 implies a median mass ratio of 0.23, or a 1:4 merger.
There is also indirect evidence that the progenitors of the remnants were typically major mergers rather than low mass accretion events. Simulations by Johnston, Sackett, & Bullock (2001) show that surface brightness levels $`28<\mu <33`$ are typical for the debris of small satellites such as the Local Group dwarfs. Similarly, the average surface brightness of the giant stream of M31 is $`\mu _V30`$ mag arcsec<sup>-2</sup> (Ibata et al. 2001) and that of the debris from the Sagittarius dwarf $`\mu _V31`$ mag arcsec<sup>-2</sup> (Johnston et al. 2001). All these values are well beyond the detection limit of our survey. Furthermore, the fraction of tidally disturbed galaxies is much lower among disk-dominated galaxies than among bulge-dominated galaxies. Ignoring the ongoing mergers only 8 % of disk-dominated galaxies show tidal features compared to 62 % of bulge-dominated galaxies. This difference is consistent with the idea that the events responsible for the tidal features were usually sufficiently strong to disturb any dominant disk component. Finally, the fact that the features are typically broad and red suggests that the progenitors were dynamically hot systems with old stellar populations, consistent with the properties of the galaxies in the ongoing mergers. We also note that the fraction of light associated with the features ($`5`$ %) is very similar to the fraction of light associated with the tidal debris in the 1:3 merger simulation shown in Fig. 1.
We conclude that approximately 35 % of bulge-dominated red galaxies experienced a major merger with mass ratio $`>1:4`$ in the time window probed by our observations. This result is direct observational confirmation of the hierarchical assembly of massive galaxies.
### 6.3 Merger Rate and Mass Accretion Rate
Up to this point the analysis did not require an estimate of the timescale of the mergers. Such estimates are obviously uncertain, but they are necessary for turning the merger fraction into a merger rate and a mass accretion rate. These numbers are more easily compared to models and other observational studies, and are needed for extrapolating the results to higher redshifts.
The merger rate can be defined in a variety of ways (see, e.g., Patton et al. 2002). Here it is expressed as the number of remnants that are formed per Gyr within our selection area:
$$R=\frac{f_\mathrm{m}}{t_\mathrm{m}}\mathrm{Gyr}^1,$$
(6)
with $`f_\mathrm{m}`$ the fraction of the galaxy population involved in a merger โ with pairs counted as single objects โ and $`t_\mathrm{m}`$ a characteristic timescale for the mergers. We restrict the analysis to the sample of nineteen ongoing mergers, as simulations of the fading of tidal debris around the remnants of dry mergers have not yet been done in a systematic way.
Following Patton et al. (2000) we assume that the timescale of the mergers can be approximated by the dynamical friction timescale, given by
$$T_{\mathrm{fric}}=\frac{2.64\times 10^5r^2v_c}{M\mathrm{ln}\mathrm{\Lambda }}\mathrm{Gyr},$$
(7)
where $`r`$ is the physical separation of the pairs, $`v_c`$ is the circular velocity, $`M`$ is the mass of the lowest mass galaxy, and $`\mathrm{ln}\mathrm{\Lambda }`$ is the Coulomb logarithm (see Binney & Tremaine 1987; Patton et al. 2000). The median projected separation of the paired galaxies is $`8\stackrel{}{\mathrm{.}}6`$, or $`16\pm 3`$ kpc at $`z=0.10\pm 0.02`$. Assuming random orientations this corresponds to a median physical separation $`r=20\pm 4`$ kpc. To obtain an estimate of $`v_c`$ we assume that the pairs have similar line-of-sight velocity differences as the seven elliptical-elliptical mergers discussed in Combes et al. (1995). The mean velocity difference of the Combes et al. pairs $`\mathrm{\Delta }v=296\pm 91`$ km s<sup>-1</sup>(with the error determined by the jackknife method), implying $`v_c=\sqrt{3}\mathrm{\Delta }v=513\pm 158`$ km s<sup>-1</sup> for an isotropic velocity distribution. The median mass of the companions is $`(7\pm 3)\times 10^{10}M_{}`$, where we used $`M/L_R=4.6`$ in Solar units to convert luminosity to mass (van der Marel 1991) and assumed $`z=0.10\pm 0.02`$. Taking $`\mathrm{ln}\mathrm{\Lambda }2`$ (following Dubinski et al. 1999 and Patton et al. 2000) we obtain $`T_{\mathrm{fric}}=0.4\pm 0.2`$ Gyr. With $`f_m=19/122=0.16\pm 0.03`$ we obtain $`R=0.4\pm 0.2`$ Gyr<sup>-1</sup>.
The effect of the mergers on the mass evolution of red galaxies not only depends on the merger rate but also on the mass change resulting from individual mergers. The mass accretion rate can be approximated by
$$\mathrm{\Delta }M/M=R\times \overline{M_2/M_1}\mathrm{Gyr}^1,$$
(8)
with $`\overline{M_2/M_1}`$ the median mass ratio of the mergers. As shown in ยง 6.2 this ratio is approximately 0.23. With $`R_\mathrm{m}=0.4`$ we obtain $`\mathrm{\Delta }M/M=0.09\pm 0.04`$ Gyr<sup>-1</sup>, i.e., merging increases the masses of galaxies on the red sequence by $`10`$ % every $`10^9`$ years.
For comparison to other studies it is also of interest to consider the major merger fraction within a projected separation of 20 kpc. There are seven red pairs with luminosity ratio $`>0.3`$ and projected separation $`<11\mathrm{}`$, corresponding to a fraction of $`0.06\pm 0.02`$. It should be stressed that this number refers to mergers within the red sequence, not to the merger fraction within the full sample of $`R<17`$ galaxies. The colors of the well-separated pairs show that red galaxies โpreferโ to merge with other red galaxies. We have not examined the prevalence of mergers among luminous blue galaxies, but as discussed in ยง 1 the stellar populations of ellipticals rule out widespread major mergers among this population. Therefore, the major merger rate in the full sample of red and blue galaxies is presumably much lower than that within the restricted sample of red galaxies. A very rough estimate of the merger fraction in the full sample is $`0.06\times N_{\mathrm{red}}/N_{\mathrm{total}}0.02`$, in reasonable agreement with previous studies of close pairs (see, e.g., Patton et al. 2002, Lin et al. 2004, and references therein).
### 6.4 Effects on the Evolution of the Luminosity Density
At redshifts $`z<1`$, the observed evolution of the luminosity function of red galaxies reflects passive evolution of the stellar populations and possible changes in the underlying mass function. As discussed in, e.g., McIntosh et al. (2005) these changes can be due to mergers, galaxies entering the red sample due to changes in their star formation rate, or other effects. The best available constraints on the evolution of the luminosity function of red galaxies were derived by Bell et al. (2004), using the COMBO-17 survey. They find that the luminosity density of luminous red galaxies is approximately constant out to $`z1`$, which is surprising given the expected evolution of a factor of $`34`$ in the $`M/L`$ ratios of the galaxies (e.g., van Dokkum et al. 1998a, Treu et al. 2005, van der Wel et al. 2005). A possible explanation is that the underlying stellar mass density evolves as well, compensating for passive evolution of the stellar populations (Bell et al. 2004).
We first determine the effect of the observed red mergers only, i.e., the 52 % of red galaxies that are merger remnants or a merger pair. Dry mergers have no effect on the total luminosity density, but they have a strong effect on the luminosity density of galaxies brighter than a fixed magnitude. The effect of a single generation of mergers can be approximated by
$$j(z)\left(1f_\mathrm{m}\frac{L_2}{L_1+L_2}\right)j(0),$$
(9)
with $`j(z)`$ the luminosity density of luminous galaxies before the mergers and $`j(0)`$ the luminosity density after the mergers. For $`f_\mathrm{m}=0.52`$ and $`L_2/L_1=0.3`$ we find $`j(z)0.88j(0)`$. We use Monte Carlo simulations to test this approximation for a fiducial luminosity function with $`\alpha =0.6`$ and $`M_{}=19.9`$ (Bell et al. 2004). The luminosity function is evolved backward in time by breaking 50 % of the galaxies into pieces with luminosity ratio 0.3. Calculating the luminosity density for $`M<19`$ gives $`j(z)=0.89j(0)`$, in very good agreement with the simple estimate given above. The conclusion is that the effect of the observed mergers and remnants on the luminosity density of bright galaxies is small, of order 10 %.
The observed interactions only probe a relatively short period of at most a few Gyr. In order to extrapolate the effect of the mergers back in time we assume the following: 1) the mass accretion rate at $`z=0.1`$ is $`\mathrm{\Delta }M/M=0.09\pm 0.04`$ Gyr<sup>-1</sup> (see ยง 6.3); 2) the change in luminosity density due to mergers is proportional to the accreted luminosity; 3) all red galaxies are equally likely to undergo mergers; and 4) the mass accretion rate evolves as $`(1+z)^m`$. The value of $`m`$ is treated as a free parameter: observational constraints on the evolution of the merger rate may not be quite consistent (see, e.g., Patton et al. 2002, Concelice et al. 2003, Lin et al. 2004), and no studies have specifically considered the evolution of the pair fraction among galaxies on the red sequence.
The dashed lines in Fig. 8 show the predicted merger-driven evolution of the luminosity density of bright red galaxies with these assumptions, for three values of $`m`$. The model with $`m=0`$ has a constant accretion rate, and in the model with $`m=1.5`$ the accretion rate is $`3\times `$ higher at $`z=1`$ than it is at $`z=0`$. The dotted lines show the evolution of the $`M/L_B`$ ratio of field early-type galaxies, as measured by van der Wel et al. (2005). These authors find $`\mathrm{\Delta }\mathrm{ln}M/L_B=(1.20\pm 0.18)z`$ for galaxies with $`M>2\times 10^{11}M_{}`$ (appropriate for our sample), which is consistent with the independent measurement by Treu et al. (2005). The solid line shows the predicted evolution of $`j_B`$ when both mergers and $`M/L`$ evolution are taken into account. Grey bands indicate the combined uncertainties in $`\mathrm{\Delta }\mathrm{ln}M/L_B`$ and the merger rate. The predicted evolution of $`j_B`$ depends rather strongly on the assumed evolution of the mass accretion rate: it is positive for a non-evolving accretion rate, constant for $`m=1`$, and negative for $`m>1`$.
Solid points show the luminosity density in luminous red galaxies as measured by Bell et al. (2004) (their Fig. 5). As discussed extensively by these authors the data are inconsistent with passive evolution alone (dotted curves). However, as can be seen in Fig. 8 the data are fully consistent with models that include the mass accretion rate measured here. The uncertainties in the datapoints, the measured $`M/L`$ evolution, and the merger rate are too large to distinguish between models with different values of $`m`$, although the data at $`z>0.8`$ seem to favor models with $`m>0`$.
We infer that the cumulative effects of dry mergers may be the dominant cause of the constant luminosity density of luminous red sequence galaxies in the COMBO-17 survey. There are several caveats: both the datapoints in Fig. 8 and the $`z=0.1`$ mass accretion rate have substantial errors, leading to a wide range of allowed models; although the most recent generation of mergers appears to be (nearly) dissipationless, this may no longer hold at $`z0.5`$; and the fact that disk-dominated red galaxies probably evolve differently from bulge-dominated ones is ignored. The lack of evolution in the luminosity density of bright red galaxies, if confirmed, is probably due to a combination of effects with dry merging a significant, but not the only, contributor (see also McIntosh et al. 2005).
## 7 Conclusions
From an analysis of tidal features associated with bright red galaxies we find that $`70`$ % of bulge-dominated galaxies experienced a merger with median mass ratio 1:4 in the recent past. Expressed in other ways, $`35`$ % of bulge-dominated galaxies experienced a major merger involving $`>20`$ % of its final mass, and the current mass accretion rate of galaxies on the red sequence $`\mathrm{\Delta }M/M=0.09\pm 0.04`$ Gyr<sup>-1</sup>. Assuming a constant or increasing mass accretion rate with redshift it is inferred that the stellar mass density in luminous red galaxies has increased by a factor of $`2`$ over the redshift range $`0<z<1`$.
Neither of the two standard paradigms for elliptical galaxy formation appears to be consistent with our results: โmonolithicโ assembly at high redshift or late assembly via mergers of gas-rich disk galaxies. Instead, elliptical galaxies appear to have been assembled in mergers of bulge-dominated, red galaxies. This โdryโ form of merging is qualitatively consistent with the high central densities of ellipticals (Gao et al. 2004), their red colors and uniform properties, the existence of the $`M_{}\sigma `$ relation (Wyithe & Loeb 2005), and their specific frequency of globular clusters.
It remains to be seen whether widespread dry mergers are consistent with the slope and scatter of the color-magnitude relation (see, e.g., Bower et al. 1998, Ciotti & van Albada 2001). The median merger in our sample makes the brightest galaxy more luminous by $`30`$ % and bluer by only $`0.02`$ mag in $`BR`$, and its remnant will lie within $`0.03`$ mag of the color-magnitude relation. However, it is doubtful whether multiple generations of such mergers can be accommodated, unless there is a strong correlation between the masses of the progenitors (see Peebles 2002). It will also be interesting to see whether the evolution is consistent with the number density of massive galaxies at high redshift. Again, the observed mergers only have a $`10`$ % effect on the mass function of red galaxies, and the extrapolation to higher redshifts is obviously still very uncertain. Furthermore, the evolution of the mass function of all galaxies may be different from that of the subset of galaxies on the red sequence.
The high merger rate confirms predictions from hierarchical galaxy formation models in a $`\mathrm{\Lambda }`$CDM universe (e.g., Kauffmann et al. 1993; Kauffmann 1996; Cole et al. 2000; Somerville, Primack, & Faber 2001; Murali et al. 2002). Furthermore, semi-analytical models have predicted that gas-poor mergers between bulge-dominated systems, rather than mergers of disk systems, are responsible for the formation of the most massive ellipticals (Kauffmann & Haehnelt 2000; and, in particular, Khochfar & Burkert 2003). Quantitative comparisons of mass growth are difficult as current models do not naturally produce red field galaxies without star formation: the colors of field ellipticals in the simulations are too blue and their $`M/L`$ ratios too low (e.g., Kauffmann 1996; van Dokkum et al. 2001a). Additional mechanisms, such as heating by active nuclei, appear to be required to halt gas cooling and star formation in massive galaxies (e.g., Binney 2004, Somerville 2004, Dekel & Birnboim 2005, Keres et al. 2005). Our results imply that revised models which address these issues should not only reproduce the โred and deadโ nature of ellipticals today but also of their immediate progenitors โ which may occur naturally if the progenitors have masses greater than some critical mass (see Cooray & Milosavljeviฤ 2005).
The main uncertainty in the observed merger fraction is the possibility that the Northern NDWFS field, which contains over 90 % of the sample, is special in its frequency of tidally disturbed objects. This seems unlikely given its area of $`400`$ Mpc<sup>2</sup> at $`z=0.1`$, but the issue of field-to-field variations in the merger fraction will only be settled conclusively when independent fields of similar size are studied in the same way. The main uncertainty in the merger rate and mass accretion rate is the timescale of the mergers. Although our estimates broadly agree with other studies (see, e.g., Lin et al. 2004, and references therein), this may simply reflect the fact that similar assumptions lead to similar results. Modeling of red, gas-poor mergers has not been done in a systematic way using modern techniques, and it will be interesting to see what the timescales are for the initial coalescence and the subsequent surface brightness evolution of tidal debris. Specifically, modeling of the 19 merging systems and 44 remnants presented here would provide much better constraints on the merger rate and mass accretion rate, particularly when more complete redshift information is available.
More detailed observational studies of the mergers and their remnants may help answer the question why the mergers are red, i.e., what made the progenitors lose their gas? If active nuclei prevent the cooling of gas above some critical mass at early times they may play the same role during mergers, and it will be interesting to compare the degree of nuclear activity in undisturbed galaxies, ongoing mergers, and remnants. Also, sensitive diagnostics of young populations (e.g., H$`\delta `$ line strengths and ultra-violet photometry) can provide better constraints on the star formation histories of the mergers and remnants. Finally, studies with higher spatial resolution can provide information on the detailed isophotal shapes (boxy or disky) of the remnants and their progenitors, and on possible correlations between large scale smooth distortions and the sharp ripples and shells that have been reported in $`z0`$ ellipticals (see Hernquist & Spergel 1995).
Our study focuses on events that we can see today, and it will be very interesting to push the analysis to higher redshifts. Although the most recent generation of mergers could be largely โdryโ, previous generations likely involved blue galaxies and/or were accompanied by strong star formation (see, e.g., Sanders et al. 1988). Unfortunately it will be difficult to identify the broad red tidal features that we see here at significantly higher redshift due to the $`(1+z)^4`$ cosmological surface brightness dimming. The limiting depth of the MUSYC and NDWFS images is $`29`$ mag arcsec<sup>-2</sup> ($`1\sigma `$, AB). An equivalent survey at $`z1`$ should cover an area of $`1`$ square degree and reach levels of $`31.5`$ mag arcsec<sup>-2</sup> at $`z=1`$ and $`33.5`$ mag arcsec<sup>-2</sup> at $`z=2`$. Even when (unfavorable) $`K`$-corrections are ignored these requirements are well beyond the capabilities of current ground- or space-based telescopes. A more viable technique is to focus on the fraction of red galaxies in pairs. Although pair statistics require large corrections due to the short timescale of the mergers, pairs are easily detectable out to high redshift with the Hubble Space Telescope (see van Dokkum et al. 1999). Based on the work presented here the merger fraction among galaxies on the red sequence should be $`(0.06\pm 0.02)\times (1+z)^m`$ for separations $`<20`$ kpc and luminosity ratios $`0.3`$.
This paper was made possible by the dedicated efforts of all the individuals behind the NOAO Deep Wide-Field Survey and the Multi-wavelength Survey by Yale-Chile. Particular thanks go to Buell Januzzi and Arjun Dey for initiating and executing the NDWFS and to Eric Gawiser, who is the driving force behind MUSYC. David Herrera is responsible for reducing most of the MUSYC optical imaging data. Marijn Franx and Jeff Kenney are thanked for useful discussions. Comments of the anonymous referee improved the paper significantly. The NDWFS is supported by the National Optical Astronomy Observatory (NOAO). NOAO is operated by AURA, Inc., under a cooperative agreement with the National Science Foundation.
## Appendix A Catalog and Atlas
Here we provide coordinates, magnitudes, colors, morphologies, and tidal classifications for all 126 red galaxies. We also present images of all objects in the sample. These small images at fixed contrast level do not do the data justice, but space limitations prohibit an atlas on the scale of Figs. 3 and 4.
The atlas is available at www.astro.yale.edu/dokkum/mergers/ along with a version of the paper with much higher quality figures.
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# The K-group of Substitutional Systems
## 1.1. Bratteli diagram.
A Bratteli diagram is an infinite directed graph $`(V,E)`$, where $`V`$ is the vertex set and $`E`$ is the edge set. Both $`V`$ and $`E`$ are partitioned into non-empty disjoint finite sets
$$V=V_0V_1V_2\mathrm{}\mathrm{and}E=E_1E_2\mathrm{}$$
There are two maps $`r,s:EV`$ the range and source maps. The following properties hold:
1. $`V_0=\{v_0\}`$ consists of a single point, referred to as the โtop vertexโ of the Bratteli diagram
2. $`r(E_n)V_n,s(E_n)V_{n1},n=1,2,\mathrm{}`$. Also $`s^1(v)\varphi `$ $`vV`$ and $`r^1(v)\varphi `$ for all $`vV_1,V_2,\mathrm{}`$.
Maps between Bratteli diagrams are assumed to preserve gradings and intertwine the range and source maps. If $`vV_n`$ and $`wV_m`$, where $`m>n`$, then a path from $`v`$ to $`w`$ is a sequence of edges $`(e_{n+1},\mathrm{},e_m)`$ such that $`s(e_{n+1})=v,r(e_m)=w`$ and $`s(e_{j+1})=r(e_j)`$. Infinite paths from $`v_0V_0`$ are defined similarly. The Bratteli diagram is called simple if for any $`n=0,1,2,\mathrm{}`$ there exists $`m>n`$ such that every vertex of $`V_n`$ can be joined to every vertex of $`V_m`$ by a path.
## 1.2. Ordered Bratteli diagram.
An ordered Bratteli diagram $`(V,E,)`$ is a Bratteli diagram $`(V,E)`$ together with a linear order on $`r^1(v),vV\{v_0\}=V_1V_2V_3\mathrm{}`$. We say that an edge $`eE_n`$ is a maximal edge (resp. minimal edge) if $`e`$ is maximal (resp. minimal) with respect to the linear order in $`r^1(r(e))`$.
Given $`vV_n`$, it is easy to see that there exists a unique path $`(e_1,e_2,\mathrm{},e_n)`$ from $`v_0`$ to $`v`$ such that each $`e_i`$ is maximal (resp. minimal).
Note that if $`m>n`$, then for any $`wV_m`$, the set of paths starting from $`V_n`$ and ending at $`w`$ obtains an induced (lexicographic) linear order:
$$(e_{n+1},e_{n+2},\mathrm{},e_m)>(f_{n+1},f_{n+2},\mathrm{},f_m)$$
if for some $`i`$ with $`n+1im,e_j=f_j`$ for $`1<jm`$ and $`e_i>f_i`$.
## 1.3. Proper order.
A properly ordered Bratteli diagram is a simple ordered Bratteli diagram $`(V,E)`$ which possesses a unique infinite path $`x_{\mathrm{max}}=(e_1,e_2,\mathrm{})`$ such that each $`e_i`$ is a maximal edge and a unique infinite path $`x_{\mathrm{min}}=(f_1,f_2,\mathrm{})`$ such that each $`f_i`$ is a minimal edge.
Given a properly ordered Bratteli diagram $`B=(V,E,)`$ we denote by $`X_B`$ its infinite path space. So
$$X_B=\{(e_1,e_2,\mathrm{})e_iE_i,r(e_i)=s(e_{i+1}),i=1,2,\mathrm{}\}$$
For an initial segment $`(e_1,e_2,\mathrm{},e_n)`$ we define the cylinder sets
$$U(e_1,e_2,\mathrm{}e_n)=\{(f_1,f_2,\mathrm{})X_Bf_i=e_i,1in\}.$$
By taking cylinder sets to be a basis for open sets $`X_B`$ becomes a topological space. We exclude trivial cases (where $`X_B`$ is finite, or has isolated points). Thus, $`X_B`$ is a Cantor set. $`X_B`$ is a metric space, where for two paths $`x,y`$ whose initial segments to level $`m`$ agree but not to level $`m+1,d(x,y)=1/m+1`$.
## 1.4. Vershik map for a properly ordered Bratteli diagram.
If $`x=(e_1,e_2,\mathrm{}e_n,\mathrm{})X_B`$ and if at least one $`e_i`$ is not maximal define
$$V_B(x)=y=(f_1,f_2,\mathrm{},f_j,e_{j+1},e_{j+2},\mathrm{})X_B$$
where $`e_1,e_2,\mathrm{},e_{j1}`$ are maximal, $`e_j`$ is not maximal and has $`f_j`$ as successor in the linearly ordered set $`r^1(r(e_j))`$ and $`(f_1,f_2,\mathrm{},f_{j1})`$ is the minimal path from $`v_0`$ to $`s(f_j)`$. Extend the above $`V_B`$ to all of $`X_B`$ by setting $`V_B(x_{\mathrm{max}})=x_{\mathrm{min}}`$. Then $`(X_B,V_B)`$ is a Cantor minimal dynamical system.
Next, we describe the construction of a dynamical system associated to a non-properly ordered Bratteli diagram. The Bratteli diagram need not be simple. To motivate this construction, it is perhaps worthwhile to begin by indicating how it works in the case of an ordered Bratteli diagram associated to a nested sequence of Kakutani-Rohlin partitions of a Cantor dynamical system $`(X,T)`$.
## 1.5. K-R partition.
A Kakutani-Rohlin partition of the Cantor minimal system $`(X,T)`$ is a clopen partition $`๐ซ`$ of the kind
$$๐ซ=\{T^jZ_kkA\mathrm{and}0j<h_k\}$$
where $`A`$ is a finite set and $`h_k`$ is a positive integer. The $`k^{th}`$ tower $`๐ฎ_k`$ of $`๐ซ`$ is $`\{T^jZ_k0j<h_k\}`$ ; its floors are $`T^jZ_k,(0j<h_k)`$. The base of $`๐ซ`$ is the set $`Z=_{kA}Z_k`$.
Let $`\{๐ซ_n\},(n\text{N})`$ be a sequence of Kakutani-Rohlin partitions
$$๐ซ_n=\{T^jZ_{n,k}kA_n,\text{and }0j<h_{n,k}\},$$
with $`๐ซ_0=\{X\}`$ and with base $`Z_n=_{kA_n}Z_{n,k}`$. We say that this sequence is nested if, for each $`n`$,
1. $`Z_{n+1}Z_n`$
2. $`๐ซ_{n+1}`$ refines the partition $`๐ซ_n`$.
For the Bratteli-Vershik system $`(X_B,V_B)`$ of sections 1.3-1.4, one obtains a Kakutani-Rohlin partition $`๐ซ_n`$ for each $`n`$ by taking the sets in the partition to be the cylinder sets $`U(e_1,e_2,\mathrm{}e_n)`$ of section 1.3 and taking as the base of the partition the union $`U(e_1,e_2,\mathrm{}e_n)`$ over minimal paths (i.e., each $`e_i`$ is a minimal edge). This is a nested sequence.
1.6. To any nested sequence $`\{๐ซ_n\},(n\text{N})`$ of Kakutani-Rohlin partitions we associate an ordered Bratteli diagram $`B=(V,E,)`$ as follows (see \[DHS, section 2.3\]): the $`A_n`$ towers in $`๐ซ_n`$ are in $`11`$ correspondence with $`V_n`$, the set of vertices at level $`n`$. Let $`v_{n,k}V_n`$ correspond to the tower $`๐ฎ_{n,k}=\{T^jZ_{n,k}0j<h_{n,k}\}`$ in $`๐ซ_n`$. We refer to $`T^jZ_{n,k},0j<h_{n,k}`$ as floors of the tower $`๐ฎ_{n,k}`$ and to $`h_{n,k}`$ as the height of the tower. We will exclude nested sequences of K-R partitions where the infimum (over $`k`$ for fixed $`n`$) of the height $`h_{n,k}`$ does not go to infinity with $`n`$. Let us view the tower $`๐ฎ_{n,k}`$ against the partition $`๐ซ_{n1}=\{T^jZ_{n1,k}kA_{n1}`$, and $`0j<h_{n1,k}\}`$. As the floors of $`๐ฎ_{n,k}`$ rise from level $`j=0`$ to level $`j=h_{n,k}1`$, $`๐ฎ_{n,k}`$ will start traversing a tower $`๐ฎ_{n1,i_1}`$ from the bottom to the top floor, then another tower $`๐ฎ_{n1,i_2}`$ from the bottom to the top floor, then another tower $`๐ฎ_{n1,i_3}`$ likewise and so on till a final segment of $`๐ฎ_{n,k}`$ traverses a tower $`๐ฎ_{n1,i_m}`$ from the bottom to the top. Note that in this final step the top floor $`T^jZ_{n,k}`$ for $`j=h_{n,k}1`$ of $`๐ฎ_{n,k}`$ reaches the top floor $`T^qZ_{n1,i_m}`$ for $`q=h_{n1,i_m}1`$ of $`๐ฎ_{n1,i_m}`$ as a consequence of the assumption $`Z_nZ_{n1}`$ and the fact that $`T^1`$ (union of bottom floors) = union of top floors. Bearing in mind this order in which $`๐ฎ_{n,k}`$ traverses $`๐ฎ_{n1,i_1},๐ฎ_{n1,i_2},\mathrm{},๐ฎ_{n1,i_m}`$ we associate $`m`$ edges, ordered as $`e_{1,k}<e_{2,k}<\mathrm{}<e_{m,k}`$ and we set the range and source maps for edges by $`r(e_{j,k})=v_{n,k}`$ and $`s(e_{j,k})=v_{n1,i_j}`$. Note that $`m`$ depends on the index $`kA_n`$ (and that by convention the indexing sets $`A_n`$ are disjoint). $`E_n`$ is the disjoint union over $`kA_n`$ of the edges having range in $`V_n`$.
1.7. For $`xX`$, we define $`x_n๐ซ_n^\text{Z},n\text{N}`$ as follows: $`x_n=(x_{n,i})_{i\text{Z}}`$, where $`x_{n,i}๐ซ_n`$ is the unique floor in $`๐ซ_n`$ to which $`T^i(x)`$ belongs. If $`m>n`$, let $`j_{m,n}:๐ซ_m๐ซ_n`$ be the unique map defined by $`j_{m,n}(F)=F^{}`$ if $`FF^{}`$. (By abuse of notation, we use the same symbol $`F`$ to denote a point of the finite set $`๐ซ_m`$ and also to denote the subset of $`X`$, in the partition $`๐ซ_n`$, which $`F`$ represents). An important property of the map
$$X\underset{n}{}(๐ซ_n^\text{Z}),x(x_1,x_2,\mathrm{}),x_n=(x_{n,i})_{i\text{Z}},$$
defined above is the following:
1.8. If $`F`$ and $`TF`$ are two successive floors of a $`๐ซ_n`$-tower and if $`x_{n,i}=F`$ then $`x_{n,i+1}=TF`$. If $`x_{n,i}`$ is the top floor of a $`๐ซ_n`$-tower, then $`x_{n,i+1}`$ is the bottom floor of a $`๐ซ_n`$-tower. More importantly, given integers $`K`$ and $`n`$, there exist $`m>n`$ and a single tower $`๐ฎ_{m,k}`$ of level $`m`$ such that the finite sequence $`(x_{n,i})_{KiK}`$ is an interval segment contained in
$$\{j_{m,n}(T^{\mathrm{}}(Z_{m,k}))0\mathrm{}<h_{m,k}\}.$$
This is a consequence of the assumption that the infimum of the heights of level-$`n`$ towers goes to infinity. It is true that $`x_{n,i}=j_{m,n}(x_{m,i})`$, but the sequence $`(x_{m,i})_{KiK}`$ need not be an interval segment of $`\{T^{\mathrm{}}(Z_{m,k})0\mathrm{}<h_{m,k}\}`$.
The foregoing observations in the case of an ordered Bratteli diagram associated to a nested sequence of Kakutani-Rohlin partitions gives us the hint to define a dynamical system $`(X_B,T_B)`$ of a non properly ordered Bratteli diagram $`B=(V,E,)`$ as follows:
## 1.9. Definition.
For each $`n`$ define $`\varpi _n=`$ the set of paths from $`V_0`$ to $`V_n`$. There is an obvious truncation map $`j_{m,n}:\varpi _m\varpi _n`$ which truncates paths from $`V_0`$ to $`V_m`$ to the initial segment ending in $`V_n`$. For each $`vV_n`$, the set $`\varpi (v)`$ of paths from $`\{\}V_0`$ ending at $`v`$ will be called a โ$`\varpi _n`$-tower parametrised by $`v`$โ. Each tower is a linearly ordered set (whose elements may be referred to as floors of the tower) since paths from $`v_0`$ to $`v`$ acquire a linear order (cf. 1.2).We will exclude unusual examples of ordered Bratteli diagram where the infimum of the height of level-$`n`$ towers does not go to infinity, with $`n`$ (for example like \[HPS, Example 3.2\]). Now, we define
1.10. Definition. $`X_B=\{x=(x_1,x_2,\mathrm{},x_n,\mathrm{})\}`$ where
1. $`x_n=(x_{n,i})_{i\text{Z}}\varpi _n^\text{Z}`$,
2. $`j_{m,n}(x_{m,i})=x_{n,i}`$ for $`m>n`$ and $`i\text{Z}`$ and
3. given $`n`$ and $`K`$ there exists $`m`$ such that $`m>n`$ and a vertex $`vV_m`$, such that the interval segment $`x_n[K,K]:=(x_{n,K},x_{n,K+1},\mathrm{},x_{n,K})`$ is obtained by applying $`j_{m,n}`$ to an interval segment of the linearly ordered set of paths from $`v_0`$ to $`v`$.
The condition (iii) is the crucial part of the definition. Without it what one gets is an inverse system.
The condition (iii) implies that a property similar to (1.8) holds. Since each $`\varpi _n`$ is a finite set $`\varpi _n^\text{Z}`$ has a product topology which makes it a compact set - in fact a Cantor set. Likewise, $`_n(\varpi _n^\text{Z})`$ is again a Cantor set. Thus, $`X_B_n(\varpi _n^\text{Z})`$ has an induced topology. The lemma below and the following proposition are analogous to corresponding facts for the Vershik model associated to properly ordered Bratteli diagrams.
The following results (1.11) and (1.12) are proved in \[EP\].
## 1.11. Lemma.
The topological space $`X_B`$ is compact.
Denote by $`T_B`$ the restriction of the shift operator to $`X_B`$. So, if $`x=(x_1,x_2,\mathrm{},x_n,\mathrm{}),`$ where $`x_n=(x_{n,i})_{i\text{Z}}\varpi _n^\text{Z}`$, then $`T_B(x)=(x_1^{},x_2^{},\mathrm{},x_n^{},\mathrm{})`$, where $`x_n^{}=(x_{n,i}^{})_{i\text{Z}}\varpi _n^\text{Z}`$ and $`x_{n,i}^{}=x_{n,i+1}`$.
$`(X_B,T_B)`$ will be called the dynamical system associated to $`B=(V,E,)`$.
## 1.12. Proposition.
If $`B=(V,E,)`$ is a simple ordered Bratteli diagram, then $`(X_B,T_B)`$ is a Cantor minimal dynamical system.
## 1.13.
In (1.7), given a nested sequence of Kakutani-Rohlin partitions of $`(X,T)`$, we defined a map from $`(X,T)`$ to the dynamical system $`(X_B,T_B)`$ of the associated ordered Bratteli diagram. It follows that if $`(X,T)`$ is minimal, and if the Bratteli diagram of the nested sequence of K-R partitions is a simple Bratteli diagram, then $`(X,T)(X_B,T_B)`$ is onto. If the topology of $`(X,T)`$ is spanned by the collection of the clopen sets belonging to the K-R partitions then clearly the map $`(X,T)(X_B,T_B)`$ is injective. In particular, if the Bratteli diagram is properly ordered then the Bratteli-Vershik system is naturally isomorphic to the system given by our construction in 1.10.
## 1.14.
Note that the same term โtowersโ has been used to denote two different but related objects \[in (1.5) and (1.9)\]. For $`vV_n`$, let $`y`$ be a path from $`\{\}`$ to $`v`$ in $`(V,E,)`$. So, $`y`$ is a โfloorโ (consisting of the single element $`y`$) belonging to the $`\varpi _n`$\- tower $`\varpi (v)`$ (a finite set) parametrized by $`vV_n`$ \- all in the sense of $`(1.9)`$. Here, $`\varpi (v)`$= all paths from $`\{\}`$ to $`v`$. Put $`_y=\{x=(x_1,x_2,\mathrm{},x_n,\mathrm{})X_Bx_{n,0}=y\}`$. $`_y`$ is a clopen set of the Cantor set $`X_B`$. Put $`๐ซ_n=\{_yy\varpi (v),vV_n\}.`$ Then, in the sense of $`(1.5)`$ $`๐ซ_n`$ is a K-R partition of $`X_B`$ whose base is the union of $`_y,(y\text{ minimal }\varpi (v),vV_n)`$. Its towers $`S_v`$ are parametrized by $`vV_n`$: $`S_v=\{_yy\varpi (v)\}`$. $`_y,(y\varpi (v))`$ are the floors of the tower $`S_v`$. (We encountered this K-R partition earlier in the case of the Bratteli-Vershik system at the end of 1.5). The ordered Bratteli diagram obtained from $`\{_yy\varpi (v),vV_n\}`$ is $`(V,E,)`$.
2. The Bratteli diagram $`(V^๐ฌ,E^๐ฌ,)`$
2.1. We will now define two nested sequences of K-R partitions of $`X`$. For $`vV_n`$, let $`y`$ be a path from $`\{\}`$ to $`v`$ in $`(V,E,)`$. So, $`y`$ is a โfloorโ belonging to the $`\varpi _n`$\- tower $`\varpi (v)`$ parametrized by $`v`$. Put $`_y=\{x=(x_1,x_2,\mathrm{},x_n,\mathrm{})Xx_{n,0}=y\}`$
$$๐ซ_n=\{_yy\varpi (v),vV_n\}.$$
Then $`\{๐ซ_n\}_n`$ is a nested sequence of K-R partitions of $`X`$ . But, the topology of $`X`$ need not be spanned by the collection of clopen sets $`\{_y\},(y\varpi (v),vV_n,n\text{N})`$. In contrast, the topology of $`X`$ is indeed spanned by the collection of clopen sets in another nested sequence $`\{๐ฌ_n\}_n`$ of K-R partitions, defined below. Let $`\varpi =\varpi (u),\varpi ^{}=\varpi (v),\varpi ^{\prime \prime }=\varpi (w)`$ be three $`\varpi _n`$-towers and $`y`$ a floor of $`\varpi ^{}`$. For any $`xX`$ and for any $`n`$ if $`x_{n,i}`$ is a floor of a $`\varpi _n`$-tower $`\overline{\varpi }`$, then for some $`a,b\text{Z}`$ such that $`aib`$, the segment $`x_n[a,b]`$ is just the sequence of floors in $`\overline{\varpi }`$. We define
$`(\varpi ,\varpi ^{},\varpi ^{\prime \prime };y)=`$ the clopen subset of $`_y`$ consisting of the elements $`x=(x_1,x_2,\mathrm{},x_n,\mathrm{})`$ with the property that for some $`a_1<a_20<a_3<a_4\text{Z}`$, the segment $`x_n[a_1,a_21]`$ is the sequence of floors of $`\varpi `$, the segment $`x_n[a_2,a_31]`$ is the sequence of floors of $`\varpi ^{}`$ and the segment $`x_n[a_3,a_4]`$ is the sequence of floors of $`\varpi ^{\prime \prime }`$. Some of the sets $`(\varpi ,\varpi ^{},\varpi ^{\prime \prime };y)`$ may be empty, but the non-empty sets $`(\varpi ,\varpi ^{},\varpi ^{\prime \prime };y)`$ form a K-R partition which we denote by $`๐ฌ_n`$. For fixed $`\varpi ,\varpi ^{},\varpi ^{\prime \prime }`$ the subcollection $`\{(\varpi ,\varpi ^{},\varpi ^{\prime \prime };y)\}`$ as $`y`$ varies through the floors of $`\varpi ^{}`$, is a $`๐ฌ_n`$-tower parametrized by $`[u,v,w]`$. We denote this $`๐ฌ_n`$-tower by $`๐ฎ_{(\varpi ,\varpi ^{},\varpi ^{\prime \prime })}`$. The floors of the tower $`๐ฎ_{(\varpi ,\varpi ^{},\varpi ^{\prime \prime })}`$ are $`\{(\varpi ,\varpi ^{},\varpi ^{\prime \prime };y)\}`$ as $`y`$ runs through the sequence of floors of $`\varpi ^{}`$.
2.2. The tripling of $`(V,E,)`$. Let $`(V,E,)`$ be an arbitrary simple, ordered Bratteli diagram. Define $`(V^๐ฌ,E^๐ฌ,)`$ as follows: $`V_0^๐ฌ=\{\}`$, a single point.
$`V_n^๐ฌ`$ consists of triples $`(u,v,w)V_n\times V_n\times V_n`$ such that for some $`yV_m`$ where $`m>n`$, the level-$`m`$ tower $`\varpi (y)`$ passes successively through the level-$`n`$ tower $`\varpi (u)`$, then $`\varpi (v)`$ and then $`\varpi (w)`$. An edge $`\stackrel{~}{e}E_n^๐ฌ`$ is a triple $`(u,e,w)`$ such that $`e`$ is an edge of $`(V,E)`$ and $`(u,r(e),w)V_n^๐ฌ`$. Let
1. $`\{e_1,e_2,\mathrm{},e_k\}`$ be all the edges in $`r^1(r(e))`$,
2. $`\{f_1,f_2,\mathrm{},f_{\mathrm{}}\}`$ be all the edges in $`r^1(u)`$ and
3. $`\{g_1,g_2,\mathrm{},g_m\}`$ be all the edges in $`r^1(w)`$.
The sources of $`(u,e_1,w),(u,e_2,w),\mathrm{},(u,e_k,w)`$ are defined to be
$`(s(f_{\mathrm{}}),s(e_1),s(e_2)),(s(e_1),s(e_2),s(e_3)),\mathrm{},(s(e_{k1}),s(e_k),s(g_1))`$
respectively. The range of $`(u,e,w)`$ is of course $`(u,r(e),w)`$. If $`r^1(r(e))`$ is ordered as $`\{e_1,e_2,\mathrm{},e_k\}`$, we declare the ordering of $`r^1(r(u,e,w))`$ to be $`\{(u,e_1,w),(u,e_2,w),\mathrm{},(u,e_k,w)\}`$. The ordered Bratteli diagram $`(V^๐ฌ,E^๐ฌ,)`$ thus defined will be called the tripling of $`(V,E,)`$.
The map $`\pi :(V^๐ฌ,E^๐ฌ,)(V,E,)`$ given by $`(u,v,w)v,(u,e,w)e`$ enjoys the โunique path liftingโ property in the following sense. If $`m>n1`$, and $`(e_n,e_{n+1},\mathrm{},e_m)`$ is a path in $`(V,E)`$ from $`V_{n1}`$ to $`V_m`$ with $`r(e_m)=v`$ then for any $`(u,v,w)V_m^๐ฌ`$, there is a unique path $`(\stackrel{~}{e}_n,\stackrel{~}{e}_{n+1},\mathrm{},\stackrel{~}{e}_m)`$ in $`(V^๐ฌ,E^๐ฌ)`$ which maps onto $`(e_n,e_{n+1},\mathrm{},e_m)`$ under $`\pi `$ and such that $`r(\stackrel{~}{e}_m)=(u,v,w)`$. It is quite elementary to check that the map $`\pi :(V^๐ฌ,E^๐ฌ,)(V,E,)`$ induces an isomorphism between the corresponding dynamical systems given by 1.10 (\[EP, 2.20\]).
\[ Two different edges on the left with the same source may map into the same edge on the right. Two different edges on the left with the same range cannot map to the same edge on the right.\]
Let $`\{n_k\}_{k=0}^{\mathrm{}}`$ be a subsequence of $`\{0,1,2,\mathrm{}\}`$ where we assume $`n_0=0`$. A Bratteli diagram $`(V^{},E^{})`$ is called a โtelescopingโ of $`(V,E)`$ if $`V_k^{}=V_{n_k}`$ and $`E_k^{}`$ consists of paths $`(e_{n_{k1}+1},\mathrm{},e_{n_k})`$ from $`V_{n_{k1}}`$ to $`V_{n_k}`$ in $`(V,E)`$, the range and source maps being the obvious ones. It is easy to see that tripling is compatible with telescoping.
2.3. Stationary Bratteli diagrams. A Bratteli diagram is stationary if the diagram repeats itself after level $`1`$. (One may relax by allowing a period from some level onwards; but, a telescoping will be stationary in the above restricted sense.) If $`(V,E,)`$ is an ordered Bratteli diagram and the diagram together with the order repeats itself after level $`1`$, then $`(V,E,)`$ will be called a stationary ordered Bratteli diagram. We refer the reader to \[DHS, section (3.3)\] for the usual definition of a substitutional system and how they give rise to stationary Bratteli diagrams. Some details are recalled below. Let $`(V,E,)`$ be as above and suppose moreover that it is a simple Bratteli diagram. We have
1. an enumeration $`\{v_{n,1},v_{n,2},\mathrm{},v_{n,L}\}`$ of $`V_n,n1`$,
2. for $`n>1`$ and $`1jL`$ an enumeration $`\{e_{n,j,1},e_{n,j,2},\mathrm{},e_{n,j,a_j}\}`$ of $`r^1(e_{n,j})`$ which is assumed to be listed in the linear order in $`r^1(v_{n,j})`$,
3. in the enumerations above, $`L`$ does not depend on $`n`$ and $`a_j`$ depends only on $`j`$ and not on $`n`$. Moreover, the ordering in $`r^1(e_{n,j})`$ is stationary, i.e, if $`n,m>1`$, if $`1jL,1kL,1ia_j`$, then $`\mathrm{`}\mathrm{`}s(e_{n,j,i})=v_{n1,k}`$$`\mathrm{`}\mathrm{`}s(e_{m,j,i})=v_{m1,k}`$โ.
2.4. Substitutional systems. Let $`A`$ be an alphabet set. Write $`A^+`$ for the set of words of finite length in the alphabets of $`A`$. Let $`\sigma :AA^+`$ be a primitive aperiodic non-proper substitution, written, $`\sigma (a)=\alpha \beta \gamma \mathrm{}`$. The stationary ordered Bratteli diagram $`B=(V,E,)`$ associated to $`(A,\sigma )`$ (cf. \[DHS, section 3.3\] can be described as
$$V_n=A,\text{ }n1,V_0=\{\}$$
$$E_n=\{(a,k,b)a,bA,k\text{N},\text{ }a\text{ is the }k^{th}\text{ alphabet in the word }\sigma (b)\}.$$
(The reader who prefers a more carefully evolved notation can consider intoducing an extra factor โ$`\times \{n\}`$โ so that vertices and edges at different levels are seen to be disjoint). The source and range maps $`s`$ and $`r`$ are defined by $`s(a,k,b)=a,r(a,k,b)=b`$. In the linear order in $`r^1(b)`$, $`(a,k,b)`$ is the $`k^{th}`$ edge.
To the stationary ordered Bratteli diagram $`B`$ of $`(A,\sigma )`$ (which may not be properly ordered unless $`\sigma `$ is a primitive, aperiodic, proper substitution, โ see \[DHS, section 3\]) we can associate a dynamical system $`X_B`$ following the construction of 1.10; this is naturally isomorphic to the substitutional dynamical system $`(X_\sigma ,T_\sigma )`$ associated to $`(A,\sigma )`$ defined for example in \[DHS, section 3.3.1\]. (See \[EP, section 2.5\].)
2.5. Tripling for a substitutional system $`(A,\sigma )`$. Let $`(A,\sigma )`$ be a substitutional system and suppose $`B=(V,E,)`$ is the stationary ordered Bratteli diagram associated to $`(A,\sigma )`$. Define $`A^๐ฌ=\{(a,b,c)A\times A\times Aabc`$ occurs as a subword of $`\sigma ^n(d)`$ for some $`dA`$ and some $`n\}`$ . Define
$$\sigma ^๐ฌ:A^๐ฌ(A^๐ฌ)^+$$
by $`\sigma ^๐ฌ[(a,b,c)]=(a_m,b_1,b_2)(b_1,b_2,b_3)\mathrm{}(b_{n2},b_{n1},b_n)(b_{n1},b_n,c_1)`$, where $`\sigma (b)=b_1b_2\mathrm{}b_n`$, and $`a_m`$ is the last alphabet in $`\sigma (a)`$, while $`c_1`$ is the first alphabet in $`\sigma (c)`$. Then $`(V^๐ฌ,E^๐ฌ,)`$ is the stationary ordered Bratteli diagram associated to $`(A^๐ฌ,\sigma ^๐ฌ)`$.
3. The groups $`K^0(X,T),K_{}0(V,E,)`$ and $`K_0(V,E)`$.
3.1.Definition. Let $`(X,T)`$ be a Cantor minimal system. Let $`C(X,\text{Z})`$ be the space of integer valued continuous functions on $`X`$. Let
$$K^0(X,T)=C(X,\text{Z})/_TC(X,\text{Z})$$
where $`_T:C(X,\text{Z})C(X,\text{Z})`$ denotes the coboundary operator $`_T(f)=ffT`$. A function of the form $`ffT`$ is called a coboundary. Define the positive cone
$$K^0(X,T)^+=\{[f]fC(X,\text{Z}^+)\}$$
where $`[f]`$ denotes the projection modulo coboundaries. The ordered group $`(K^0(X,T),K^0(X,T)^+)`$ has a distinguished order unit, namely , the projection of the constant function 1.
Let $`(V,E)`$ be a Bratteli diagram and $`(V,E,)`$ the same thing equipped with a linear order on edges which makes it an ordered Bratteli diagram. As usual the dimension group $`K_0(V,E)`$ is defined to be the inductive limit of the system of ordered groups
$$\text{Z}^{|V_0|}\stackrel{A_0}{}\text{Z}^{|V_1|}\stackrel{A_1}{}\text{Z}^{|V_2|}\stackrel{A_2}{}\text{Z}^{|V_3|}\stackrel{A_3}{}\mathrm{}$$
where the positive homomorphism $`A_n`$ is given by matrix multiplication with the incidence matrix between levels $`n1`$ and $`n`$. The inductive limit $`K_0(V,E)`$ is endowed with the induced order, the positive cone being denoted by $`K_0(V,E)^+`$. The image of $`1\text{Z}^{|V_0|}`$ in $`(K_0(V,E),K_0(V,E)^+)`$ is an order unit.
On the other hand we define the group $`K_{}0(V,E,)`$ in the following way. Whenever we have $`mn`$ and two paths $`\varpi _1`$ and $`\varpi _2,(\varpi _1\varpi _2)`$ from $`V_n`$ to $`V_m`$ with the same range $`uV_m`$ define $`[\varpi _1,\varpi _2)`$ to be the set consisting of all paths from $`V_n`$ to $`V_m`$ lying between $`\varpi _1\text{ (included) and }\varpi _2\text{ (excluded) }`$ranging at $`u`$. Put $`B\text{Z}^{|V_n|}=\{\overline{m}=(m_k)_{kV_n}\text{Z}^{|V_n|}\mathrm{\Sigma }_{\varpi [\varpi _1,\varpi _2)}m_{s(\varpi )}=0\text{ for all }m\text{ and all }\varpi _1,\varpi _2`$ as above with the same source and same range$`\}`$. Observe that $`A_n(B\text{Z}^{|V_n|})B\text{Z}^{|V_{n+1}|}`$. Moreover, suppose $`\overline{p},\overline{q}\text{Z}^{|V_n|,+},\overline{m}B\text{Z}^{|V_n|}`$ and $`\overline{p}=\overline{q}+\overline{m}`$. Then, $`\overline{m}=\overline{p}+\overline{q}\text{Z}^{|V_n|,+}`$, which forces $`\overline{m}`$ to be zero because of the defining conditions of $`B\text{Z}^{|V_n|}`$. Thus, the natural order in $`\text{Z}^{|V_n|}`$ induces an order in the quotient group $`\text{Z}^{|V_n|}/B\text{Z}^{|V_n|}`$ making it an ordered group. Define $`(K_{}0(V,E,),K_{}0(V,E,)^+)`$ to be the inductive limit of the system of ordered groups
$$\frac{\text{Z}^{|V_0|}}{B\text{Z}^{|V_0|}}\stackrel{A_0}{}\frac{\text{Z}^{|V_1|}}{B\text{Z}^{|V_1|}}\stackrel{A_1}{}\frac{\text{Z}^{|V_2|}}{B\text{Z}^{|V_2|}}\stackrel{A_2}{}\frac{\text{Z}^{|V_3|}}{B\text{Z}^{|V_3|}}\stackrel{A_3}{}\mathrm{}$$
Observe that $`B\text{Z}^{|V_0|}=0`$. The image of $`1\text{Z}^{|V_0|}`$ in $`(K_{}0(V,E,),K_{}0(V,E,)^+)`$ is an order unit.
3.2. Theorem. For $`B=(V,E,)`$ let $`(X_B,T_B)`$ be defined as in (1.10). Write $`(X,T)=(X_B,T_B)`$. Define the tripling $`B^๐ฌ=(V^๐ฌ,E^๐ฌ,)`$ as in 2.2. Then $`K^0(X,T)`$ is naturally order isomorphic to $`K_{}0(V^๐ฌ,E^๐ฌ,),`$ preserving distinguished order units.
Proof. We recall the notation introduced in 2.1. Given $`fC(X,\text{Z})`$, choose $`n`$ sufficiently large such that $`f,_T(f)`$ are both constant on the sets of the partition $`๐ฌ_n`$. The vertices of the Bratteli diagram $`(V^๐ฌ,E^๐ฌ,)`$ correspond to towers $`๐ฎ_{(\varpi ,\varpi ^{},\varpi ^{\prime \prime })}`$ of a K-R partition which in turn are partitioned into floors $`\{(\varpi ,\varpi ^{},\varpi ^{\prime \prime };y)\}`$ as $`y`$ varies through the floors of $`\varpi ^{}`$. For $`f`$ as above, define $`\gamma _n(f)\text{Z}^{|V_n^๐ฌ|}`$ by $`\gamma _n(f)(\varpi ,\varpi ^{},\varpi ^{\prime \prime })=f(x)+f(Tx)+f(T^2x)+\mathrm{}+f(T^{h1}x)`$ where $`x`$ belongs to the lowest floor of $`๐ฎ_{(\varpi ,\varpi ^{},\varpi ^{\prime \prime })}`$ and $`h`$ is the height of the tower $`๐ฎ_{(\varpi ,\varpi ^{},\varpi ^{\prime \prime })}`$. Then $`A_n^๐ฌ(\gamma _n(f))=\gamma _{n+1}(f)`$ and $`\gamma _n(_T(f))B\text{Z}^{|V_n^๐ฌ|}`$. This gives rise to a map
$$\gamma :K^0(X_B,T_B)K_{}0(V^๐ฌ,E^๐ฌ,).$$
3.3. Lemma. Let $`fC(X,\text{Z})`$ and suppose that $`f`$ is constant on the sets of the partition $`๐ฌ_n`$. Suppose that $`\gamma _n(f)B\text{Z}^{|V_n^๐ฌ|}`$. Then, $`f=_T(g)`$, for some $`gC(X,\text{Z})`$.
3.4. Lemma. With the same assumptions as in 3.3 suppose that $`x,y(X)`$ both lie in the same floor of a $`๐ฌ_n`$-tower $`๐ฎ_{(\zeta ,\zeta ^{},\zeta ^{\prime \prime })}`$. Furthermore, suppose that for some positive integers $`k,\mathrm{}`$ both $`T^kx`$ and $`T^{\mathrm{}}y`$ lie in the same floor of a $`๐ฌ_n`$-tower $`๐ฎ_{(\vartheta ,\vartheta ^{},\vartheta ^{\prime \prime })}`$. Then,
$$f(x)+f(Tx)+\mathrm{}+f(T^kx)=f(y)+f(Ty)+\mathrm{}+f(T^{\mathrm{}}y).$$
Proof of 3.4. Let $`mn`$. Let $`U_y`$ be a neighborhood of $`y`$ such that $`zU_y`$ and for $`i[0,\mathrm{}],T^iy`$ and $`T^iz`$ belong to the same floor of the $`๐ฌ_n`$ partition. Since the orbit of $`T^kx`$ by iterations of $`T`$ is dense $`j`$ such that $`T^{j+k}xU_y`$. For sufficiently large $`mn,`$ a $`๐ฌ_m`$-tower $`๐ฎ`$ such that $`x,T^kx,T^{k+j}x,T^{k+j+\mathrm{}}x`$ belong to different floors of $`๐ฎ`$.
Let $`u`$ be the vertex of $`V_m^๐ฌ`$ represented by $`๐ฎ`$. The floors of the $`๐ฌ_m`$-tower $`๐ฎ`$ are linearly ordered reflecting the linear order in the set of paths in $`(V^๐ฌ,E^๐ฌ,)`$ from the top vertex to $`u`$. Similarly, the paths from $`V_n^๐ฌ`$ to $`uV_m^๐ฌ`$ are linearly ordered reflecting the order in which $`๐ฎ`$ traverses the level-$`n`$ towers of $`(V^๐ฌ,E^๐ฌ,)`$. Denote by $`\varpi _1,\varpi _2,\varpi _3,\mathrm{},\varpi _L`$ the paths from $`V_n^๐ฌ`$ to $`uV_m^๐ฌ`$ in their linear order. Write $`s(\varpi _1),s(\varpi _2),s(\varpi _3),\mathrm{},s(\varpi _L)V_n^๐ฌ`$ for their sources and $`๐ฎ_{s(\varpi _1)},๐ฎ_{s(\varpi _2)},๐ฎ_{s(\varpi _3)},\mathrm{},๐ฎ_{s(\varpi _L)}`$ for the $`๐ฌ_n`$-towers represented by these sources. Thus, $`๐ฎ`$ traverses $`๐ฌ_n`$-towers in the order $`๐ฎ_{s(\varpi _1)},๐ฎ_{s(\varpi _2)},๐ฎ_{s(\varpi _3)},\mathrm{},๐ฎ_{s(\varpi _L)}`$. Choose $`1a<b<c<dL`$ such that $`x`$,(resp.$`T^kx`$, resp.$`T^{k+j}x`$, resp.$`T^{k+j+\mathrm{}}x`$), is picked up by $`๐ฎ`$ at the $`a^{th}`$(resp.$`b^{th}`$, resp.$`c^{th}`$, resp.$`d^{th}`$) instance of $`๐ฎ`$ traversing through a $`๐ฌ_n`$-tower, namely, $`๐ฎ_{s(\varpi _a)}`$,(resp.$`๐ฎ_{s(\varpi _b)}`$, resp.$`๐ฎ_{s(\varpi _c)}`$, resp.$`๐ฎ_{s(\varpi _d)}`$). In particular, observe that $`s(\varpi _a)=s(\varpi _c)`$ and $`s(\varpi _b)=s(\varpi _d)`$.
Since $`\gamma _n(f)B\text{Z}^{|V_n^๐ฌ|}`$, we have
$$\gamma _n(f)(s(\varpi _a))+\gamma _n(f)(s(\varpi _{a+1}))+\mathrm{}+\gamma _n(f)(s(\varpi _{c1}))=0$$
and similarly,
$$\gamma _n(f)(s(\varpi _b))+\gamma _n(f)(s(\varpi _{b+1}))+\mathrm{}+\gamma _n(f)(s(\varpi _{d1}))=0.$$
These two equations imply that
$$f(x)+f(Tx)+\mathrm{}+f(T^{k+j1}x)=0$$
and
$$f(T^{k+1}x)+f(T^{k+2}x)+\mathrm{}+f(T^{k+j+\mathrm{}}x)=0.$$
Hence,
$$f(x)+f(Tx)+\mathrm{}+f(T^kx)=$$
$$f(x)+f(Tx)+\mathrm{}+f(T^kx)+f(T^{k+1}x)+\mathrm{}+f(T^{k+j}x)+\mathrm{}+f(T^{k+j+\mathrm{}}x)=$$
$$f(T^{k+j}x)+f(T^{k+j+1}x)+\mathrm{}+f(T^{k+j+\mathrm{}}x)=$$
$$f(y)+f(Ty)+\mathrm{}+f(T^{\mathrm{}}y).$$
This ends the proof of 3.4. $`\mathrm{}`$
Proof of 3.3. Choose $`x_0X`$. Now, for any $`zX`$ choose $`k\text{Z}^+`$ such that $`T^kx_0`$ and $`z`$ belong to the same $`๐ฌ_n`$-floor. Define $`gC(X,\text{Z})`$ by $`g(z)=f(x_0)+f(Tx_0)+\mathrm{}+f(T^kx_0)`$. Then, 3.4 implies that $`g`$ is well defined and $`_T(g)=fT`$. So, $`_T(gT^1)=f`$.
From 3.3 one can immediately deduce that the map $`\gamma :K^0(X_B,T_B)K_{}0(V^๐ฌ,E^๐ฌ,)`$ defined just before the statement of lemma 3.3 is an isomorphism.
This completes the proof of Theorem 3.2. $`\mathrm{}`$
3.5. A subgroup of $`B\text{Z}^{|V_n^๐ฌ|}`$. In practice it is quite tedious to determine whether a given element $`\overline{p}`$ of $`\text{Z}^{|V_n^๐ฌ|}`$ lies in $`B\text{Z}^{|V_n^๐ฌ|}`$. We now begin to describe a subgroup $`\mathrm{\Delta }\text{Z}^{|V_n^๐ฌ|}B\text{Z}^{|V_n^๐ฌ|}`$, which is more easily identifiable than $`B\text{Z}^{|V_n^๐ฌ|}`$. Eventhough, in general, this inclusion is proper we will later see that the distinction disappears when one takes inductive limits. As a consequence, we are able to obtain theorem 3.9, which yields a feasible method to compute $`K_0`$ effectively. Clearly, $`\text{Z}^{|V_n^๐ฌ|}`$ is the space of integral valued functions on the set $`V_n^๐ฌ`$. For a function $`\phi :V_n\times V_n\text{Z}`$ define $`\delta (\phi )\text{Z}^{|V_n^๐ฌ|}`$ by $`\delta (\phi )(a,b,c)=\phi (b,c)\phi (a,b)`$.
Lemma 3.6. $`\delta (\phi )B\text{Z}^{|V_n^๐ฌ|}`$.
Proof. Let $`\overline{p}\text{Z}^{|V_n^๐ฌ|}`$. Write $`\overline{p}=\{p_{(u,v,w)}\}_{(u,v,w)V_n^๐ฌ}`$. Take two paths from $`V_n^๐ฌ`$ to $`V_m^๐ฌ(m>n)`$ with the same source in $`V_n^๐ฌ`$ and same range in $`V_m^๐ฌ`$. The sequence of sources of paths lying between the above two paths is of the form $`\{(u_1,v_1,w_1),(u_2,v_2,w_2),\mathrm{},(u_j,v_j,w_j)\}`$ where
1. $`(u_1,v_1,w_1)=(u_j,v_j,w_j)`$,
2. $`u_{i+1}=v_i`$ and
3. $`v_{i+1}=w_i`$, for $`i=1,2,\mathrm{},j1`$.
If $`\overline{p}=\delta (\phi )`$, the sum $`p_{(u_1,v_1,w_1)}+p_{(u_2,v_2,w_2)}+\mathrm{}+p_{(u_{j1},v_{j1},w_{j1})}`$ equals
$`\{\phi (v_1,w_1)\phi (u_1,v_1)\}+\{\phi (v_2,w_2)\phi (u_2,v_2)\}+\mathrm{}+\{\phi (v_{j1},w_{j1})\phi (u_{j1},v_{j1})\}`$
$$\begin{array}{ccc}=& \phi (u_1,v_1)+\phi (v_{j1},w_{j1})& (\text{in view of (ii) and (iii) above})\\ =& \phi (u_1,v_1)+\phi (u_j,v_j)& (\text{in view of (ii) and (iii) above})\\ =& 0& (\text{in view of (i) above}).\end{array}$$
The condition for $`\overline{p}`$ to belong to $`B\text{Z}^{|V_n^๐ฌ|}`$ is precisely that the sums of the type $`p_{(u_1,v_1,w_1)}+p_{(u_2,v_2,w_2)}+\mathrm{}+p_{(u_{j1},v_{j1},w_{j1})}`$ as above should all vanish. As the foregoing calculation shows this holds whenever $`\overline{p}=\delta (\phi )`$ for some $`\phi :V_n\times V_n\text{Z}`$.
$`\mathrm{}`$
##
We define $`\mathrm{\Delta }\text{Z}^{|V_n^๐ฌ|}`$ to be the subgroup $`\delta (\text{Z}^{|V_n\times V_n|})`$ of $`\text{Z}^{|V_n^๐ฌ|}`$.
##
Recall the map $`A_n^๐ฌ:\text{Z}^{|V_n^๐ฌ|}\text{Z}^{|V_{n+1}^๐ฌ|}`$ given by matrix multiplication by the incidence matrix between the levels $`V_n^๐ฌ`$ and $`V_{n+1}^๐ฌ`$. For a function $`\phi :V_n\times V_n\text{Z}`$ define $`\phi ^{}:V_{n+1}\times V_{n+1}\text{Z}`$ by $`\phi ^{}(u^{},v^{})=\phi (u,v)`$ where $`u`$(resp.$`v`$) is the source of the last(resp.first) edge ranging at $`u^{}`$(resp.$`v^{}`$).
Lemma 3.7. With notation as above, $`A_n^๐ฌ(\delta (\phi ))=\delta (\phi ^{})`$. In particular the identity endomorphism of $`\text{Z}^{|V_n^๐ฌ|}`$ induces a map
$$\frac{\text{Z}^{|V_n^๐ฌ|}}{\mathrm{\Delta }\text{Z}^{|V_n^๐ฌ|}}\stackrel{\rho _n}{}\frac{\text{Z}^{|V_n^๐ฌ|}}{B\text{Z}^{|V_n^๐ฌ|}}$$
and
$$\begin{array}{ccc}\frac{\text{Z}^{|V_n^๐ฌ|}}{B\text{Z}^{|V_n^๐ฌ|}}& \stackrel{A_n^๐ฌ}{}& \frac{\text{Z}^{|V_{n+1}^๐ฌ|}}{B\text{Z}^{|V_{n+1}^๐ฌ|}}\\ & & \\ \rho _n& & \rho _{n+1}\\ & & \\ \frac{\text{Z}^{|V_n^๐ฌ|}}{\mathrm{\Delta }\text{Z}^{|V_n^๐ฌ|}}& \stackrel{A_n^๐ฌ}{}& \frac{\text{Z}^{|V_{n+1}^๐ฌ|}}{\mathrm{\Delta }\text{Z}^{|V_{n+1}^๐ฌ|}}\end{array}$$
is commutative.
Proof. The proof of the assertion $`A_n^๐ฌ(\delta (\phi ))=\delta (\phi ^{})`$ is a straightforward calculation using definitions and notation. The rest follows immediately.
$`\mathrm{}`$
We might wish to ask whether every element $`\overline{p}`$ of $`B\text{Z}^{|V_n^๐ฌ|}`$ is of the form $`\delta (\phi )`$ for some function $`\phi :V_n\times V_n\text{Z}`$. The proposition 3.8 below shows that after applying a finite iteration $`A_{n+i}^๐ฌ\mathrm{}A_{n+1}^๐ฌA_n^๐ฌ`$ to $`\overline{p}`$ it will indeed be so.
Let $`\overline{p}B\text{Z}^{|V_n^๐ฌ|}`$. Let $`gC(X,\text{Z})`$ be chosen as in Lemma 3.3 so that
1. $`_T(g)`$ is constant on the sets of the partition $`๐ฌ_n`$ and moreover, $`\overline{p}=\gamma _n(_T(g))`$.
2. $`g`$ itself is constant on the sets of the partition $`๐ฌ_{n+i}`$ for some positive integer $`i`$.
Choose a positive integer $`j`$ such that any $`๐ฌ_{n+i+j}`$-tower traverses through at least two $`๐ฌ_{n+i}`$-towers. Then, of course, any $`๐ซ_{n+i+j}`$-tower traverses through at least two $`๐ซ_{n+i}`$-towers.
For $`u^{},v^{}V_{n+i+j}`$ let $`u_a`$(resp.$`u_{a1}`$, resp.$`v_1,`$ resp.$`v_2`$) be the source of the last(resp.last but one, resp.first, resp.second) path from $`V_{n+i}`$ to $`V_{n+i+j}`$ ranging at $`u^{}`$(resp.$`u^{}`$,resp.$`v^{}`$,resp.$`v^{}`$). If there exists $`xX`$ such that
1. $`x`$ bottom floor of the $`Q_{n+i}`$-tower represented by $`(u_a,v_1,v_2)`$,
2. $`T^1x`$ top floor of the $`Q_{n+i}`$-tower represented by $`(u_{a1},u_a,v_1)`$
define $`\phi ^{}(u^{},v^{})=g(x)`$; then, $`\phi ^{}(u^{},v^{})`$ is independent of $`x`$. For given $`u^{},v^{}V_{n+i+j}`$ if no such $`x`$ exists, define $`\phi ^{}(u^{},v^{})`$ arbitrarily.
Proposition 3.8. With notation as above $`A_{n+i+j1}^๐ฌ\mathrm{}A_{n+1}^๐ฌA_n^๐ฌ(\overline{p})=\delta (\phi ^{})`$.
Proof. $`\delta (\varphi ^{})(u^{},v^{},w^{})=\varphi ^{}(v^{},w^{})\varphi ^{}(u^{},v^{})=g(T^hy)g(y)`$, if $`y`$ lies in the lowest floor of the $`๐ฌ_{n+i+j}`$-tower of height $`h`$ represented by $`(u^{},v^{},w^{})`$. Also, for $`(u,v,w)V_n^๐ฌ`$, $`\overline{p}(u,v,w)=\gamma _n(_T(g))(u,v,w)=`$ the sum $`(gTg)(z)+(gTg)(Tz)+(gTg)(T^2z)+\mathrm{}+(gTg)(T^{k1}z)`$, where $`k`$ is the height of the $`๐ฌ_n`$-tower represented by $`(u,v,w)`$ and $`z`$ lies in its lowest floor. Thus, $`A_{n+i+j1}^๐ฌ\mathrm{}A_{n+1}^๐ฌA_n^๐ฌ(\overline{p})(u^{},v^{},w^{})`$ is the sum of $`gTg`$ taken over all the floors of the $`๐ฌ_{n+i+j}`$-tower represented by $`(u^{},v^{},w^{})`$. This sum also equals $`gT^h(y)g(y)`$.
$`\mathrm{}`$
We can therefore give an alternate description of $`K^0(X_B,T_B)`$ which is more elegant than the description in Theorem 3.2. For the same reasons as in the case of $`{\displaystyle \frac{\text{Z}^{|V_n^๐ฌ|}}{B\text{Z}^{|V_n^๐ฌ|}}}`$, we see that the natural order in $`\text{Z}^{|V_n^๐ฌ|}`$ induces an order in $`{\displaystyle \frac{\text{Z}^{|V_n^๐ฌ|}}{\mathrm{\Delta }\text{Z}^{|V_n^๐ฌ|}}}`$.
As we already observed, $`A_n^๐ฌ[\delta (\text{Z}^{|V_n\times V_n|})]\delta (\text{Z}^{|V_{n+1}\times V_{n+1}|})`$.
3.9. Theorem. For $`B=(V,E,)`$ let $`(X_B,T_B)`$ be defined as in (1.10). Write $`(X,T)=(X_B,T_B)`$. Define the tripling $`B^๐ฌ=(V^๐ฌ,E^๐ฌ,)`$ as in 2.2. Then the map induced between the inductive limits of the two (horizontal) systems of ordered groups in the following diagram is an isomorphism.
$$\begin{array}{ccccccccc}\frac{\text{Z}^{|V_0^๐ฌ|}}{B\text{Z}^{|V_0^๐ฌ|}}& \stackrel{A_0^๐ฌ}{}& \frac{\text{Z}^{|V_1^๐ฌ|}}{B\text{Z}^{|V_1^๐ฌ|}}& \stackrel{A_1^๐ฌ}{}& \frac{\text{Z}^{|V_2^๐ฌ|}}{B\text{Z}^{|V_2^๐ฌ|}}& \stackrel{A_2^๐ฌ}{}& \frac{\text{Z}^{|V_3^๐ฌ|}}{B\text{Z}^{|V_3^๐ฌ|}}& \stackrel{A_3^๐ฌ}{}& \mathrm{}\\ & & & & & & & & \\ \rho _0& & \rho _1& & \rho _2& & \rho _3& & \\ & & & & & & & & \\ \frac{\text{Z}^{|V_0^๐ฌ|}}{\mathrm{\Delta }\text{Z}^{|V_0^๐ฌ|}}& \stackrel{A_0^๐ฌ}{}& \frac{\text{Z}^{|V_1^๐ฌ|}}{\mathrm{\Delta }\text{Z}^{|V_1^๐ฌ|}}& \stackrel{A_1^๐ฌ}{}& \frac{\text{Z}^{|V_2^๐ฌ|}}{\mathrm{\Delta }\text{Z}^{|V_2^๐ฌ|}}& \stackrel{A_2^๐ฌ}{}& \frac{\text{Z}^{|V_3^๐ฌ|}}{\mathrm{\Delta }\text{Z}^{|V_3^๐ฌ|}}& \stackrel{A_3^๐ฌ}{}& \mathrm{}\end{array}$$
Furthermore, the two inductive limits are both isomorphic to $`K^0(X,T)`$.
\[To avoid messy notation and display, we have hidden $`G^+`$ while referring to the ordered group $`(G,G^+`$).\]
Proof. By lemma 3.6, $`\mathrm{\Delta }\text{Z}^{|V_n^๐ฌ|}B\text{Z}^{|V_n^๐ฌ|}`$. By lemma 3.8, for sufficiently large $`K`$, $`A_{n+K1}^๐ฌ\mathrm{}A_{n+1}^๐ฌA_n^๐ฌ(B\text{Z}^{|V_n^๐ฌ|})\mathrm{\Delta }\text{Z}^{|V_{n+K}^๐ฌ|}`$. Hence, the induced map between the inductive limits is an isomorphism. That $`K^0(X,T)`$ is isomorphic to the inductive limit of the top horizontals was already proved in theorem 3.2. Observe that $`\text{Z}^{|V_0^๐ฌ|}=\text{Z}`$ and $`B\text{Z}^{|V_0^๐ฌ|}=\mathrm{\Delta }\text{Z}^{|V_0^๐ฌ|}=0`$. The image of $`1\text{Z}^{|V_0^๐ฌ|}`$ in the inductive limit maps to the order unit $`u`$ in $`K^0(X,T)`$ corresponding to the image of the constant function $`1C(X,\text{Z})`$.
$`\mathrm{}`$
Remark. Since it is known that $`K^0(X,T)`$ is isomorphic to the $`K_0`$-group of the associated $`C^{}`$-crossed product $`C(X)_T\text{Z}`$, we see that as a corollary to theorem 3.9 , we can effectively compute $`K_0(C(X_B)_{T_B}\text{Z})`$.
Finally, we should point out how these descriptions simplify further for properly ordered Bratteli diagrams and yield the isomorphism $`K^0(X,T)K_0(V,E)`$, (see 3.1), proved by Hermann, Putnam and Skau \[HPS, Theorem 5.4 and Corollary 6.3\].
3.10. Let $`(V,E,)`$ be a properly ordered Bratteli diagram. Telescoping if necessary, assume that every level $`n+1`$-tower traverses through at least two level $`n`$-towers. Telescoping further if necessary, (see \[HPS, Proposition 2.8\]), we can assume that any two maximal edges of $`E_n`$ have the same source. Similarly, we can assume that any two minimal edges of $`E_n`$ have the same source. For the rest of the paper we assume that these properties hold. Then for any $`๐ฌ_{n+2}`$-tower $`๐ฎ(u,v,w)`$ the first $`๐ฌ_n`$-tower traversed by $`๐ฎ(u,v,w)`$ is independent of $`(u,v,w)V_{n+2}^๐ฌ`$. Thus one sees from the definition of $`B\text{Z}^{|V_n^๐ฌ|}`$ that $`A_{n+1}^๐ฌA_n^๐ฌ(B\text{Z}^{|V_n^๐ฌ|})=0`$. As a consequence, the map induced between the inductive limits of the top two horizontals in the following diagram is an isomorphism.
$$\begin{array}{cccccccccc}\text{Z}& \stackrel{A_0^๐ฌ}{}& \frac{\text{Z}^{|V_1^๐ฌ|}}{B\text{Z}^{|V_1^๐ฌ|}}& \stackrel{A_1^๐ฌ}{}& \frac{\text{Z}^{|V_2^๐ฌ|}}{B\text{Z}^{|V_2^๐ฌ|}}& \stackrel{A_2^๐ฌ}{}& \mathrm{}& \frac{\text{Z}^{|V_n^๐ฌ|}}{B\text{Z}^{|V_n^๐ฌ|}}& & \mathrm{}\\ & & & & & & & & & \\ 1& & \tau & & \tau & & & \tau & & \\ & & & & & & & & & \\ \text{Z}& \stackrel{A_0^๐ฌ}{}& \text{Z}^{|V_1^๐ฌ|}& \stackrel{A_1^๐ฌ}{}& \text{Z}^{|V_2^๐ฌ|}& \stackrel{A_2^๐ฌ}{}& \mathrm{}& \text{Z}^{|V_n^๐ฌ|}& & \mathrm{}\\ & & & & & & & & & \\ 1& & \pi ^{}& & \pi ^{}& & & \pi ^{}& & \\ & & & & & & & & & \\ \text{Z}& \stackrel{A_0}{}& \text{Z}^{|V_1|}& \stackrel{A_1}{}& \text{Z}^{|V_2|}& \stackrel{A_2}{}& \mathrm{}& \text{Z}^{|V_n|}& & \mathrm{}\end{array}$$
The map $`\pi ^{}:\text{Z}^{|V_n|}\text{Z}^{|V_n^๐ฌ|}`$ is induced by the map $`\pi :V_n^๐ฌV_n`$ given by $`(u,v,w)v`$ and of course commutes with multiplication by the respective incidence matrices (i.e, $`A^n,A_n^๐ฌ`$). Recall that after doing necessary telescoping we have arranged so that the properly ordered Bratteli diagram $`(V,E,)`$ has the properties described in the beginning of 3.10. Now let $`e`$ be an edge in $`(V,E,)`$ with range $`vV_{n+1}`$. Let $`(u,v,w),(u^{},v,w^{})V_{n+1}^๐ฌ`$. Let $`\stackrel{~}{e},\stackrel{~}{e}^{}`$ be the (unique) lifts of $`e`$ to $`(V^๐ฌ,E^๐ฌ,)`$ with ranges $`(u,v,w),(u^{},v,w^{})`$ respectively. Then, from the description in 2.2, $`\stackrel{~}{e},\stackrel{~}{e}^{}`$ have the same sources in $`V_n^๐ฌ`$. From this it follows that for $`\overline{p}\text{Z}^{|V_n^๐ฌ|},A_n^๐ฌ(\overline{p})(u,v,w)=A_n^๐ฌ(\overline{p})(u^{},v,w^{})`$; in other words, $`A_n^๐ฌ(\overline{p})\pi ^{}(\text{Z}^{|V_{n+1}|})`$. Thus the map induced between the inductive limits of the two bottom horizontals in the above diagram is also an isomorphism.
3.11. Specialization of Theorem 3.9 to substitutional systems. We recall the notation from 2.5. Let $`(A,\sigma )`$ be a primitive aperiodic non-proper substitutional system. Let $`B=(V,E,)`$ be the stationary ordered Bratteli diagram associated to $`(A,\sigma )`$. Define $`A^๐ฌ=\{(a,b,c)A\times A\times Aabc\text{ occurs as a subword of }\sigma ^n(d)\text{ for some }dA\text{ and some }n\}`$. Define
$$\sigma ^๐ฌ:A^๐ฌ(A^๐ฌ)^+$$
by $`\sigma ^๐ฌ[(a,b,c)]=(a_m,b_1,b_2)(b_1,b_2,b_3)\mathrm{}(b_{n2},b_{n1},b_n)(b_{n1},b_n,c_1)`$, where $`\sigma (b)=b_1b_2\mathrm{}b_n`$, and $`a_m`$ is the last alphabet in $`\sigma (a)`$, while $`c_1`$ is the first alphabet in $`\sigma (c)`$. For a function $`\phi :A\times A\text{Z}`$ define $`\delta (\phi )\text{Z}^{|A^๐ฌ|}`$ by $`\delta (\phi )(a,b,c)=\phi (b,c)\phi (a,b)`$. Let $`\mathrm{\Delta }\text{Z}^{|A^๐ฌ|}`$ be the subgroup $`\delta (\text{Z}^{|A\times A|})`$ of $`\text{Z}^{|A^๐ฌ|}`$. Suppose $`\overline{p},\overline{q}\text{Z}^{|A^๐ฌ|}`$ and take values in $`\text{Z}^+`$ and further that $`\overline{p}=\overline{q}`$mod $`\mathrm{\Delta }\text{Z}^{|A^๐ฌ|}`$. Then $`\overline{p}=\overline{q}=0`$. The natural order in $`\text{Z}^{|A^๐ฌ|}`$ induces an order in $`\text{Z}^{|A^๐ฌ|}/\mathrm{\Delta }\text{Z}^{|A^๐ฌ|}`$ making it an ordered group.
##
Let $`\beta ^๐ฌ:\text{Z}^{|A^๐ฌ|}\text{Z}^{|A^๐ฌ|}`$ be given by matrix multiplication by the incidence matrix of the substitution $`\sigma ^๐ฌ`$. For a function $`\phi :A\times A\text{Z}`$ define $`\phi ^{}:A\times A\text{Z}`$ by $`\phi ^{}(u^{},v^{})=\phi (u,v)`$ where $`u`$(resp.$`v`$) is the last(resp.first) alphabet in the substitution $`\sigma (u^{})`$(resp.$`\sigma (v^{})`$). Then $`\beta ^๐ฌ(\delta (\phi ))=\delta (\phi ^{})`$; thus, $`\beta ^๐ฌ[\delta (\text{Z}^{|A\times A|})]\delta (\text{Z}^{|A\times A|})`$. Hence, $`\beta ^๐ฌ`$ induces a homomorphism of ordered groups
$$\frac{\text{Z}^{|A^๐ฌ|}}{\mathrm{\Delta }\text{Z}^{|A^๐ฌ|}}\frac{\text{Z}^{|A^๐ฌ|}}{\mathrm{\Delta }\text{Z}^{|A^๐ฌ|}}$$
still denoted by $`\beta ^๐ฌ`$.
From theorem 3.9 one deduces immediately
Theorem 3.12. The inductive limit of the system of ordered groups
$$\begin{array}{ccccccccc}\frac{\text{Z}^{|A^๐ฌ|}}{\mathrm{\Delta }\text{Z}^{|A^๐ฌ|}}& \stackrel{\beta ^๐ฌ}{}& \frac{\text{Z}^{|A^๐ฌ|}}{\mathrm{\Delta }\text{Z}^{|A^๐ฌ|}}& \stackrel{\beta ^๐ฌ}{}& \frac{\text{Z}^{|A^๐ฌ|}}{\mathrm{\Delta }\text{Z}^{|A^๐ฌ|}}& \stackrel{\beta ^๐ฌ}{}& \frac{\text{Z}^{|A^๐ฌ|}}{\mathrm{\Delta }\text{Z}^{|A^๐ฌ|}}& \stackrel{\beta ^๐ฌ}{}& \mathrm{}\end{array}$$
is isomorphic to the dimension group $`K^0(X_\sigma ,T_\sigma )`$ of the substitution system associated to $`(A,\sigma )`$.
References
1. \] F. Durand, B. Host and C. Skau. Substitutional dynamical systems, Bratteli diagrams and dimension groups, Ergodic Th. and Dynam. Sys. 19, (1999) 953-993.
2. \] A. El Kacimi, R. Parthasarathy. Skew-product for group-valued edge labellings of Bratteli diagrams, arXiv.math.DS/0506304.
3. \] R.H. Herman, I.F. Putnam and C.F. Skau. Ordered Bratteli diagrams, dimension groups and topological dynamics, Internat. J. Math. 3, (1992) 827-864.
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# The Extended Chandra Deep Field-South Survey. Chandra Point-Source Catalogs
## 1 Introduction
Deep and wide X-ray surveys indicate that the cosmic X-ray background is largely due to accretion onto supermassive black holes (SMBHs) integrated over cosmic time (e.g., see Brandt & Hasinger 2005 for a review). Follow-up studies of deep-survey sources with 8โ10 m optical telescopes as well as multiwavelength correlative studies have shown that most of the X-ray sources are active galactic nuclei (AGNs), many of which are obscured (e.g., Bauer et al. 2004; Szokoly et al. 2004; Barger et al. 2005). X-ray surveys have found the highest density of AGNs on the sky (up to $``$7200 deg<sup>-2</sup>). In addition to AGNs, the deepest X-ray surveys have also detected respectable numbers of starburst and normal galaxies out to cosmologically interesting distances ($`z1`$; e.g., Hornschemeier et al. 2003; Bauer et al. 2004; Norman et al. 2004).
Presently, the two deepest X-ray surveys are the $`2`$ Ms Chandra Deep Field-North (CDF-N; Brandt et al. 2001, hereafter B01; Alexander et al. 2003, hereafter A03) and the $`1`$ Ms Chandra Deep Field-South (CDF-S; Giacconi et al. 2002, hereafter G02). These $``$400 arcmin<sup>2</sup> surveys have been performed in regions of sky with extensive multiwavelength coverage. They have provided 50โ250 times the sensitivity of surveys by previous X-ray missions, detecting large numbers of point sources (584 for the CDF-N and 346 for the CDF-S; G02; A03) and about a dozen extended groups and poor clusters (Bauer et al. 2002; G02).
The X-ray surveys performed to date have explored an impressive amount of the sensitivity versus solid angle โdiscovery spaceโ (see Figure 1 and Brandt & Hasinger 2005). However, one limitation of the present surveys is that there is only a relatively small amount of sky probed to 0.5โ2 keV flux levels of (2โ50)$`\times 10^{17}`$ erg cm<sup>-2</sup> s<sup>-1</sup>, a flux regime where many obscured AGNs are observed (e.g., Bauer et al. 2004). As a result, our understanding of the X-ray universe at these faint fluxes suffers from limited source statistics and field-to-field variance. To mitigate this limitation, the Extended Chandra Deep Field-South (E-CDF-S) survey was undertaken as part of the Chandra Cycle 5 guest observer program. The E-CDF-S is composed of four 250 ks Chandra ACIS-I pointings flanking the original CDF-S; these are arranged in a contiguous two-by-two pattern and cover a total solid angle of $``$1100 arcmin<sup>2</sup>.<sup>1</sup><sup>1</sup>1The $``$1 Ms CDF-S data cover $`35`$% of the E-CDF-S; much of this coverage, however, has limited sensitivity due to point spread function (PSF) broadening and vignetting at large off-axis angles (see ยง 3 for details). The same effects limit the sensitivity and positions derived from the XMM-Newton data (Streblyanska et al. 2004) extending outside the region with Chandra coverage. The pointings have sufficient sensitivity to detect the X-ray emission from moderate-luminosity AGNs ($`L_\mathrm{X}=10^{43}`$$`10^{44}`$ erg s<sup>-1</sup>) to $`z`$ 3โ6 as well as X-ray luminous starburst galaxies to $`z1`$. The E-CDF-S therefore can significantly improve understanding of SMBH accretion at high redshift where the source statistics are still limited. The contiguous nature of the E-CDF-S will allow wider field studies of the remarkable AGN clustering already found in the CDF-S (e.g., Gilli et al. 2003, 2005), and comparisons with other surveys of comparable depth (e.g., Stern et al. 2002; Harrison et al. 2003; Wang et al. 2004a,b; Nandra et al. 2005) will allow further assessment of the field-to-field variance of X-ray source populations.
The E-CDF-S field was selected for this program primarily due to its superb and growing multiwavelength coverage over a $``$900 arcmin<sup>2</sup> area, which ensures that it will remain a prime survey field in coming decades (see Figure 2). For example, the E-CDF-S has been imaged intensively with the HST Advanced Camera for Surveys (ACS) via the Galaxy Evolution from Morphology and Spectral Energy Distributions (GEMS; Rix et al. 2004; 117 HST orbits) and Great Observatories Origins Deep Survey (GOODS; Giavalisco et al. 2004; 199 HST orbits) projects. Excellent ground-based imaging is also available (e.g., Arnouts et al. 2001; Renzini et al. 2003; Giavalisco et al. 2004; Wolf et al. 2004; Gawiser et al. 2005), and several spectroscopic campaigns are underway to identify sources in the E-CDF-S, most notably with the Very Large Telescope (VLT; e.g., Le Fevre et al. 2004; Szokoly et al. 2004; Vanzella et al. 2005). The E-CDF-S has been targeted by Spitzer via the GOODS (M. Dickinson et al., in preparation), the Spitzer Wide-Area Infrared Extragalactic Survey (SWIRE; Lonsdale et al. 2003), guaranteed time (e.g., Papovich et al. 2004), and guest observer (PI: P. van Dokkum) programs. Radio observations of the E-CDF-S have been made with the Australia Telescope Compact Array (ATCA; J. Afonso et al., in preparation) and the Very Large Array.
In this paper, we present Chandra point-source catalogs and data products derived from the E-CDF-S data set along with details of the observations, data reduction, and technical analysis. The observational procedures and data processing were similar in nature to those presented in B01 and A03. Detailed follow-up investigations and scientific interpretation of the E-CDF-S sources will be presented in subsequent papers.
The Galactic column density along the line of sight to the E-CDF-S is remarkably low: $`N_\mathrm{H}=8.8\times 10^{19}`$ cm<sup>-2</sup> (e.g., Stark et al. 1992). The coordinates throughout this paper are J2000. Cosmological parameters of $`H_0=70`$ km s<sup>-1</sup> Mpc<sup>-1</sup>, $`\mathrm{\Omega }_\mathrm{M}=0.3`$, and $`\mathrm{\Omega }_\mathrm{\Lambda }=0.7`$ are adopted.
## 2 Observations and Data Reduction
### 2.1 Instrumentation and Observations
The Advanced CCD Imaging Spectrometer (ACIS; Garmire et al. 2003) was used for all of the Chandra observations.<sup>2</sup><sup>2</sup>2For additional information on ACIS and Chandra see the Chandra Proposersโ Observatory Guide at http://cxc.harvard.edu/proposer/CfP/. ACIS is composed of ten CCDs (each $`1024\times 1024`$ pixels) designed for efficient X-ray source detection and spectroscopy. ACIS-I consists of four CCDs (CCDs I0โI3) arranged in a $`2\times 2`$ array with each CCD tipped slightly to approximate the curved focal surface of the Chandra High Resolution Mirror Assembly (HRMA). The aim-point of ACIS-I lies on CCD I3. The remaining six CCDs (ACIS-S; CCDs S0โS5) reside in a linear array and are tipped to approximate the Rowland circle of the objective gratings that can be inserted behind the HRMA.
The ACIS-I full field of view is $`16\stackrel{}{\mathrm{.}}9\times 16\stackrel{}{\mathrm{.}}9`$ ($``$285 arcmin<sup>2</sup>), and the sky-projected ACIS pixel size is $``$$`0\stackrel{}{\mathrm{.}}492`$. The PSF is smallest at the lowest photon energies and for sources at small off-axis angles. For example, the 95% encircled-energy radius at 1.5 keV for off-axis angles of $`0^{}`$$`8^{}`$ is $`1\stackrel{}{\mathrm{.}}8`$$`7\stackrel{}{\mathrm{.}}5`$ (Feigelson, Broos, & Gaffney 2000; Jerius et al. 2000).<sup>3</sup><sup>3</sup>3Feigelson et al. (2000) is available at http://www.astro.psu.edu/xray/acis/memos/memoindex.html. The PSF is approximately circular at small off-axis angles, broadens and elongates at intermediate off-axis angles, and becomes complex at large off-axis angles.
The entire Chandra observation program consisted of nine separate Chandra observations taken between 2004 February 29 and 2004 November 20 and is described in Table 1. The four ACIS-I CCDs were operated in all of the observations; the ACIS-S CCD S2 was in operation for observations 5019โ5022 and 6164. Due to the large off-axis angle of ACIS-S, and consequently its low sensitivity, these data were not used in this analysis. All observations were taken in Very Faint mode to improve the screening of background events and thus increase the sensitivity of ACIS in detecting faint X-ray sources.<sup>4</sup><sup>4</sup>4For more information on the Very Faint mode see http://cxc.harvard.edu/cal/Acis/Cal\_prods/vfbkgrnd/ and Vikhlinin (2001). The observations were made in four distinct observational fields (hereafter, fields 1, 2, 3, and 4; see Table 1 for more observational details) and cover a total solid angle of 1128.4 arcmin<sup>2</sup>. The focal-plane temperature was kept at $`120^{}`$ C for all of the nine observations.
Background light curves for all nine observations were inspected using event browser in the Tools for ACIS Real-time Analysis (tara; Broos et al. 2000) software package.<sup>5</sup><sup>5</sup>5tara is available at http://www.astro.psu.edu/xray/docs. All but two are free from significant flaring and are stable to within $``$20%. The two observations with significant flaring are 5015 and 5017. The background was $`\stackrel{>}{}`$1.5 times higher than nominal for two $``$1 ks intervals of observation 5015, and during observation 5017 the background rose to $`\stackrel{>}{}`$1.5 times the nominal rate and remained above this level for $``$10 ks near the end of the observation. Intervals with flaring were retained because the flaring strengths were not strong enough to have significant negative effects on our analyses.
### 2.2 Data Reduction
Chandra X-ray Center (hereafter CXC) pipeline software was used for basic data processing, and the pipeline versions are listed in Table 1. The reduction and analysis of the data used Chandra Interactive Analysis of Observations (ciao) Version 3.2 tools whenever possible;<sup>6</sup><sup>6</sup>6See http://cxc.harvard.edu/ciao/ for details on ciao. however, custom software, including the tara package, was also used extensively.
All data were corrected for the radiation damage sustained by the CCDs during the first few months of Chandra operations using the Charge Transfer Inefficiency (CTI) correction procedure of Townsley et al. (2000, 2002).<sup>7</sup><sup>7</sup>7The software associated with the correction method of Townsley et al. (2000, 2002) is available at http://www.astro.psu.edu/users/townsley/cti/. In addition to correcting partially for the positionally dependent grade distribution due to CTI effects, this procedure also partially corrects for quantum efficiency losses (see Townsley et al. 2000, 2002 for further details).
All bad columns, bad pixels, and cosmic ray afterglows were removed using the โstatusโ information in the event files, and we only used data taken during times within the CXC-generated good-time intervals. The ciao tool acis\_process\_events was used to remove the standard pixel randomization.
## 3 Production of the Point-Source Catalogs
By design, the four 250 ks E-CDF-S observations have their regions of highest sensitivity located where the sensitivity of the original $``$1 Ms CDF-S observation is poorest (see Table 1 and Figures 2 and 17). The loss of sensitivity is due to the combination of substantial degradation of the Chandra PSF at large off-axis angles and vignetting. In fact, most of the area where the E-CDF-S observations have their highest sensitivity lack any Chandra coverage in the $``$1 Ms CDF-S. We experimented with source searching utilizing the addition of the $``$1 Ms CDF-S and the 250 ks E-CDF-S images; such searching was done with wavdetect (Freeman et al. 2002) runs that did not utilize detector-specific PSF information (i.e., the โDETNAMโ keyword in the image files was deleted). However, such searching did not find a substantial number of new sources compared to those presented below combined with those from the original $``$1 Ms CDF-S (G02; A03); these results were verified via inspection of adaptively smoothed images. Therefore, our basic approach here is to present just the sources detected in the new 250 ks E-CDF-S observations. The X-ray sources in the $``$1 Ms CDF-S catalog of A03 were processed using the same techniques presented here with two main differences:
1. Our main Chandra catalog includes sources detected by running wavdetect at a false-positive probability threshold of 10<sup>-6</sup>, somewhat less conservative than the 10<sup>-7</sup> value adopted by A03; see $`\mathrm{\S }`$ 3.2 for details.
2. The E-CDF-S consists of four ACIS-I observational fields and subtends a larger solid angle than the fields presented in A03. Therefore, our main Chandra catalog of the entire E-CDF-S exposure was generated by merging sub-catalogs created in each of the four observational fields; see $`\mathrm{\S }`$ 3.2, Table 1, and Figure 2 for details.
### 3.1 Image and Exposure-Map Creation
We constructed images of each of the four E-CDF-S fields using the standard ASCA grade set (ASCA grades 0, 2, 3, 4, 6) for three standard bands (i.e., 12 images in total): 0.5โ8.0 keV (full band; FB), 0.5โ2.0 keV (soft band; SB), and 2โ8 keV (hard band; HB). These images have 0$`\stackrel{}{\mathrm{.}}`$492 per pixel. For each of the standard bands, the images from all four observational fields were merged into a single image using the ciao script merge\_all.<sup>8</sup><sup>8</sup>8See http://cxc.harvard.edu/ciao/threads/merge\_all/ In Figures 3 and 4 we display the full-band raw and exposure-corrected adaptively smoothed images (see discussion below), respectively.<sup>9</sup><sup>9</sup>9Raw and adaptively smoothed images for all three standard bands are available at the E-CDF-S website (see http://www.astro.psu.edu/users/niel/ecdfs/ecdfs-chandra.html). Furthermore, equivalent images obtained by merging the E-CDF-S and CDF-S are also available at the E-CDF-S website. Our point-source detection analyses have been restricted to the raw images constructed for each of the four observational fields so that the Chandra PSF is accounted for correctly (see $`\mathrm{\S }`$ 3.2).
We constructed exposure maps for each of the four observational fields in the three standard bands. These were created following the basic procedure outlined in $`\mathrm{\S }`$ 3.2 of Hornschemeier et al. (2001) and are normalized to the effective exposures of sources located at the aim points. Briefly, this procedure takes into account the effects of vignetting, gaps between the CCDs, bad column filtering, and bad pixel filtering. Also, with the release of ciao version 3.2, the spatially dependent degradation in quantum efficiency due to contamination on the ACIS optical blocking filters is now incorporated into the generation of exposure maps.<sup>10</sup><sup>10</sup>10See http://cxc.harvard.edu/ciao/why/acisqedeg.html A photon index of $`\mathrm{\Gamma }=1.4`$, the slope of the X-ray background in the 0.5โ8.0 keV band (e.g., Marshall et al. 1980; Gendreau et al. 1995), was assumed in creating the exposure maps. For each standard band, a total exposure map, covering the entire E-CDF-S, was constructed by merging the exposure maps of the four observational fields using the ciao script dmregrid. The resulting full-band exposure map is shown in Figure 5. Figure 6 displays the survey solid angle as a function of full-band effective exposure for both the total E-CDF-S exposure (Figure 6a) and the four individual observational fields (Figure 6b). Each observational field has comparable coverage with the majority of the solid angle coverage ($``$900 arcmin<sup>2</sup>) having at least 200 ks of effective exposure.
Using the exposure maps and adaptively smoothed images discussed above, we produced exposure-corrected images following the prescription outlined in $`\mathrm{\S }`$ 3.3 of Baganoff et al. (2003). Figure 7 shows a โfalse-colorโ composite image made using exposure-corrected adaptively smoothed 0.5โ2.0 keV (red), 2โ4 keV (green), and 4โ8 keV (blue) images.
### 3.2 Point-Source Detection
Point-source detection was performed in each band with wavdetect using a โ$`\sqrt{2}`$ sequenceโ of wavelet scales (i.e., 1, $`\sqrt{2}`$, 2, $`2\sqrt{2}`$, 4, $`4\sqrt{2}`$, and 8 pixels). Our key criterion for source detection, and inclusion in the main Chandra catalog, is that a source must be found with a given false-positive probability threshold in at least one of the three standard bands. The false-positive probability threshold in each band was set to $`1\times 10^6`$; a total of 762 distinct sources met this criterion. We also ran wavdetect using false-positive probability thresholds of $`1\times 10^7`$ and $`1\times 10^8`$ to evaluate the significance of each detected source.
If we conservatively treat the 12 images (i.e., the three standard bands over the four observational fields) as being independent, it appears that $``$50 (i.e., $``$6%) false sources are expected in our total Chandra source catalog for the case of a uniform background over $``$5.0 $`\times `$ 10<sup>7</sup> pixels. However, since wavdetect suppresses fluctuations on scales smaller than the PSF, a single pixel usually should not be considered a source detection cell, particularly at large off-axis angles. Hence, our false-source estimates are conservative. As quantified in $`\mathrm{\S }`$ 3.4.1 of A03 and by new source-detection simulations (P. E. Freeman 2005, private communication), the number of false-sources is likely $``$2โ3 times less than our conservative estimate, leaving only $``$15โ25 (i.e., $`\stackrel{<}{}`$3%) false sources. In $`\mathrm{\S }`$ 3.3.1 below we provide additional source-significance information that a user can utilize to perform more conservative source screening if desired.
### 3.3 Point-Source Catalogs
#### 3.3.1 Main Chandra Source Catalog
We ran wavdetect with a false-positive probability threshold of $`1\times 10^6`$ on all of the 12 images. The resulting source lists were then merged to create the point-source catalog given in Table 2. For cross-band matching, a matching radius of $`2\stackrel{}{\mathrm{.}}5`$ was used for sources within $`6^{}`$ of the average aim point. For larger off-axis angles, a matching radius of $`4\stackrel{}{\mathrm{.}}0`$ was used. These matching radii were chosen based on inspection of histograms showing the number of matches obtained as a function of angular separation (e.g., see ยง2 of Boller et al. 1998); with these radii the mismatch probability is $`\stackrel{<}{}`$1% over the entire field.
We improved the wavdetect source positions using a matched-filter technique (A03). This technique convolves the full-band image in the vicinity of each source with a combined PSF. The combined PSF is automatically generated as part of the acis\_extract procedure (Broos et al. 2002) within tara (see Footnote 5) and is produced by combining the โlibraryโ PSF of a source for each observation, weighted by the number of detected counts.<sup>11</sup><sup>11</sup>11acis\_extract can be accessed from http://www.astro.psu.edu/xray/docs/TARA/ae\_users\_guide.html. The PSFs are taken from the CXC PSF library; see http://cxc.harvard.edu/ciao/dictionary/psflib.html. This technique takes into account the fact that, due to the complex PSF at large off-axis angles, the X-ray source position is not always located at the peak of the X-ray emission. The matched-filter technique provides a small improvement ($``$0$`\stackrel{}{\mathrm{.}}`$1 on average) in the positional accuracy for sources further than 6 from the average aim-point. For sources with off-axis angles ($`\theta `$$`<`$ 6, we found that the off-axis-angle-weighted combination of centroid and matched-filter positions returned the most significant improvement to source positions. Algebraically, this can be written as:
$$(6{}_{}{}^{}\theta )/6{}_{}{}^{}\times \text{ centroid position }+\theta /6{}_{}{}^{}\times \text{ matched-filter position }$$
(1)
This method is similar to that employed by A03.
Manual correction of the source properties was required in some special cases: (1) There were 11 close doubles (i.e., sources with overlapping PSFs) and one close triple. These sources incur large photometric errors due to the difficulty of the separation process. (2) A total of eight sources were located close to bright sources, in regions of high background, in regions with strong gradients in exposure time, or partially outside of an observational field. The properties of these sources have been adjusted manually and are flagged in column 39 of Table 2 (see below).
For each observational field, we refined the absolute X-ray source positions by matching X-ray sources from the main point-source catalog to $`R`$-band optical source positions from deep observations ($`R_{\mathrm{lim},6\sigma }`$ $``$ 27 \[AB\] over the entire E-CDF-S) obtained with the Wide Field Imager (WFI) of the MPG/ESO telescope at LaSilla (see $`\mathrm{\S }`$ 2 of Giavalisco et al. 2004). X-ray sources from each of the four observational fields were matched to optical sources using a 2$`\stackrel{}{\mathrm{.}}`$5 matching radius. Using this matching radius, a small number of sources were observed to have more than one optical match; the brightest of these sources was selected as the most probable counterpart. Under these criteria, 640 ($``$84%) X-ray sources have optical counterparts. We also note that in a small number of cases the X-ray source may be offset from the center of the optical source even though both are associated with the same galaxy (e.g., a galaxy with bright optical emission from starlight that also has an off-nuclear ultraluminous X-ray binary with $`L_\text{X}`$ $``$ 10<sup>38-40</sup> erg s<sup>-1</sup>; see e.g., Hornschemeier et al. 2004). The accuracy of the X-ray source positions was improved by centering the distribution of offsets in right ascension and declination between the optical and X-ray source positions; this resulted in small ($`<`$ 1$`\stackrel{}{\mathrm{.}}`$0) field-dependent astrometric shifts for all sources in each field. We also checked for systematic offsets as a function of right ascension and declination that may arise from differing โplate scalesโ and rotations between the X-ray and optical images. These investigations were performed by plotting the right ascension and declination offsets (between optical and X-ray sources) as functions of right ascension and declination; no obvious systematic offsets were found.
Figure 8 shows the positional offset between the X-ray and optical sources versus off-axis angle after applying the positional corrections discussed above. Here, the off-axis angles are computed for each observational field appropriately; this allows for the consistent analysis of Chandra positional uncertainties as a function of off-axis angle. The median offset is $``$0$`\stackrel{}{\mathrm{.}}`$35; however, there are clear off-axis angle and source-count dependencies. The off-axis angle dependence is due to the HRMA PSF becoming broad at large off-axis angles, while the count dependency is due to the difficulty of centroiding a faint X-ray source. The median offset of the bright X-ray sources ($``$ 50 full-band counts) is only $``$0$`\stackrel{}{\mathrm{.}}`$25, while the median offset of the faint X-ray sources ($`<`$ 50 full-band counts) is $``$0$`\stackrel{}{\mathrm{.}}`$47. The positional uncertainty of each source is estimated following equations 2 and 3.
The main Chandra source catalog is presented in Table 2, and the details of the columns are given below.
* Column 1 gives the source number. Sources are listed in order of increasing RA.
* Columns 2 and 3 give the RA and Dec of the X-ray source, respectively. Note that more accurate positions are available for sources detected near the aim-point of the $``$1 Ms CDF-S through the catalogs presented in A03; see columns 19โ21. To avoid truncation error, we quote the positions to higher precision than in the International Astronomical Union (IAU) registered names beginning with the acronym โCXO ECDFSโ for โChandra X-ray Observatory Extended Chandra Deep Field-South.โ The IAU names should be truncated after the tenths of seconds in RA and after the arcseconds in Dec.
* Column 4 gives the positional uncertainty. As shown above, the positional uncertainty is dependent on off-axis angle and the number of detected counts. For the brighter X-ray sources ($``$ 50 full-band counts) the positional uncertainties are given by the empirically determined equation:
$$\mathrm{\Delta }=\{\begin{array}{cc}0.6\hfill & \theta <5^{}\hfill \\ & \\ 0.6+\left(\frac{\theta 5^{}}{20^{}}\right)\hfill & \theta 5^{}\hfill \end{array}$$
(2)
where $`\mathrm{\Delta }`$ is the positional uncertainty in arcseconds and $`\theta `$ is the off-axis angle in arcminutes (compare with Figure 8).
For the fainter X-ray sources ($`<`$ 50 full-band counts) the positional uncertainties are given by the empirically determined equation:
$$\mathrm{\Delta }=\{\begin{array}{cc}0.85\hfill & \theta <5^{}\hfill \\ & \\ 0.85+\left(\frac{\theta 5^{}}{4^{}}\right)\hfill & \theta 5^{}\hfill \end{array}$$
(3)
The stated positional uncertainties are somewhat conservative, corresponding to the $``$80โ90% confidence level.
* Column 5 gives the off-axis angle for each source in arcminutes. This is calculated using the source position given in columns 2 and 3 and the aim point (see Table 1) for the corresponding field in which it was detected (column 37).
* Columns 6โ14 give the source counts and the corresponding $`1\sigma `$ upper and lower statistical errors (from Gehrels 1986), respectively, for the three standard bands. All values are for the standard ASCA grade set, and they have not been corrected for vignetting. Source counts and statistical errors have been calculated using circular aperture photometry; extensive testing has shown that this method is more reliable than the wavdetect photometry. The circular aperture was centered at the position given in columns 2 and 3 for all bands.
The local background is determined in an annulus outside of the source-extraction region. The mean number of background counts per pixel is calculated from a Poisson model using $`\frac{n_1}{n_0}`$, where $`n_0`$ is the number of pixels with 0 counts and $`n_1`$ is the number of pixels with 1 count. Although only the numbers of pixels with 0 and 1 counts are measured, this technique directly provides the mean background even when $`n_1n_0`$. Furthermore, by ignoring all pixels with more than 1 count, this technique guards against background contamination from sources. We note that relatively bright nearby sources may contribute counts to nearby pixels where the background is estimated. Since the number density of relatively bright sources in the E-CDF-S is low, we estimate that only $``$10โ20 of these sources are thereby contaminated; the majority of these sources have been corrected via manual photometry of close doubles (see above). The principal requirement for using this Poisson-model technique is that the background is low and follows a Poisson distribution; in ยง4.2 of A03 it has been shown that the ACIS-I background matches this criterion for exposures as long as $``$2 Ms. The total background for each source is calculated and subtracted to give the net number of source counts.
For sources with fewer than 1000 full-band counts, we have chosen the aperture radii based on the encircled-energy function of the Chandra PSF as determined using the CXCโs mkpsf software (Feigelson et al. 2000; Jerius et al. 2000). In the soft band, where the background is lowest, the aperture radius was set to the 95% encircled-energy radius of the PSF. In the other bands, the 90% encircled-energy radius of the PSF was used. Appropriate aperture corrections were applied to the source counts by dividing the extracted source counts by the encircled-energy fraction for which the counts were extracted.
For sources with more than 1000 full-band counts, systematic errors in the aperture corrections often exceed the expected errors from photon statistics when the apertures described in the previous paragraph are used. Therefore, for such sources we used larger apertures to minimize the importance of the aperture corrections; this is appropriate since these bright sources dominate over the background. We set the aperture radii to be twice those used in the previous paragraph and inspected these sources to verify that the measurements were not contaminated by neighboring objects.
We have performed several consistency tests to verify the quality of the photometry. For example, we have checked that the sum of the counts measured in the soft and hard bands does not differ from the counts measured in the full band by an amount larger than that expected from measurement error. Systematic errors that arise from differing full-band counts and soft-band plus hard-band counts are estimated to be $`\stackrel{<}{}`$4%.
When a source is not detected in a given band, an upper limit is calculated; upper limits are indicated as a โ$``$1โ in the error columns. All upper limits are determined using the circular apertures described above. When the number of counts in the aperture is $`10`$, the upper limit is calculated using the Bayesian method of Kraft, Burrows, & Nousek (1991) for 99% confidence. The uniform prior used by these authors results in fairly conservative upper limits (see Bickel 1992), and other reasonable choices of priors do not materially change our scientific results. For larger numbers of counts in the aperture, upper limits are calculated at the $`3\sigma `$ level for Gaussian statistics.
* Columns 15 and 16 give the RA and Dec of the optical source centroid, which was obtained by matching our X-ray source positions (columns 2 and 3) to WFI $`R`$-band positions using a matching radius of 1.5 times the positional uncertainty quoted in column 4. For a small number of sources more than one optical match was found, and for these sources the brightest match was selected as the most probable counterpart. Using these criteria, 594 ($``$78%) of the sources have optical counterparts. Note that the matching criterion used here is more conservative than that used in the derivation of our positional errors discussed in $`\mathrm{\S }`$ 3.3.1. Sources with no optical counterparts have RA and Dec values set to โ00 00 00.00โ and โ+00 00 00.0โ.
* Column 17 gives the measured offset between the optical and X-ray sources (i.e., $`OX`$) in arcseconds. Sources with no optical counterparts have a value set to โ0โ.
* Column 18 gives the $`R`$-band magnitude (AB) of each X-ray source. Sources with no optical counterparts have a value set to โ0โ.
* Column 19 gives the $``$1 Ms CDF-S source number from the main Chandra catalog presented in A03 (see column 1 of Table 3a in A03) for E-CDF-S sources that were matched to A03 counterparts. We used a matching radius of 1.5 times the sum of the positional errors of the E-CDF-S and A03 source positions. We note that for each matched source only one match was observed; E-CDF-S sources with no A03 match have a value of โ0โ.
* Columns 20 and 21 give the RA and Dec of the corresponding $``$1 Ms CDF-S A03 source indicated in column 19. Sources with no A03 match have RA and Dec values set to โ00 00 00.00โ and โ+00 00 00.0โ.
* Column 22 gives the $``$1 Ms CDF-S source number from the main Chandra catalog presented in G02 (see โIDโ column of Table 2 in G02) for E-CDF-S sources that were matched to G02 counterparts. When matching our E-CDF-S source positions with G02 counterparts, we removed noted offsets to the G02 positions of $`1\stackrel{}{\mathrm{.}}2`$ in RA and $`+0\stackrel{}{\mathrm{.}}8`$ in Dec (see $`\mathrm{\S }`$ A3 of A03); these positions are corrected in the quoted source positions in columns 23 and 24. We used a matching radius of 1.5 times the E-CDF-S positional error plus the G02-quoted positional error for each source position. We note that for each matched source only one match was observed; E-CDF-S sources with no G02 match have a value of โ0โ.
* Columns 23 and 24 give the RA and Dec of the corresponding $``$1 Ms CDF-S G02 source indicated in column 22. Note that the quoted positions have been corrected by the noted offsets described in column 22 (see $`\mathrm{\S }`$ A3 of A03). Sources with no G02 match have RA and Dec values set to โ00 00 00.00โ and โ+00 00 00.0โ.
* Columns 25โ27 give the effective exposure times derived from the standard-band exposure maps (see ยง3.1 for details on the exposure maps). Dividing the counts listed in columns 6โ14 by the corresponding effective exposures will provide vignetting-corrected and quantum efficiency degradation-corrected count rates.
* Columns 28โ30 give the band ratio, defined as the ratio of counts between the hard and soft bands, and the corresponding upper and lower errors, respectively. Quoted band ratios have been corrected for differential vignetting between the hard band and soft band using the appropriate exposure maps. Errors for this quantity are calculated following the โnumerical methodโ described in ยง1.7.3 of Lyons (1991); this avoids the failure of the standard approximate variance formula when the number of counts is small (see ยง2.4.5 of Eadie et al. 1971). Note that the error distribution is not Gaussian when the number of counts is small. Upper limits are calculated for sources detected in the soft band but not the hard band and lower limits are calculated for sources detected in the hard band but not the soft band. For these sources, the upper and lower errors are set to the computed band ratio.
* Columns 31โ33 give the effective photon index ($`\mathrm{\Gamma }`$) with upper and lower errors, respectively, for a power-law model with the Galactic column density. The effective photon index has been calculated based on the band ratio in column 28 when the number of counts is not low.
A source with a low number of counts is defined as being (1) detected in the soft band with $`<30`$ counts and not detected in the hard band, (2) detected in the hard band with $`<15`$ counts and not detected in the soft band, (3) detected in both the soft and hard bands, but with $`<15`$ counts in each, or (4) detected only in the full band. When the number of counts is low, the photon index is poorly constrained and is set to $`\mathrm{\Gamma }=1.4`$, a representative value for faint sources that should give reasonable fluxes. Upper and lower limits are indicated by setting the upper and lower errors to the computed effective photon index.
* Columns 34โ36 give observed-frame fluxes in the three standard bands; quoted fluxes are in units of $`10^{15}`$ erg cm<sup>-2</sup> s<sup>-1</sup>. Fluxes have been computed using the counts in columns 6, 9, and 12, the appropriate exposure maps (columns 25โ27), and the spectral slopes given in column 31. The fluxes have not been corrected for absorption by the Galaxy or material intrinsic to the source. For a power-law model with $`\mathrm{\Gamma }=1.4`$, the soft-band and hard-band Galactic absorption corrections are $``$2.1% and $`0.1`$%, respectively. More accurate fluxes for these sources would require direct fitting of the X-ray spectra for each observation, which is beyond the scope of this paper.
* Column 37 gives the observational field number corresponding to the detected source. The observational fields overlap in a few areas (see Figures 5 and 6a) over $``$50 arcmin<sup>2</sup>, which allowed for duplicate detections of a single source. Fourteen sources in the Chandra catalog were detected in more than one observational field; these sources are flagged in column 39 (see below). The data from the observation that produced the greatest number of full-band counts for these sources is included here; properties derived from the cross-field observations are provided in Table 3.
* Column 38 gives the logarithm of the minimum false-positive probability run with wavdetect in which each source was detected (see $`\mathrm{\S }`$ 3.2). A lower false-positive probability indicates a more significant source detection. Note that 655 ($``$86%) and 596 ($``$78%) of our sources are detected with false-positive probability thresholds of 1 $`\times `$ 10<sup>-7</sup> and 1 $`\times `$ 10<sup>-8</sup>, respectively.
* Column 39 gives notes on the sources. โDโ denotes a source detected in more than one of the four observational fields. โUโ denotes objects lying in the UDF (see Figure 2). โGโ denotes objects that were identified as Galactic stars through the optical spectrophotometric COMBO-17 survey (Wolf et al. 2004). โOโ refers to objects that have large cross-band (i.e., between the three standard bands) positional offsets ($`>2^{\prime \prime }`$); all of these sources lie at off-axis angles of $`>8^{}`$. โMโ refers to sources where the photometry was performed manually. โSโ refers to close-double or close-triple sources where manual separation was required. โCโ refers to sources detected within the boundary of the $``$1 Ms CDF-S exposure that have no A03 or G02 counterparts. Several of these sources are located in low-sensitivity regions of the $``$1 Ms CDF-S, and a few of these sources may be variable. For further explanation of many of these notes, see the above text in this section on manual correction of the wavdetect results.
In Table 3 we summarize the cross-field source properties of the 14 sources detected in more than one observational field; none of these sources were detected in more than two fields. These properties were derived from the observation not included in the main Chandra catalog (see columns 37 and 39 of Table 2) and are included here for comparison. The columns of Table 3 are the same as those in Table 2; the source number for each source corresponds to its duplicate listed in the main Chandra catalog. In Table 4 we summarize the source detections in the three standard bands for each of the observational fields and the main Chandra catalog. In total 776 point sources are detected in one or more of the three standard bands; 14 of these sources are detected in more than one of the four observational fields (see columns 37 and 39 of Tables 2 and 3) leaving a total of 762 distinct point sources. Out of these 762 distinct point sources, we find that 173 are coincident with sources included in the main Chandra catalog for the $``$1 Ms CDF-S presented in A03 (see columns 19โ24 of Table 2). For these sources, we find reasonable agreement between the derived X-ray properties presented here and in A03. A total of 589 new point sources are thus detected here, which brings the total number of E-CDF-S plus $``$1 Ms CDF-S sources to 915.
In Table 5 we summarize the number of sources detected in one band but not another. All but two of the detected sources are detected in either the soft or full bands. From Tables 4 and 5, the fraction of hard-band sources not detected in the soft band is $`96/45321\%`$. The fraction is somewhat higher than for the Chandra Deep Fields, where it is $`14\%`$. Some of this difference is likely due to differing methods of cross-band matching (i.e., compare $`\mathrm{\S }`$ 3.4.1 of A03 with our $`\mathrm{\S }`$ 3.3.1). Furthermore, this fraction is physically expected to vary somewhat with sensitivity limit. We have also attempted comparisons with X-ray surveys of comparable depth (Stern et al. 2002; Harrison et al. 2003; Wang et al. 2004a,b; Nandra et al. 2005) to the E-CDF-S. Such comparisons are not entirely straightforward due to varying energy bands utilized, source-selection techniques, and source-searching methods. However, the โhard-band but not soft-bandโ fractions for these surveys appear plausibly consistent ($``$15โ25%) with that for the E-CDF-S.
In Figure 9 we show the distributions of detected counts in the three standard bands. There are 154 sources with $`>`$ 100 full-band counts, for which basic spectral analyses are possible; there are eight sources with $`>`$ 1000 full-band counts. Figure 10 shows the distribution of effective exposure time for the three standard bands. The median effective exposure times for the soft and hard bands are $``$216 ks and $``$212 ks, respectively. In Figure 11 we show the distributions of X-ray flux in the three standard bands. The X-ray fluxes in this survey span roughly four orders of magnitude with $``$50% of the sources having soft band fluxes less than 50 $`\times `$ 10<sup>-17</sup> erg cm<sup>-2</sup> s<sup>-1</sup>, a flux regime that few X-ray surveys have probed with significant areal coverage.
In Figure 12 we show โpostage-stampโ images from the WFI $`R`$-band image with adaptively-smoothed full-band contours overlaid for sources included in the main Chandra catalog. The wide range of X-ray source sizes observed in these images is largely due to PSF broadening with off-axis angle. In Figure 13 we plot the positions of sources detected in the main Chandra catalog. Sources that are also included in the A03 CDF-S source catalog are indicated as open circles, and new X-ray sources detected in this survey are indicated as filled circles; the circle sizes depend upon the most significant false-positive probability run with wavdetect for which each source was detected (see column 38 of Table 2). The majority of the sources lie in the vicinities of the aim points where the fields are most sensitive. In Figure 14 we show the band ratio as a function of full-band count rate for sources in the main Chandra catalog. This plot shows that the mean band ratio for sources detected in both the soft and hard bands hardens for fainter fluxes, a trend observed in other studies (e.g., della Ceca et al. 1999; Ueda et al. 1999; Mushotzky et al. 2000; B01; Tozzi et al. 2001; A03). This trend is due to the detection of more absorbed AGNs at low flux levels, and it has been shown that AGNs will dominate the number counts down to 0.5โ2.0 keV fluxes of $``$$`\times `$ 10<sup>-17</sup> erg cm<sup>-2</sup> s<sup>-1</sup> (e.g., Bauer et al. 2004). Figure 15a shows the $`R`$-band magnitude versus the soft band flux for sources included in the main Chandra catalog. The approximate X-ray to $`R`$-band flux ratios for AGNs and galaxies (e.g., Maccaccaro et al. 1998; Stocke et al. 1991; Hornschemeier et al. 2001; Bauer et al. 2004) are indicated with dark and light shading, respectively. The majority of the sources in this survey appear to be AGNs. Sixty-one of the sources were reliably classified as AGNs and seventeen sources have been identified as Galactic stars (see column 39 of Table 2) in the COMBO-17 survey (Wolf et al. 2004). A significant minority of the sources appear to have X-ray-to-optical flux ratios characteristic of normal or starburst galaxies.
#### 3.3.2 Supplementary Optically Bright Chandra Source Catalog
The density of optically bright ($`R<23`$) sources on the sky is comparatively low. Therefore, we can search for X-ray counterparts to optically bright sources at a lower X-ray significance threshold than that used in the main catalog without introducing many false sources (see ยง5.3 of Richards et al. 1998 for a similar technique applied at radio wavelengths). We ran wavdetect with a false-positive probability threshold of $`1\times 10^5`$ on images created in the three standard bands. A basic lower significance Chandra catalog was produced containing 323 X-ray sources not present in the main Chandra source catalog.
In our matching of these lower significance Chandra sources to optically bright sources, we used the WFI $`R`$-band source catalog described in ยง3.3. We searched for X-ray counterparts to these optical sources using a matching radius of $`1\stackrel{}{\mathrm{.}}3`$. Based upon offset tests, as described below, we found empirically that we could match to sources as faint as $`R=23`$ without introducing an unacceptable number of false matches; this $`R`$-band cutoff provides an appropriate balance between the number of detected sources and the expected number of false sources.
In total 26 optically bright X-ray sources were found via our matching. We estimated the expected number of false matches by artificially offsetting the X-ray source coordinates in RA and Dec by both $`5^{\prime \prime }`$ and $`10^{\prime \prime }`$ (using both positive and negative shifts) and then re-correlating with the optical sources. On average $``$3 matches were found with these tests, demonstrating that the majority of the 26 X-ray matches are real X-ray sources; only about 12% of these sources are expected to be spurious matches.
We also included seven $`R<21`$ sources where the X-ray source lay $`1\stackrel{}{\mathrm{.}}3`$$`10\stackrel{}{\mathrm{.}}0`$ from the centroid of the optical source but was still within the extent of the optical emission. Using optical spectrophotometric redshift information from COMBO-17, we required that our off-nuclear sources have 0.5โ2.0 keV luminosities of $`\stackrel{<}{}`$$`10^{40}`$ erg s<sup>-1</sup>. This restriction was intended to remove obvious sources not associated with their host galaxies; this led to the removal of one candidate source (J$`033210.9280230`$) when forming our sample of seven plausible off-nuclear X-ray sources. Of the seven selected off-nuclear sources included in the supplementary catalog, we found that six of the host galaxies have COMBO-17 redshift information available. These galaxies were found to have $`z0.100.25`$ and 0.5โ2.0 keV luminosities in the range of $``$$`10^{3940}`$ erg s<sup>-1</sup>. These derived luminosities are consistent with these sources being off-nuclear X-ray binaries or star-forming regions associated with bright host galaxies. Since these seven sources were identified in a somewhat subjective manner, it is not meaningful to determine a false-matching probability for them. These sources are indicated in column 33 of Table 6. Thus, in total, the supplementary optically bright Chandra source catalog contains 33 sources.
The format of Table 6 is similar to that of Table 2. Details of the columns in Table 6 are given below.
* Column 1 gives the source number (see column 1 of Table 2 for details).
* Columns 2 and 3 give the RA and Dec of the X-ray source, respectively. The wavdetect positions are given for these faint X-ray sources. Whenever possible, we quote the position determined in the full band; when a source is not detected in the full band we use, in order of priority, the soft-band position and then the hard-band position. The priority ordering of position choices above was designed to maximize the signal-to-noise of the data being used for positional determination.
* Column 4 gives the positional uncertainty in arcseconds. For these faint X-ray sources, the positional uncertainty is take to be $`1\stackrel{}{\mathrm{.}}2`$, the 90th percentile of the average optical-X-ray positional offsets given in column 17.
* Column 5 gives the off-axis angle for each source in arcminutes (see column 5 of Table 2 for details).
* Columns 6โ14 give the counts and the corresponding 1$`\sigma `$ upper and lower statistical errors (using Gehrels 1986), respectively, for the three standard bands. The photometry is taken directly from wavdetect for these faint X-ray sources.
* Columns 15 and 16 give the RA and Dec of the optical source centroid, respectively.
* Column 17 gives the measured offset between the optical and X-ray sources (i.e., $`OX`$) in arcseconds.
* Column 18 gives the $`R`$-band magnitude (AB) of the optical source.
* Column 19 gives the $``$1 Ms CDF-S source number from the main Chandra catalog presented in A03 (see column 1 of Table 3a in A03) for supplementary sources that were matched to A03 counterparts. We used a matching radius of 1.5 times the sum of the positional errors of the E-CDF-S and A03 source positions. We note that for each matched source only one match was observed; supplementary sources with no A03 match have a value of โ0โ.
* Columns 20 and 21 give the RA and Dec of the corresponding $``$1 Ms CDF-S A03 source indicated in column 19. Sources with no A03 match have RA and Dec values set to โ00 00 00.00โ and โ+00 00 00.0โ.
* Column 22 gives the $``$1 Ms CDF-S source number from the main Chandra catalog presented in G02 (see โIDโ column of Table 2 in G02) for supplementary sources that were matched to G02 counterparts. When matching our supplementary source positions with G02 counterparts, we removed noted offsets to the G02 positions of $`1\stackrel{}{\mathrm{.}}2`$ in RA and $`+0\stackrel{}{\mathrm{.}}8`$ in Dec (see $`\mathrm{\S }`$ A3 of A03); these positions are corrected in the quoted source positions in columns 23 and 24. We used a matching radius of 1.5 times the E-CDF-S positional error plus the G02-quoted positional error for each source position. We note that for each matched source only one match was observed; supplementary sources with no G02 match have a value of โ0โ.
* Columns 23 and 24 give the RA and Dec of the corresponding $``$1 Ms CDF-S G02 source indicated in column 22. Note that the quoted positions have been corrected by the noted offsets described in column 22 (see $`\mathrm{\S }`$ A3 of A03). Sources with no G02 match have RA and Dec values set to โ00 00 00.00โ and โ+00 00 00.0โ.
* Columns 25โ27 give the effective exposure times derived from the standard-band exposure maps (see columns 25โ27 of Table 2 for details).
* Column 28 gives the photon index used to calculate source fluxes (columns 26โ28). We used a constant photon index of $`\mathrm{\Gamma }=2.0`$ since our source-selection technique preferentially selects objects with flux-ratios $`f_{0.52.0\mathrm{keV}}/f_R<0.1`$, which are observed to have effective photon indices of $`\mathrm{\Gamma }2`$ (e.g., $`\mathrm{\S }`$ 4.1.1 of Bauer et al. 2004).
* Column 29โ31 give observed-frame fluxes in the three standard bands; quoted fluxes are in units of $`10^{15}`$ erg cm<sup>-2</sup> s<sup>-1</sup> and have been calculated assuming $`\mathrm{\Gamma }=2.0`$. The fluxes have not been corrected for absorption by the Galaxy or material intrinsic to the sources (see columns 34โ36 of Table 2 for details).
* Column 32 gives the observational field number corresponding to the detected source (see column 37 of Table 2 for details).
* Column 33 gives notes on the sources. With the exception of the additional note given below, the key for these notes is given in column 39 of Table 2. โLโ refers to objects where the X-ray source lies $`>1\stackrel{}{\mathrm{.}}3`$ from the centroid of the optical source but is still within the extent of the optical emission (see the text above for further discussion).
The $`R`$-band magnitudes of the supplementary sources span $`R=`$ 16.6โ22.9. In Figure 15b we show the $`R`$-band magnitude versus soft-band flux. All of the sources lie in the region expected for starburst and normal galaxies. Three of these sources have been classified as Galactic stars via optical classifications from COMBO-17 (see column 33 of Table 6). Some of these sources may be low-luminosity AGNs; the small number of hard-band detections ($``$6$`\%`$) indicates that few of these are absorbed AGNs. Due to the low number of counts detected for these sources, we do not provide postage-stamp images as we did for sources in the main catalog (i.e., Figure 12).
The addition of the optically bright supplementary sources increases the number of extragalactic objects in the E-CDF-S with $`f_{0.52.0\mathrm{keV}}/f_R<0.1`$ and $`f_{0.52.0\mathrm{keV}}/f_R<0.01`$ by $``$20% and $``$50%, respectively. However, the optically bright supplementary sources are not representative of the faintest X-ray sources as a whole since our selection criteria preferentially select optically bright and X-ray faint non-AGNs (e.g., A03; Hornschemeier et al. 2003).
## 4 Background and Sensitivity Analysis
The faintest sources in the main Chandra catalog have $``$4 counts (see Table 4). For a $`\mathrm{\Gamma }=`$ 1.4 power law with Galactic absorption, the corresponding soft-band and hard-band fluxes at the aim points are $``$8.9 $`\times `$ 10<sup>-17</sup> erg cm<sup>-2</sup> s<sup>-1</sup> and $``$4.4 $`\times `$ 10<sup>-16</sup> erg cm<sup>-2</sup> s<sup>-1</sup>, respectively. This gives a measure of the ultimate sensitivity of this survey, however, these numbers are only relevant for a small region close to the aim point. To determine the sensitivity across the field it is necessary to take into account the broadening of the PSF with off-axis angle, as well as changes in the effective exposure and background rate across the field. We estimated the sensitivity across the field by employing a Poisson model, which was calibrated for sources detected in the main Chandra catalog. Our resulting relation is
$$\mathrm{log}(N)=\alpha +\beta \mathrm{log}(b)+\gamma [\mathrm{log}(b)]^2+\delta [\mathrm{log}(b)]^3$$
(4)
where $`N`$ is the required number of source counts for detection, and $`b`$ is the number of background counts in a source cell; $`\alpha =0.967`$, $`\beta =0.414`$, $`\gamma =0.0822`$, and $`\delta =0.0051`$ are fitting constants. The only component within this equation that we need to measure is the background. For the sensitivity calculations here we measured the background in a source cell using the background maps described below and assumed an aperture size of 70% of the encircled-energy radius of the PSF; the 70% encircled energy-radius was chosen as a compromise between having too few source counts and too many background counts. The total background includes contributions from the unresolved cosmic background, particle background, and instrumental background (e.g., Markevitch 2001; Markevitch et al. 2003). For our analyses we are only interested in the total background and do not distinguish between these different components.
To create background maps for all of the twelve images, we first masked out the point sources from the main Chandra catalog using apertures with radii twice that of the $``$ 90% PSF encircled-energy radii. The resultant images should include minimal contributions from detected point sources. They will, however, include contributions from extended sources (e.g., Bauer et al. 2002; see $`\mathrm{\S }`$ 6), which will cause a slight overestimation of the measured background close to extended sources. Extensive testing of background-count distributions in all three standard bands has shown that the X-ray background follows a nearly Poisson count distribution (see $`\mathrm{\S }`$ 4.2 of A03). We filled in the masked regions for each source with a local background estimate by making a probability distribution of counts using an annulus with inner and outer radii of 2 and 4 times the $``$90% PSF encircled-energy radius, respectively. The background properties are summarized in Table 7, and the full-band background map is shown in Figure 16. The majority of the pixels have no background counts (e.g., in the full band $``$97% of the pixels are zero) and the mean background count rates for these observations are broadly consistent with those presented in A03.
Following equation 4, we generated sensitivity maps using the background and exposure maps; we assumed a $`\mathrm{\Gamma }=`$ 1.4 power-law model with Galactic absorption. In Figure 17 we show the full-band sensitivity map, and in Figure 18 we show plots of flux limit versus solid angle for the full, soft, and hard bands. The $``$1 arcmin<sup>2</sup> regions at the aim points has average 0.5โ2.0 keV and 2โ8 keV sensitivity limits of $``$1.1 $`\times `$ 10<sup>-16</sup> erg cm<sup>-2</sup> s<sup>-1</sup> and $``$6.7 $`\times `$ 10<sup>-16</sup> erg cm<sup>-2</sup> s<sup>-1</sup>, respectively. Since we do not filter out detected sources with our sensitivity maps, a number of sources have fluxes below these sensitivity limits (4 sources in the soft band and 17 sources in the hard band). Approximately 800 arcmin<sup>2</sup> of the field (i.e., $``$ 3 times the size of a single ACIS-I field) has a soft-band sensitivity limit of $`\stackrel{<}{}3\times 10^{16}`$ erg cm<sup>-2</sup> s<sup>-1</sup>, well into the flux range where few X-ray surveys have probed (see Figure 1).
## 5 Number Counts for Main Chandra Catalog
We have calculated cumulative number counts, $`N(>S)`$, for the soft and hard bands using sources presented in our main Chandra catalog (see Table 2). We restricted our analyses to flux levels where we expect to be mostly complete based on our sensitivity maps and simulations performed by Bauer et al. (2004); this also helps to guard against Eddington bias at low flux levels. We empirically set our minimum flux levels to 3.0 $`\times `$ 10<sup>-16</sup> erg cm<sup>-2</sup> s<sup>-1</sup> in the soft band and 1.2 $`\times `$ 10<sup>-15</sup> erg cm<sup>-2</sup> s<sup>-1</sup> in the hard band, which correspond to the minimum detected fluxes for sources with $`\stackrel{>}{}`$15 counts in each respective band.
Assuming completeness to the flux levels quoted above, the cumulative number of sources, $`N(>S)`$, brighter than a given flux, $`S`$, weighted by the appropriate aerial coverage, is
$$N(>S)=\underset{S_i>S}{}\mathrm{\Omega }_i^1$$
(5)
where $`\mathrm{\Omega }_i`$ is the maximum solid angle for which a source with measured flux, $`S_i`$, could be detected. Each maximum solid angle was computed using the profiles presented in Figure 18. In Figure 19 we show the cumulative number counts for the main Chandra catalog. Number counts derived for the $``$1 Ms CDF-S from Rosati et al. (2002) have been plotted for comparison. The E-CDF-S number counts appear to be consistent with those from the $``$1 Ms CDF-S to within $``$$`\sigma `$ over the overlapping flux ranges. We note, however, that the hard-band number counts appear to be generally elevated with respect to those from the $``$1 Ms CDF-S; this effect is likely a signature of field-to-field variance from the $``$1 Ms CDF-S where a smaller solid angle is surveyed. Even with the conservative flux constraints used in our number-counts analysis we reach source densities exceeding $``$2000 deg<sup>-2</sup>; as noted in $`\mathrm{\S }`$ 3.3.1, a large majority of these sources are AGNs. For comparison, the number density of COMBO-17 sources with reliable AGN identifications is $``$300 deg<sup>-2</sup> (Wolf et al. 2004).
## 6 Extended Sources
We searched the standard-band images for extended sources using the Voronoi Tessellation and Percolation algorithm vtpdetect (Ebeling & Wiedenmann 1993; Dobrzycki et al. 2002). In our vtpdetect searching, we adopted a false-positive probability threshold of 1 $`\times `$ 10<sup>-7</sup> and a โcoarseโ parameter of 50. Following the source-detection criteria presented in Bauer et al. (2002), we further required that vtpdetect-detected sources have (1) average vtpdetect radii (i.e., average of the 3$`\sigma `$ major and minor axes estimated by vtpdetect) $``$ three times the 95% encircled-energy radius of a point source at the given position and (2) visible evidence for extended emission in the adaptively smoothed, exposure-corrected images (see $`\mathrm{\S }`$ 3.3.1 and Figure 4). Application of these somewhat conservative selection criteria yielded three extended X-ray sources, all of which are detected only in the soft band. The soft emission from the most significant of these three extended sources, CXOECDFS J033320.3$``$274836, is clearly visible as an extended red โglowโ near the left-hand side of Figure 7.
The X-ray properties of these three extended sources are presented in Table 8; our analysis was limited to the soft band where we find all of our detections. The counts for extended sources were determined using manual aperture photometry; point sources were masked out using circular apertures with radii of twice the 95% encircled-energy radii (see Footnote 3). We extracted extended-source counts using elliptical apertures with sizes and orientations that most closely matched the apparent extent of X-ray emission ($`>`$ 10% above the background level) as observed in the adaptively smoothed images (see Table 8). The local background was estimated using elliptical annuli with inner and outer sizes of 1 and 2.5 times those used for extracting source counts. In order to calculate properly the expected numbers of background counts in our source extraction regions, we extracted total exposure times from the source and background regions (with point sources removed) and normalized the extracted background counts to the source exposure times. That is, using the number of background counts $`b_\mathrm{m}`$ and total background exposure time $`T_\mathrm{m}`$ as measured from the elliptical annuli, we calculated the expected number of background counts $`b_\mathrm{s}`$ in a source extraction region with total exposure time $`T_\mathrm{s}`$ as being $`b_\mathrm{s}=b_\mathrm{m}T_\mathrm{s}/T_\mathrm{m}`$. This technique was used to account for extended emission from sources that are spatially distributed over more than one observational field, which was the case for CXOECDFS J033320.3$``$274836.
Figure 20 shows WFI $`R`$-band images of the extended sources with adaptively smoothed soft-band contours overlaid. Inspection of the spectrophotometric redshifts of optical sources (from COMBO-17) in these regions suggests that the extended X-ray emission for all three sources may originate from low-to-moderate redshift groups or poor clusters. The most conspicuous of these is the apparent clustering of galaxies at $`z0.1`$ in the area of CXOECDFS J033320.3$``$274836, an $``$20 arcmin<sup>2</sup> extended X-ray source. CXOECDFS J033209.6$``$274242 lies in the $``$1 Ms CDF-S and was previously detected as an extended source by G02. Optical spectroscopic follow-up observations using the VLT have shown that this source is associated with a galaxy cluster at $`z=0.73`$ (Szokoly et al. 2004). Suggestive evidence for clustering at $`z0.7`$ is also observed for CXOECDFS J033257.9$``$280155; this may be an extension of the large-scale structures observed in the $``$1 Ms CDF-S (Gilli et al. 2003, 2005). Under the assumption that these sources are indeed groups or poor clusters at the discussed redshifts, we computed the expected soft-band fluxes and luminosities assuming a Raymond-Smith thermal plasma (Raymond & Smith 1977) with $`kT=1.0`$ keV (see Table 8). We find that these sources would have rest-frame 0.5โ2.0 keV luminosities of $``$1โ5 $`\times 10^{42}`$ erg s<sup>-1</sup>. Further optical spectroscopic observations and analyses beyond the scope of this paper are required for confirmation of the nature of these sources.
## 7 Summary
We have presented catalogs and basic analyses of point sources detected in the 250 ks $``$0.3 deg<sup>2</sup> Extended Chandra Deep Field-South (E-CDF-S). The survey area consists of four observational fields, of similar exposure, with average on-axis flux limits in the 0.5โ2.0 keV and 2โ8 keV bandpasses of $``$1.1 $`\times `$ 10<sup>-16</sup> erg cm<sup>-2</sup> s<sup>-1</sup> and $``$6.7 $`\times `$ 10<sup>-16</sup> erg cm<sup>-2</sup> s<sup>-1</sup>, respectively. We have presented two catalogs: a main Chandra catalog of 762 sources (589 of these are new), which was generated by running wavdetect with a false-positive probability threshold of 1 $`\times `$ 10<sup>-6</sup>, and a supplementary catalog of 33 lower-significance (false-positive probability threshold of 1 $`\times `$ 10<sup>-5</sup>) X-ray sources with optically bright ($`R`$ $`<`$ 23) counterparts. The X-ray spectral properties and optical fluxes of sources in our main Chandra catalog indicate a variety of source types, most of which are absorbed AGNs that dominate at lower X-ray fluxes. The X-ray and optical properties of sources in the supplementary optically bright Chandra catalog are mostly consistent with those expected for starburst and normal galaxies. We have presented basic number-count results for point sources in our main Chandra catalog and find overall consistency with number counts derived for the $``$1 Ms CDF-S in both the 0.5โ2.0 keV and 2โ8 keV bandpasses. We have also presented three 0.5โ2.0 keV extended sources, which were detected using a conservative detection criterion. These sources are likely associated with groups or poor clusters at $`z0.10.7`$ with $`L_\text{X}`$ $``$ 1โ5 $`\times 10^{42}`$ erg s<sup>-1</sup>.
We thank G. Chartas, P.E. Freeman, A.T. Steffen, and L.K. Townsley for helpful discussions. We thank the anonymous referee for useful comments that improved the manuscript. Support for this work was provided by NASA through Chandra Award Number GO4-5157 (BDL, WNB, DPS, RG, AMK, TM) issued by the Chandra X-ray Observatory Center, which is operated by the Smithsonian Astrophysical Observatory under NASA contract NAS8-03060. We also acknowledge the financial support of NSF CAREER award AST-9983783 (BDL, WNB), the Royal Society (DMA), the Chandra Fellowship program (FEB), NSF AST 03-07582 (DPS), NASA LTSA Grant NAG5-10875 (TM), and MIUR COFIN grant 03-02-23 (CV).
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# Loops of Bol-Moufang type with a subgroup of index two
## 1. Introduction
Due to the specialized nature of this paper we assume that the reader is already familiar with the theory of quasigroups and loops. We therefore omit basic definitions and results (see Bruck (1971), Pflugfelder (1990)).
In a sense, a nonassociative loop is closest to a group when it contains a subgroup of index two. Such loops proved useful in the study of Moufang loops, and it is our opinion that they will also prove useful in the study of other varieties of loops.
Here is the well-known construction of Moufang loops with a subgroup of index two:
###### Theorem 1.1 (Chein Chein (1978)).
Let $`G`$ be a group, $`g_0Z(G)`$, and $``$ an involutory antiautomorphism of $`G`$ such that $`g_0^{}=g_0`$, $`gg^{}Z(G)`$ for every $`gG`$. For an indeterminate $`u`$, define multiplication $``$ on $`GGu`$ by
$$gh=gh,g(hu)=(hg)u,guh=(gh^{})u,guhu=g_0h^{}g,$$
(1)
where $`g`$, $`hG`$. Then $`L=(GGu,)`$ is a Moufang loop. Moreover, $`L`$ is associative if and only if $`G`$ is commutative.
It has been shown in Vojtฤchovskรฝ (2003) that (1) is the only construction of its kind for Moufang loops. (This statement will be clarified later.) In Vojtฤchovskรฝ (2004), all constructions similar to (1) were determined for Bol loops.
The purpose of this paper is to give a complete list of all constructions similar to (1) for all loops of Bol-Moufang type. A groupoid identity is of *Bol-Moufang type* if it has three distinct variables, two of the variables occur once on each side, the third variable occurs twice on each side, and the variables occur in the same order on both sides. A loop is of *Bol-Moufang type* if it belongs to a variety of loops defined by a single identity of Bol-Moufang type. Figure 1 shows all varieties of loops of Bol-Moufang type and all inclusions among them (cf. Fenyves (1969), Phillips and Vojtฤchovskรฝ (2005)). Some varieties of Figure 1 can be defined equivalently by other identities of Bol-Moufang type. For instance, Moufang loops are equivalently defined by the identity $`x(y(xz))=((xy)x)z`$. See Phillips and Vojtฤchovskรฝ (2005) for all such equivalences. Furthermore, although some defining identities of Figure 1 do not appear to be of Bol-Moufang type, they are in fact equivalent to some Bol-Moufang identity. For instance, the flexible law $`x(yx)=(xy)x`$ is equivalent to the Bol-Moufang identity $`(x(yx))z=((xy)x)z`$ in any variety of loops.
As we shall see, the computational complexity of our programme is overwhelming (for humans). We therefore first carefully define what we mean by a construction similar to (1) (see Section 2), and then identify situations in which two given constructions are โthe sameโ (see Sections 3, 4, 5). Upon showing which constructions yield loops, we work out one construction by hand (see Section 6), and then switch to a computer search, described in Section 7. The results of the computer search are summarized in Section 8.
## 2. Similar Constructions
Throughout the paper, we assume that $`G`$ is a finite group, $`g_0Z(G)`$, and $``$ is an involutory automorphism of $`G`$ such that $`g_0^{}=g_0`$ and $`gg^{}Z(G)`$ for every $`gG`$.
The following property of $``$ will be used without reference:
###### Lemma 2.1.
Let $`G`$ be a group and $`:GG`$ an involutory map such that $`gg^{}Z(G)`$ for every $`gG`$. Then $`g^{}g=gg^{}Z(G)`$ for every $`gG`$.
###### Proof.
For $`gG`$, we have $`g^{}g=g^{}(g^{})^{}Z(G)`$. Then $`(g^{}g)g^{}=g^{}(g^{}g)`$, and $`gg^{}=g^{}g`$ follows upon cancelling $`g^{}`$ on the left. โ
Consider the following eight bijections of $`G\times G`$:
$$\begin{array}{cccc}\theta _{xy}(g,h)=(g,h),\hfill & \theta _{xy^{}}(g,h)=(g,h^{}),\hfill & \theta _{x^{}y}(g,h)=(g^{},h),\hfill & \theta _{x^{}y^{}}(g,h)=(g^{},h^{}),\hfill \\ \theta _{yx}(g,h)=(h,g),\hfill & \theta _{yx^{}}(g,h)=(h,g^{}),\hfill & \theta _{y^{}x}(g,h)=(h^{},g),\hfill & \theta _{y^{}x^{}}(g,h)=(h^{},g^{}).\hfill \end{array}$$
They form a group $`\mathrm{\Theta }`$ under composition, isomorphic to the dihedral group $`D_8`$. It is generated by $`\{\theta _{yx},\theta _{xy^{}}\}`$, say. Let $`\mathrm{\Theta }_0`$ be the group generated by $`\mathrm{\Theta }`$ and $`\theta _{g_0}`$, where $`\theta _{g_0}(g,h)=(g_0g,h)`$.
Let $`\mathrm{\Delta }:G\times GG`$ be the evaluation map $`\mathrm{\Delta }(g,h)=gh`$, and $`u`$ an indeterminate. Given $`\alpha `$, $`\beta `$, $`\gamma `$, $`\delta \mathrm{\Theta }_0`$, define multiplication $``$ on $`GGu`$ by
$$gh=\mathrm{\Delta }\alpha (g,h),ghu=(\mathrm{\Delta }\beta (g,h))u,guh=(\mathrm{\Delta }\gamma (g,h))u,guhu=\mathrm{\Delta }\delta (g,h),$$
where $`g`$, $`hG`$. The resulting groupoid $`(GGu,)`$ will be denoted by
$$Q(G,,g_0,\alpha ,\beta ,\gamma ,\delta ),$$
or by $`Q(G,\alpha ,\beta ,\gamma ,\delta )`$, when $`g_0`$, $``$ are known from the context or if they are not important. It is easy to check that $`Q(G,,g_0,\alpha ,\beta ,\gamma ,\delta )`$ is a quasigroup.
We also define
$$๐ฌ(G,,g_0)=\{Q(G,,g_0,\alpha ,\beta ,\gamma ,\delta );\alpha ,\beta ,\gamma ,\delta \mathrm{\Theta }_0\},$$
and
$$๐ฌ(G)=\underset{,g_0}{}๐ฌ(G,,g_0),$$
where the union is taken over all involutory antiautomorphisms $``$ satisfying $`gg^{}Z(G)`$ for every $`gG`$, and over all elements $`g_0`$ such that $`g_0^{}=g_0Z(G)`$. By definition, we call elements of $`๐ฌ(G)`$ *quasigroups obtained from $`G`$ by a construction similar to* (1).
## 3. Reductions
The goal of this section is to show that one does not have to take all elements of $`\mathrm{\Theta }_0`$ into consideration in order to determine $`๐ฌ(G,,g_0)`$.
Note that $`g_0^n=(g_0^n)^{}Z(G)`$ for every integer $`n`$. Therefore
$$g_0^n\mathrm{\Delta }\theta _0(g,h)=\mathrm{\Delta }\theta _{g_0}^n\theta _0(g,h)=\mathrm{\Delta }\theta _0\theta _{g_0}^n(g,h)$$
(2)
for every $`\theta _0\mathrm{\Theta }_0`$ and every $`g`$, $`hG`$.
###### Lemma 3.1.
For every integer $`n`$, the quasigroup $`Q(G,\theta _{g_0}^n\alpha ,\theta _{g_0}^n\beta ,\theta _{g_0}^n\gamma ,\theta _{g_0}^n\delta )`$ is isomorphic to $`Q(G,\alpha ,\beta ,\gamma ,\delta )`$.
###### Proof.
We use (2) freely in this proof. Let $`t=g_0^n`$. Denote by $``$ the multiplication in $`Q(G,\alpha ,\beta ,\gamma ,\delta )`$, and by $``$ the multiplication in $`Q(G,\theta _{g_0}^n\alpha `$, $`\theta _{g_0}^n\beta `$, $`\theta _{g_0}^n\gamma `$, $`\theta _{g_0}^n\delta )`$. Let $`f`$ be the bijection of $`GGu`$ defined by $`gt^1g`$, $`gu(t^1g)u`$, for $`gG`$. Then for $`g`$, $`hG`$, we have
$`f(gh)=t^1\mathrm{\Delta }\alpha (g,h)=t\mathrm{\Delta }\alpha (t^1g,t^1h)=t^1gt^1h=f(g)f(h),`$
$`f(ghu)=t^1\mathrm{\Delta }\beta (g,h)u=t\mathrm{\Delta }\beta (t^1g,t^1h)u=t^1g(t^1h)u=f(g)f(hu),`$
and similarly for $`\gamma `$, $`\delta `$. Hence $`f`$ is the desired isomorphism. โ
Therefore, if we only count the quasigroups in $`๐ฌ(G,,g_0)`$ up to isomorphism, we can assume that $`๐ฌ(G,,g_0)=\{Q(G,,g_0,\alpha ,\beta ,\gamma ,\delta );\alpha \mathrm{\Theta }`$, and $`\beta `$, $`\gamma `$, $`\delta `$ are of the form $`\theta \theta _{g_0}^n`$ for some $`n`$ and $`\theta \mathrm{\Theta }\}`$.
Given a groupoid $`(A,)`$, the *opposite groupoid* $`(A,^{\mathrm{op}})`$ is defined by $`x^{\mathrm{op}}y=yx`$.
###### Lemma 3.2.
The quasigroups $`Q(G,\alpha ,\beta ,\gamma ,\delta )`$ and $`Q(G,\theta _{yx}\alpha `$, $`\theta _{yx}\gamma `$, $`\theta _{yx}\beta `$, $`\theta _{yx}\delta )`$ are opposite to each other.
###### Proof.
Let $``$ denote the multiplication in $`Q(G,\alpha ,\beta ,\gamma ,\delta )`$, and $``$ the multiplication in $`Q(G,\theta _{yx}\alpha `$, $`\theta _{yx}\gamma `$, $`\theta _{yx}\beta `$, $`\theta _{yx}\delta \}`$. For $`g`$, $`hG`$ we have
$`gh=\mathrm{\Delta }\alpha (g,h)=\mathrm{\Delta }\theta _{yx}\alpha (h,g)=hg,`$
$`ghu=\mathrm{\Delta }\beta (g,h)u=\mathrm{\Delta }\theta _{yx}\beta (h,g)u=hug,`$
$`guh=\mathrm{\Delta }\gamma (g,h)u=\mathrm{\Delta }\theta _{yx}\gamma (h,g)u=hgu,`$
$`guhu=\mathrm{\Delta }\delta (g,h)=\mathrm{\Delta }\theta _{yx}\delta (h,g)=hugu.`$
Therefore, if we only count the quasigroups in $`๐ฌ(G,,g_0)`$ up to isomorphism and opposites, we can assume that $`๐ฌ(G,,g_0)=\{Q(G,,g_0,\alpha ,\beta ,\gamma ,\delta );\alpha \{\theta _{xy}`$, $`\theta _{xy^{}}`$, $`\theta _{x^{}y}`$, $`\theta _{x^{}y^{}}\}`$, and $`\beta `$, $`\gamma `$, $`\delta `$ are of the form $`\theta \theta _{g_0}^n`$ for some $`n`$ and $`\theta \mathrm{\Theta }\}`$.
###### Assumption 3.3.
From now on we assume that $`\alpha \{\theta _{xy}`$, $`\theta _{xy^{}}`$, $`\theta _{x^{}y}`$, $`\theta _{x^{}y^{}}\}`$, and that $`\beta `$, $`\gamma `$, $`\delta `$ are of the form $`\theta \theta _{g_0}^n`$ for some $`n`$ and $`\theta \mathrm{\Theta }`$.
## 4. When $``$ is identical on $`G`$
Assume for a while that $`g=g^{}`$ for every $`gG`$. Then $`gh=(gh)^{}=h^{}g^{}=hg`$ shows that $`G`$ is commutative. In particular, $`\mathrm{\Theta }=\{\theta _{xy}\}`$, and $`\mathrm{\Theta }_0=_n\theta _{g_0}^n`$. We show in this section that loops $`Q(G,,g_0,\alpha ,\beta ,\gamma ,\delta )`$ obtained with identical $``$ are not interesting.
Let $`\psi `$ be a groupoid identity, and let $`\mathrm{var}\psi `$ be all the variables appearing in $`\psi `$. Assume that for every $`x\mathrm{var}\psi `$ a decision has been made whether $`x`$ is to be taken from $`G`$ or from $`Gu`$. Then, while evaluating each side of the identity $`\psi `$ in $`GGu`$, we have to use the multiplications $`\alpha `$, $`\beta `$, $`\gamma `$ and $`\delta `$ certain number of times.
###### Example 4.1.
Consider the left alternative law $`x(xy)=(xx)y`$. With $`xG`$, $`yGu`$, we see that we need $`\beta `$ twice to evaluate $`x(xy)`$, while we need $`\alpha `$ once and $`\beta `$ once to evaluate $`(xx)y`$.
A groupoid identity is said to be *strictly balanced* if the same variables appear on both sides of the identity the same number of times and in the same order. For instance $`(x(y(xz)))(yx)=((xy)x)(z(yx))`$ is strictly balanced.
The above example shows that the same multiplications do not have to be used the same number of times even while evaluating a strictly balanced identity. However:
###### Lemma 4.2.
Let $`\psi `$ be a strictly balanced identity. Assume that for $`x\mathrm{var}\psi `$ a decision has been made whether $`xG`$ or $`xGu`$. Then, while evaluating $`\psi `$ in $`Q(G,,g_0,\alpha ,\beta ,\gamma ,\delta )`$, $`\delta `$ is used the same number of times on both sides of $`\psi `$.
###### Proof.
Let $`k`$ be the number of variables on each side of $`\psi `$, with repetitions, whose value is assigned to be in $`Gu`$. The number $`k`$ is well-defined since $`\psi `$ is strictly balanced.
While evaluating the identity $`\psi `$, each multiplication reduces the number of factors by $`1`$. However, only $`\delta `$ reduces the number of factors from $`Gu`$ (by two). Since the coset multiplication in $`GGu`$ modulo $`G`$ is associative, and since $`\psi `$ is strictly balanced, either both evaluated sides of $`\psi `$ will end up in $`G`$ (in which case $`\delta `$ is applied $`k/2`$ times on each side), or both evaluated sides of $`\psi `$ will end up in $`Gu`$ (in which case $`\delta `$ is applied $`(k1)/2`$ times on each side). โ
###### Lemma 4.3.
If $`\alpha \mathrm{\Theta }`$ and $`L=Q(G,,g_0,\alpha ,\beta ,\gamma ,\delta )`$ is a loop, then the neutral element of $`Q`$ coincides with the neutral element of $`G`$.
###### Proof.
Let $`e`$ be the neutral element of $`L`$ and $`1`$ the neutral element of $`G`$. Since $`1=1^{}`$, we have $`11=\mathrm{\Delta }\alpha (1,1)=1=1e`$, and the result follows from the fact that $`L`$ is a quasigroup. โ
###### Proposition 4.4.
Assume that $`g^{}=g`$ for every $`gG`$, and let $`\alpha `$, $`\beta `$, $`\gamma `$, $`\delta \mathrm{\Theta }_0`$. If $`L=Q(G,,g_0,\alpha ,\beta ,\gamma ,\delta )`$ happens to be a loop, then every strictly balanced identity holds in $`L`$. In particular, $`L`$ is an abelian group.
###### Proof.
Since $``$ is identical on $`G`$, we have $`\mathrm{\Theta }_0=\{\theta _{g_0}^n;n\}`$. By Assumption 3.3, we have $`\alpha =\theta _{xy}`$. Then by Lemma 4.3, $`L`$ has neutral element $`1`$. Assume that $`\beta =\theta _{g_0}^n`$ for some $`n`$. Then $`gu=1gu=(\mathrm{\Delta }\beta (1,g))u=(g_0^ng)u`$, which means that $`n=0`$. Similarly, if $`\gamma =\theta _{g_0}^m`$ then $`m=0`$.
Let $`\delta =\theta _{g_0}^k`$. Let $`\psi `$ be a strictly balanced identity. For every $`x\mathrm{var}\psi `$, decide if $`xG`$ or $`xGu`$. By Lemma 4.2, while evaluating $`\psi `$ in $`L`$, the multiplication $`\delta `$ is used the same number of times on the left and on the right, say $`t`$ times. Since $`\alpha =\beta =\gamma =\theta _{xy}`$, we conclude that $`\psi `$ reduces to $`g_0^{kt}z=g_0^{kt}z`$, for some $`zGGu`$.
Since the associative law is strictly balanced, $`L`$ is associative. We have already noticed that identical $``$ forces $`G`$ to be abelian. Then $`L`$ is abelian too, as $`guh=(gh)u=(hg)u=hgu`$ and $`guhu=g_0^kgh=g_0^khg=hugu`$ for every $`g`$, $`hG`$. โ
We have just seen that if $`g=g^{}`$ for every $`gG`$ then our constructions do not yield nonassociative loops. Therefore:
###### Assumption 4.5.
From now on, we assume that there exists $`gG`$ such that $`g^{}g`$.
## 5. Loops
In this section we further narrow the choices of $`\alpha `$, $`\beta `$, $`\gamma `$, $`\delta `$ when $`Q(G,\alpha ,\beta ,\gamma ,\delta )`$ is supposed to be a loop.
###### Proposition 5.1.
Let $`L=Q(G,,g_0,\alpha ,\beta ,\gamma ,\delta )`$. Then $`L`$ is a loop if and only if $`\alpha =\theta _{xy}`$, $`\beta \{\theta _{xy}`$, $`\theta _{x^{}y}`$, $`\theta _{yx}`$, $`\theta _{yx^{}}\}`$, $`\gamma \{\theta _{xy}`$, $`\theta _{xy^{}}`$, $`\theta _{yx}`$, $`\theta _{y^{}x}\}`$, and $`\delta `$ is of the form $`\theta \theta _{g_0}^n`$ for some integer $`n`$ and $`g_0G`$.
###### Proof.
If $`L`$ is a loop then $`\alpha \{\theta _{xy}`$, $`\theta _{xy^{}}`$, $`\theta _{x^{}y}`$, $`\theta _{x^{}y^{}}\}`$ and Lemma 4.3 imply that $`1`$ is the neutral element of $`L`$.
The equation $`g=1g`$ holds for every $`gG`$ if and only if $`\mathrm{\Delta }\alpha (1,g)=g`$ for every $`gG`$, which happens if and only if $`\alpha \{\theta _{xy}`$, $`\theta _{x^{}y}\}`$. (Note that we use Assumption 4.5 here.) Similarly, $`g=g1`$ holds for every $`gG`$ if and only if $`\mathrm{\Delta }\alpha (g,1)=g`$ for every $`gG`$, which happens if and only if $`\alpha \{\theta _{xy}`$, $`\theta _{xy^{}}\}`$. Therefore $`g=1g=g1`$ holds for every $`gG`$ if and only if $`\alpha =\theta _{xy}`$.
Now, $`gu=1gu`$ holds for every $`gG`$ if and only if $`\mathrm{\Delta }\beta (1,g)=g`$ for every $`gG`$, which happens if and only if $`\beta \{\theta _{xy}`$, $`\theta _{x^{}y}`$, $`\theta _{yx}`$, $`\theta _{yx^{}}\}`$. Similarly, $`gu=gu1`$ holds for every $`gG`$ if and only if $`\mathrm{\Delta }\gamma (g,1)=g`$ for every $`gG`$, which happens if and only if $`\gamma \{\theta _{xy}`$, $`\theta _{xy^{}}`$, $`\theta _{yx}`$, $`\theta _{y^{}x}\}`$. โ
We are only interested in loops, and we have already noted that $`(g_0^n)^{}=g_0^nZ(G)`$. Since we allow $`g_0=1`$, we can agree on:
###### Assumption 5.2.
From now on, we assume that $`\alpha =\theta _{xy}`$, $`\beta \{\theta _{xy}`$, $`\theta _{x^{}y}`$, $`\theta _{yx}`$, $`\theta _{yx^{}}\}`$, $`\gamma \{\theta _{xy}`$, $`\theta _{xy^{}}`$, $`\theta _{yx}`$, $`\theta _{y^{}x}\}`$, and $`\delta \theta _{g_0}\mathrm{\Theta }`$.
Our last reduction concerns the maps $`\beta `$ and $`\gamma `$.
###### Lemma 5.3.
We have $`\mathrm{\Delta }\theta _{x^{}y^{}}\theta _0=\mathrm{\Delta }\theta _0\theta _{x^{}y^{}}`$ for every $`\theta _0\mathrm{\Theta }_0`$.
###### Proof.
The group $`\mathrm{\Theta }_0`$ is generated by $`\theta _{yx}`$, $`\theta _{xy^{}}`$ and $`\theta _{g_0}`$. It therefore suffices to check that $`\mathrm{\Delta }\theta _{x^{}y^{}}\theta _0=\mathrm{\Delta }\theta _0\theta _{x^{}y^{}}`$ holds for $`\theta _0\{\theta _{yx}`$, $`\theta _{xy^{}}`$, $`\theta _{g_0}\}`$, which follows by straightforward calculation. โ
###### Lemma 5.4.
The quasigroups $`Q(G,,g_0,\alpha ,\beta ,\gamma ,\delta )`$, $`Q(G,,g_0,\alpha ,\beta ^{},\gamma ^{},\theta _{x^{}y^{}}\delta )`$ are isomorphic if
$`(\beta ,\beta ^{})`$ $``$ $`\{(\theta _{xy},\theta _{yx^{}}),(\theta _{yx},\theta _{x^{}y}),(\theta _{x^{}y},\theta _{yx}),(\theta _{yx^{}},\theta _{xy})\},`$
$`(\gamma ,\gamma ^{})`$ $``$ $`\{(\theta _{xy},\theta _{y^{}x}),(\theta _{yx},\theta _{xy^{}}),(\theta _{xy^{}},\theta _{yx}),(\theta _{y^{}x},\theta _{xy})\}.`$
###### Proof.
Let $``$ denote the multiplication in $`Q(G,,g_0,\alpha ,\beta ,\gamma ,\delta )`$, and $``$ the multiplication in $`Q(G,,g_0,\alpha ,\beta ^{},\gamma ^{},\theta _{x^{}y^{}}\delta )`$. Consider the permutation $`f`$ of $`G`$ defined by $`f(g)=g`$, $`f(gu)=g^{}u`$, for $`gG`$.
We show that $`f`$ is an isomorphism of $`(GGu,)`$ onto $`(GGu,)`$ if and only if
$$(\mathrm{\Delta }\beta (g,h))^{}=\mathrm{\Delta }\beta ^{}(g,h^{}),(\mathrm{\Delta }\gamma (g,h))^{}=\mathrm{\Delta }\gamma ^{}(g^{},h).$$
(3)
Once we establish this fact, the proof is finished by checking that the pairs $`(\beta ,\beta ^{})`$, $`(\gamma ,\gamma ^{})`$ in the statement of the Lemma satisfy $`(\text{3})`$.
Let $`g`$, $`hG`$. Then
$`f(gh)`$ $`=`$ $`f(\mathrm{\Delta }\alpha (g,h))=\mathrm{\Delta }\alpha (g,h),`$
$`f(ghu)`$ $`=`$ $`f(\mathrm{\Delta }\beta (g,h)u)=(\mathrm{\Delta }\beta (g,h))^{}u,`$
$`f(guh)`$ $`=`$ $`f(\mathrm{\Delta }\gamma (g,h)u)=(\mathrm{\Delta }\gamma (g,h))^{}u,`$
$`f(guhu)`$ $`=`$ $`f(\mathrm{\Delta }\delta (g,h))=\mathrm{\Delta }\delta (g,h),`$
while
$`f(g)f(h)`$ $`=`$ $`gh=\mathrm{\Delta }\alpha (g,h),`$
$`f(g)f(hu)`$ $`=`$ $`gh^{}u=\mathrm{\Delta }\beta ^{}(g,h^{})u,`$
$`f(gu)f(h)`$ $`=`$ $`g^{}uh=\mathrm{\Delta }\gamma ^{}(g^{},h)u,`$
$`f(gu)f(hu)`$ $`=`$ $`g^{}uh^{}u=\mathrm{\Delta }\theta _{g^{}h^{}}\delta (g^{},h^{}).`$
We see that $`f(gh)=f(g)f(h)`$ always holds. By Lemma 5.3, $`f(guhu)=f(gu)f(hu)`$ always holds. Finally, $`f(ghu)=f(g)f(hu)`$, $`f(guh)=f(gu)f(h)`$ hold if and only if $`(\beta ,\beta ^{})`$, $`(\gamma ,\gamma ^{})`$ satisfy $`(\text{3})`$. โ
Assume that $`Q(G,,g_0,\theta _{xy},\beta ,\gamma ,\delta )`$ is a loop (satisfying Assumption 5.2). Then Lemma 5.4 provides an isomorphism of $`Q(G,,g_0,\theta _{xy},\beta ,\gamma ,\delta )`$ onto some loop$`Q(G,,g_0,\theta _{xy},\beta ^{},\gamma ^{},\delta ^{})`$ such that if $`\gamma =\theta _{xy^{}}`$ then $`\gamma ^{}=\theta _{yx}`$, and if $`\gamma =\theta _{y^{}x}`$ then $`\gamma ^{}=\theta _{xy}`$. We can therefore assume:
###### Assumption 5.5.
From now on, we assume that $`\alpha =\theta _{xy}`$, $`\beta \{\theta _{xy}`$, $`\theta _{x^{}y}`$, $`\theta _{yx}`$, $`\theta _{yx^{}}\}`$, $`\gamma \{\theta _{xy}`$, $`\theta _{yx}\}`$, and $`\delta \theta _{g_0}\mathrm{\Theta }`$.
It is easy to see how much calculation is needed to find all loops $`Q(G,,g_0,\alpha ,\beta ,\gamma ,\delta )`$ that satisfy a given groupoid identity $`\psi `$. We have $`1428=64`$ choices for $`(\alpha ,\beta ,\gamma ,\delta )`$. (To appreciate the reductions, compare this with the unrestricted case $`\alpha `$, $`\beta `$, $`\gamma `$, $`\delta \mathrm{\Theta }_0`$.) Once $`(\alpha ,\beta ,\gamma ,\delta )`$ is chosen, we must verify $`2^k`$ equations in $`G`$, where $`k`$ is the number of variables in $`\psi `$ (since each variable can be assigned value in $`G`$ or in $`Gu`$).
We work out the calculation for one identity $`\psi `$ and one choice of multiplication $`(\alpha ,\beta ,\gamma ,\delta )`$. After seing the routine nature of the calculations, we gladly switch to a computer search.
## 6. C-loops arising from the construction of de Barros and Juriaans
C-loops are loops satisfying the identity $`((xy)y)z=x(y(yz))`$. In de Barros and Juriaans (1998), de Barros and Juriaans used a construction similar to (1) to obtain loops whose loop algebras are flexible. In our systematic notation, their construction is
$$Q(G,,g_0,\theta _{xy},\theta _{xy},\theta _{y^{}x},\theta _{g_0}\theta _{xy^{}}),$$
(4)
with the usual conventions on $`g_0`$ and $``$. The construction (4) violates Assumption 5.5 but, by Lemma 5.4, it is isomorphic to
$$Q(G,,g_0,\theta _{xy},\theta _{yx^{}},\theta _{xy},\theta _{g_0}\theta _{x^{}y}),$$
which complies with all assumptions we have made.
###### Theorem 6.1.
Let $`G`$ be a group and let $`L`$ be the loop defined by $`(\text{4})`$. Then $`L`$ is a flexible loop, and the following conditions are equivalent:
1. $`L`$ is associative,
2. $`L`$ is Moufang,
3. $`G`$ is commutative.
Furthermore, $`L`$ is a C-loop if and only if $`G/Z(G)`$ is an elementary abelian $`2`$-group. When $`L`$ is a C-loop, it is diassociative.
###### Proof.
Throughout the proof, we use $`g_0=g_0^{}Z(G)`$, $`gg^{}=g^{}gZ(G)`$, $`(g^{})^{}=g`$ and $`(gh)^{}=h^{}g^{}`$ without warning.
By Proposition 5.1, $`L`$ is a loop.
Flexibility. For $`x`$, $`yG`$ we have:
$`(xy)x=(xy)x=x(yx)=x(yx),`$
$`(xyu)x=(xy)ux=x^{}xyu=xx^{}yu=xx^{}yu=x(yux),`$
$`(xuy)xu=y^{}xuxu=g_0y^{}xx^{}=g_0xx^{}y^{}=xu(yx)u=xu(yxu),`$
$`(xuyu)xu=g_0xy^{}xu=g_0xy^{}xu=xug_0yx^{}=xu(yuxu).`$
Thus $`L`$ is flexible.
Associativity. For $`x`$, $`y`$, $`zG`$ we have:
$`x(yz)=x(yz)=(xy)z=(xy)z,`$
$`x(yzu)=x(yz)u=(xy)zu=(xy)zu,`$
$`xu(yz)=xuyz=z^{}y^{}xu=y^{}xuz=(xuy)z,`$
$`x(yuzu)=xg_0yz^{}=g_0xyz^{}=xyuzu=(xyu)zu,`$
$`xu(yuz)=xuz^{}yu=g_0xy^{}z=g_0xy^{}z=(xuyu)z.`$
Furthermore,
$`x(yuz)=xz^{}yu=xz^{}yu,(xyu)z=xyuz=z^{}xyu,`$
$`xu(yzu)=xuyzu=g_0xz^{}y^{},(xuy)zu=y^{}xuzu=g_0y^{}xz^{},`$
$`xu(yuzu)=xug_0yz^{}=g_0zy^{}xu,(xuyu)zu=g_0xy^{}zu=g_0xy^{}zu.`$
Thus $`L`$ is associative if and only if $`G`$ is commutative. (Sufficiency is obvious. For necessity, note that $``$ is onto, and substitute $`1`$ for one of $`x`$, $`y`$, $`z`$ if needed.)
Moufang property. Let $`x`$, $`y`$, $`zG`$. Then
$`x(yu(xz))=x(yuxz)=xz^{}x^{}yu=xz^{}x^{}yu,`$
$`((xyu)x)z=(xyux)z=x^{}xyuz=z^{}x^{}xyu.`$
Therefore, this particular form of the Moufang identity holds if and only if $`xz^{}x^{}=z^{}x^{}x`$. Now, given $`x`$, $`yG`$, there is $`zG`$ such that $`z^{}x^{}=y`$. Therefore $`xz^{}x^{}=z^{}x^{}x`$ holds in $`G`$ if and only if $`G`$ is commutative. However, when $`G`$ is commutative, then $`L`$ is associative, and we have proved the equivalence of (i), (ii), (iii).
C property. Let $`x`$, $`y`$, $`zG`$. Then
$`x(y(yz))=x(y(yz))=((xy)y)z=((xy)y)z,`$
$`x(y(yzu))=(x(y(yz))u=((xy)y)z)u=((xy)y)zu,`$
$`x(yu(yuz))=x(yuz^{}yu)=xg_0yy^{}z=g_0xyy^{}z=g_0xyy^{}z`$
$`=(xyuyu)z=((xyu)yu)z,`$
$`xu(y(yz))=xuyyz=z^{}y^{}y^{}xu=y^{}y^{}xuz=(y^{}xuy)z`$
$`=((xuy)y)z,`$
$`x(yu(yuzu))=x(yug_0yz^{})=xg_0zy^{}yu=g_0xzy^{}yu`$
$`=g_0xyy^{}zu=g_0xyy^{}zu=(xyuyu)zu=((xyu)yu)zu,`$
$`xu(yu(yuz))=xu(yuz^{}yu)=xug_0yy^{}z=g_0z^{}yy^{}xu=g_0z^{}xy^{}yu`$
$`=g_0xy^{}yuz=(g_0xy^{}yu)z=((xuyu)yu)z,`$
$`xu(yu(yuzu))=xu(yug_0yz^{})=xug_0zy^{}yu=g_0^2xy^{}yz^{}=g_0xy^{}yuzu`$
$`=(g_0xy^{}yu)zu=((xuyu)yu)zu.`$
While verifying the remaining form of the C identity, we obtain
$`xu(y(yzu))=xuyyzu=g_0xz^{}y^{}y^{},`$
$`((xuy)y)zu=(y^{}xuy)zu=y^{}y^{}xuzu=g_0y^{}y^{}xz^{}.`$
The identity therefore holds if and only if $`y^{}y^{}`$ commutes with all elements of $`G`$, which happens if and only if $`G/Z(G)`$ is an elementary abelian $`2`$-group.
Finally, by Lemma 4.4 of Phillips and Vojtฤchovskรฝ (2004), flexible C-loops are diassociative. โ
## 7. The Algorithm
### 7.1. Collecting Identities
Let $`G`$ be a group, $`\psi `$ a groupoid identity and $`(\alpha ,\beta ,\gamma ,\delta )`$ a multiplication. Then the following algorithm will output a set $`\mathrm{\Psi }`$ of group identities such that $`Q(G,,g_0,\alpha ,\beta ,\gamma ,\delta )`$ satisfies $`\psi `$ if and only if $`G`$ satisfies all identities of $`\mathrm{\Psi }`$:
1. Let $`f:\mathrm{var}\psi \{0,1\}`$ be a function that decides whether $`x\mathrm{var}\psi `$ is to be taken from $`G`$ or from $`Gu`$.
2. Upon assigning the variables of $`\psi `$ according to $`f`$, let $`\psi _f=(u,v)`$ be the identity $`\psi `$ evaluated in $`Q(G,,g_0,\alpha ,\beta ,\gamma ,\delta )`$.
3. Let $`\mathrm{\Psi }=\{\psi _f;f:\mathrm{var}\psi \{0,1\}\}`$.
This algorithm is straightforward but not very useful, since it typically outputs a large number of complicated group identities.
### 7.2. Understanding the identities in the Bol-Moufang case
We managed to decipher the meaning of $`\mathrm{\Psi }`$ for all multiplications $`(\alpha ,\beta ,\gamma ,\delta )`$ and for all identities of Bol-Moufang type by another algotihm. First, we reduced the identity $`\psi _f=(u,v)`$ to a canonical form as follows:
1. replace $`g_0^{}`$ by $`g_0`$,
2. move all $`g_0`$ to the very left,
3. replace $`x^{}x`$ by $`xx^{}`$,
4. move all substrings $`xx^{}`$ immediately to the right of the power $`g_0^m`$, and order the substrings $`xx^{}`$, $`yy^{}`$, $`\mathrm{}`$ lexicographically,
5. cancel as much as possible on the left and on the right of the resulting identity.
Then we used Lemmas 7.17.5 to understand what the canonical identities collected in $`\mathrm{\Psi }`$ say about the group $`G`$:
###### Lemma 7.1.
If an identity of $`\mathrm{\Psi }`$ reduces to $`x^{}=x`$ then it does not hold in any group.
###### Proof.
Since we assume that $``$ is not identical on $`G`$. โ
###### Lemma 7.2.
The following conditions are equivalent:
1. $`G/Z(G)`$ is an elementary abelian $`2`$-group,
2. $`xxy=yxx`$,
3. $`xyx^{}=x^{}yx`$.
###### Proof.
We have $`xyx^{}=x^{}yx`$ if and only if $`x^{}xyx^{}x=x^{}x^{}yxx`$. Since $`x^{}xZ(G)`$, the latter identity is equivalent to $`x^{}xx^{}xy=x^{}x^{}yxx`$. Since $`xx^{}=x^{}x`$, we can rewrite it equivalently as $`x^{}x^{}xxy=x^{}x^{}yxx`$, which is by cancellation equivalent to $`xxy=yxx`$. โ
###### Lemma 7.3.
The following conditions are equivalent:
1. $`G`$ is commutative,
2. $`xx^{}y=x^{}yx`$.
###### Proof.
If $`xx^{}y=x^{}yx`$ then $`x^{}xy=x^{}yx`$ and so $`xy=yx`$. โ
###### Lemma 7.4.
If $`\psi `$ is a strictly balanced identity that reduces to $`xy=yx`$ upon substituting $`1`$ for some of the variables of $`\psi `$, then $`\psi `$ is equivalent to commutativity.
###### Proof.
$`\psi `$ implies commutativity. Once commutativity holds, we can rearrange the variables of $`\psi `$ so that both sides of $`\psi `$ are the same, because $`\psi `$ is strictly balanced. โ
###### Lemma 7.5.
The following conditions are equivalent:
1. $`xxy=yx^{}x^{}`$ holds in $`G`$,
2. $`(xx)^{}=xx`$ and $`G/Z(G)`$ is an elementary abelian $`2`$-group.
###### Proof.
Condition (ii) clearly implies (i). If (i) holds, we have $`xx=x^{}x^{}`$ (with $`y=1`$) and so $`(xx)^{}=xx`$. Also $`xxy=yx^{}x^{}=yxx`$. โ
### 7.3. What the identities mean in the Bol-Moufang case
Lemmas 7.17.5 are carefully tailored to loops of Bol-Moufang type, and we discovered them upon studying the canonical identities $`\mathrm{\Psi }`$ obtained by the computer search.
It just so happens that every identity $`\psi _f`$ of $`\mathrm{\Psi }`$ is equivalent to a combination of the following properties of $`G`$:
1. No group satisfies $`\psi _f`$.
2. All groups satisfy $`\psi _f`$.
3. $`G`$ is commutative.
4. $`G/Z(G)`$ is an elementary abelian $`2`$-group.
5. $`(gg)^{}=gg`$ for every $`gG`$.
A prominent example of $``$ is the inverse operation <sup>-1</sup> in $`G`$. Then (PB) says that $`G`$ is of exponent $`4`$, and it is therefore not difficult to obtain examples of groups satisfying any possible combination of (PN), (PA), (PC), (PB) and (PS).
We have implemented the algorithm in GAP GAP (1999), and made it available online at
http://www.math.du.edu/~petr
in section Research. The algorithm is not safe for identities that are not strictly balanced.
## 8. Results
We now present the results of the computer search. In order to organize the results, observe that if $`L=Q(G,,g_0,\alpha ,\beta ,\gamma ,\delta )`$ is associative, it satisfies all identities of Bol-Moufang type. Since we do not want to list the multiplications and properties of $`G`$ repeatedly, we first describe all cases when $`L`$ is associative, then all cases when $`L`$ is an extra loop, then all cases when $`L`$ is a Moufang loop, etc., guided by the inclusions of Figure 1.
All results of this section are computer generated. To avoid errors in transcribing, the source of the statements of the results is also computer generated. In the statements, we write $`xy`$ instead of $`\theta _{xy}`$, $`g_0yx^{}`$ instead of $`\theta _{g_0}\theta _{yx^{}}`$, etc., in order to save space and improve legibility. Some results are mirror versions of others (cf. Theorem 8.5 versus Theorem 8.6), but we decided to include them anyway for quicker future reference. Finally, when $`G`$ is commutative, $`\mathrm{\Delta }(\mathrm{\Theta }\theta _{g_0}\mathrm{\Theta })`$ coincides with $`\mathrm{\Delta }(S\theta _{g_0}S)`$, where $`S=\{\theta _{xy}`$, $`\theta _{xy^{}}`$, $`\theta _{x^{}y}`$, $`\theta _{x^{}y^{}}\}`$. We therefore report only maps $`\alpha `$, $`\beta `$, $`\gamma `$, $`\delta `$ from $`S\theta _{g_0}S`$ in the commutative case.
In Theorems 8.18.14, $`G`$ is a group, $``$ is a nonidentical involutory antiautomorphism of $`G`$ satisfying $`gg^{}Z(G)`$ for every $`gG`$, the element $`g_0Z(G)`$ satisfies $`g_0^{}=g_0`$, and the maps $`\alpha `$, $`\beta `$, $`\gamma `$, $`\delta `$ are as in Assumption 5.5.
###### Theorem 8.1.
The loop $`Q(G,,g_0,\theta _{xy},\beta ,\gamma ,\delta )`$ is associative iff the following conditions are satisfied: $`(\beta ,\gamma ,\delta )`$ is equal to
$`(xy,xy,g_0xy)`$, or
$`G`$ is commutative and $`(\beta ,\gamma ,\delta )`$ is equal to $`(x^{}y,xy,g_0x^{}y)`$.
###### Theorem 8.2.
The loop $`Q(G,,g_0,\theta _{xy},\beta ,\gamma ,\delta )`$ is extra iff it is associative or if the following conditions are satisfied:
$`G/Z(G)`$ is an elementary abelian $`2`$-group and $`(\beta ,\gamma ,\delta )`$ is equal to
$`(x^{}y,yx,g_0yx^{})`$.
###### Theorem 8.3.
The loop $`Q(G,,g_0,\theta _{xy},\beta ,\gamma ,\delta )`$ is Moufang iff it is extra or if the following conditions are satisfied: $`(\beta ,\gamma ,\delta )`$ is equal to
$`(x^{}y,yx,g_0yx^{})`$.
###### Theorem 8.4.
The loop $`Q(G,,g_0,\theta _{xy},\beta ,\gamma ,\delta )`$ is a C-loop iff it is extra or if the following conditions are satisfied:
$`G/Z(G)`$ is an elementary abelian $`2`$-group and $`(\beta ,\gamma ,\delta )`$ is among
$`(yx,yx,g_0yx)`$, $`(yx^{},xy,g_0x^{}y)`$.
###### Theorem 8.5.
The loop $`Q(G,,g_0,\theta _{xy},\beta ,\gamma ,\delta )`$ is left Bol iff it is Moufang or if the following conditions are satisfied:
$`G/Z(G)`$ is an elementary abelian $`2`$-group and $`(\beta ,\gamma ,\delta )`$ is among
$`(xy,yx,g_0yx)`$, $`(x^{}y,xy,g_0x^{}y)`$, or
$`G`$ is commutative, $`(xx)^{}=xx`$ for every $`xG`$ and $`(\beta ,\gamma ,\delta )`$ is among
$`(xy,xy,g_0x^{}y)`$, $`(x^{}y,xy,g_0xy)`$, or
$`G/Z(G)`$ is an elementary abelian $`2`$-group, $`(xx)^{}=xx`$ for every $`xG`$ and $`(\beta ,\gamma ,\delta )`$ is among
$`(xy,xy,g_0x^{}y)`$, $`(xy,yx,g_0yx^{})`$, $`(x^{}y,xy,g_0xy)`$, $`(x^{}y,yx,g_0yx)`$.
###### Theorem 8.6.
The loop $`Q(G,,g_0,\theta _{xy},\beta ,\gamma ,\delta )`$ is right Bol iff it is Moufang or if the following conditions are satisfied:
$`G/Z(G)`$ is an elementary abelian $`2`$-group and $`(\beta ,\gamma ,\delta )`$ is among
$`(yx,xy,g_0yx)`$, $`(yx^{},yx,g_0x^{}y)`$, or
$`G`$ is commutative, $`(xx)^{}=xx`$ for every $`xG`$ and $`(\beta ,\gamma ,\delta )`$ is among
$`(xy,xy,g_0xy^{})`$, $`(x^{}y,xy,g_0x^{}y^{})`$, or
$`G/Z(G)`$ is an elementary abelian $`2`$-group, $`(xx)^{}=xx`$ for every $`xG`$ and $`(\beta ,\gamma ,\delta )`$ is among
$`(xy,xy,g_0xy^{})`$, $`(x^{}y,yx,g_0y^{}x^{})`$, $`(yx,xy,g_0y^{}x)`$, $`(yx^{},yx,g_0x^{}y^{})`$.
###### Theorem 8.7.
The loop $`Q(G,,g_0,\theta _{xy},\beta ,\gamma ,\delta )`$ is an LC-loop iff it is a C-loop or if the following conditions are satisfied:
$`G/Z(G)`$ is an elementary abelian $`2`$-group and $`(\beta ,\gamma ,\delta )`$ is among
$`(xy,yx,g_0yx)`$, $`(x^{}y,xy,g_0x^{}y)`$, $`(yx,xy,g_0xy)`$, $`(yx^{},yx,g_0yx^{})`$, or
$`G`$ is commutative, $`(xx)^{}=xx`$ for every $`xG`$ and $`(\beta ,\gamma ,\delta )`$ is among
$`(xy,xy,g_0x^{}y)`$, $`(x^{}y,xy,g_0xy)`$, or
$`G/Z(G)`$ is an elementary abelian $`2`$-group, $`(xx)^{}=xx`$ for every $`xG`$ and $`(\beta ,\gamma ,\delta )`$ is among
$`(xy,xy,g_0x^{}y)`$, $`(xy,yx,g_0yx^{})`$, $`(x^{}y,xy,g_0xy)`$, $`(x^{}y,yx,g_0yx)`$,
$`(yx,xy,g_0x^{}y)`$, $`(yx,yx,g_0yx^{})`$, $`(yx^{},xy,g_0xy)`$, $`(yx^{},yx,g_0yx)`$.
###### Theorem 8.8.
The loop $`Q(G,,g_0,\theta _{xy},\beta ,\gamma ,\delta )`$ is an RC-loop iff it is a C-loop or if the following conditions are satisfied:
$`G/Z(G)`$ is an elementary abelian $`2`$-group and $`(\beta ,\gamma ,\delta )`$ is among
$`(xy,yx,g_0xy)`$, $`(x^{}y,xy,g_0yx^{})`$, $`(yx,xy,g_0yx)`$, $`(yx^{},yx,g_0x^{}y)`$, or
$`G`$ is commutative, $`(xx)^{}=xx`$ for every $`xG`$ and $`(\beta ,\gamma ,\delta )`$ is among
$`(xy,xy,g_0xy^{})`$, $`(x^{}y,xy,g_0x^{}y^{})`$, or
$`G/Z(G)`$ is an elementary abelian $`2`$-group, $`(xx)^{}=xx`$ for every $`xG`$ and $`(\beta ,\gamma ,\delta )`$ is among
$`(xy,xy,g_0xy^{})`$, $`(xy,yx,g_0xy^{})`$, $`(x^{}y,xy,g_0y^{}x^{})`$, $`(x^{}y,yx,g_0y^{}x^{})`$,
$`(yx,xy,g_0y^{}x)`$, $`(yx,yx,g_0y^{}x)`$, $`(yx^{},xy,g_0x^{}y^{})`$, $`(yx^{},yx,g_0x^{}y^{})`$.
###### Theorem 8.9.
The loop $`Q(G,,g_0,\theta _{xy},\beta ,\gamma ,\delta )`$ is flexible iff it is Moufang or if the following conditions are satisfied: $`(\beta ,\gamma ,\delta )`$ is among
$`(xy,xy,g_0y^{}x^{})`$, $`(x^{}y,yx,g_0xy^{})`$, $`(x^{}y,yx,g_0x^{}y)`$, $`(x^{}y,yx,g_0y^{}x)`$,
$`(yx,yx,g_0x^{}y^{})`$, $`(yx,yx,g_0yx)`$, $`(yx^{},xy,g_0xy^{})`$, $`(yx^{},xy,g_0x^{}y)`$,
$`(yx^{},xy,g_0yx^{})`$, $`(yx^{},xy,g_0y^{}x)`$, or
$`G/Z(G)`$ is an elementary abelian $`2`$-group and $`(\beta ,\gamma ,\delta )`$ is among
$`(xy,xy,g_0x^{}y^{})`$, $`(xy,xy,g_0yx)`$, $`(yx,yx,g_0xy)`$, $`(yx,yx,g_0y^{}x^{})`$.
###### Theorem 8.10.
The loop $`Q(G,,g_0,\theta _{xy},\beta ,\gamma ,\delta )`$ is left alternative iff it is left Bol or an LC-loop or if the following conditions are satisfied: $`(\beta ,\gamma ,\delta )`$ is among
$`(xy,xy,g_0x^{}y)`$, $`(xy,yx,g_0yx^{})`$, $`(x^{}y,xy,g_0x^{}y)`$, $`(yx,xy,g_0x^{}y)`$,
$`(yx,yx,g_0yx^{})`$, $`(yx^{},xy,g_0x^{}y)`$, $`(yx^{},yx,g_0yx^{})`$, or
$`(xx)^{}=xx`$ for every $`xG`$ and $`(\beta ,\gamma ,\delta )`$ is equal to
$`(x^{}y,xy,g_0xy)`$.
###### Theorem 8.11.
The loop $`Q(G,,g_0,\theta _{xy},\beta ,\gamma ,\delta )`$ is right alternative iff it is right Bol or an RC-loop or if the following conditions are satisfied: $`(\beta ,\gamma ,\delta )`$ is among
$`(xy,xy,g_0xy^{})`$, $`(xy,yx,g_0xy^{})`$, $`(x^{}y,xy,g_0yx^{})`$, $`(yx,xy,g_0y^{}x)`$,
$`(yx,yx,g_0y^{}x)`$, $`(yx^{},xy,g_0x^{}y)`$, $`(yx^{},yx,g_0x^{}y)`$, or
$`(xx)^{}=xx`$ for every $`xG`$ and $`(\beta ,\gamma ,\delta )`$ is equal to
$`(yx^{},yx,g_0x^{}y^{})`$.
###### Theorem 8.12.
The loop $`Q(G,,g_0,\theta _{xy},\beta ,\gamma ,\delta )`$ is a left nuclear square loop iff it is an LC-loop or if the following conditions are satisfied: $`(\beta ,\gamma ,\delta )`$ is among
$`(xy,xy,g_0xy^{})`$, $`(yx^{},yx,g_0x^{}y)`$, $`(yx^{},yx,g_0x^{}y^{})`$, or
$`G/Z(G)`$ is an elementary abelian $`2`$-group and $`(\beta ,\gamma ,\delta )`$ is among
$`(xy,xy,g_0yx)`$, $`(xy,xy,g_0y^{}x)`$, $`(xy,yx,g_0xy)`$, $`(xy,yx,g_0xy^{})`$,
$`(xy,yx,g_0y^{}x)`$, $`(x^{}y,xy,g_0x^{}y^{})`$, $`(x^{}y,xy,g_0yx^{})`$, $`(x^{}y,xy,g_0y^{}x^{})`$,
$`(x^{}y,yx,g_0x^{}y)`$, $`(x^{}y,yx,g_0x^{}y^{})`$, $`(x^{}y,yx,g_0y^{}x^{})`$, $`(yx,xy,g_0xy^{})`$,
$`(yx,xy,g_0yx)`$, $`(yx,xy,g_0y^{}x)`$, $`(yx,yx,g_0xy)`$, $`(yx,yx,g_0xy^{})`$,
$`(yx,yx,g_0y^{}x)`$, $`(yx^{},xy,g_0x^{}y^{})`$, $`(yx^{},xy,g_0yx^{})`$, $`(yx^{},xy,g_0y^{}x^{})`$,
$`(yx^{},yx,g_0y^{}x^{})`$, or
$`G/Z(G)`$ is an elementary abelian $`2`$-group, $`(xx)^{}=xx`$ for every $`xG`$ and $`(\beta ,\gamma ,\delta )`$ is among
$`(xy,xy,g_0x^{}y^{})`$, $`(xy,xy,g_0yx^{})`$, $`(xy,xy,g_0y^{}x^{})`$, $`(xy,yx,g_0x^{}y)`$,
$`(xy,yx,g_0x^{}y^{})`$, $`(xy,yx,g_0y^{}x^{})`$, $`(x^{}y,xy,g_0xy^{})`$, $`(x^{}y,xy,g_0yx)`$,
$`(x^{}y,xy,g_0y^{}x)`$, $`(x^{}y,yx,g_0xy)`$, $`(x^{}y,yx,g_0xy^{})`$, $`(x^{}y,yx,g_0y^{}x)`$,
$`(yx,xy,g_0x^{}y^{})`$, $`(yx,xy,g_0yx^{})`$, $`(yx,xy,g_0y^{}x^{})`$, $`(yx,yx,g_0x^{}y)`$,
$`(yx,yx,g_0x^{}y^{})`$, $`(yx,yx,g_0y^{}x^{})`$, $`(yx^{},xy,g_0xy^{})`$, $`(yx^{},xy,g_0yx)`$,
$`(yx^{},xy,g_0y^{}x)`$, $`(yx^{},yx,g_0xy)`$, $`(yx^{},yx,g_0xy^{})`$, $`(yx^{},yx,g_0y^{}x)`$.
###### Theorem 8.13.
The loop $`Q(G,,g_0,\theta _{xy},\beta ,\gamma ,\delta )`$ is a middle nuclear square loop iff it is an LC-loop or an RC-loop or if the following conditions are satisfied: $`(\beta ,\gamma ,\delta )`$ is among
$`(xy,xy,g_0y^{}x^{})`$, $`(yx^{},xy,g_0xy^{})`$, $`(yx^{},xy,g_0yx^{})`$, or
$`G/Z(G)`$ is an elementary abelian $`2`$-group and $`(\beta ,\gamma ,\delta )`$ is among
$`(xy,xy,g_0x^{}y^{})`$, $`(xy,xy,g_0yx)`$, $`(xy,yx,g_0x^{}y^{})`$, $`(xy,yx,g_0y^{}x^{})`$,
$`(x^{}y,xy,g_0xy^{})`$, $`(x^{}y,xy,g_0y^{}x)`$, $`(x^{}y,yx,g_0xy^{})`$, $`(x^{}y,yx,g_0x^{}y)`$,
$`(x^{}y,yx,g_0y^{}x)`$, $`(yx,xy,g_0x^{}y^{})`$, $`(yx,xy,g_0y^{}x^{})`$, $`(yx,yx,g_0xy)`$,
$`(yx,yx,g_0x^{}y^{})`$, $`(yx,yx,g_0y^{}x^{})`$, $`(yx^{},xy,g_0y^{}x)`$, $`(yx^{},yx,g_0xy^{})`$,
$`(yx^{},yx,g_0y^{}x)`$, or
$`G/Z(G)`$ is an elementary abelian $`2`$-group, $`(xx)^{}=xx`$ for every $`xG`$ and $`(\beta ,\gamma ,\delta )`$ is among
$`(xy,xy,g_0yx^{})`$, $`(xy,xy,g_0y^{}x)`$, $`(xy,yx,g_0x^{}y)`$, $`(xy,yx,g_0y^{}x)`$,
$`(x^{}y,xy,g_0x^{}y^{})`$, $`(x^{}y,xy,g_0yx)`$, $`(x^{}y,yx,g_0xy)`$, $`(x^{}y,yx,g_0x^{}y^{})`$,
$`(yx,xy,g_0xy^{})`$, $`(yx,xy,g_0yx^{})`$, $`(yx,yx,g_0xy^{})`$, $`(yx,yx,g_0x^{}y)`$,
$`(yx^{},xy,g_0yx)`$, $`(yx^{},xy,g_0y^{}x^{})`$, $`(yx^{},yx,g_0xy)`$, $`(yx^{},yx,g_0y^{}x^{})`$.
###### Theorem 8.14.
The loop $`Q(G,,g_0,\theta _{xy},\beta ,\gamma ,\delta )`$ is a right nuclear square loop iff it is an RC-loop or if the following conditions are satisfied: $`(\beta ,\gamma ,\delta )`$ is among
$`(xy,xy,g_0x^{}y)`$, $`(x^{}y,xy,g_0xy)`$, $`(x^{}y,xy,g_0x^{}y)`$, or
$`G/Z(G)`$ is an elementary abelian $`2`$-group and $`(\beta ,\gamma ,\delta )`$ is among
$`(xy,xy,g_0yx)`$, $`(xy,xy,g_0yx^{})`$, $`(xy,yx,g_0x^{}y)`$, $`(xy,yx,g_0yx)`$,
$`(xy,yx,g_0yx^{})`$, $`(x^{}y,xy,g_0yx)`$, $`(x^{}y,yx,g_0xy)`$, $`(x^{}y,yx,g_0x^{}y)`$,
$`(x^{}y,yx,g_0yx)`$, $`(yx,xy,g_0xy)`$, $`(yx,xy,g_0x^{}y)`$, $`(yx,xy,g_0yx^{})`$,
$`(yx,yx,g_0xy)`$, $`(yx,yx,g_0x^{}y)`$, $`(yx,yx,g_0yx^{})`$, $`(yx^{},xy,g_0xy)`$,
$`(yx^{},xy,g_0yx)`$, $`(yx^{},xy,g_0yx^{})`$, $`(yx^{},yx,g_0xy)`$, $`(yx^{},yx,g_0yx)`$,
$`(yx^{},yx,g_0yx^{})`$, or
$`G/Z(G)`$ is an elementary abelian $`2`$-group, $`(xx)^{}=xx`$ for every $`xG`$ and $`(\beta ,\gamma ,\delta )`$ is among
$`(xy,xy,g_0x^{}y^{})`$, $`(xy,xy,g_0y^{}x)`$, $`(xy,xy,g_0y^{}x^{})`$, $`(xy,yx,g_0x^{}y^{})`$,
$`(xy,yx,g_0y^{}x)`$, $`(xy,yx,g_0y^{}x^{})`$, $`(x^{}y,xy,g_0xy^{})`$, $`(x^{}y,xy,g_0x^{}y^{})`$,
$`(x^{}y,xy,g_0y^{}x)`$, $`(x^{}y,yx,g_0xy^{})`$, $`(x^{}y,yx,g_0x^{}y^{})`$, $`(x^{}y,yx,g_0y^{}x)`$,
$`(yx,xy,g_0xy^{})`$, $`(yx,xy,g_0x^{}y^{})`$, $`(yx,xy,g_0y^{}x^{})`$, $`(yx,yx,g_0xy^{})`$,
$`(yx,yx,g_0x^{}y^{})`$, $`(yx,yx,g_0y^{}x^{})`$, $`(yx^{},xy,g_0xy^{})`$, $`(yx^{},xy,g_0y^{}x)`$,
$`(yx^{},xy,g_0y^{}x^{})`$, $`(yx^{},yx,g_0xy^{})`$, $`(yx^{},yx,g_0y^{}x)`$, $`(yx^{},yx,g_0y^{}x^{})`$.
## 9. Concluding remarks
(I) Figure 1 and Theorems 8.18.14 taken together tell us more than if we consider them separately. For instance, Theorem 8.1 and Theorem 8.3 plus the fact that every group is a Moufang loop imply that the construction of Theorem 8.3 yields a nonassociative loop if and only if the group $`G`$ is not commutative. In other words, the two theorems encompass Theorem 1.1, and, in addition, show that Cheinโs construction is unique for Moufang loops.
(II) Note that we have also recovered (an isomorphic copy of) the construction (4) of de Barros and Juriaans. Our results on Bol loops agree with those of Vojtฤchovskรฝ (2004), obtained by hand.
(III) To illustrate how the algorithm works for loops that are not of Bol-Moufang type, we show the output for nonassociative RIF loops. A loop is an *RIF loop* if it satisfies $`(xy)(z(xy))=((x(yz))x)y`$.
###### Theorem 9.1.
The loop $`Q(G,,g_0,\theta _{xy},\beta ,\gamma ,\delta )`$ is RIF iff it is associative or if the following conditions are satisfied: $`(\beta ,\gamma ,\delta )`$ is among
$`(x^{}y,yx,g_0yx^{})`$, $`(yx^{},xy,g_0x^{}y)`$, or
$`(\beta ,\gamma ,\delta )`$ and $`G`$ are as in the following list:
$`(yx,yx,g_0yx)`$ and $`xyzxy=yxzyx`$.
Note that the algorithm did not manage to decipher the meaning of the group identity $`xyzxy=yxzyx`$, so it simply listed it.
(IV) We conclude the paper with the following observation:
###### Lemma 9.2.
Let $`L=Q(G,,g_0,\alpha ,\beta ,\gamma ,\delta )`$ be a loop. Then $`L`$ has two-sided inverses.
###### Proof.
Let $`gG`$. Since $`g^{}(g^1)^{}=(g^1g)^{}=1^{}=1`$, we have $`(g^{})^1=(g^1)^{}`$, and the antiautomorphisms <sup>-1</sup> and $``$ commute. Let us denote $`(g^1)^{}=(g^{})^1`$ by $`g^{}`$.
We show that for every $`\alpha \mathrm{\Theta }_0`$ and $`gG`$, there is $`hG`$ such that $`\mathrm{\Delta }\alpha (g,h)=gh=1=hg=\mathrm{\Delta }\alpha (h,g)`$. The proof for $`guGu`$ is similar.
Assume that $`\alpha \{\theta _{xy}`$, $`\theta _{xy^{}}`$, $`\theta _{x^{}y}`$, $`\theta _{x^{}y^{}}\}`$. Then
$`\mathrm{\Delta }\theta _{xy}(g,g^1)=gg^1=1=g^1g=\mathrm{\Delta }\theta _{xy}(g^1,g),`$
$`\mathrm{\Delta }\theta _{xy^{}}(g,g^{})=g(g^{})^{}=1=g^{}g^{}=\mathrm{\Delta }\theta _{xy^{}}(g^{},g),`$
$`\mathrm{\Delta }\theta _{x^{}y}(g,g^{})=g^{}g^{}=1=(g^{})^{}g=\mathrm{\Delta }\theta _{x^{}y}(g^{},g),`$
$`\mathrm{\Delta }\theta _{x^{}y^{}}(g,g^1)=g^{}g^{}=1=g^{}g^{}=\mathrm{\Delta }\theta _{x^{}y^{}}(g^1,g)`$
show that the two-sided inverse $`h`$ exists. The case $`\alpha \{\theta _{yx}`$, $`\theta _{y^{}x}`$, $`\theta _{yx^{}}`$, $`\theta _{y^{}x^{}}\}`$ is similar. The general case $`\alpha \mathrm{\Theta }_0`$ then follows thanks to $`g_0=g_0^{}Z(G)`$. โ
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# Rotating Resonator-Oscillator Experiments to Test Lorentz Invariance in Electrodynamics
## 1 Introduction
The Einstein Equivalence Principle (EEP) is a founding principle of relativity Will . One of the constituent elements of EEP is Local Lorentz Invariance (LLI), which postulates that the outcome of a local experiment is independent of the velocity and orientation of the apparatus. The central importance of this postulate has motivated tremendous work to experimentally test LLI. Also, a number of unification theories suggest a violation of LLI at some level. However, to test for violations it is necessary to have an alternative theory to allow interpretation of experiments Will , and many have been developed Robertson ; MaS ; LightLee ; Ni ; Kosto1 ; KM . The kinematical frameworks (RMS) Robertson ; MaS postulate a simple parameterization of the Lorentz transformations with experiments setting limits on the deviation of those parameters from their values in special relativity (SR). Because of their simplicity they have been widely used to interpret many experiments Brillet ; Wolf ; Muller ; WolfGRG . More recently, a general Lorentz violating extension of the standard model of particle physics (SME) has been developed Kosto1 whose Lagrangian includes all parameterized Lorentz violating terms that can be formed from known fields.
This work analyses rotating laboratory Lorentz invariance experiments that compare precisely the resonant frequencies of two high-Q factor (or high finesse) cavity resonators. High stability electromagnetic oscillatory fields are generated by implementing state of the art frequency stabilization systems with the narrow line width of the resonators. Previous non-rotating experiments Lipa ; Muller ; Wolf04 relied on the rotation of the Earth to modulate putative Lorentz violating effects. This is not optimal for two reasons. Firstly, the sensitivity to Lorentz violations is proportional to the noise of the oscillators at the modulation frequency, typically best for periods between 10 and 100 seconds. Secondly, the sensitivity is proportional to the square root of the number of periods of the modulation signal, therefore taking a relatively long time to acquire sufficient data. Thus, by rotating the experiment the data integration rate is increased and the relevant signals are translated to the optimal operating regime Mike .
In this work we outline the two most commonly used test theories (RMS and SME) for testing LLI of the photon. Then we develop the general frame-work of applying these test theories to resonator experiments with an emphasis on rotating experiments in the laboratory. We compare the inherent sensitivity factors of common experiments and propose some new configurations. Finally we apply the test theories to the rotating cryogenic experiment at the University of Western Australia, which recently set new limits in both the RMS and SME frameworks PRL . Note added: Two other concurrent experiments have set some similar limits schil ; achim .
## 2 Common Test Theories to Characterize Lorentz Invariance
The most famous test of LLI (or the constancy of the speed of light) was that conducted by Michelson and Morley in 1887 MM with a rotating table and a Michelson interferometer. In actual fact, the theoretical framework used by Michelson and Morley was not a test of LLI, since the concept did not exist at the time, but that of an aether drift. The relative motion of the apparatus through the aether was thought to induce a phase difference between the arms of the interferometer (and hence an interference pattern) depending on the orientation. Thus, as the Earth moved from one end of its orbit to the opposite end, the change in its velocity should be a detectable value. Michelson and Morley found no fringe shifts due to Earth motion around the sun and reported a null result. Since the Michelson Morley experiment, there have been many other types of experiments devised to test the validity of SR and the constancy of light. However, to interpret these experiments one must formulate an alternative test theory, and in this section we outline two of the most commonly used.
### 2.1 Robertson, Mansouri, Sexl Framework
A simple kinematic test theory that has been widely used is that of Robertson, Mansouri and Sexl (RMS)Robertson ; MaS , where time standards (โclocksโ) and length standards (โrodsโ) are considered without taking into account their underlying structure. This framework postulates a preferred frame $`\mathrm{\Sigma }(T,๐)`$ which satisfies LLI, and a moving frame $`S(t,๐ฑ)`$, which does not. The prime candidate for the preferred frame is taken as the Cosmic Microwave Background (CMB), since any anticipated non-symmetries are expected to arise from Planck-scale effects during the creation of the universe. In this framework we analyse the Poynting vector direction of the electromagnetic signal with respect to the velocity of the lab through the CMB.
The normal Lorentz Transformations for a boost in the $`x`$ direction are expressed in a special form below (where $`c`$ is the speed of light in the $`\mathrm{\Sigma }`$ frame):
$$dT=\frac{1}{a}\left(dt+\frac{vdx}{c^2}\right);dX=\frac{dx}{b}+\frac{v}{a}\left(dt+\frac{vdx}{c^2}\right);dY=\frac{dy}{d};dZ=\frac{dz}{d};$$
(1)
Here we take a Taylor expansion for $`a`$, $`b`$ and $`d`$ of the form: $`a1+\alpha v^2/c^2+๐ช(c^4)`$; $`b1+\beta v^2/c^2+๐ช(c^4)`$; $`d1+\delta v^2/c^2+๐ช(c^4)`$. Recalling that $`\gamma =\frac{1}{\sqrt{1v^2/c^2}}`$ from Special Relativity (SR), we see that SR predicts $`\alpha =1/2`$ and $`\beta =1/2`$. Since SR predicts no contraction in directions orthogonal to a boost, it also predicts that $`\delta =0`$. Thus, the RMS parameterizes a possible Lorentz violation by a deviation of the parameters ($`\alpha ,\beta ,\delta `$) from the SR values ($`\frac{1}{2},\frac{1}{2},0`$).
By manipulating equation (1) to form the infinitesimals in the $`S`$ frame, we can separate the equation into a boost term $`(\beta \alpha 1)`$, anisotropy term $`(\delta \beta +\frac{1}{2})`$ and time dilation parameter $`\alpha +\frac{1}{2}`$. Thus, a complete verification of LLI in the RMS framework Robertson ; MaS requires a test of (i) the isotropy of the speed of light (measuring $`P_{MM}=\delta \beta +\frac{1}{2}`$), a Michelson-Morley (MM) experiment MM , (ii) the boost dependence of the speed of light (measuring $`P_{KT}=\beta \alpha 1`$), a Kennedy-Thorndike (KT) experiment KT and (iii) the time dilation parameter (measuring $`P_{IS}=\alpha +\frac{1}{2}`$), an Ives-Stillwell (IS) experiment IS ; Saat . Rotating experiments may be considered Michelson-Morley experiments and only measure $`P_{MM}`$, so in this section we restrict ourselves to these types of measurements.
Assuming only a MM type Lorentz violation, and setting $`ds^2=c^2dT^2dX^2dY^2dZ^2=0`$ in $`\mathrm{\Sigma }`$, and transforming according to Eqn. (1) we find the coordinate travel time of a light signal in $`S`$ becomes;
$$dt=\frac{dl}{c}\left(P_{MM}\times \mathrm{sin}^2\theta \frac{v^2}{c^2}\right)+๐ช(c^4)$$
(2)
where $`dl=\sqrt{dx^2+dy^2+dz^2}`$ and $`\theta `$ is the angle between the Poynting vector and the velocity v of $`S`$ in $`\mathrm{\Sigma }`$. For a modern MM experiment that measures the difference frequency between two resonant cavities, the fractional frequency difference may be calculated from (2) in a similar way to WolfGRG to give:
$$\frac{\mathrm{\Delta }\nu _0}{\nu _0}=\frac{P_{MM}}{2\pi c^2}[(๐ฏ.\widehat{๐}_๐(q_a))^2dq_a(๐ฏ.\widehat{๐}_๐(q_b))^2dq_b]$$
(3)
Where $`\widehat{๐}_j(q_j)`$ is the unit vector in the direction of light propagation (Poynting vector) of each resonator (labeled by subscripts $`a`$ and $`b`$), and $`q_j`$ is the variable of integration around the closed path coordinates of the Poynting vector of each resonator.
To calculate the relevant time dependent expressions for $`๐ฏ`$, velocities are transformed to a geocentric non-rotating (with respect to distant stars) reference frame (denoted as the MM-frame) centred at the centre of mass of the Earth with its $`z`$-axis perpendicular to the equator, pointing north, the x-axis pointing towards 11.2h right ascension (aligned with the equatorial projection of $`๐ฎ`$ defined below). A pictorial representation of the frame is shown in Fig. 1.
Classical (Galilean) transformations for the velocities and $`\mathrm{sin}^2\theta `$ are sufficient as relativistic terms are of order $`๐ช(c^2)`$ and therefore give rise in Eqn. (2) to terms of order $`๐ช(c^5)`$. We consider two velocities, the velocity of the sun with respect to the CMB $`๐ฎ`$ (declination -6.4, right ascension 11.2 h) and the orbital velocity of the Earth $`๐ฏ_o`$. Velocities due to the spinning of the Earth and laboratory are much smaller and do not impact on the calculations and may be ignored. Thus, the sum of the two provide the velocity of the laboratory in the universal frame to be inserted in Eqn. (3). In the MM-Earth frame, the CMB velocity is:
$`๐ฎ=u\left(\begin{array}{ccc}\mathrm{cos}\varphi _\mu & & \\ 0& & \\ \mathrm{sin}\varphi _\mu & & \end{array}\right)`$ (7)
where $`u377\mathrm{k}\mathrm{m}/\mathrm{s}`$ and $`\varphi _\mu 6.4^{}`$. To calculate the orbital velocity we first consider the Earth in a barycentric non-rotating frame (BRS) with the $`z`$-axis perpendicular to the Earthโs orbital plane and the $`x`$-axis pointing towards $`0^o`$ right ascension (pointing from the Sun to the Earth at the moment of the autumn equinox).
$`๐ฏ_o^{BRS}=v_o\left(\begin{array}{ccc}\mathrm{sin}\lambda _0& & \\ \mathrm{cos}\lambda _0& & \\ 0& & \end{array}\right)`$ (11)
where $`v_o29.78\mathrm{km}/\mathrm{s}`$ is the orbital speed of the Earth, and $`\lambda _0=\mathrm{\Omega }_{}(tt_o)`$ with $`\mathrm{\Omega }_{}\mathrm{2.0\; 10}^7\mathrm{rad}/\mathrm{s}`$ the angular orbital velocity and $`tt_0`$ the time since the autumnal equinox. We first transform to a geocentric frame (GRS) that has its $`x`$-axis aligned with the BRS one but its $`z`$-axis perpendicular to the equatorial plane of the Earth
$`๐ฏ_o^{GRS}=\left(\begin{array}{ccc}v_{}^{x}{}_{o}{}^{BRS}& & \\ v_{}^{y}{}_{o}{}^{BRS}\mathrm{cos}\epsilon v_{}^{z}{}_{o}{}^{BRS}\mathrm{sin}\epsilon & & \\ v_{}^{y}{}_{o}{}^{BRS}\mathrm{sin}\epsilon +v_{}^{z}{}_{o}{}^{BRS}\mathrm{cos}\epsilon & & \end{array}\right)`$ (15)
where $`\epsilon 23.27^o`$ is the angle between the equatorial and orbital planes of the Earth. We then transform to the MM-Earth frame:
$`๐ฏ_o=\left(\begin{array}{ccc}v_{}^{x}{}_{o}{}^{GRS}\mathrm{cos}\alpha _\mu +v_{}^{y}{}_{o}{}^{GRS}\mathrm{sin}\alpha _\mu & & \\ v_{}^{x}{}_{o}{}^{GRS}\mathrm{sin}\alpha _\mu +v_{}^{y}{}_{o}{}^{GRS}\mathrm{cos}\alpha _\mu & & \\ v_{}^{z}{}_{o}{}^{GRS}& & \end{array}\right)`$ (19)
where $`\alpha _\mu 167.9^o`$ is the right ascension of $`๐ฎ`$. Summing the two velocities from Eqns. (7) and (19) we obtain the velocity of the lab with respect to the โuniverse rest frameโ, transformed to the MM-Earth frame
$`๐ฏ=\left(\begin{array}{ccc}u\mathrm{cos}\varphi _\mu +v_o(\mathrm{sin}\lambda _0\mathrm{cos}\alpha _\mu +\mathrm{cos}\lambda _0\mathrm{sin}\alpha _\mu \mathrm{cos}\epsilon )& & \\ v_o(\mathrm{sin}\lambda _0\mathrm{sin}\alpha _\mu +\mathrm{cos}\lambda _0\mathrm{cos}\alpha _\mu \mathrm{cos}\epsilon )& & \\ u\mathrm{sin}\varphi _\mu +v_o\mathrm{cos}\lambda _0\mathrm{sin}\epsilon & & \end{array}\right).`$ (23)
Substituting in the numeric values gives an orbital velocity of (in $`m/s`$);
$`๐ฏ=\left(\begin{array}{ccc}374651+5735\mathrm{cos}(\lambda _0)+29118\mathrm{sin}(\lambda _0)& & \\ 26750\mathrm{cos}(\lambda _0)+6242\mathrm{sin}(\lambda _0)& & \\ 42024+11765\mathrm{cos}(\lambda _0)& & \end{array}\right).`$ (27)
The last calculation to make is the time dependence of the the unit vector $`\widehat{๐ฅ}`$ along the direction of light propagation, which will depend on the configuration of the experiment, including the type of resonator and whether it is rotating in the laboratory or not. In section 5 we calculate this dependence for a specific experiment, which uses WG modes rotating in the laboratory.
### 2.2 Standard Model Extension
The SME Kosto1 conglomerates all possible Lorentz-Violating terms and incorporates them in a framework, which is an extension of the Standard Model of Particle Physics. There are numerous Lorentz-violating terms per particle sector (i.e. fermions, bosons and photons). However in this work we are restricted to the so called minimal โphoton-sectorโ, which only includes 19 terms. The SME adds additional terms to the Lagrangian of the Standard Model for photons. Where as the standard Lagrangian was simply:
$$=\frac{1}{4}F^{\mu \nu }F_{\mu \nu }$$
(28)
Under the SME, it becomes KM :
$$=\frac{1}{4}F^{\mu \nu }F_{\mu \nu }\frac{1}{4}(k_F)_{\kappa \lambda \mu \nu }F^{\kappa \lambda }F^{\mu \nu }+\frac{1}{2}(k_{AF})^\kappa ฯต_{\kappa \lambda \mu \nu }A^\lambda F^{\mu \nu }$$
(29)
where $`A^\lambda `$ is the 4-potential. The $`(k_{AF})^\kappa `$ terms have the dimensions of mass, and are the CPT odd terms cptodd . It is argued in Kosto1 that these should be zero because they induce instabilities as they are non-negative in the Lagrangian. There are also astronomical measurements KM which place stringent limits on $`k_{AF}`$. From here on these terms are set to zero.
On the other hand, the $`(k_F)_{\kappa \lambda \mu \nu }`$ terms are CPT even, dimensionless and have 19 independent terms out of the 256 possible combinations of $`\kappa `$, $`\lambda `$, $`\mu `$ and $`\nu `$. Out of these independent Lorentz violating terms, 10 combinations have been analysed using astrophysical polarisation tests and have an upper-limit of $`2\times 10^{32}`$ KM . This limit is many orders of magnitude less than what is expected from laboratory experiments, so these terms are set to zero to simplify the calculations and to remain consistent with previous results.
We can derive the equations of motion for this system by minimising the action given by (29), using variational techniques and the definition $`F^{\mu \nu }_\mu A_\nu _\nu A_\mu `$ and $`A_\mu (\varphi ,๐)`$. These equations are similar to those of a Maxwellian model in anisotropic media instead of a vacuum. In order to express these in a convenient form, we form linear combinations of the CPT even term. These are given below KM :
$`\kappa _{DE}^{}{}_{}{}^{jk}=2(k_F)^{0j0k};\kappa _{HB}^{}{}_{}{}^{jk}={\displaystyle \frac{1}{2}}ฯต^{jpq}ฯต^{krs}(k_F)^{pqrs};`$
$`\kappa _{DB}^{}{}_{}{}^{jk}=\kappa _{HE}^{}{}_{}{}^{kj}=(k_F)^{0jpq}ฯต^{kpq}.`$ (30)
The dynamics of the model can be described in terms of equivalent $`๐`$, $`๐`$, $`๐`$ and $`๐`$ fields KM ; WolfGRG in a vacuum using the matrices in (2.2):
$$\left(\begin{array}{c}๐\\ ๐\end{array}\right)=\left(\begin{array}{cc}ฯต_0(1+\kappa _{DE})& \sqrt{\frac{ฯต_0}{\mu _0}}\kappa _{DB}\\ \sqrt{\frac{ฯต_0}{\mu _0}}\kappa _{HE}& \mu _0^1(1+\kappa _{HB})\end{array}\right)\left(\begin{array}{c}๐\\ ๐\end{array}\right)$$
(31)
Note that (31) is rank 6, as the $`\kappa `$ matrices are rank 3 as defined in (2.2). The standard Maxwell equations in a vacuum are recovered if these $`\kappa `$ matrices are set to zero.
Thus the effect of the SME in the photon-sector can be interpreted as introducing medium-like properties to the vacuum. In the full SME, this is considered as an effect from Planck-scale physics in the early universe. The $`\kappa `$ matrices are all position dependent and thus act as โvaluesโ positioned throughout space. If one or more of these values are zero, it does not imply the rest are also zero as there is no relation between each of the independent components. However, there is a linear combination of these components which allows us to separate them into birefringent KM and non-birefringent terms. By eliminating those values which have been constrained beyond what we hope to achieve in this experiment, these terms can be simply written as in (32).
$`(\stackrel{~}{\kappa }_{e+})^{jk}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\kappa _{DE}^{}{}_{}{}^{jk}+\kappa _{HB}^{}{}_{}{}^{jk};(\stackrel{~}{\kappa }_e)^{jk}={\displaystyle \frac{1}{2}}\kappa _{DE}^{}{}_{}{}^{jk}\kappa _{HB}^{}{}_{}{}^{jk}{\displaystyle \frac{1}{3}}\delta ^{jk}(\kappa _{DE})^{ll}`$
$`(\stackrel{~}{\kappa }_{o+})^{jk}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\kappa _{DB}^{}{}_{}{}^{jk}+\kappa _{HE}^{}{}_{}{}^{jk};\stackrel{~}{\kappa }_o^{jk}={\displaystyle \frac{1}{2}}\kappa _{DB}^{}{}_{}{}^{jk}\kappa _{HE}^{}{}_{}{}^{jk}`$ (32)
$`\stackrel{~}{\kappa }_{tr}`$ $`=`$ $`{\displaystyle \frac{1}{3}}\delta ^{jk}(\kappa _{DE})^{ll}`$
In the above definitions, $`\stackrel{~}{\kappa }_{e+}`$, $`\stackrel{~}{\kappa }_e`$ and $`\stackrel{~}{\kappa }_{tr}`$ are parity-even matrices, while $`\stackrel{~}{\kappa }_{o+}`$ and $`\stackrel{~}{\kappa }_o`$ are the parity-odd matrices. As mentioned in KM ; Kost01 , the $`\stackrel{~}{\kappa }_{e+}`$ and $`\stackrel{~}{\kappa }_o`$ are constrained such that $`\left|\stackrel{~}{\kappa }^{JK}\right|2\times 10^{32}`$. Thus, both $`\stackrel{~}{\kappa }_{e+}`$ and $`\stackrel{~}{\kappa }_o`$ are set to zero each time they appear in our equations. Also, $`\stackrel{~}{\kappa }_{tr}`$ is a scalar, which resonator experiments are usually insensitive to TobarPRD , and is not considered in this work.
The standard reference frame that we use is the Sun-Centred Celestial Equatorial Frame (SCCEF), which is shown in Fig. 2. This is the frame in which the sun is at the center, and is inertial with respect to the CMB to first order. The axes in this frame are labeled $`X`$, $`Y`$and $`Z`$. The $`Z`$ axis is defined (KM, , pg. 6)(Bluhm, , pg. 3) to be parallel to the Earthโs north pole, or $`90^{}`$ declination. The $`X`$ axis points from the sun towards the Earth at the moment of the autumn equinox, or $`0^{}`$ right ascension (RA) and $`0^{}`$ declination, while the $`Y`$ axis is at $`90^{}`$ RA and also at $`0^{}`$ declination, usually taken in the J2000.0 frame.
The convention described in (KM, , pg. 18), which has the raised capital indicies $`(J,K)`$ in the SCCEF, has been used. Local coordinates $`x`$, $`y`$ and $`z`$ are defined on the Earthโs surface (at the point of the experiment). The $`z`$ axis is defined as being locally normal to the ground, vertically upwards. The $`x`$ axis points south and the $`y`$ axis points east. These coordinates are denoted by the lowered capital indicies $`(j,k)`$ and they rotate with sidereal period $`\mathrm{\Delta }T_{}=\frac{1}{\omega _{}}23`$ h $`56`$ min. There is a relation between these two coordinates which is given by the following rotation matrix:
$$R^{jJ}=\left(\begin{array}{ccc}\mathrm{cos}\chi \mathrm{cos}\omega _{}T_{}& \mathrm{cos}\chi \mathrm{sin}\omega _{}T_{}& \mathrm{sin}\chi \\ \mathrm{sin}\omega _{}T_{}& \mathrm{cos}\omega _{}T_{}& 0\\ \mathrm{sin}\chi \mathrm{cos}\omega _{}T_{}& \mathrm{sin}\chi \mathrm{sin}\omega _{}T_{}& \mathrm{sin}\chi \end{array}\right)$$
(33)
Here $`\chi `$ is the co-latitude of the laboratory from the north pole, and the $`T_{}`$ is the time coordinate that is related to the sidereal frequency of the Earth. The time $`T_{}`$ is defined in KM as any time when the $`y`$ axis and the $`Y`$ axis align. This has been taken to be the first time this occurs after the vernal equinox, which points along the negative $`x`$-axis.
When searching for leading order violations it is only necessary to consider the spinning of the Earth about itself. However, the orbit of the Earth about the sun may also be considered, since it induces Lorentz boosts and we may calculate proportional terms to these. Since the Earth moves relatively slowly around the Sun compared to the speed of light, the boost terms will be suppressed by the velocity with respect to the speed of light ($`\beta _{}=\frac{v_{}}{c}10^4`$). The boost velocity of a point on the Earthโs surface is given by the following relation:
$$\stackrel{}{\beta }=\beta _{}\left(\begin{array}{c}\mathrm{sin}\mathrm{\Omega }_{}T\\ \mathrm{cos}\eta \mathrm{cos}\mathrm{\Omega }_{}T\\ \mathrm{sin}\eta \mathrm{cos}\mathrm{\Omega }_{}T\end{array}\right)+\beta _L\left(\begin{array}{c}\mathrm{sin}\omega _{}T\\ \mathrm{cos}\omega _{}T\\ 0\end{array}\right)$$
(34)
Here $`\beta _{}`$ is the value for the boost speed of the orbital motion of the Earth and $`\beta _L`$ is the boost speed of the lab at the surface of the Earth due to its spin motion. The latter is location dependent, but is less than $`1.5\times 10^6`$ and is zero at the poles ($`\eta `$ is as defined in Fig. 2). The Lorentz matrix, $`\mathrm{\Lambda }_\nu ^\mu `$, that implements the transformation from the SCCEF to the laboratory frame with the sidereal rotation $`R^{jJ}`$ and a boost $`\stackrel{}{\beta }`$ is given by,
$$\mathrm{\Lambda }_\nu ^\mu =\left(\begin{array}{cccc}1& \beta ^1& \beta ^2& \beta ^3\\ (R\stackrel{}{\beta })^1& R^{11}& R^{12}& R^{13}\\ (R\stackrel{}{\beta })^2& R^{21}& R^{22}& R^{23}\\ (R\stackrel{}{\beta })^3& R^{31}& R^{32}& R^{33}\end{array}\right)$$
(35)
After some calculation KM the $`\kappa `$ matrices from the SCCEF (indexed by J and K) can be express in terms of the values in the laboratory frame (indexed by j and k).
$`(\kappa _{DE})_{lab}^{jk}`$ $`=`$ $`T_0^{jkJK}(\kappa _{DE})^{JK}T_1^{kjJK}(\kappa _{DB})^{JK}T_1^{jkJK}(\kappa _{DB})^{JK}`$ (36)
$`(\kappa _{HB})_{lab}^{jk}`$ $`=`$ $`T_0^{jkJK}(\kappa _{DE})^{JK}T_1^{(kjKJ}(\kappa _{DB})^{JK}T_1^{jkKJ}(\kappa _{DB})^{JK}`$ (37)
$`(\kappa _{DB})_{lab}^{jk}`$ $`=`$ $`T_0^{jkJK}(\kappa _{DB})^{JK}+T_1^{kjJK}(\kappa _{DB})^{JK}+T_1^{jkJK}(\kappa _{HB})^{JK}`$ (38)
Where:
$`T_0^{jkJK}`$ $`=`$ $`R^{jJ}R^{kK}`$ (39)
$`T_1^{jkJK}`$ $`=`$ $`R^{jP}R^{kJ}ฯต^{KPQ}\beta ^Q`$ (40)
Here $`ฯต`$ is the standard anti-symmetric tensor.
In the section 3 we apply the above to calculate the sensitivity of typical resonator experiment. To do this the sensitivity to the components given in Eqn. (2.2) are derived.
## 3 Applying the SME to Resonator Experiments
The modified Lagrangian of the SME introduces perturbations of the electric and magnetic fields in a vacuum. The unperturbed fields are denoted by a zero subscript to distinguish them from the Lorentz-violating fields. Putative Lorentz violations are produced by motion with respect to a preferred frame, which perturbs the fields generating an observable signal. The general framework KM for denoting the sensitivity of this observable signal in the laboratory frame is a linear expression as follows:
$$\delta ๐ช=(_{DE})_{lab}^{jk}(\kappa _{DE})_{lab}^{jk}+(_{HB})_{lab}^{jk}(\kappa _{HB})_{lab}^{jk}+(_{DB})_{lab}^{jk}(\kappa _{DB})_{lab}^{jk}$$
(41)
The summation over the indicies is implied, and the components of the $`_{lab}^{jk}`$ matricies are in general a function of time. The observable is dependent on the type of experiment, and in the case of a resonant cavity experiments it is the resonance frequency deviation, $`\frac{\delta \nu }{\nu _0}`$. Since the laboratory frame and the resonator frame do not necessarily coincide, we first consider $`_{res}^{jk}`$ coefficients in the resonator frame and later relate it to the laboratory and sun-centred frame.
In general, resonators may be constructed from dielectric and magnetic materials. To calculate the $`_{res}^{jk}`$ matricies for such structures a more general form of (31) must be considered, which includes the properties of the medium, $`\mu _r`$ (permeability) and $`ฯต_r`$ (permittivity), which are in general second order tensors.
$$\left(\begin{array}{c}๐\\ ๐\end{array}\right)=\left(\begin{array}{cc}ฯต_0(ฯต_r+\kappa _{DE})& \sqrt{\frac{ฯต_0}{\mu _0}}\kappa _{DB}\\ \sqrt{\frac{ฯต_0}{\mu _0}}\kappa _{HE}& \mu _0^1(\mu _r^1+\kappa _{HB})\end{array}\right)\left(\begin{array}{c}๐\\ ๐\end{array}\right)$$
(42)
Here as was derived in KM , we assume that the fractional frequency shift due to Lorentz violations is given by:
$`{\displaystyle \frac{\mathrm{\Delta }\nu }{\nu _0}}={\displaystyle \frac{1}{4U}}\times `$ (43)
$`{\displaystyle _V}d^3x\left(ฯต_0๐_0^{}\kappa _{DE}๐_0\mu _0^1๐_0^{}\kappa _{HB}๐_0+2Re\left(\sqrt{{\displaystyle \frac{ฯต_0}{\mu _0}}}๐_0^{}\kappa _{DB}๐_0\right)\right)`$
Here $`U`$ is the energy stored in the field and is given by the standard electrodynamic integral.
$$U=\frac{1}{4}_Vd^3x(๐_0๐_0^{}+๐_0๐_0^{})$$
(44)
In Maxwellian electrodynamics the balance of magnetic and electrical energy in a resonator is equal, so the following is true:
$$U=\frac{1}{2}_Vd^3x๐_0๐_0^{}=\frac{1}{2}_Vd^3x๐_0๐_0^{}$$
(45)
This reduces $`U`$ to an effective normalisation factor for either an electric or magnetic filling factor. Also, since the $`\kappa `$ terms from the integral of (43) are only time dependent rather than spatially dependent, the $`\kappa `$ terms can be removed from the integral. Thus, the final term in (43) will be zero since the electric and magnetic fields are orthogonal in a resonant structure. By equating (41) and (43) in the resonator frame the $`(_{DB})_{res}^{jk}`$ coefficients are calculated to be zero, eliminating the possibility of making a measurement of $`\kappa _{tr}`$TobarPRD ; Bail . Assuming the resonator permeability and permittivity have no off-diagonal coefficients (i.e. non-gyrotropic) such that;
$$ฯต_r=ฯต_0\left(\begin{array}{ccc}ฯต_x& 0& 0\\ 0& ฯต_y& 0\\ 0& 0& ฯต_z\end{array}\right)\mu _r=\mu _0\left(\begin{array}{ccc}\mu _x& 0& 0\\ 0& \mu _y& 0\\ 0& 0& \mu _z\end{array}\right)$$
(46)
the only non-zero coefficients may then be calculated to be
$$(_{DE})_{res}^{jj}=\frac{1}{ฯต_j}\frac{_Vd^3x\left|E_0^j\right|^2}{2_Vd^3x๐_{0}^{}{}_{}{}^{}๐_0}=\frac{Pe_j}{2ฯต_j}$$
(47)
$$(_{HB})_{res}^{jj}=\mu _j\frac{_Vd^3x\left|H_0^j\right|^2}{2_Vd^3x๐_{0}^{}{}_{}{}^{}๐_0}=\mu _j\frac{Pm_j}{2}$$
(48)
Thus the $`_{DE}`$ and $`_{HB}`$ matricies are diagonal and simply related to the electric and magnetic energy filling factors, $`Pe_j`$ and $`Pm_j`$ respectively WolfGRG . In general a resonator may consist of more than one material, and may include vacuum. In this case (47) and (48) may be written more generally ($`s`$ is the number of different materials including vacuum).
$`(_{DE})_{res}^{jj}`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{s}{}}}{\displaystyle \frac{Pe_j^i}{2ฯต_j^i}}`$ (49)
$`(_{HB})_{res}^{jj}`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{s}{}}}{\displaystyle \frac{\mu _j^iPm_j^i}{2}}`$ (50)
To measure the resonant frequency it is necessary excite electromagnetic fields inside the resonator and then compare it against a similar frequency. To be sensitive to violations of LLI, the comparison frequency must be generated by a source which exhibits a different dependence on Lorentz violations in the photon sector. For example, an atomic standard (such as a hydrogen maser) may operate in a mode which is not sensitive to Lorentz violations WolfGRG ; Wolf04 . Alternatively, the resonant frequency may be compared against another resonator designed to have a different dependence. The latter can be achieved by orientating two identical resonators orthogonally Muller , or by exciting two modes in a matter filled resonator with orthogonal polarizations. In both cases the field components must be considered with respect to the laboratory frame and not the resonator. For such an experiment the observable becomes the frequency difference (between a resonator labeled $`a`$ and $`b`$) such that;
$$\delta ๐ช=\frac{\delta \nu _a}{\nu _a}\frac{\delta \nu _b}{\nu _b}$$
(51)
Thus, with respect to the laboratory frame, the effective $`(_{DE})_{lab}`$ and $`(_{HB})_{lab}`$ matricies consistent with (41) become:
$`(_{DE})_{ab}=`$ (52)
$`\left(\begin{array}{ccc}(_{DE})_a^{xx}(_{DE})_b^{xx}& 0& 0\\ 0& (_{DE})_a^{yy}(_{DE})_b^{yy}& 0\\ 0& 0& (_{DE})_a^{zz}(_{DE})_b^{zz}\end{array}\right)`$
$`(_{HB})_{ab}=`$ (56)
$`\left(\begin{array}{ccc}(_{HB})_a^{xx}(_{HB})_b^{xx}& 0& 0\\ 0& (_{HB})_a^{yy}(_{HB})_b^{yy}& 0\\ 0& 0& (_{HB})_a^{zz}(_{HB})_b^{zz}\end{array}\right)`$
These equations are general for any resonator experiments, including Fabry-Perot and microwave cavity experiments, and simplify the analysis for complex resonator configurations, such as whispering gallery mode resonators. Only the electric and magnetic filling factors need to be calculated to determine the sensitivity coefficients to the observable, which is possible using standard numerical techniques krupka .
To determine the sensitivity of stationary laboratory experiments one calculates the time dependence of (51) due to the sidereal and orbital motion of the Earth around the Sun in terms of the Sun-centred coefficients given in (32) and (36). This calculation has already been done in KM ; WolfGRG ; Muller2 and will not be repeated here. In the following subsection we generalize this analysis to rotating experiments.
### 3.1 Rotation in the Laboratory Frame
Non-rotating experiments Lipa ; Muller ; Wolf04 that rely on Earth rotation to modulate a Lorentz violating effect are not optimal for two reasons. Firstly, the sensitivity is proportional to the noise in the system at the modulation frequency, typically best for microwave resonator-oscillators and Fabry-Perot stabilized lasers for periods between 10 to 100 seconds. Secondly, the sensitivity is proportional to the square root of the number of periods of the modulation signal, therefore taking a relatively long time to acquire sufficient data. Thus, by rotating the experiment the data integration rate is increased and the relevant signals are translated to the optimal operating regime Mike . For rotation in the laboratory frame the $`()_{lab}^{jk}`$ coefficients become a function of time and depend on the axis of rotation. In the laboratory it is most practical to rotate around the axis of the gravitational field to reduce gravity induced perturbation of the experiment. Thus, our analysis includes rotation about the laboratory $`z`$-axis. If we set the time, $`t=0`$, to be defined when the experiment and laboratory axes are aligned, and we only consider the time varying components (i.e. the most sensitive ones induced by rotation), then for clock-wise rotation of $`\omega _R`$ rads/sec, (52) and (56)become:
$`(_{DE})_{lab}`$ $`=`$ $`\left(\begin{array}{ccc}๐ฎ_{DE}\mathrm{cos}(2\omega _Rt)& ๐ฎ_{DE}\mathrm{sin}(2\omega _Rt)& 0\\ ๐ฎ_{DE}\mathrm{sin}(2\omega _Rt)& ๐ฎ_{DE}\mathrm{cos}(2\omega _Rt)& 0\\ 0& 0& 0\end{array}\right)`$ (63)
$`๐ฎ_{DE}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left((_{DE})_a^{xx}(_{DE})_a^{yy}(_{DE})_b^{xx}+(_{DE})_b^{yy}\right)`$ (64)
$`(_{HB})_{lab}`$ $`=`$ $`\left(\begin{array}{ccc}๐ฎ_{HB}\mathrm{cos}(2\omega _Rt)& ๐ฎ_{HB}\mathrm{sin}(2\omega _Rt)& 0\\ ๐ฎ_{HB}\mathrm{sin}(2\omega _Rt)& ๐ฎ_{HB}\mathrm{cos}(2\omega _Rt)& 0\\ 0& 0& 0\end{array}\right)`$ (68)
$`๐ฎ_{HB}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left((_{HB})_a^{xx}(_{HB})_a^{yy}(_{HB})_b^{xx}+(_{HB})_b^{yy}\right)`$ (69)
Note that if one resonator is tested with respect to a stationary generated frequency, then the $`()_i^{jj}`$ coefficients in the definition of $`๐ฎ_{HB}`$ and $`๐ฎ_{HB}`$ pertaining to that frequency must be set to zero.
To determine the time dependence of the observable (51) we follow the same procedure as presented in subsection 2.2 to transform the $`\stackrel{~}{\kappa }`$ matricies given in (32) to the $`\kappa `$ matricies in the laboratory given by (36), (37) and (38). We then substitute (64) to (69) into (41) to calculate the time dependence of the observable of (51). This is a tedious process and the details are omitted. Essentially, because the ($`_{DE})_{lab}`$ and ($`_{HB})_{lab}`$ matricies are time dependent at $`2\omega _R`$, the observable signals are at frequencies close to this value and are summarized in Table 1. Here the frequency of Earth rotation is defined as $`\omega _{}`$, and orbit around the sun as $`\mathrm{\Omega }_{}`$. The $`\omega _{}`$ is commonly referred to as the sidereal frequency, while the $`\mathrm{\Omega }_{}`$ is referred to as the annual frequency. We also define the sensitivity factor, $`๐ฎ`$ of the experiment as:
$$๐ฎ=๐ฎ_{HB}๐ฎ_{DE}$$
(70)
To decorrelate all side bands, more than one year of data is necessary. In this case we have eight unknown $`\stackrel{~}{\kappa }`$ coefficients and thirty possible individual measurements listed in Table 1, which is an over parameterization. For short data sets (less than a year) we do not have enough information to satisfy the Nyquist condition to distinguish between frequencies that differ by the annual offset (collected in the same blocks). Thus, to make a short data set approximation, we collect the sidebands together (see Fig. 3).
The short data set approximation is achieved by knowing the angle of the orbit, $`\mathrm{\Phi }=\mathrm{\Omega }_{}t`$, in the sun-centered frame with respect to the negative $`X`$-axis (which occurs at the vernal equinox as shown in Fig. 2), and then taking a Taylor series expansion around that angle. Here we define the phase of the combined rotational and sidereal term as $`\theta `$ and $`\mathrm{\Phi }_0`$ as the value of $`\mathrm{\Phi }`$ when a short data set is taken. Since $`\delta \mathrm{\Phi }\mathrm{\Phi }\mathrm{\Phi }_0`$ is small with respect to $`2\pi `$, via the double angle rule we can derive the following relationships:
$`\mathrm{sin}(\theta \pm (\delta \mathrm{\Phi }+\mathrm{\Phi }_0))`$ $`=`$ $`\mathrm{sin}(\theta \pm \delta \mathrm{\Phi })\mathrm{cos}(\mathrm{\Phi }_0)\pm \mathrm{cos}(\theta \pm \delta \mathrm{\Phi })\mathrm{sin}(\mathrm{\Phi }_0)`$
$``$ $`\mathrm{sin}(\theta )\mathrm{cos}(\mathrm{\Phi }_0)\pm \mathrm{cos}(\theta )\mathrm{sin}(\mathrm{\Phi }_0)`$
$`\mathrm{cos}(\theta \pm (\delta \mathrm{\Phi }+\mathrm{\Phi }_0))`$ $`=`$ $`\mathrm{cos}(\theta \pm \delta \mathrm{\Phi })\mathrm{cos}(\mathrm{\Phi }_0)\mathrm{sin}(\theta \pm \delta \mathrm{\Phi })\mathrm{sin}(\mathrm{\Phi }_0)`$
$``$ $`\mathrm{cos}(\theta )\mathrm{cos}(\mathrm{\Phi }_0)\mathrm{sin}(\theta )\mathrm{sin}(\mathrm{\Phi }_0)`$
Now we can combine the sidebands as shown in Fig. 3 by applying the above relationships to eliminate the dependence on $`\mathrm{\Omega }_{}`$. In this case the components from Table 1 decompose to those listed in Table 2.
The first feature to notice is that the $`2\omega _R\pm 2\omega _{}`$ sidebands are redundant. One might also expect the $`2\omega _R\pm \omega _{}`$ sidebands to be redundant as well. The only reason they are not is because we have taken into account the velocity of the laboratory due to the Earth spinning on its axis, $`\beta _L`$. In fact it turns out that it is not useful to keep this term because $`\beta _{}`$ is two orders of magnitude larger and when one applies the data analysis procedures the sensitivities will be degraded if the analysis depends on the $`\beta _L`$ terms for the uniqueness of the solution. Thus, since it makes no practical sense to keep these terms, we set them to zero. For this case the coefficients are listed in Table 3.
For data sets of less than one year, the components in Table 3 may be used to set upper limits on the $`\stackrel{~}{\kappa }`$ coefficients in the SME. Since there are only five possible independent components, to set limits on eight coefficients we use the same technique as adopted by Lipa et. al. Lipa . The $`\stackrel{~}{\kappa }_{o+}`$ boost coefficients are set to zero to calculate limits on the $`\stackrel{~}{\kappa }_e`$ isotropy coefficients and vice versa. This technique assumes no correlation between the isotropy and boost coefficients. It would be unlikely that a cancelation of Lorentz violating effects would occur, as this would necessitate a fortuitous relationship between the coefficients of the same order of value as the boost suppression coefficient (i.e. orbit velocity, $`\beta _{}`$), and consistent with the correct linear combinations as presented in Tab 3.
Another practical point is that the largest systematic effect occurs at $`2\omega _R`$. Thus, when setting the limits on the three $`\stackrel{~}{\kappa }_{o+}^{XY}`$, $`\stackrel{~}{\kappa }_{o+}^{XZ}`$ and $`\stackrel{~}{\kappa }_{o+}^{YZ}`$ coefficients we only use the data collected at $`2\omega _R\pm \omega _{}`$ and $`2\omega _R\pm 2\omega _{}`$ frequencies. Likewise for the $`\stackrel{~}{\kappa }_e^{XY}`$, $`\stackrel{~}{\kappa }_e^{XZ}`$, $`\stackrel{~}{\kappa }_e^{YZ}`$ and $`(\stackrel{~}{\kappa }_e^{XX}\stackrel{~}{\kappa }_e^{YY})`$ coefficients. These are the same coefficients that have had limits set by the non-rotating experiments Muller Lipa Wolf04 . The remaining coefficient $`\stackrel{~}{\kappa }_e^{ZZ}(\stackrel{~}{\kappa }_e^{XX}+\stackrel{~}{\kappa }_e^{YY})`$ can only be set amongst a systematic signal at $`2\omega _R`$, which is in general much greater than the statistical uncertainties at the other frequencies. In this case we can assume all coefficients are zero except for the $`\stackrel{~}{\kappa }_e^{ZZ}`$ coefficient. However, it is not straight forward to set a limit on any putative Lorentz violation amongst a large systematic as one can not be sure if the systematic signal actually cancels an effect. Since the signal at $`2\omega _R`$ is dominated by systematic effects, it is likely that its phase and amplitude will vary across different data sets. In this case the systematic signal from multiple data sets can be treated statistically to place an upper limit on $`\stackrel{~}{\kappa }_e^{ZZ}`$. In our experiment we use this technique to set an upper of $`2.1(5.7)\times 10^{14}`$ PRL (see section 6).
### 3.2 Phase with respect to the SCCEF
To extract the $`\kappa `$ components of the SME out of our observed signal we first need to determine the relevant $`C_{2\omega _i}`$ and $`S_{2\omega _i}`$ coefficients listed in Table 1. This in turn requires us to know the phase of the experimentโs orientation with respect to the SCCEF. In this section we will derive an expression for this phase in terms of the time origins of the experimentโs rotation, the Earthโs sidereal rotation, and the orbit of the Earth around the Sun.
In general, we are interested in the frequency components
$$2\omega _{[a,b]}=2\omega _R+a\omega _{}+b\mathrm{\Omega }_{}$$
(71)
where $`a`$ and $`b`$ take on values in the domains
$$a[2,2],b[1,1]$$
(72)
Thus to determine the $`C_{2\omega _{[a,b]}}`$ coefficient we fit the data with a model of the form
$$\mathrm{cos}(2\omega _RT_R+a\omega _{}T_{}+b\mathrm{\Omega }_{}T)$$
(73)
where $`T_R`$ is the experimentโs rotation time, $`T_{}`$ is the sidereal time, and $`T`$ is the time since the vernal equinox.
To simplify our analysis we aim to transform this expression to the form $`\mathrm{cos}(\alpha t+\varphi )`$. To achieve this we note that the difference $`\delta _R`$ between the experimentโs rotation time $`T_R`$ and the time since the vernal equinox $`T`$ is constant over the course of the measurement, as determined by the initial configuration of the experiment, and we may write,
$$\delta _R=T_RT.$$
(74)
Similarly the sidereal time and the time since the vernal equinox are related by $`\delta _{}`$,
$$\delta _{}=T_{}T$$
(75)
By combining (73), (74) and (75) we arrive at an expression of the desired form.
$`\mathrm{cos}(2\omega _RT_R+a\omega _{}T_{}+b\mathrm{\Omega }_{}T)`$ $`=`$ $`\mathrm{cos}(2\omega _R(\delta _R+T)+a\omega _{}(\delta _{}+T)+b\mathrm{\Omega }_{}T)`$
$`=`$ $`\mathrm{cos}((2\omega _R+a\omega _{}+b\mathrm{\Omega }_{})T`$
$`+2\omega _R\delta _R+a\omega _{}\delta _{})`$
Thus we can account for the phase of the experiment relative to the SCCEF by determining $`\delta _R`$ and $`\delta _{}`$. The origin of the experimentโs rotation time $`T_R`$ is defined to be the instant at which the axis of symmetry of the first resonator (resonator $`a`$) is aligned with the local $`y`$ axis. Our experiment has been designed such that the time origin of the data acquisition coincides with the same event, rendering $`\delta _R=0`$ in our case.
We also need to obtain $`\delta _{}`$ for the sidereal rotation. We define $`T_{}=0`$ as in KM to be the instant the local $`y`$ axis and the SCCEF $`Y`$ axis are aligned (noon) in the laboratory (see figure 2). Let us define $`T_v`$ to be the time in seconds after midnight UTC+0, at which the vernal equinox has occurred in the J2000.0 frame KM . For convenience we also define our longitude $`T_L`$ in terms of sidereal seconds from midnight (in the case of our laboratory $`T_l=115.826^{}\times \frac{23hr56min}{360^{}}=27721sec`$). There exists a special location whose meridian is at noon at the vernal equinox. For this special location (during the vernal equinox), $`\delta _{}=0`$ since the time when the $`y`$ and $`Y`$ axes align and the vernal equinox are the same. We see geometrically that any longitude greater than this meridian will have positive $`\delta _{}`$, otherwise if the longitude is less than this meridian it would have negative $`\delta _{}`$. As shown in Fig. 4, we can now derive an expression for $`\delta _{}`$.
$$\delta _{}=T_L+T_v\frac{23hr56min}{2}$$
(76)
Hence we are able to determine the phase of the experimentโs orientation relative to the SCCEF.
## 4 Comparison of sensitivity of various resonator experiments in the SME
In this section we show how the general analysis may be applied to some common resonator configurations for testing LLI. Also, we propose some new configurations based on exciting two modes in matter filled resonators. The comparison is made by calculating the sensitivity parameter $`๐ฎ`$ of the resonator using Eqns. (49) to (70). Note that the sign of the $`๐ฎ`$ factor depends on the definition of the first resonator. Practically this will need to be the resonator that exhibits the largest value of frequency. In this work, where appropriate, we assume the first resonator is aligned along the $`y`$-axis.
### 4.1 Fabry-Perot resonators
Experiments based on laser stabilized Fabry-Perot resonators typically use either one Brillet or two Muller cavities placed with the lengths orthogonal to the laboratory $`z`$-axis. In a vacuum filled cavity it is easy to show that $`|๐ฎ|=\frac{1}{2}`$ for the configuration in Fig .5. In contrast, when one rotating cavity is compared to a stationary one the value is reduced by a factor of 2, to $`|๐ฎ|=\frac{1}{4}`$.
It is also interesting to consider the sensitivity of matter filled cavities in the photon sector. Here, for simplicity we assume the relative permeability and permittivity are scalars of $`\mu _r`$ and $`ฯต_r`$ respectively. It is straight forward to add anisotropy and only modifies the sensitivity slightly, so for brevity is not considered here. If similar configurations to Fig. 5 are constructed from solid material the sensitivity factor, $`๐ฎ`$, becomes dependent on polarization. This effect also allows for a sensitive experiment by exciting two modes of different polarization inside one cavity (Dual-Mode), of which some examples are shown in Fig . 6. Such cavities have been built previously to measure birefringence Hall . High finesse matter cavities can be made by using low-loss crystalline dielectric materials at optical frequencies FSC ; Schiller . The sensitivity for different configurations are compared in Table 4.
For a low-loss dielectric material with $`E_z`$ polarization in the two orthogonal cavities (Fig .5) the sensitive factor, $`๐ฎ`$, is the same as the vacuum cavity, while for the circularly polarized case, the sensitivity is close to that of the single vacuum cavity resonator. In contrast the same experiment with $`H_z`$ polarization has reduced sensitivity of the order of the permittivity of the material. The sensitivity of the two Dual-Mode resonators gives the possibility of realizing a similar sensitivity to dual cavity experiments, but within the same cavity. The configuration should have a large degree of common mode rejection, and will be much more insensitive to external effects like temperature, vibration etc. and other systematics, and may be worth pursuing for these reasons. Note that Mรผller has recently completed an analysis of conventional cavity configurations in the electron (due to dispersion changes) and photon sector MulMat . In our analysis we have only considered the photon sector and we have proposed some new unconventional configurations. It may be interesting to analyze these configurations in the electron sector. In the next subsection we consider similar configurations for Whispering Gallery (WG) modes.
### 4.2 Whispering gallery mode resonators
In this subsection we consider โpureโ WG modes, with the electric and magnetic fields propagating around with cylindrical symmetry at a discontinuity, with the direction of the Poynting vector ($`๐\times ๐`$) as shown in figure 7. Thus, it is natural to analyse such modes in cylindrical coordinates $`\{r,\varphi ,z\}`$.
For an actual WG mode the wave is reflected off an electromagnetic discontinuity, and the fields mainly lie within the radius of the discontinuity and a smaller inner causticWolfGRG . However, by taking the limit as the azimuthal mode number $`m`$ tends to infinity, the inner caustic converges to the radius of the discontinuity and the fields are reduced to a Dirac delta function. There are two possible polarizations, WGE with dominant $`H_z`$ and $`E_r`$ fields and WGH with dominant $`E_z`$ and $`H_r`$. For โpureโ WG modes, WGE have non-zero electric and magnetic filling factors of $`Pe_r=1`$ and $`Pm_z=1`$, and WGH have electric and magnetic filling factors of $`Pe_z=1`$ and $`Pm_r=1`$, in cylindrical coordinates. The electric and magnetic filling factors may be converted from cylindrical to cartesian symmetry by (the $`z`$ component of the filling factor need not be transformed):
$$Pe_x=Pe_y=\frac{Pe_r+Pe_\varphi }{2}:Pm_x=Pm_y=\frac{Pm_r+Pm_\varphi }{2}$$
(77)
We can now do a similar analysis to subsection 4.1 for configurations shown in Fig. 8 and 9 for the WG case, with the computed sensitivities listed in Table 5. In vacuum the $`๐ฎ`$ factor is half that of the FP cavities in subsection 4.1, and the Dual-Mode resonator is insensitive. However, in a low loss dielectric the $`๐ฎ`$ factor approaches the same value for WGE modes as the FP cavity experiments, but the WGH modes remain about a factor of two less sensitive. The value of the $`๐ฎ`$ factor for the Dual-Mode resonator is the mean value of the WGE and WGH modes.
We have shown that similar sensitivities can be achieved with FP and WG cavity resonators. At UWA we have developed an experiment that uses low loss sapphire crystals, which exhibit a small uniaxial dielectric anistropy. The calculations of the sensitivity are presented section 6.
## 5 Applying the RMS to Whispering Gallery Mode Resonator Experiments
In this section we restrict ourselves to analysis of whispering gallery mode resonator experiments, as the analysis has been well described for Fabry-Perot resonators previously Muller2 . For the whispering gallery mode experiment as shown in Fig. 8, the variable of integration around the path of the resonator is naturally chosen as the azimuthal angle, $`\varphi _j`$, relative to the cylindrical co-ordinates of each resonator. Thus, from Eqn. (3) a frequency shift due to a putative Lorentz violation in the RMS framework is given by,
$$\frac{\mathrm{\Delta }\nu _0}{\nu _0}=\frac{P_{MM}}{2\pi c^2}[(๐ฏ.\widehat{๐}_๐(\varphi _๐))^2d\varphi _a(๐ฏ.\widehat{๐}_๐(\varphi _๐))^2d\varphi _b]$$
(78)
The dominant components of the velocity vector $`๐ฏ`$ were already calculated in section 2.1, so to complete the calculation the time dependence of $`\widehat{๐}_a`$ and $`\widehat{๐}_b`$ must be calculated with respect to the MM-Earth frame. This of course depends on the sidereal and semi-sidereal frequencies, as well as the rotation frequency of the experiment. To start the calculation we define the time, $`t=0`$ when the axis of the two WG resonators are aligned as shown in Fig. 8 (i.e. the resonators align with the laboratory frame). Then from this time we assume the resonator is rotated in a anti-clockwise direction of frequency $`\omega _s`$, so the angle of rotation is $`\gamma =\omega _s(tt_s)`$. Also, the longitudinal angle of the experiment is $`\lambda `$, which is dependent on the sidereal frequency and given by $`\lambda =\omega _{}(tt_l)`$. Then we define the resonator with its cylinder axis in the $`y`$ direction as resonator $`a`$, and the resonator with its cylinder axis in the $`x`$ direction as resonator $`b`$. We also assume the WG modes are oscillating in a clockwise direction. In actual fact the calculation has been verified to be independent of the WG mode direction, and in most experiments is usually a standing wave (depending on the excitation) WGTW . Thus in the laboratory frame at $`t=0`$ the unit vectors in the direction of the Poynting vector are;
$`๐_a(\varphi _a)=\left(\begin{array}{ccc}\mathrm{sin}\varphi _a& & \\ 0& & \\ \mathrm{cos}\varphi _a& & \end{array}\right)๐_b(\varphi _b)=\left(\begin{array}{ccc}0& & \\ \mathrm{cos}\varphi _b& & \\ \mathrm{sin}\varphi _b& & \end{array}\right)`$ (85)
Now if we transform from the resonator to the laboratory, then to the MM-Earth frame the unit vectors become.
$`๐_{Earth:a}=\left(\begin{array}{ccc}\mathrm{sin}\varphi _a(\mathrm{cos}\lambda \mathrm{cos}\chi \mathrm{cos}\gamma \mathrm{sin}\gamma \mathrm{sin}\lambda )+\mathrm{cos}\varphi _a\mathrm{cos}\lambda \mathrm{sin}\chi & & \\ \mathrm{sin}\varphi _a(\mathrm{cos}\lambda \mathrm{sin}\gamma +\mathrm{cos}\gamma \mathrm{cos}\chi \mathrm{sin}\lambda )+\mathrm{cos}\varphi _a\mathrm{sin}\lambda \mathrm{sin}\chi & & \\ \mathrm{cos}\chi \mathrm{cos}\varphi _a+\mathrm{sin}\chi \mathrm{cos}\gamma \mathrm{sin}\varphi _a& & \end{array}\right)`$ (89)
$`๐_{Earth:b}=\left(\begin{array}{ccc}\mathrm{sin}\varphi _b(\mathrm{cos}\lambda \mathrm{cos}\chi \mathrm{sin}\gamma +\mathrm{cos}\gamma \mathrm{sin}\lambda )+\mathrm{cos}\varphi _b\mathrm{cos}\lambda \mathrm{sin}\chi & & \\ \mathrm{sin}\varphi _b(\mathrm{cos}\lambda \mathrm{cos}\gamma \mathrm{sin}\gamma \mathrm{cos}\chi \mathrm{sin}\lambda )+\mathrm{cos}\varphi _b\mathrm{sin}\lambda \mathrm{sin}\chi & & \\ \mathrm{cos}\chi \mathrm{cos}\varphi _b+\mathrm{sin}\chi \mathrm{sin}\gamma \mathrm{sin}\varphi _b& & \end{array}\right)`$ (93)
Here as in the previous sections $`\chi `$ is the angle from the north pole (co-latitude).
The next step is to substitute (89), (93) and (27) into (78). However, to be consistent with the SME analysis the phase should be calculated with respect to the vernal equinox, so that $`\lambda _0=\mathrm{\Phi }_0+\pi `$ is substituted into Eqn. (27) before we substitute it into (78) to calculate the frequency shift. Also, because we defined the rotation to be clockwise in the SME, to be consistent we define $`\gamma _R=\omega _R(tt_s)`$ where $`\omega _R=\omega _s`$. In this case the frequency components, which experience a frequency shift are given in Table 6.
From the results of the calculation we note that perturbations due to Lorentz violations occur at the same frequencies as the SME (see subsection 7.2). Fortunately, it is not necessary to consider perturbations at exactly twice the spin frequency, $`2\omega _s`$, that are primarily due to the larger systematic effects associated with the rotation, as we only need to put a limit on one parameter. Also, the cosine components $`(Cu_{\omega _i})`$ with respect to the CMB are the most sensitive, so we need not consider the sine components.
### 5.1 Phase with respect to the CMB
To extract the $`P_{MM}`$ term from our data we must first determine the phase of our experiment with respect to the CMB. Thus, in similar way to the reasoning for the SME (see subsection 3.2) we require $`\delta _R`$, the difference between the experimentโs rotation time and the time since the vernal equinox, and $`\delta _{}`$, the difference between the sidereal time and the time since the vernal equinox.
As was the case for the SME, $`\delta _R=0`$ since the axis of symmetry of the first resonator, $`a`$, is aligned with the local $`y`$-axis at $`T_R=0`$. However, $`\delta _{}`$ will be different since in the case of the RMS it is measured with respect to the CMB (or MM-Earth frame), not the SCCEF. The CMB is oriented at 11.2 h right ascension, 6.4 degrees declination relative to the equatorial plane. Let us define $`T_v`$ to be the time in seconds after midnight UTC+0, at which the vernal equinox has occurred in the J2000.0 frame KM . $`T_u`$ is the direction of the CMB (11.2h). For convenience we also define our longitude $`T_L`$ in terms of sidereal seconds from midnight (in the case of our laboratory $`T_l=115.826^{}\times \frac{23hr56min}{360^{}}=27721sec`$). As shown in Fig. 10, we now have an expression for $`\delta _{}`$.
$$\delta _{}=T_L+T_v(T_u+\frac{23hr56min}{2})$$
(94)
Hence we are able to determine the phase of the experiment orientation relative to the CMB.
## 6 The University of Western Australia Rotating Experiment
Our experiment consists of two cylindrical sapphire resonators of 3 cm diameter and height supported by spindles at either end within superconducting niobium cavities Giles , which are oriented with their cylindrical axes orthogonal to each other in the horizontal plane (see Fig. 11).
Whispering gallery modes wgmode are excited close to 10 GHz, with a difference frequency of 226 kHz. The frequencies are stabilized using a Pound locking scheme, and amplitude variations are suppressed using an additional control circuit. A detailed description of the cryogenic oscillators can be found in Mann ; Hartnett , and a schematic of the experimental setup is shown in Fig. 12. The resonators are mounted in a common copper block, which provides common mode rejection of temperature fluctuations due to high thermal conductivity at cryogenic temperatures. The structure is in turn mounted inside two successive stainless steel vacuum cylinders from a copper post, which provides the thermal connection between the cavities and the liquid helium bath. A stainless steel section within the copper post provides thermal filtering of bath temperature fluctuations. A foil heater and carbon-glass temperature sensor attached to the copper post controls the temperature set point to 6 K with mK stability. Two stages of vacuum isolation are used to avoid contamination of the sapphire resonators from the microwave and temperature control devices located in the cryogenic environment.
A schematic of the rotation system is shown in Fig.13. A cryogenic dewar containing the resonators, along with the room temperature oscillator circuits and control electronics, is suspended within a ring bearing. A multiple โVโ shaped suspension made from loops of elastic shock cord avoids high Q-factor pendulum modes by ensuring that the cord has to stretch and shrink (providing damping losses) for horizontal motion as well as vertical. The rotation system is driven by a microprocessor controlled stepper motor. A commercial 18 conductor slip ring connector, with a hollow through bore, transfers power and various signals to and from the rotating experiment. A mercury based rotating coaxial connector transmits the difference frequency to a stationary frequency counter referenced to an Oscilloquartz oscillator. The data acquisition system logs the difference frequency as a function of orientation, as well as monitoring systematic effects including the temperature of the resonators, liquid helium bath level, ambient room temperature, oscillator control signals, tilt, and helium return line pressure.
Inside the sapphire crystals standing waves are set up with the dominant electric and magnetic fields in the axial and radial directions respectively, corresponding to a propagation (Poynting) vector around the circumference. The observable of the experiment is the difference frequency, and to test for Lorentz violations the perturbation of the observable with respect to an alternative test theory must be derived. The mode which we excite is a Whispering Gallery mode we have a choice of WGE<sub>m,n,p</sub> or WGH<sub>m,n,p</sub> modes, the first subscript, $`m`$, gives the azimuthal mode number, while $`n`$ and $`p`$ give the number of zero crossings in the radial and z-direction respectively. Typically the so called fundamental mode families WGE<sub>m,0,0</sub> or WGH<sub>m,0,0</sub> as they have the highest Q-factors. To calculate the sensitivity in the RMS we use the technique presented in section 5, while in the following subsection we numerically compute the sensitivity in the SME.
### 6.1 Sensitivity in the SME
In this subsection we calculate the sensitivity of the fundamental WG mode families, WGE<sub>m,0,0</sub> and WGH<sub>m,0,0</sub> to putative Lorentz violation in the SME, and compare it with the โpureโ WG approximation given in Fig. 8. For a proper analysis of the sapphire loaded cavity resonators two regions of space need to be taken into account: the anisotropic crystal and the cavity free space surrounding it (see Fig. 14). The latter has a relative permittivity of 1, while both have relative permeability of 1 in all directions. The calculations proceed by splitting up $`V`$ into $`V_1`$ (the crystal) and $`V_2`$ (the freespace), so we may sum the components of the $``$ matricies over the two volumes (see Eqns. (49) and (50)). The resonator operates close to liquid helium temperatures (6 Kelvin), where the permittivity of sapphire is, $`ฯต_{}=9.272`$ and $`ฯต_{||}=11.349`$.
To determine the sensitivity, we need to just calculate the experiments $`๐ฎ`$ factor in a similar way to the calculation for the โpureโ WG modes in subsection 4.2. In this case the electric and magnetic filling factors must be calculated using a numeric technique such as finite element analysis, method of lines or separation of variablesWolfGRG . In this work we have chosen to use method of lines developed at IRCOM at the University of Limoges MoL . The calculated electric and magnetic field densities for the chosen mode $`(WGH_{8,0,0})`$ of operation at 10 GHz is shown in figure 15, and the $`๐ฎ`$ factor is calculated to be 0.19575.
The actual WG modes have all field components in both regions of the crystal. This modifies the sensitivity slightly, but approaches the limit of the โpureโ WG mode as $`m\mathrm{}`$. The magnitude of the $`๐ฎ`$ factor for the fundamental $`WGE`$ and $`WGH`$ modes at X-Band (8GHz-12GHz) are plotted in Fig. 16. The WGH modes seem converge nicely towards the predicted โpureโ WGH mode sensitivity, while the WGE modes have a dip in sensitivity. This can be explained by an intersection with another mode of the same $`m`$ number, resulting in a spurious mode interactionwhispering . This does not occur in WGH modes since they are the lowest frequency modes for the mode number $`m`$ (refer to figure 2 of whispering ). It is important to note that about a factor of two in sensitivity can be gained if we use a WGE mode rather than a WGH mode. However because we are using a 3 cm crystal rather than a 5 cm crystal, the Q-factor of WGE modes are degraded due to radiation and wall losses. In the future we can markedly improve the sensitivity by employing the typical 5 cm cavities that operate in WGE modes, as were used in the non-rotating experiments WolfGRG .
## 7 Data Analysis and Interpretation of Results
Fig.17 shows typical fractional frequency instability of the 226 kHz difference with respect to 10 GHz, and compares the instability when rotating and stationary. A minimum of $`1.6\times 10^{14}`$ is recorded at 40s. Rotation induced systematic effects degrade the stability up to 18s due to signals at the rotation frequency of $`0.056Hz`$ and its harmonics. We have determined that tilt variations dominate the systematic by measuring the magnitude of the fractional frequency dependence on tilt and the variation in tilt at twice the rotation frequency, $`2\omega _R(0.11Hz)`$, as the experiment rotates. We minimize the effect of tilt by manually setting the rotation bearing until our tilt sensor reads a minimum at $`2\omega _R`$. The latter data sets were up to an order of magnitude reduced in amplitude as we became more experienced at this process. The remaining systematic signal is due to the residual tilt variations, which could be further annulled with an automatic tilt control system. It is still possible to be sensitive to Lorentz violations in the presence of these systematics by measuring the sidereal, $`\omega _{}`$, and semi-sidereal, $`2\omega _{}`$, sidebands about $`2\omega _R`$, as was done in Brillet . The amplitude and phase of a Lorentz violating signal is determined by fitting the parameters of Eq. (95) to the data, with the phase of the fit adjusted according to the test theory used.
$$\frac{\mathrm{\Delta }\nu _0}{\nu _0}=A+Bt+\underset{i}{}C_{\omega _i}\mathrm{cos}(\omega _it+\phi _i)+S_{\omega _i}\mathrm{sin}(\omega _it+\phi _i)$$
(95)
Here $`\nu _0`$ is the average unperturbed frequency of the two sapphire resonators, and $`\mathrm{\Delta }\nu _0`$ is the perturbation of the 226 kHz difference frequency, $`A`$ and $`B`$ determine the frequency offset and drift, and $`C_{\omega _i}`$ and $`S_{\omega _i}`$ are the amplitudes of a cosine and sine at frequency $`\omega _i`$ respectively. In the final analysis we fit 5 frequencies to the data, $`\omega _i=(2\omega _R,2\omega _R\pm \omega _{},2\omega _R\pm 2\omega _{})`$, as well as the frequency offset and drift. The correlation coefficients between the fitted parameters are all between $`10^2`$ to $`10^5`$. Since the residuals exhibit a significantly non-white behavior, the optimal regression method is weighted least squares (WLS) Wolf04 . WLS involves pre-multiplying both the experimental data and the model matrix by a whitening matrix determined by the noise type of the residuals of an ordinary least squares analysis.
We have acquired 5 sets of data over a period of 3 months beginning December 2004, totaling 18 days. The length of the sets (in days) and size of the systematic are ($`3.6,2.3\times 10^{14}`$), ($`2.4,2.1\times 10^{14}`$), ($`1.9,2.6\times 10^{14}`$), ($`4.7,1.4\times 10^{15}`$), and ($`6.1,8.8\times 10^{15}`$) respectively. We have observed leakage of the systematic into the neighboring side bands due to aliasing when the data set is not long enough or the systematic is too large. Fig.18 shows the total amplitude resulting from a WLS fit to 2 of the data sets over a range of frequencies about $`2\omega _R`$. It is evident that the systematic of data set 1 at $`2\omega _R`$ is affecting the fitted amplitude of the sidereal sidebands $`2\omega _R\pm \omega _{}`$ due to its relatively short length and large systematics. By analyzing all five data sets simultaneously using WLS the effective length of the data is increased, reducing the width of the systematic sufficiently as to not contribute significantly to the sidereal and semi-sidereal sidebands.
### 7.1 Standard Model Extension Framework
In the photon sector of the SME 10 independent components of $`\stackrel{~}{\kappa }_{e+}`$ and $`\stackrel{~}{\kappa }_o`$ have been constrained by astronomical measurements to $`<2\times 10^{32}`$ KM ; Kost01 . Seven components of $`\stackrel{~}{\kappa }_e`$ and $`\stackrel{~}{\kappa }_{o+}`$ have been constrained in optical and microwave cavity experiments Muller ; Wolf04 at the $`10^{15}`$ and $`10^{11}`$ level respectively, while the scalar $`\stackrel{~}{\kappa }_{tr}`$ component recently had an upper limit set of $`<10^4`$ TobarPRD . The remaining $`\stackrel{~}{\kappa }_e^{ZZ}`$ component could not be previously constrained in non-rotating experiments Muller ; Wolf04 . In contrast, our rotating experiment is sensitive to $`\stackrel{~}{\kappa }_e^{ZZ}`$. However, it appears only at $`2\omega _R`$, which is dominated by systematic effects. By using the formulas derived in Table 3 for short data sets and the $`๐ฎ`$ factor for the $`WGH_{8,0,0}`$ mode in Fig. 16, the resulting numerical relation between the parameters of the SME and the $`C_{\omega _i}`$ and $`S_{\omega _i}`$ coefficients were calculated and are given in Table 7.
From our combined analysis of all data sets, and using the relation to $`\stackrel{~}{\kappa }_e^{ZZ}`$ given in Table 7, we determine a value for $`\stackrel{~}{\kappa }_e^{ZZ}`$ of $`4.1(0.5)\times 10^{15}`$. However, since we do not know if the systematic has canceled a Lorentz violating signal at $`2\omega _R`$, we cannot reasonably claim this as an upper limit. Since we have five individual data sets, a limit can be set by treating the $`C_{2\omega _R}`$ coefficient as a statistic. The phase of the systematic depends on the initial experimental conditions, and is random across the data sets. Thus, we have five values of $`C_{2\omega _R}`$, ($`\{4.2,11.4,21.4,1.3,8.1\}`$ in $`10^{15}`$), two are negative coefficients and three are positive. If we take the mean of these coefficients, the systematic signal will cancel if the phase is random, but the possible Lorentz violating signal will not, since the phase is constant. Thus a limit can be set by taking the mean and standard deviation of the five coefficient of $`C_{2\omega _R}`$. This gives a more conservative bound of $`2.1(5.7)\times 10^{14}`$, which includes zero. Our experiment is also sensitive to all other seven components of $`\stackrel{~}{\kappa }_e`$ and $`\stackrel{~}{\kappa }_{o+}`$ (see Table 7) and improves present limits by up to a factor of 7, as shown in Table 8.
### 7.2 Robertson, Mansouri, Sexl Framework
From Eqn. (78), the dominant coefficients are calculated to be only due to the cosine terms with respect to the CMB right ascension, $`Cu_{\omega _i}`$, and the theory predicts no perturbations in the quadrature term. Since our experiment is rotating clock wise we can substitute $`2\omega _s=2\omega _R`$, and once we perform the integral and substitute all the numeric values. Following this method we calculate the coefficients as shown in Table 9.
The same five data sets were then re-analysed in the correct quadrature with respect to the CMB, with the results listed with the coefficients in Table 9. The measured and statistical uncertainty of $`P_{MM}`$ is determined to be $`0.9(2.0)\times 10^{10}`$, which represents a factor of 7.5 improvement over previous results $`2.2(1.5)\times 10^9`$Muller .
## 8 Summary
Rotating resonator experiments are emerging as one of the most sensitive types of Local Lorentz Invariance tests in electrodynamics (also see other contributions within these proceedings). In this work we have analysed in detail such experiments to putative Lorentz violation in both the RMS and SME frameworks. In the RMS, rotating experiments only enhance the sensitivity to the Michelson-Morley parameter, $`P_{MM}`$, and are not sensitive to the Kennedy-Thorndike , $`P_{KT}`$, or Ives-Stilwell , $`P_{IS}`$, parameters. In the SME non-rotating resonator experiments in the laboratory test for four components of the $`\stackrel{~}{\kappa }_e`$ tensor and three components of the $`\stackrel{~}{\kappa }_{o+}`$ tensor, with the scalar coefficient $`\stackrel{~}{\kappa }_{tr}`$ and $`\stackrel{~}{\kappa }_e^{ZZ}`$ unmeasurable. Rotation in the SME enhances the sensitivity to the seven components and also allows the determination of the $`\stackrel{~}{\kappa }_e^{ZZ}`$ component. We have shown that all resonator experiment exhibit the same relative frequency spectrum to the putative signal to within a multiplicative sensitivity factor, $`๐ฎ`$. This was utilized to compare the sensitivity of different FP and WG resonator configurations, leading to the proposal of some new dual-mode resonator experiments.
We applied the above analysis to our experiment at the University of Western Australia, which is based on rotating cryogenic sapphire whispering gallery mode microwave oscillators. In summary, we presented the first results of the experiment, which we set bounds on 7 components of the SME photon sector (Table 8) and $`P_{MM}`$ (Table 9) of the RMS framework up to a factor of 7.5 more stringent than those obtained from previous experiments. We also set an upper limit ($`2.1(5.7)\times 10^{14}`$) on the previously unmeasured SME component $`\stackrel{~}{\kappa }_e^{ZZ}`$. To further improve these results, tilt and environmental controls will be implemented to reduce systematic effects. To remove the assumption that the $`\stackrel{~}{\kappa }_{o+}`$ and $`\stackrel{~}{\kappa }_e`$ do not cancel each other, data integration will continue for more than a year.
## 9 Acknowledgment
This work was funded by the Australian Research Council.
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# Decay of Correlations for the RauzyโVeechโZorich Induction Map on the Space of Interval Exchange Transformations and the Central Limit Theorem for the Teichmรผller Flow on the Moduli Space of Abelian Differentials.
## 1 Introduction
The aim of this paper is to prove a stretched-exponential bound for the decay of correlations for the Rauzy-Veech-Zorich induction map on the space of interval exchange transformations (Theorem 4). A Corollary is the Central Limit Theorem for the Teichmรผller flow (Theorem 10).
The proof of Theorem 4 proceeds by approximating the induction map by a Markov chain satisfying the Doeblin condition, the method of Sinai and BunimovichโSinai . The main โloss of memoryโ estimate is Lemma 4.
### 1.1 Interval exchange transformations.
Let $`m`$ be a positive integer. Let $`\pi `$ be a permutation on $`m`$ symbols. The permutation $`\pi `$ will always be assumed irreducible, which means that $`\pi \{1,\mathrm{},k\}=\{1,\mathrm{},k\}`$ only if $`k=m`$.
Let $`\lambda `$ be a vector in $`_+^m`$, $`\lambda =(\lambda _1,\mathrm{},\lambda _m)`$, $`\lambda _i>0`$ for all $`i`$. Denote
$$|\lambda |=\underset{i=1}{\overset{m}{}}\lambda _i.$$
Consider the half-open interval $`[0,|\lambda |)`$. Consider the points $`\beta _i=_{j<i}\lambda _j`$, $`\beta _i^\pi =_{j<i}\lambda _{\pi ^1j}`$.
Denote $`I_i=[\beta _i,\beta _{i+1})`$, $`I_i^\pi =[\beta _i^\pi ,\beta _{i+1}^\pi )`$. The length of $`I_i`$ is $`\lambda _i`$, whereas the length of $`I_i^\pi `$ is $`\lambda _{\pi ^1i}`$.
Set
$$T_{(\lambda ,\pi )}(x)=x+\beta _{\pi i}^\pi \beta _i\mathrm{for}xI_i.$$
The map $`T_{(\lambda ,\pi )}`$ is called an interval exchange transformation corresponding to $`(\lambda ,\pi )`$.
The map $`T_{(\lambda ,\pi )}`$ is an order-preserving isometry from $`I_i`$ onto $`I_{\pi (i)}^\pi `$.
We say that $`\lambda `$ is irrational if there are no rational relations between $`|\lambda |`$, $`\lambda _1`$,$`\lambda _2`$, โฆ$`\lambda _{m1}`$.
###### Theorem 1 (Oseledets(), Keane())
Let $`\pi `$ be irreducible and $`\lambda `$ irrational. Then for any $`x[0,_{i=1}^m\lambda _i)`$, the set $`\{T_{(\lambda ,\pi )}^nx,n0\}`$ is dense in $`[0,_{i=1}^m\lambda _i).`$
### 1.2 Rauzy operations $`a`$ and $`b`$.
Let $`(\lambda ,\pi )`$ be an interval exchange. Assume that $`\pi `$ is irreducible and $`\lambda `$ is irrational.
Following Rauzy , consider the induced map of $`(\lambda ,\pi )`$ on the interval $`[0,|\lambda |min(\lambda _m,\lambda _{\pi ^1(m)}))`$. The induced map is again an interval exchange of $`m`$ intervals. For $`i,j=1,\mathrm{},m`$, denote by $`E_{ij}`$ an $`m\times m`$ matrix of which the $`i,j`$-th element is equal to $`1`$, all others to $`0`$. Let $`E`$ be the $`m\times m`$-identity matrix.
#### 1.2.1 Case $`a`$: $`\lambda _{\pi ^1m}>\lambda _m`$.
Define
$$A(a,\pi )=\underset{i=1}{\overset{\pi ^1(m)}{}}E_{ii}+E_{m,\pi ^1m+1}+\underset{i=\pi ^1m+1}{\overset{m}{}}E_{i,i+1}$$
$$a\pi (j)=\{\begin{array}{cc}\pi j,\hfill & \text{if }j\pi ^1m\text{;}\hfill \\ \pi m,\hfill & \text{if }j=\pi ^1m+1\text{;}\hfill \\ \pi (j1),\hfill & \text{ other }j\text{.}\hfill \end{array}$$
If $`\lambda _{\pi ^1m}>\lambda _m`$, then the induced interval exchange of $`T_{(\lambda ,\pi )}`$ on the interval $`[0,_{im}\lambda _i)`$ is $`T_{(\lambda ^{},\pi ^{})}`$, where $`\lambda ^{}=A(a,\pi )^1\lambda `$ and $`\pi ^{}=a\pi `$.
#### 1.2.2 Case $`b`$: $`\lambda _m>\lambda _{\pi ^1m}`$.
Define
$$A(b,\pi )=E+E_{m,\pi ^1m}$$
$$b\pi (j)=\{\begin{array}{cc}\pi j,\hfill & \text{if }\pi j\pi m\text{;}\hfill \\ \pi j+1,\hfill & \text{if }\pi m<\pi j<m\text{;}\hfill \\ \pi m+1,\hfill & \text{ if }\pi j=m\text{.}\hfill \end{array}$$
If $`\lambda _{\pi ^1m}<\lambda _m`$, then the induced interval exchange of $`T_{(\lambda ,\pi )}`$ on the interval $`[0,_{i\mathrm{\Pi }^1m}\lambda _i)`$ is $`T_{(\lambda ^{},\pi ^{})}`$, where $`\lambda ^{}=A(b,\pi )^1\lambda `$ and $`\pi ^{}=b\pi `$.
Note that operations $`a`$ and $`b`$ are invertible on the space of permutations, namely, we have:
$$a^1\pi (j)=\{\begin{array}{cc}\pi (j),\hfill & \text{if }j\pi ^1(m)\text{;}\hfill \\ \pi (j+1),\hfill & \text{if }\pi ^1(m)+1<j<m\text{;}\hfill \\ \pi (\pi ^1(\pi (m)+1),\hfill & \text{if }j=m\text{.}\hfill \end{array}$$
$$b^1\pi (j)=\{\begin{array}{cc}\pi (j),\hfill & \text{ if }\pi (j)\pi (m)\hfill \\ m,\hfill & \text{if }j=\pi ^1(\pi (m)+1)\text{;}\hfill \\ \pi (j)1,\hfill & \text{if }\pi (j)>\pi (m)+1\text{.}\hfill \end{array}$$
For $`(\lambda ,\pi )\mathrm{\Delta }()`$, denote
$$T_{a^1}(\lambda ,\pi )=(A(a^1\pi ,a)\lambda ,a^1\pi ),T_{b^1}(\lambda ,\pi )=(A(b^1\pi ,b)\lambda ,b^1\pi ).$$
(1)
The interval exchange $`T_{a^1}(\lambda ,\pi )`$ is the preimage of $`(\lambda ,\pi )`$ under the operation $`a`$, and the interval exchange $`T_{b^1}(\lambda ,\pi )`$ is the preimage of $`(\lambda ,\pi )`$ under the operation $`b`$.
Normalize (dividing by $`|\lambda |=\lambda _1+\mathrm{}+\lambda _m`$) and set:
$$t_{a^1}(\lambda ,\pi )=(\frac{A(a^1\pi ,a)\lambda }{|A(a^1\pi ,a)\lambda |},a^1\pi ),t_{b^1}(\lambda ,\pi )=(\frac{A(b^1\pi ,b)\lambda }{|A(b^1\pi ,b)\lambda |},b^1\pi ).$$
(2)
### 1.3 Rauzy class and Rauzy graph.
If $`\pi `$ is an irreducible permutation, then its Rauzy class is the set of all permutations that can be obtained from $`\pi `$ by applying repeatedly the operations $`a`$ and $`b`$; the Rauzy class of the permutation $`\pi `$ is denoted $`(\pi )`$. Rauzy class has a natural structure of an oriented labelled graph: namely, the permutations of the Rauzy class are the vertices of the graph, and if $`\pi =a\pi ^{}`$ then we draw an edge from $`\pi `$ to $`\pi ^{}`$ and label it by $`a`$, and if $`\pi =b\pi ^{}`$ then we draw an edge from $`\pi `$ to $`\pi ^{}`$ and label it by $`b`$. This labelled graph will be called the Rauzy graph of the permutation $`\pi `$.
For example, the Rauzy graph of the permutation $`(4321)`$ is
For a permutation $`\pi `$, consider the set $`\{a^n\pi ,n0\}`$. This set forms a cycle in the Rauzy graph which will be called the $`a`$-cycle of $`\pi `$. Similarly, the set $`\{b^n\pi ,n0\}`$ will be called the $`b`$-cycle of $`\pi `$.
### 1.4 The Rauzy-Veech-Zorich induction.
Denote
$$\mathrm{\Delta }_{m1}=\{\lambda _+^m:|\lambda |=1\},$$
$$\mathrm{\Delta }_\pi ^+=\{\lambda \mathrm{\Delta }_{m1},\lambda _{\pi ^1m}>\lambda _m\},\mathrm{\Delta }_\pi ^{}=\{\lambda \mathrm{\Delta }_{m1},\lambda _m>\lambda _{\pi ^1m}\},$$
$$\mathrm{\Delta }()=\mathrm{\Delta }_{m1}\times (\pi ).$$
Define a map
$$๐ฏ:\mathrm{\Delta }()\mathrm{\Delta }()$$
by
$$๐ฏ(\lambda ,\pi )=\{\begin{array}{cc}(\frac{A(\pi ,a)^1\lambda }{|A(\pi ,a)^1\lambda |},a\pi ),\hfill & \text{if }\lambda \mathrm{\Delta }_\pi ^+\text{;}\hfill \\ (\frac{A(\pi ,b)^1\lambda }{|A(\pi ,b)^1\lambda |},b\pi ),\hfill & \text{if }\lambda \mathrm{\Delta }_\pi ^{}\text{.}\hfill \end{array}$$
Each $`(\lambda ,\pi )\mathrm{\Delta }()`$ has exactly two preimages under the map $`๐ฏ`$, namely, $`t_{a^1}(\lambda ,\pi )`$ and $`t_{b^1}(\lambda ,\pi )`$ (2).
The set $`\mathrm{\Delta }()`$ is a finite union of simplices. Let $`๐ฆ`$ be the Lebesgue measure on $`\mathrm{\Delta }()`$ normalized in such a way that $`๐ฆ(\mathrm{\Delta }())=1`$.
###### Theorem 2 (Veech)
The map $`๐ฏ`$ has an infinite conservative ergodic invariant measure, absolutely continuous with respect to Lebesgue measure on $`\mathrm{\Delta }()`$.
From this result Veech derives that almost all (with respect to m) interval exchange transformations are uniquely ergodic.
Denote
$$\mathrm{\Delta }^+=_{\pi ^{}(\pi )}\mathrm{\Delta }_\pi ^{}^+,\mathrm{\Delta }^{}=_{\pi ^{}(\pi )}\mathrm{\Delta }_\pi ^{}^{}.$$
Following Zorich , we define the function $`n(\lambda ,\pi )`$ in the following way.
$$n(\lambda ,\pi )=\{\begin{array}{cc}inf\{k>0:๐ฏ^k(\lambda ,\pi )\mathrm{\Delta }^{}\},\hfill & \text{if }\lambda \mathrm{\Delta }_\pi ^+\text{;}\hfill \\ inf\{k>0:๐ฏ^k(\lambda ,\pi )\mathrm{\Delta }^+\},\hfill & \text{if }\lambda \mathrm{\Delta }_\pi ^{}\text{.}\hfill \end{array}$$
Define
$$๐ข(\lambda ,\pi )=๐ฏ^{n(\lambda ,\pi )}(\lambda ,\pi ).$$
The map $`๐ข`$ will be referred to as the Rauzy-Veech-Zorich induction map .
For $`(\lambda ,\pi )\mathrm{\Delta }()`$, denote
$$t_{a^n}(\lambda ,\pi )=t_{a^1}^n(\lambda ,\pi ),t_{b^n}(\lambda ,\pi )=t_{b^1}^n(\lambda ,\pi ),T_{a^n}(\lambda ,\pi )=T_{a^1}^n(\lambda ,\pi ),T_{b^n}(\lambda ,\pi )=T_{b^1}^n(\lambda ,\pi ).$$
Under the map $`๐ข`$, each interval exchange $`(\lambda ,\pi )`$ has countably many preimages:
$$๐ข^1(\lambda ,\pi )=\{\begin{array}{cc}\{t_{a^n}(\lambda ,\pi ),n\},\hfill & \text{if }(\lambda ,\pi )\mathrm{\Delta }^+\text{;}\hfill \\ \{t_{b^n}(\lambda ,\pi ),n\},\hfill & \text{if }(\lambda ,\pi )\mathrm{\Delta }^{}\text{.}\hfill \end{array}$$
###### Theorem 3 (Zorich)
The map $`๐ข`$ has an ergodic invariant probability measure, absolutely continuous with respect to Lebesgue on $`\mathrm{\Delta }()`$.
Denote this invariant measure by $`\nu `$; the probability with respect to $`\nu `$ will be denoted by $``$.
Let $`\rho (\lambda ,\pi )`$ be the density of $`\nu `$ with respect to the Lebesgue measure $`๐ฆ`$. Zorich showed that for any $`\pi `$ there exist two positive rational homogeneous of degree $`m`$ functions $`\rho _\pi ^+`$, $`\rho _\pi ^{}`$ such that
$$\rho (\lambda ,\pi )=\{\begin{array}{cc}\rho _\pi ^+(\lambda ),\hfill & \text{if }\lambda \mathrm{\Delta }_\pi ^+\text{;}\hfill \\ \rho _\pi ^{}(\lambda ),\hfill & \text{if }\lambda \mathrm{\Delta }_\pi ^{}\text{.}\hfill \end{array}$$
(3)
Remark. In particular, the invariant density is bounded from below: there exists a positive constant $`C()`$, depending on the Rauzy class only and such that $`\rho (\lambda ,\pi )C()`$ for any $`(\lambda ,\pi )\mathrm{\Delta }()`$.
The map $`๐ข`$ is not mixing: indeed, from the definition of $`๐ข`$, we have
$$๐ข(\mathrm{\Delta }^+)=\mathrm{\Delta }^{},๐ข(\mathrm{\Delta }^{})=\mathrm{\Delta }^+.$$
Let $``$ be the Borel $`\sigma `$-algebra on $`\mathrm{\Delta }()`$, and let $`_n=๐ข^n`$. We have $`_{n+2}_n`$. Recall that exactness of the map $`๐ข^2`$ means, by definition, that the $`\sigma `$-algebra $`_{n=1}^{\mathrm{}}_{2n}`$ is trivial (in other words, that Kolmogorovโs $`01`$ law holds for the map $`๐ข^2`$.)
###### Proposition 1
The map $`๐ข^2:\mathrm{\Delta }^+\mathrm{\Delta }^+`$ is exact with respect to $`\nu |_{\mathrm{\Delta }^+}`$.
This Proposition is proven in Section 4; it implies strong mixing for the map $`๐ข^2`$.
### 1.5 The main result
Introduce a metric on $`\mathrm{\Delta }_{m1}`$ by setting
$$d(\lambda ,\lambda ^{})=\mathrm{log}\frac{\mathrm{max}_i\frac{\lambda _i}{\lambda _i^{}}}{\mathrm{min}_i\frac{\lambda _i}{\lambda _i^{}}}.$$
(4)
Now introduce a metric on $`\mathrm{\Delta }()`$ by setting
$$d((\lambda ,\pi ),(\lambda ^{},\pi ^{}))=\{\begin{array}{cc}2+d(\lambda ,\lambda ^{}),\hfill & \text{if }\pi \pi ^{}\text{;}\hfill \\ d(\lambda ,\lambda ^{}),\hfill & \text{if }\pi =\pi ^{}\text{.}\hfill \end{array}$$
For $`\alpha >0`$, let $`H_\alpha `$ be the space of functions $`\varphi :\mathrm{\Delta }()`$ such that if $`d((\lambda ,\pi ),(\lambda ^{},\pi ^{})1`$, then $`|\varphi (\lambda ,\pi )\varphi (\lambda ^{},\pi ^{})|Cd((\lambda ,\pi ),(\lambda ^{},\pi ^{}))^\alpha `$ for some constant $`C`$.
Define
$$C_{H_\alpha }(\varphi )=\underset{d((\lambda ,\pi ),(\lambda ^{},\pi ^{}))1}{\mathrm{max}}\frac{|\varphi (\lambda ,\pi )\varphi (\lambda ^{},\pi ^{})|}{d((\lambda ,\pi ),(\lambda ^{},\pi ^{}))^\alpha },$$
The main result of this paper is
###### Theorem 4
Let $`๐ข:\mathrm{\Delta }()\mathrm{\Delta }()`$ be the Rauzy-Veech-Zorich induction map and let $`\nu `$ be the absolutely continuous invariant measure.
Let $`p>2`$. Then, for any $`\alpha >0`$, there exist positive constants $`C,\delta `$ such that for any $`\varphi H_\alpha L_p(\mathrm{\Delta }^+(),\nu )`$ and $`\psi L_2(\mathrm{\Delta }^+(),\nu )`$ we have
$$|\varphi \times \psi ๐ข^{2n}๐\nu \varphi ๐\nu \psi ๐\nu |C\mathrm{exp}(\delta n^{1/6})(C_{H_\alpha }(\varphi )+|\varphi |_{L_p})(|\psi |_{L_2}).$$
Denote by $`๐ฉ(0,\sigma )`$ the Gaussian distribution with mean $`0`$ and variance $`\sigma `$. By , we have
###### Corollary 1
Let $`\varphi H_\alpha L_p(\mathrm{\Delta }()^+,\nu )`$, $`\varphi ๐\nu =0`$. Assume that there does not exist $`\psi L_2(\mathrm{\Delta }()^+,\nu )`$ such that $`\varphi =\psi ๐ข^2\psi `$. Then there exists $`\sigma >0`$ such that
$$\frac{1}{\sqrt{N}}\underset{n=0}{\overset{N1}{}}\varphi ๐ข^{2n}\stackrel{๐}{}๐ฉ(0,\sigma )\mathrm{as}N\mathrm{}.$$
### 1.6 Veechโs space of zippered rectangles
A zippered rectangle associated to the Rauzy class $``$ is a quadruple $`(\lambda ,h,a,\pi )`$, where $`\lambda _+^m,h_+^m,a^m,\pi `$, and the vectors $`h`$ and $`a`$ satisfy the following equations and inequalities (one introduces auxiliary components $`a_0=h_0=a_{m+1}=h_{m+1}=0`$, and sets $`\pi (0)=0`$, $`\pi ^1(m+1)=m+1`$.):
$$h_ia_i=h_{\pi ^1(\pi (i)+1)}a_{\pi ^1(\pi (i)+1)1},i=0,\mathrm{},m$$
$$h_i0,i=1,\mathrm{},m,a_i0,i=1,\mathrm{},m1,$$
$$a_i\mathrm{min}(h_i,h_{i+1})\mathrm{for}im,i\pi ^1m,$$
$$a_mh_m,a_mh_{\pi ^1m},a_{\pi ^1m}h_{\pi ^1m+1}$$
The area of a zippered rectangle is given by the expression $`\lambda _1h_1+\mathrm{}+\lambda _mh_m`$. Following Veech, we denote by $`\mathrm{\Omega }()`$ the space of all zippered rectangles, corresponding to a given Rauzy class $``$ and satisfying the condition
$$\lambda _1h_1+\mathrm{}+\lambda _mh_m=1.$$
We shall denote by $`x`$ an individual zippered rectangle.
Veech further defines a map $`๐ฐ`$ and a flow $`P^t`$ on the space of zippered rectangles in the following way:
$$P^t(\lambda ,h,a,\pi )=(e^t\lambda ,e^th,e^ta,\pi ).$$
$$๐ฐ(\lambda ,h,a,\pi )=\{\begin{array}{cc}(A^1(a,\pi )\lambda ,A^t(a,\pi )h,a^{},a\pi ),\hfill & \text{if }(\lambda ,\pi )\mathrm{\Delta }^{}\hfill \\ (A^1(b,\pi )\lambda ,A^t(b,\pi )h,a^{\prime \prime },b\pi ),\hfill & \text{if }(\lambda ,\pi )\mathrm{\Delta }^+,\hfill \end{array}$$
where
$$a_i^{}=\{\begin{array}{cc}a_i,\hfill & \text{if }j<\pi ^1m\text{,}\hfill \\ h_{\pi ^1m}+a_{m1},\hfill & \text{if }i=\pi ^1m\text{,}\hfill \\ a_{i1},\hfill & \text{other }i.\hfill \end{array}$$
$$a_i^{\prime \prime }=\{\begin{array}{cc}a_i,\hfill & \text{if }j<m\text{,}\hfill \\ h_{\pi ^1m}+a_{\pi ^1m1},\hfill & \text{if }i=m\text{.}\hfill \end{array}$$
The map $`๐ฐ`$ is invertible; $`๐ฐ`$ and $`P^t`$ commute ().
Denote
$$\tau (\lambda ,\pi )=(\mathrm{log}(|\lambda |\mathrm{min}(\lambda _m,\lambda _{\pi ^1m})),$$
and for $`x\mathrm{\Omega }()`$, $`x=(\lambda ,h,a,\pi )`$, write
$$\tau (x)=\tau (\lambda ,\pi ).$$
Now define
$$๐ด()=\{x\mathrm{\Omega }():|\lambda |=1\}.$$
and
$$\mathrm{\Omega }_0()=\underset{x๐ด(),0t\tau (x)}{}P^tx.$$
$`\mathrm{\Omega }_0()`$ is a fundamental domain for $`๐ฐ`$ and, identifying the points $`x`$ and $`๐ฐx`$ in $`\mathrm{\Omega }_0()`$, we obtain a natural flow, also denoted by $`P^t`$, on $`\mathrm{\Omega }_0()`$.
The space $`\mathrm{\Omega }()`$ has a natural Lebesgue measure class and so does the transversal $`๐ด()`$. Veech has proved the following Theorem.
###### Theorem 5
There exists a measure $`\mu _{}`$ on $`\mathrm{\Omega }()`$, absolutely continuous with respect to Lebesgue, preserved by both the map $`๐ฐ`$ and the flow $`P^t`$ and such that $`\mu _{}(\mathrm{\Omega }_0())<\mathrm{}`$.
For $`x๐ด()`$, define
$$๐ฎ(x)=๐ฐP^{\tau (x)}(x).$$
The map $`๐ฎ`$ is a lift of $`๐ฏ`$ to the space of zippered rectangles: indeed, if
$$๐ฎ(\lambda ,h,a,\pi )=(\lambda ^{},h^{},a^{},\pi ^{}),$$
then $`(\lambda ^{},\pi ^{})=๐ฏ(\lambda ^{},\pi ^{})`$.
Since $`๐ด()`$ is a transversal to the flow, the measure $`\mu _{}`$ induces an absolutely continuous measure $`\mu _{}^{(1)}`$ on $`๐ด()`$; since $`\mu _{}`$ is both $`๐ฐ`$ and $`P^t`$-invariant, the measure $`\mu _{}^{(1)}`$ is $`๐ฎ`$-invariant. Since $`\mu _{}(\mathrm{\Omega }_0())<\mathrm{}`$, the measure $`\mu _{}^{(1)}`$ is conservative; it is, however, infinite (Veech ).
Zorich constructed a different section for the flow $`P^t`$, for which the restricted measure has finite total mass.
Following Zorich , define
$$\mathrm{\Omega }^+()=\{x=(\lambda ,h,a,\pi ):(\lambda ,\pi )\mathrm{\Delta }^+,a_m0\}.$$
$$\mathrm{\Omega }^{}()=\{x=(\lambda ,h,a,\pi ):(\lambda ,\pi )\mathrm{\Delta }^{},a_m0\},$$
$$๐ด^+()=๐ด()\mathrm{\Omega }^+(),๐ด^{}()=๐ด()\mathrm{\Omega }^{}(),๐ด^\pm ()=๐ด^+()๐ด^{}().$$
Take $`x๐ด^\pm ()`$, $`x=(\lambda ,h,a,\pi )`$, and define
$$(x)=๐ฎ^{n(\lambda ,\pi )}x.$$
The map $``$ is a lift of the map $`๐ข`$ to the space of zippered rectangles: if
$$(\lambda ,h,a,\pi )=(\lambda ^{},h^{},a^{},\pi ^{}),$$
then $`(\lambda ^{},\pi ^{})=๐ข(\lambda ^{},\pi ^{})`$.
We shall see, moreover, that the map $``$ can be almost surely (with respect to Lebesgue) identified with the natural extension of the map $`๐ข`$ (Section 3).
If $`x๐ด^+`$, then $`(x)๐ด^{}`$, and if $`x๐ด^{}`$, then $`(x)๐ด^+`$. The map $``$ is the induced map of $`๐ฎ`$ to the subset $`๐ด^\pm ()`$.
Since $`๐ด^\pm ()`$ is a transversal to the flow $`P^t`$, the measure $`\mu _{}`$ naturally induces an absolutely continuous measure $`\overline{\nu }`$ on $`๐ด^\pm ()`$; since $`\mu _{}`$ is both $`๐ฐ`$ and $`P^t`$-invariant, the measure $`\overline{\nu }`$ is $``$-invariant.
Zorich proved
###### Theorem 6
The measure $`\overline{\nu }`$ is finite and ergodic for $``$.
Since the map $`๐ข`$ is exact (as is shown in Section 4), the map $``$ satisfies the $`K`$-property of Kolmogorov, and, in particular, is strongly mixing. Decay of correlations is proven for the map $``$ as well.
Introduce a metric on the space of zippered rectangles in the following way. Take two zippered rectangles $`x=(\lambda ,h,a,\pi )`$ and $`x^{}=(\lambda ^{},h^{},a^{},\pi ^{})`$. Write
$$d((\lambda ,h,a),(\lambda ^{},h^{},a^{}))=\mathrm{log}\frac{\mathrm{max}_i\frac{\lambda _i}{\lambda _i^{}},\frac{h_i}{h_i^{}},\frac{|a_i|}{|a_i^{}|},\frac{|h_ia_i|}{|h_i^{}a_i^{}|}}{\mathrm{min}_i\frac{\lambda _i}{\lambda _i^{}},\frac{h_i}{h_i^{}},\frac{|a_i|}{|a_i^{}|},\frac{|h_ia_i|}{|h_i^{}a_i^{}|}}.$$
Define the metric on $`\mathrm{\Omega }()`$ by
$$d(x,x^{})=\{\begin{array}{cc}d((\lambda ,h,a),(\lambda ^{},h^{},a^{})\hfill & \text{if }\pi =\pi ^{}\text{ and }\frac{a_m}{a_m^{}}>0\text{;}\hfill \\ 2+d((\lambda ,h,a),\lambda ^{},h^{},a^{}),\hfill & \text{otherwise}.\hfill \end{array}$$
As above, for $`\alpha >0`$, let $`H_\alpha `$ be the space of functions $`\varphi :๐ด^\pm ()`$ such that if $`d(x,x^{})1`$, then $`|\varphi (x)\varphi (x^{})|Cd(x,x^{})^\alpha `$ for some constant $`C`$.
Note that the distance $`d(x,x^{})`$ is not defined if $`a_i=0`$ or $`a_i^{}=0`$ for some $`i=1,\mathrm{},m`$; nothing, therefore, is said about the values of a function from $`H_\alpha `$ at such points. This does not represent a problem, however, since we only need the space $`H_\alpha `$ for the Central Limit Theorem, and for for such a result we may deal with functions defined almost everywhere.
Define
$$C_{H_\alpha }(\varphi )=\underset{d(x,x^{})1}{\mathrm{max}}\frac{|\varphi (x)\varphi (x^{})|}{d(x,x^{})^\alpha }.$$
###### Theorem 7
Let $`:๐ด^\pm ()๐ด^\pm ()`$ be the Rauzy-Veech-Zorich induction map on the space of zippered rectangles and let $`\overline{\nu }_{}`$ be the absolutely continuous invariant probability measure. Let $`p>2`$. Then, for any $`\alpha >0`$, there exist positive constants $`C,\delta `$ such that for any $`\varphi ,\psi H_\alpha L_p(๐ด(),\overline{\nu }_{})`$ we have
$$|\varphi \times \psi ^{2n}๐\overline{\nu }_{}\varphi ๐\overline{\nu }_{}\psi ๐\overline{\nu }_{}|C\mathrm{exp}(\delta n^{1/6})(C_{H_\alpha }(\varphi )+|\varphi |_{L_p})(C_{H_\alpha }(\psi )+|\psi |_{L_p})$$
Theorem 7 will be established simultaneosuly with the Theorem 4. Indeed, the map $``$ can be almost surely identified with the natural extension of the map $`๐ข`$, and the method of Markov approximations of of Sinai and BunimovichโSinai allows to obtain the decay of correlations for the invertible case simultaneously with that for the noninvertible one.
Since the flow $`P^t`$ is a special flow over the map $``$, by the Theorem of Melbourne and Tรถrรถk , the decay of correlations for the map $``$ allows to obtain the Central Limit Theorem for the flow $`P^t`$.
Denote by $`X_t`$ the derivative with respect to the flow $`P^t`$.
###### Theorem 8
Let $`p>2`$ and let $`\varphi H_\alpha (\mathrm{\Omega }_0())L_p(\mathrm{\Omega }_0(),\mu _{})`$ satisfy $`\varphi ๐\nu =0`$. Assume that there does not exist $`\psi L_2(\mathrm{\Omega }_0(),\mu _{})`$ such that $`\varphi =X_t\psi `$. Then there exists $`\sigma >0`$ such that
$$\frac{1}{\sqrt{T}}_0^T\varphi P^t\stackrel{๐}{}๐ฉ(0,\sigma )\mathrm{as}T\mathrm{}.$$
This Theorem will be proved in Section 16.
### 1.7 Zippered rectangles and the moduli space of holomorphic differentials.
Let $`g2`$ be an integer. Take an arbitrary integer vector $`\kappa =(k_1,\mathrm{},k_\sigma )`$ such that $`k_i>0`$, $`k_1+\mathrm{}+k_\sigma =2g2`$.
Denote by $`_\kappa `$ the moduli space of Riemann surfaces of genus $`g`$ endowed with a holomorphic differential of area $`1`$ with singularities of orders $`k_1,\mathrm{},k_\sigma `$. (the stratum in the moduli space of holomorphic differentials). Denote by $`g_t`$ the Teichmรผller flow on $`_\kappa `$ (see , , , ). The flow $`g_t`$ preserves a natural absolutely continuous probability measure on $`_\kappa `$ (,, ). We denote that measure by $`\mu _\kappa `$.
A zippered rectangle naturally defines a Riemann surface endowed with a holomorphic differential of area $`1`$. The orders of the singularities of $`\omega `$ are uniquely defined by the Rauzy class of the permutation $`\pi `$ ().
For any $``$ we thus have a map
$$\pi _{}:\mathrm{\Omega }_{}_\kappa ,$$
where $`\kappa `$ is uniquely defined by $``$.
Veech proved
###### Theorem 9 (Veech)
1. The set $`\pi _0(\mathrm{\Omega }_0())`$ is a connected component of $`_\kappa `$. Any connected component of any $`_\kappa `$ has the form $`\pi _0(\mathrm{\Omega }_0())`$ for some $``$.
2. The map $`\pi _0`$ is finite-to-one and almost everywhere locally bijective.
3. $`\pi _0(๐ฐx)=\pi _0(x)`$.
4. The flow $`P^t`$ on $`\mathrm{\Omega }_0()`$ projects under $`\pi _0`$ to the Teichmรผller flow $`g_t`$ on the corresponding connected component of $`_\kappa `$.
5. $`(\pi _{})_{}\mu _\kappa =\mu _{}`$.
A detailed treatment of the relationship between Rauzy classes, zippered rectangles and connected components is given by M.Kontsevich and A.Zorich in .
Say that a function $`\psi :_\kappa `$ is Hรถlder in the sense of Veech if there exists a Hรถlder function $`\varphi :\mathrm{\Omega }_0()`$ such that $`\psi \pi _0=\varphi `$.
Remark. This definition has a natural interpretation in terms of cohomological coordinates of Hubbard and Masur . Indeed, under the map $`\pi _0`$ the Veech coordinates on the space of zippered rectangles correspond, upto a linear change of variables, to the cohomological coordinates of Hubbard and Masur. Locally, one can associate a Hilbert metric to those coordinates. A function Hรถlder in the sense of Veech if and only if it is Hรถlder with respect to that metric. Note that the thus defined local Hilbert distance between two elements in $`_\kappa `$ majorates the Teichmรผller distance between their underlying surfaces. Therefore, if a function $`\varphi :_\kappa `$ is a lift of a smooth function from the underlying moduli space $`_g`$ of compact surfaces of genus $`g`$, then $`\varphi `$ is Hรถlder in the sense of Veech.
Denote by $`๐ณ_t`$ the derivative in the direction of the flow $`g_t`$.
Theorem 8 and Theorem 9 imply the following
###### Theorem 10
Let $``$ be a connected component of $`_\kappa `$. Let $`p>2`$, and let $`\psi L_p(,\mu _\kappa )`$ be Hรถlder in the sense of Veech and satisfy $`\varphi ๐\mu _\kappa =0`$. Assume that there does not exist $`\psi L_2(,\mu _\kappa )`$ such that $`\varphi =๐ณ_t\psi `$. Then there exists $`\sigma >0`$ such that
$$\frac{1}{\sqrt{T}}_0^T\varphi g_t๐t\stackrel{๐}{}๐ฉ(0,\sigma )\mathrm{as}T\mathrm{}.$$
### 1.8 Outline of the Proof of Theorem 4.
First, one takes a subset of the space $`\mathrm{\Delta }()`$ such that the induced map of $`๐ข`$ is uniformly expanding (namely, the set of all interval exchanges such that the renormalization matrix for them is a fixed matrix all whose elements are positive, see Proposition 4; note that the return map on such a subset is an essential element in Veechโs proof of unique ergodicity ). Then one estimates the statistics of return times in this subset, in the spirit of Lai-Sang Young . After that, the method of Markov approximations, due to Sinai , Bunimovich and Sinai , is used to complete the proof.
The paper is organized as follows. In Section 2, we state auxiliary propositions about unimodular matrices. In Section 3, following Veech and Zorich , we construct symbolic dynamics for the Rauzy-Veech-Zorich induction map $`๐ข`$, compute its transition probabilities in the sense of Sinai , and identify the natural extension of $`๐ข`$ with $``$. In Section 4, we establish the exactness of $`๐ข^2`$. In Section 6, we state the main Lemma 4, whose proof takes Sections 6 โ 10. In the remainder of the paper we apply the Markov approximation method of Sinai , Bunimovich and Sinai , in order to obtain the decay of correlations for $`๐ข`$ and $``$. In the final Section, we apply the Theorem of Melbourne and Tรถrรถk to obtain the Central Limit Theorem for the Teichmรผller flow.
## 2 Matrices
Let $`A`$ be an $`m\times m`$-matrix with positive entries.
Denote
$$|A|=\underset{i,j=1}{\overset{m}{}}A_{ij}$$
$$col(A)=\underset{i,j,k}{\mathrm{max}}\frac{A_{ij}}{A_{kj}},$$
$$row(A)=\underset{i,j,k}{\mathrm{max}}\frac{A_{ij}}{A_{ik}}$$
###### Proposition 2
Let $`Q`$ be a matrix with positive entries, $`A`$ a matrix with nonnegative entries without zero columns or rows.
Then all entries of the matrices $`AQ`$ and $`QA`$ are positive, and, moreover, we have
$$row(AQ)row(Q),col(QA)col(Q)$$
###### Corollary 2
Let $`Q`$ be a matrix with positive entries, $`A`$ a matrix with nonnegative entries without zero columns or rows.
$$row(QAQ)row(Q),col(QAQ)col(Q)$$
Let $`A`$ be an $`m\times m`$ matrix with nonnegative entries and determinant $`1`$. Consider the map $`J_A:\mathrm{\Delta }_{m1}\mathrm{\Delta }_{m1}`$ given by
$$J_A(\lambda )=\frac{A\lambda }{|A\lambda |}.$$
Then
$$detDJ_A(\lambda )=\frac{1}{|A\lambda |^m}.$$
(5)
Suppose all entries of $`A`$ are positive; then, for any $`\lambda ,\lambda ^{}\mathrm{\Delta }_{m1}`$, we have
$$row(A)^m\frac{detDJ_A(\lambda )}{detDJ_A(\lambda ^{})}row(A)^m,$$
(6)
whence we have the following
###### Proposition 3
Let $`C\mathrm{\Delta }_{m1}`$ and let $`A`$ be a matrix with positive entries and determinant $`1`$. Then
$$row(A)^m\frac{๐ฆ(C_1)}{๐ฆ(C_2)}\frac{๐ฆ(J_A(C_1))}{๐ฆ(J_A(C_2))}row(A)^m\frac{๐ฆ(C_1)}{๐ฆ(C_2)}.$$
We also note the following well-known Lemma (see, for example, ):
###### Lemma 1
Suppose all entries of the matrix $`A`$ are positive. Then the map $`J_A`$ is uniformly contracting with respect to the Hilbert metric.
## 3 Symbolic dynamics for $`๐ข`$.
First, following Veech and Zorich , we describe a Markov partition and a symbolic dynamics for the map $`๐ข^2`$, then we identify almost surely the induction map $``$ on the space of zippered rectangles with the natural extension of $`๐ข`$, and, finally, we compute for $`๐ข`$ its transition probabilities in the sense of Sinai .
### 3.1 The alphabet
Let $`\pi `$, and let $`n`$ be a positive integer.
Set
$$\mathrm{\Lambda }(a,n,\pi )=\{\lambda :\mathrm{there}\mathrm{exists}(\lambda ^{},\pi ^{})\mathrm{such}\mathrm{that}\lambda ^{}\mathrm{\Delta }_\pi ^{}^+\mathrm{and}(\lambda ,\pi )=t_{a^n}(\lambda ^{},\pi ^{})\}$$
$$\mathrm{\Delta }(a,n,\pi )=\{(\lambda ,\pi ),\lambda \mathrm{\Lambda }(a,n,\pi )\}$$
In other words, $`\mathrm{\Delta }(a,n,\pi )`$ is the set of interval exchange transformations such that the application of the Zorich induction results in the application of the $`a`$-operation $`n`$ times.
The sets $`\mathrm{\Delta }(a,n,\pi )`$ and $`\mathrm{\Delta }(a,n^{},\pi ^{})`$ are disjoint unless $`n=n^{}`$, $`\pi =\pi ^{}`$, and
$$\mathrm{\Delta }_\pi ^{}=_{n=1}^{\mathrm{}}\mathrm{\Delta }(a,n,\pi )$$
up to a set of measure zero (namely, a union of countably many hyperplanes on which Zorich induction is not defined).
If $`\pi ^{}=a^n\pi `$, then we have
$$๐ข\mathrm{\Delta }(a,n,\pi )=\mathrm{\Delta }_\pi ^{}^+.$$
Similarly, for $`\pi `$, and $`n`$ a positive integer, set
$$\mathrm{\Lambda }(b,n,\pi )=\{\lambda :\mathrm{there}\mathrm{exists}(\lambda ^{},\pi ^{})\mathrm{such}\mathrm{that}\lambda ^{}\mathrm{\Delta }_\pi ^{}^{}\mathrm{and}(\lambda ,\pi )=t_{b^n}(\lambda ^{},\pi ^{})\}.$$
$$\mathrm{\Delta }(b,n,\pi )=\{(\lambda ,\pi ),\lambda \mathrm{\Lambda }(b,n,\pi )\}.$$
In other words, $`\mathrm{\Delta }(b,n,\pi )`$ is the set of interval exchange transformations such that the application of the Zorich induction results in the application of the $`b`$-operation $`n`$ times.
The sets $`\mathrm{\Delta }(b,n,\pi )`$ and $`\mathrm{\Delta }(b,n^{},\pi ^{})`$ are disjoint unless $`n=n^{}`$, $`\pi =\pi ^{}`$, and
$$\mathrm{\Delta }_\pi ^+=_{n=1}^{\mathrm{}}\mathrm{\Delta }(b,n,\pi )$$
up to a set of measure zero (namely, a union of countably many hyperplanes on which the Zorich induction is not defined).
If $`\pi ^{}=b^n\pi `$, then, clearly,
$$๐ข(\mathrm{\Delta }(b,n,\pi ))=\mathrm{\Delta }_\pi ^{}^{}.$$
Note that the sets $`\mathrm{\Delta }(a,n,\pi )`$ and $`\mathrm{\Delta }(b,n^{},\pi ^{})`$ are always disjoint, since we have $`\mathrm{\Delta }(a,n,\pi )\mathrm{\Delta }_\pi ^{}`$, $`\mathrm{\Delta }(b,n^{},\pi ^{})\mathrm{\Delta }_\pi ^{}^+`$.
The sets $`\mathrm{\Delta }(a,n,\pi )`$, $`\mathrm{\Delta }(b,n,\pi )`$, for all $`n>0`$ and all $`\pi `$, form a Markov partition for $`๐ข`$.
### 3.2 Words
Consider the alphabet
$$๐=\{(c,n,\pi ),c=a\mathrm{or}b\}$$
For $`w_1๐`$, $`w_1=(c_1,n_1,\pi _1)`$, we write $`c_1=c(w_1),\pi _1=\pi (w_1),n_1=n(w_1)`$.
For $`w_1,w_2๐`$, $`w_1=(c_1,n_1,\pi _1)`$, $`w_2=(c_2,n_2,\pi _2)`$, define the function $`B(w_1,w_2)`$ in the following way: $`B(w_1,w_2)=1`$ if $`c_1^{n_1}\pi _1=\pi _2`$ and $`c_1c_2`$ and $`B(w_1,w_2)=0`$ otherwise.
Let
$$W_{๐,B}=\{w=w_1\mathrm{}w_n,w_i๐,B(w_i,w_{i+1})=1\mathrm{for}\mathrm{all}i=1,\mathrm{},n\}.$$
For $`w_1๐`$, $`w_1=(c_1,n_1,\pi _1)`$, set
$$A(w)=A(c_1,c_1^{n_1}\pi _1)\mathrm{}A(c_1,c_1^1\pi _1)A(c_1,\pi _1),$$
and for $`wW_{๐,B}`$, $`w=w_1\mathrm{}w_n`$, set
$$A(w)=A(w_1)\mathrm{}A(w_n).$$
Also, for $`w_1๐`$, $`\pi `$, set $`w_1^1\pi =c_1^{n_1}\pi `$, and for $`wW_{๐,B}`$, $`w=w_1\mathrm{}w_n`$, set
$$w^1\pi =w_1^1\mathrm{}w_n^1\pi .$$
For $`wW_{๐,B}`$, define a map $`t_w:\mathrm{\Delta }()\mathrm{\Delta }()`$ by
$$t_w(\lambda ,\pi )=(\frac{A(w)\lambda }{|A(w)\lambda |},w^1\pi )$$
Consider also the map
$$T_w(\lambda ,\pi )=(A(w)\lambda ,w^1\pi )$$
For $`w_1๐`$, $`w_1=(c_1,n_1,\pi _1)`$, we write $`\mathrm{\Delta }(w_1)=\mathrm{\Delta }(c_1,n_1\pi _1)`$.
For $`wW_{๐,B}`$, $`w=w_1\mathrm{}w_n`$, denote
$$\mathrm{\Delta }(w)=t_w(\mathrm{\Delta }()).$$
Then, by definition,
$$\mathrm{\Delta }(w)=\{(\lambda ,\pi ):(\lambda ,\pi )\mathrm{\Delta }(w_1),๐ข(\lambda ,\pi )\mathrm{\Delta }(w_2),\mathrm{},๐ข^{n1}(\lambda ,\pi )\mathrm{\Delta }(w_n)\}.$$
Say that $`w_1๐`$ is compatible with $`(\lambda ,\pi )\mathrm{\Delta }()`$ if
1. either $`\lambda \mathrm{\Delta }_\pi ^+`$, $`c_1=a`$, and $`a^{n_1}\pi _1=\pi `$
2. or $`\lambda \mathrm{\Delta }_\pi ^{}`$, $`c_1=b`$, and $`b^{n_1}\pi _1=\pi `$.
Say that a word $`wW_{๐,B}`$, $`w=w_1\mathrm{}w_n`$ is compatible with $`(\lambda ,\pi )`$ if $`w_n`$ is compatible with $`(\lambda ,\pi )`$.
We can write
$$๐ข^n(\lambda ,\pi )=\{t_w(\lambda ,\pi ):|w|=n\mathrm{and}w\mathrm{is}\mathrm{compatible}\mathrm{with}(\lambda ,\pi )\}.$$
Suppose that a word $`wW_{๐,B}`$ is compatible with both $`(\lambda ,\pi )`$ and $`(\lambda ^{},\pi )`$. Then
$$d(t_w(\lambda ,\pi ),t_w(\lambda ^{},\pi ))d((\lambda ,\pi ),(\lambda ^{},\pi ^{})).$$
If, moreover, all entries of the the matrix $`A(w)`$ are positive, then, by Lemma 1, there exists $`\alpha (w)`$, $`0<\alpha (w)<1`$, such that
$$d(t_w(\lambda ,\pi ),t_w(\lambda ^{},\pi ))\alpha (w)d((\lambda ,\pi ),(\lambda ^{},\pi ^{})).$$
We therefore have
###### Proposition 4
Let $`wW_{๐,B}`$ be such that all entries of the matrix $`A(w)`$ are positive. Then the return map of $`๐ข`$ on $`\mathrm{\Delta }(w)`$ is uniformly expanding with respect to the Hilbert metric.
### 3.3 Sequences
Now let
$$\mathrm{\Omega }_{๐,B}=\{\omega =\omega _1\mathrm{}\omega _n\mathrm{},\omega _n๐,B(\omega _n,\omega _{n+1})=1\mathrm{for}\mathrm{all}n\}$$
and
$$\mathrm{\Omega }_{๐,B}^{}=\{\omega =\mathrm{}\omega _n\mathrm{}\omega _1\mathrm{}\omega _n\mathrm{},\omega _n๐,B(\omega _n,\omega _{n+1})=1\mathrm{for}\mathrm{all}n\}$$
Denote by $`\sigma `$ the shift on both these spaces.
There is a natural map $`\mathrm{\Phi }:\mathrm{\Delta }\mathrm{\Omega }_{๐,B}`$ given by the formula
$$\mathrm{\Phi }(\lambda ,\pi )=\omega _1\mathrm{}\omega _n\mathrm{}$$
if
$$๐ข^n(\lambda ,\pi )\mathrm{\Delta }(\omega _n)$$
The measure $`\nu `$ projects under $`\mathrm{\Phi }`$ to a $`\sigma `$-invariant measure on $`\mathrm{\Omega }_{๐,B}`$; probability with respect to that measure will be denoted by $``$.
For $`wW_{๐,B}`$, $`w=w_1\mathrm{}w_n`$, let
$$C(w)=\{\omega \mathrm{\Omega }_{๐,B}:\omega _1=w_1,\mathrm{},\omega _n=w_n\}.$$
We have then
$$\mathrm{\Delta }(w)=\mathrm{\Phi }^1(C(w)).$$
W. Veech has proved the following
###### Proposition 5
The map $`\mathrm{\Phi }`$ is $`\nu `$-almost surely bijective.
We thus obtain a symbolic dynamics for the map $`๐ข`$.
### 3.4 The natural extension.
Consider the natural extension for the map $`๐ข`$.
The phase space is the space of sequences of interval exchanges; it will be convenient to number them by negative integers. We set:
$$\overline{\mathrm{\Delta }}()=$$
$$\{๐ฑ=\mathrm{}(\lambda (n),\pi (n)),\mathrm{},(\lambda (0),\pi (0))|๐ข(\lambda (n),\pi (n))=(\lambda (1n),\pi (1n)),n=1,\mathrm{}\}$$
The map $`๐ข`$ and the invariant measure $`\nu `$ are extended to $`\overline{\mathrm{\Delta }}`$ in the natural way. We shall still denote the probability with respect to the extended measure by $``$.
We extend the map $`\mathrm{\Phi }`$ to a map
$$\overline{\mathrm{\Phi }}:\overline{\mathrm{\Delta }}\mathrm{\Omega }_{๐,B}^{},$$
$$\overline{\mathrm{\Phi }}(\lambda )=\mathrm{}\omega _n\mathrm{}\omega _0\mathrm{}\omega _n\mathrm{},$$
if $`(\lambda (n),\pi (n))\mathrm{\Delta }(\omega _n)`$, and $`๐ข^n(\lambda (0),\pi (0))\mathrm{\Delta }(\omega _n).`$
Now take a zippered rectangle $`x\mathrm{\Omega }()`$, $`x=(\lambda ,h,a,\pi )`$. Set $`^n(x)=(\lambda (n),h(n),a(n),\pi (n))`$.
Consider a map
$$\stackrel{~}{\mathrm{\Phi }}:๐ด()\mathrm{\Omega }_{๐,B}^{},$$
(7)
given by
$$(\lambda ,h,a,\pi )\mathrm{}\omega _n\mathrm{}\omega _0\mathrm{}\omega _n\mathrm{},$$
where
$$(\lambda (n),\pi (n))\mathrm{\Delta }(\omega _n)$$
for all $`n`$.
Under the natural projection $`(\lambda ,h,a,\pi )(\lambda ,\pi )`$, the $``$-invariant measure $`\overline{\nu }`$ on $`๐ด^\pm ()`$ is mapped to the $`๐ข`$-invariant measure $`\nu `$ on $`\mathrm{\Delta }()`$, whence the measure $`\stackrel{~}{\mathrm{\Phi }}_{}\overline{\nu }`$ is exactly the probability measure $``$ on the space of bi-infinite sequences. To complete the identification of the spaces $`(๐ด^\pm (),\overline{\nu })`$ and $`(\mathrm{\Omega }_{๐,B}^{},)`$, it remains to show that almost surely there is at most one zippered rectangle corresponding to a given symbolic sequence.
###### Proposition 6
Let $`๐ช๐ฒ_{๐,B}`$ be such that all entries of the matrix $`A(๐ช)`$ are positive. Let $`\omega \mathrm{\Omega }_{๐,B}^{}`$ be such that the word $`๐ช`$ occurs infinitely many times in $`\omega `$. Then there exists at most one zippered rectangle corresponding to $`\omega `$.
Proof. Write
$$\omega =\mathrm{}\omega _n\mathrm{}\omega _0\mathrm{}\omega _n\mathrm{},$$
and let $`(\lambda ,h,a,\pi )`$ be a zippered rectangle corresponding to $`\omega `$; we want to show that $`(\lambda ,h,a,\pi )`$ is uniquely defiend by $`\omega `$.
First, $`(\lambda ,\pi )`$ is uniquely defined by the โfutureโ $`\omega _0\mathrm{}\omega _n\mathrm{}`$ of $`\omega `$.
Denote $`w(n)=\omega _n\mathrm{}\omega _0`$, $`(\lambda (n),h(n),a(n),\pi (n)=^n(\lambda ,h,a,\pi )`$.
For any $`n`$ , the interval exchange $`(\lambda (n),\pi (n))`$ corresponds to the symbolic sequence $`\omega _n\mathrm{}\omega _0\mathrm{}`$, and, again, is uniquely defined by that sequence.
By definition of the map $``$, we have
$$\lambda (n)=\frac{A(w(n))\lambda }{|A(w(n))\lambda |},h(n)=(A(w(n))^t)^1h|A(w(n)\lambda |.$$
Projectively, therefore, we have
$$_+hA(w(n))^t_+^m.$$
Since the subword $`๐ช`$ occurs infinitely many times, the intersection
$$\underset{n=1}{\overset{\mathrm{}}{}}A(w(n))^t_+^m$$
consists of a single line and the vector $`h`$ is therefore uniquely determined by the condition $`<\lambda ,h>=1`$.
It remains to determine the vector $`a`$.
By definition of the map $``$, for any $`n`$ there exists an orthogonal matrix $`U(n)`$, uniquely determined by $`\omega `$, and a vector $`v(n)`$, uniquely determined by the the vectors $`h(n),\mathrm{},h(0)`$ and $`\omega `$, such that
$$\frac{U(n)a(n)+v(n)}{|A(w(n)\lambda |}=a.$$
(8)
Now let $`n`$ be a moment such that all $`\lambda (n)_i>\frac{1}{100m}`$ (there are infinitely many such moments). Then $`|a(n)_i|<100m`$ for all $`i=1,\mathrm{},m`$ and, (8) since $`|A(w(n))\lambda |\mathrm{}`$ as $`n\mathrm{}`$, (8) implies that $`a`$ is also uniquely determined by $`\omega `$.
The proof is complete.
### 3.5 Transition probabilities.
Take a sequence $`c_1\mathrm{}c_n\mathrm{}\mathrm{\Omega }_{๐,B}`$. Following Sinai , consider the transition probability
$$(\omega _1=c_1|\omega _2=c_2,\mathrm{},\omega _n=c_n,\mathrm{})=\underset{n\mathrm{}}{lim}\frac{(c_1c_2\mathrm{}c_n)}{(c_2\mathrm{}c_n)}.$$
In this subsection, we give a formula for this probability in terms of $`(\lambda ,\pi )=\mathrm{\Phi }^1(c_2\mathrm{}c_n\mathrm{})`$.
Assume $`w_1๐`$ is compatible with $`(\lambda ,\pi )`$.
Denote
$$(w_1|(\lambda ,\pi ))=(((\lambda (1),\pi (1))=t_{w_1}(\lambda (0),\pi (0))|(\lambda (0),\pi (0))=(\lambda ,\pi )).$$
If $`w_1๐`$ is compatible with $`(\lambda ,\pi )`$, from the definition of $`๐ข`$ and from (5) we have
$$(w_1|(\lambda ,\pi ))=\frac{\rho (t_{w_1}(\lambda ,\pi ))}{\rho (\lambda ,\pi )|A(w_1)\lambda |^m}$$
(9)
Since the invariant density is a homogeneous function of degree $`m`$, we have
$$\rho (T_{w_1}(\lambda ,\pi ))=\frac{\rho (t_{w_1}(\lambda ,\pi ))}{|A(w_1)\lambda |^m},$$
and we can rewrite (9) as follows:
$$(w_1|(\lambda ,\pi ))=\frac{\rho (T_{w_1}(\lambda ,\pi ))}{\rho (\lambda ,\pi )}$$
(10)
Let $`w=w_1\mathrm{}w_n`$ be compatible with $`(\lambda ,\pi )`$.
Denote
$$(w|(\lambda ,\pi ))=((\lambda (k),\pi (k))=t_{w_{nk+1}}(\lambda (1k),\pi (1k)),k=1,\mathrm{},n|(\lambda (0),\pi (0))=(\lambda ,\pi )).$$
From (9), by induction, we have
$$(w|(\lambda ,\pi ))=\frac{\rho (t_w(\lambda ,\pi ))}{\rho (\lambda ,\pi )|A(w)\lambda |^m}$$
(11)
Since the invariant density is a homogeneous function of degree $`m`$, we have
$$\rho (T_w(\lambda ,\pi ))=\frac{\rho (t_w(\lambda ,\pi ))}{|A(w)\lambda |^m},$$
and we can rewrite (11) as follows:
$$(w|(\lambda ,\pi ))=\frac{\rho (T_w(\lambda ,\pi ))}{\rho (\lambda ,\pi )}$$
(12)
###### Corollary 3
There exists $`C>0`$ such that the following is true. Suppose $`wW_{๐,B}`$ is compatible with $`(\lambda ,\pi )`$. Then
$$(w|(\lambda ,\pi ))\frac{C}{\rho (\lambda ,\pi )|A(w)|^m}$$
Proof: recall that the invariant density is a positive homogeneous function of degree $`m`$ and therefore is bounded from below: there exists $`C>0`$ such that $`\rho (\lambda ,\pi )>C`$ for all $`(\lambda ,\pi )\mathrm{\Delta }()`$. In particular, $`\rho (t_w(\lambda ,\pi ))>C`$. Substituting into (11), we obtain the result.
For $`ฯต:0<ฯต<1`$, let
$$\mathrm{\Delta }_ฯต=\{(\lambda ,\pi )\mathrm{\Delta }(),\mathrm{min}|\lambda _i|ฯต\}.$$
For any $`ฯต>0`$ there exists a constant $`C(ฯต)`$ such that for any $`(\lambda ,\pi )\mathrm{\Delta }_ฯต`$ we have $`\rho (\lambda ,\pi )<C(ฯต)`$.
###### Corollary 4
For any $`ฯต>0`$ there exists $`C(ฯต)>0`$ such that if $`(\lambda ,\pi )\mathrm{\Delta }_ฯต`$, then
$$(w|(\lambda ,\pi ))\frac{C(ฯต)}{|A(w)|^m}.$$
## 4 Proof of the Exactness
First, one notes that the discrete parameter $`\pi `$ does not give rise to any period, and then the proof follows the standard pattern : since almost every point of any measurable subset is a density point, bounded distortion estimates of Proposition 3 imply that if the measure of a tail event is positive, then it must be arbitrarily close to $`1`$.
In more detail, observe that there exists an integer $`M`$ such that for any $`n>M`$ and for any $`\pi ,\pi ^{}`$ there exist $`k_1,\mathrm{},k_{2n}`$ such that $`a^{k_1}b^{k_2}\mathrm{}a^{k_{2n1}}b^{k_{2n}}\pi =\pi ^{}`$. This follows from conmnectedness of the Rauzy graph and the fact that for any $`\pi `$ there exist $`n_1,n_2`$ such that $`a^{n_1}\pi =b^{n_2}\pi =\pi `$.
Let $`\alpha _0`$ be the partition of $`\mathrm{\Delta }^+`$ into $`\mathrm{\Delta }_\pi ^+`$, $`\pi `$, and let $`\alpha _n`$ be the partition into the cylinders $`\mathrm{\Delta }(w)`$, where $`w๐ฒ_{๐,B}`$, $`|w|=2n`$.
###### Lemma 2
There exists $`k>0`$ such that the following is true. Suppose $`C\mathrm{\Delta }^+`$, and there exists $`\pi `$ such that $`\mathrm{\Delta }_\pi ^+C`$. Then $`๐ข^{2k}C=\mathrm{\Delta }^+()`$.
This implies
###### Lemma 3
There exists $`k>0`$ such that the following holds. For any $`\epsilon >0`$ there is $`\delta >0`$ such that for any $`C\mathrm{\Delta }^+()`$ satisfying $`๐ฆ(C\mathrm{\Delta }_\pi ^+)<\delta `$, we have $`๐ฆ(๐ข^{2k}C\mathrm{\Delta }^+)<\epsilon `$.
Now suppose $`C\mathrm{\Delta }^+`$ is a $`๐ข^2`$-tail event, i.e., for any $`n>0`$ there exists $`B_n`$ such that $`C=๐ข^{2n}B_n`$ and $`0<\nu (C)<1`$. Then $`\nu (B_n)=\nu (C)`$ and, by Lemma 3, we can assume that there exists $`\epsilon >0`$ such that for any $`\pi `$, we have
$$๐ฆ((\mathrm{\Delta }^+C)\mathrm{\Delta }_\pi ^+)\epsilon $$
(13)
Let $`๐ช=q_1\mathrm{}q_l`$ be a word such that the matrix $`A(๐ช)`$ is positive.
For almost any $`(\lambda ,\pi )C`$ we have
$$\underset{n\mathrm{}}{lim}\frac{๐ฆ(\alpha _n(\lambda ,\pi )C)}{๐ฆ(\alpha _n(\lambda ,\pi ))}=1$$
(14)
Now let $`n`$ be such that $`๐ข^{2n}(\lambda ,\pi )\mathrm{\Delta }(๐ช)`$. Denote $`(\lambda ^{},\pi ^{})=๐ข^{2n}(\lambda ,\pi )`$. Let $`A`$ be the corresponding renormalization matrix, that is, $`\lambda =J_A\lambda ^{}`$. Then $`A=A_1A(๐ช)`$ for some (unimodular nonnegative integer) matrix $`A_1`$. We have $`\alpha _n(\lambda ,\pi )=J_A(\mathrm{\Delta }_\pi ^{}^+)`$. By Proposition 3, from (13), we deduce that there exists $`\epsilon ^{}`$, not depending on $`n`$ such that
$$\frac{๐ฆ(\alpha _n(\lambda ,\pi )(\mathrm{\Delta }^+C))}{๐ฆ(\alpha _n(\lambda ,\pi ))}\epsilon ^{}.$$
Since, by ergodicity, for almost any $`(\lambda ,\pi )`$ we can find infinitely many $`n`$ such that $`๐ข^{2n}(\lambda ,\pi )\mathrm{\Delta }(๐ช)`$, we arrive at a contradiction with (14), which gives the exactness of $`๐ข^2`$.
## 5 The Main Lemma
We shall suppose from now on that the Rauzy class $``$ is fixed and will often suppress it from notation.
For $`ฯต:0<ฯต<1`$, define, in the same way as above,
$$\mathrm{\Delta }_ฯต=\{(\lambda ,\pi )\mathrm{\Delta }(),\mathrm{min}|\lambda _i|ฯต\}.$$
###### Lemma 4
There exist positive constants $`\gamma ,K,p`$ such that the following is true for any $`ฯต>0`$. Suppose $`(\lambda ,\pi )\mathrm{\Delta }_ฯต`$. Then
$$\{nK|\mathrm{log}ฯต|,(\lambda (n),\pi (n))\mathrm{\Delta }_\gamma |(\lambda (1),\pi (1))=(\lambda ,\pi ))\}p.$$
From Corollary 4, we obtain
###### Corollary 5
Let $`๐ชW_{๐,B}`$, $`๐ช=q_1\mathrm{}q_l`$ be such that all entries of the matrix $`A(๐ช)`$ are positive. Then there exist positive constants $`K(๐ช),p(๐ช)`$ such that the following is true for any $`ฯต>0`$. Suppose $`(\lambda ,\pi )\mathrm{\Delta }_ฯต`$. Then
$$\{nK(๐ช)|\mathrm{log}ฯต|,(\lambda (n),\pi (n))\mathrm{\Delta }(๐ช)|(\lambda (1),\pi (1))=(\lambda ,\pi ))\}p(๐ช).$$
Informally, the proof of Lemma 4 proceeds by getting rid of small intervals.
For $`\gamma >0`$, $`km`$, denote
$$\mathrm{\Delta }_{\gamma ,k}=\{(\lambda ,\pi ):i_1,\mathrm{},i_k:\lambda _{i_1},\mathrm{},\lambda _{i_k}\gamma \}.$$
and
$$\mathrm{\Delta }_{\gamma ,k,ฯต}=\{(\lambda ,\pi ):\lambda _iฯต\mathrm{for}\mathrm{all}i=1,\mathrm{},m\mathrm{and}i_1,\mathrm{},i_k:\lambda _{i_1},\mathrm{},\lambda _{i_k}\gamma \}.$$
Lemma 4 follows from
###### Lemma 5
There exist constants $`L,K,p`$, depending only on the Rauzy class, such that the following is true for any $`\gamma ,k,ฯต`$.
Assume $`(\lambda ,\pi )\mathrm{\Delta }_{\gamma ,k,ฯต}`$.
Then
$$\{nK|\mathrm{log}ฯต|:(\lambda (n),\pi (n))\mathrm{\Delta }_{\gamma /L,k+1,ฯต/L}|(\lambda (1),\pi (1))=(\lambda ,\pi ))\}p.$$
Lemma 5 is proved in the next four sections.
## 6 An estimate on the number of Rauzy operations.
Recall that, if $`(\lambda ,\pi )\mathrm{\Delta }^+`$, then the $`๐ข`$-preimages of $`(\lambda ,\pi )`$ are the exchanges $`t_{a^n}(\lambda ,\pi )`$, $`n=1,\mathrm{}`$. whereas if $`(\lambda ,\pi )\mathrm{\Delta }^{}`$, then the $`๐ข`$-preimages of $`(\lambda ,\pi )`$ are the exchanges $`t_{b^n}(\lambda ,\pi )`$, $`n=1,\mathrm{}`$.
Denote
$$๐ฉ_n(\lambda ,\pi )=\{\begin{array}{cc}((\lambda (1),\pi (1))=t_{a^n}(\lambda ,\pi )|(\lambda (0),\pi (0))=(\lambda ,\pi )),\hfill & \text{if }(\lambda ,\pi )\mathrm{\Delta }^+\text{ ;}\hfill \\ ((\lambda (1),\pi (1))=t_{b^n}(\lambda ,\pi )|(\lambda (0),\pi (0))=(\lambda ,\pi )),\hfill & \text{if }(\lambda ,\pi )\mathrm{\Delta }^{}\text{.}\hfill \end{array}$$
For $`\lambda _+^m`$, set
$$T_{a^1}^{(\pi )}(\lambda )=A(a^1\pi ,a)\lambda ,t_{a^1}^{(\pi )}(\lambda )=\frac{A(a^1\pi ,a)\lambda }{|A(a^1\pi ,a)\lambda |},$$
$$T_{b^1}^{(\pi )}(\lambda )=A(b^1\pi ,b)\lambda ,t_{b^1}^{(\pi )}(\lambda )=\frac{A(b^1\pi ,b)\lambda }{|A(b^1\pi ,b)\lambda |},$$
and
$$T_{a^n}^{(\pi )}(\lambda )=T_{a^1}^{(a^{1n}\pi )}\mathrm{}T_{a^1}^{(\pi )}\lambda ,t_{a^n}^{(\pi )}(\lambda )=t_{a^1}^{(a^{1n}\pi )}\mathrm{}t_{a^1}^{(\pi )}\lambda ,$$
$$T_{b^n}^{(\pi )}(\lambda )=T_{b^1}^{(b^{1n}\pi )}\mathrm{}T_{b^1}^{(\pi )}\lambda ,t_{b^n}^{(\pi )}(\lambda )=t_{b^1}^{(b^{1n}\pi )}\mathrm{}t_{b^1}^{(\pi )}\lambda ,$$
so that we have
$$t_{a^n}(\lambda ,\pi )=(t_{a^n}^{(\pi )}\lambda ,a^n\pi ),T_{a^n}(\lambda ,\pi )=(T_{a^n}^{(\pi )}\lambda ,a^n\pi ),$$
$$t_{b^n}(\lambda ,\pi )=(t_{b^n}^{(\pi )}\lambda ,b^n\pi ),T_{b^n}(\lambda ,\pi )=(T_{b^n}^{(\pi )}\lambda ,b^n\pi ).$$
###### Lemma 6
If $`(\lambda ,\pi )\mathrm{\Delta }^+`$, then, for any $`N1`$, we have
$$\underset{n=N+1}{\overset{\mathrm{}}{}}๐ฉ_n(\lambda ,\pi )=\frac{\rho _{a^N\pi }^+(T_{a^N}^{(\pi )}(\lambda ))}{\rho _\pi ^+(\lambda )}$$
If $`(\lambda ,\pi )\mathrm{\Delta }^{}`$, then, for any $`N1`$, we have
$$\underset{n=N+1}{\overset{\mathrm{}}{}}๐ฉ_n(\lambda ,\pi )=\frac{\rho _{b^N\pi }^{}(T_{b^N}^{(\pi )}(\lambda ))}{\rho _\pi ^{}(\lambda )}$$
Proof: We only consider the case $`(\lambda ,\pi )\mathrm{\Delta }^+`$. In this case, the formula (10) can be written as
$$๐ฉ_n(\lambda ,\pi )=\frac{\rho _{a^n\pi }^{}(T_{a^n}^{(\pi )}\lambda )}{\rho _\pi ^+(\lambda )},$$
whence we can write
$$\rho _\pi ^+(\lambda )=\underset{n=1}{\overset{\mathrm{}}{}}\rho _{a^n\pi }^{}(T_{a^n}^{(\pi )}\lambda ).$$
(15)
Note that this formula is true for any permutation $`\pi `$ and any $`\lambda `$ (i.e., even if $`\lambda \mathrm{\Delta }_\pi ^+`$, the formula, being an identity between rational functions, still holds).
Since, for any $`\lambda `$, we have
$$T_{a^{nN}}^{(\pi )}\lambda =T_{a^n}^{(a^N\pi )}(T_{a^N}^{(\pi )}\lambda ),$$
from (15) we obtain
$$\rho _{a^N\pi }^+(T_{a^N}^{(\pi )}\lambda )=\underset{n=1}{\overset{\mathrm{}}{}}\rho _{a^{nN}\pi }^{}T_{a^{nN}}^{(\pi )}\lambda =\rho _\pi ^+(\lambda )(\underset{n=N+1}{\overset{\mathrm{}}{}}๐ฉ_n(\lambda ,\pi )),$$
and the Lemma is proved.
### 6.1 Bounded growth
Let $`(\lambda ,\pi )\mathrm{\Delta }()`$.
Define
$$(\lambda ^{(n)},\pi ^{(n)})=\{\begin{array}{cc}t_{a^n}(\lambda ,\pi )),\hfill & \text{if }(\lambda ,\pi )\mathrm{\Delta }^+\text{ ;}\hfill \\ t_{b^n}(\lambda ,\pi ),\hfill & \text{if }(\lambda ,\pi )\mathrm{\Delta }^{}\text{.}\hfill \end{array}$$
$$(\mathrm{\Lambda }^{(n)},\pi ^{(n)})=\{\begin{array}{cc}T_{a^n}(\lambda ,\pi )),\hfill & \text{if }(\lambda ,\pi )\mathrm{\Delta }^+\text{ ;}\hfill \\ T_{b^n}(\lambda ,\pi ),\hfill & \text{if }(\lambda ,\pi )\mathrm{\Delta }^{}\text{.}\hfill \end{array}$$
We have
$$๐ข^1(\lambda ,\pi )=\{(\lambda ^{(n)},\pi ^{(n)}),n=1,\mathrm{}\}.$$
and
$$๐ฉ_n=((\lambda (1),\pi (1))=(\lambda ^{(n)},\pi ^{(n)})|((\lambda (0),\pi (0))=(\lambda ,\pi )).$$
For any $`n`$, there exists $`i(n)\{1,\mathrm{},m\}`$ such that
$$|\mathrm{\Lambda }^{(n)}||\mathrm{\Lambda }^{(n1)}|=\lambda _{i(n)}.$$
If $`(\lambda (1),\pi (1))`$ is a $`๐ข`$-preimage of $`(\lambda ,\pi )`$ and $`(\lambda (1),\pi (1)=t_{c^n}(\lambda ,\pi )`$, $`c=a`$ or $`b`$, then we define a vector $`\mathrm{\Lambda }(1)`$ by the relation $`(\mathrm{\Lambda }(1),\pi (1)=T_{c^n}(\lambda ,\pi )`$ (in other words, $`(\mathrm{\Lambda }(1),\pi (1))`$ is the Zorich preimage without normalization).
###### Lemma 7
There exists a constant $`C()`$, depending on the Rauzy class only, such that for any $`(\lambda ,\pi )\mathrm{\Delta }()`$ we have
$$(|\mathrm{\Lambda }(1)|>K|(\lambda (0),\pi (0)=(\lambda ,\pi ))<\frac{C()}{K2}.$$
For definiteness, assume $`\lambda \mathrm{\Delta }_\pi ^{}`$ (the proof is completely identical in the other case). Then $`๐ข`$-preimages of $`(\lambda ,\pi )`$ are $`(\lambda ^{(n)},\pi ^{(n)})=t_{b^n}(\lambda ,\pi )`$, $`n=1,2,\mathrm{}`$.
By construction , the invariant density $`\rho _\pi ^{}`$ has the form
$$\rho _\pi ^{}(\lambda )=\underset{i=1}{\overset{N}{}}\frac{1}{l_{i1}(\lambda )l_{i2}(\lambda )\mathrm{}l_{im}(\lambda )},$$
where the functions $`l_{ij}`$ are linear:
$$l_{ij}(\lambda )=a_{ij}^{(1)}\lambda _1+\mathrm{}+a_{ij}^{(m)}\lambda _m,$$
and all $`a_{ij}^{(r)}`$ are nonnegative (in fact, $`a_{ij}^{(r)}=0`$ or $`1`$, but we do not need this fact here).
Let $`l`$ be the length of the $`a`$-cycle of $`\pi `$, that is, the smallest such number that $`a^l\pi =\pi `$.
Since for any $`k>0`$ we have $`a^{kl}\pi =\pi `$, from Lemma 6 we obtain
$$\underset{n=kl+1}{\overset{\mathrm{}}{}}๐ฉ_n(\lambda ,\pi )=\frac{\rho _\pi ^{}(\mathrm{\Lambda }^{(kl)})}{\rho _\pi ^{}(\lambda )}.$$
As noted above, for any $`n>0`$ there exists $`\lambda _{i(n)}`$ such that
$$|\mathrm{\Lambda }^{(n)}||\mathrm{\Lambda }^{(n1)}|=\lambda _{i(n)},$$
and, in fact,
$$\mathrm{\Lambda }^{(n)}=(\lambda _1,\mathrm{},\lambda _{m1},\lambda _m+\lambda _{i(1)}+\mathrm{}+\lambda _{i(n)}).$$
Since
$$\underset{n=kl+1}{\overset{\mathrm{}}{}}๐ฉ_n(\lambda ,\pi )0\mathrm{as}k\mathrm{},$$
for any $`i=1,\mathrm{},N`$ there exists $`j`$ such that $`a_{ij}^{(m)}>0`$. Renumbering, if necessary, the linear forms $`l_{ij}`$, we may assume that $`a_{i1}^{(m)}>0`$ for any $`i`$. Denote $`ฯต=\mathrm{min}a_{i1}^{(m)}`$ and $`L=\mathrm{max}a_{i1}^{(r)}`$. For any $`\lambda _+^m`$ we have then
$$ฯต\lambda _ml_{i1}(\lambda )L|\lambda |,$$
whence
$$\frac{\rho _\pi ^{}(\mathrm{\Lambda }^{(kl)})}{\rho _\pi ^{}(\lambda )}\frac{L}{ฯต(\lambda _m+\lambda _{i(1)}+\mathrm{}+\lambda _{i(kl)})}.$$
(16)
Let $`N`$ be the smallest number such that $`|\mathrm{\Lambda }(N)|>K`$ and let $`s`$ be the largest such integer that $`sl<N`$. Then $`|\mathrm{\Lambda }(sl)|>K1`$ (because all $`\lambda _{i(sl+1)},\mathrm{},\lambda _{i(N)}`$ are all distinct) and $`\lambda _m+\lambda _{i(1)}+\mathrm{}\mathrm{}+\lambda _{i(sl)}>K2`$ (because $`|\mathrm{\Lambda }(sl)|=1+\lambda _{i(1)}+\mathrm{}\mathrm{}+\lambda _{i(sl)}`$.
Therefore, by (16), we obtain
$$\frac{\rho _\pi ^{}(\mathrm{\Lambda }^{(kl)})}{\rho _\pi ^{}(\lambda )}\frac{L}{ฯต}\frac{1}{K2},$$
and the Lemma is proved.
###### Lemma 8
Suppose $`(\lambda ,\pi )\mathrm{\Delta }^+`$, and let $`l`$ be the length of the $`a`$-cycle of $`\pi `$.
Then, for any $`k1`$, we have
$$\underset{n=kl+1}{\overset{\mathrm{}}{}}๐ฉ_n(\lambda ,\pi )(\frac{\lambda _{\pi ^1m}}{\lambda _{\pi ^1m}+k})^m.$$
Suppose $`(\lambda ,\pi )\mathrm{\Delta }^{}`$, and let $`l`$ be the length of the $`b`$-cycle of $`\pi `$.
Then, for any $`k1`$, we have
$$\underset{n=kl+1}{\overset{\mathrm{}}{}}๐ฉ_n(\lambda ,\pi )(\frac{\lambda _m}{\lambda _m+k})^m.$$
Proof. Again, we only consider the case $`(\lambda ,\pi )\mathrm{\Delta }^{}`$, as the proof of the other case is identical.
$$\underset{n=kl+1}{\overset{\mathrm{}}{}}๐ฉ_n(\lambda ,\pi )=\frac{\rho _\pi ^{}(\mathrm{\Lambda }^{(kl)}))}{\rho _\pi ^{}(\lambda )}.$$
Set $`\mathrm{\Lambda }^{(kl)}=(\mathrm{\Lambda }_1^{(kl)},\mathrm{},\mathrm{\Lambda }_m^{(kl)}).`$
For $`k=1`$ we have $`\mathrm{\Lambda }_i^{(l)}=\lambda _i`$ for $`i<m`$ and $`\mathrm{\Lambda }_m^{(l)}=\lambda _m+\lambda _{i(1)}+\mathrm{}+\lambda _{i(l)}`$, and for arbitrary $`k`$ by induction we obtain $`\mathrm{\Lambda }_i^{(kl)}=\lambda _i`$ for $`i<m`$ and $`\mathrm{\Lambda }_m^{(kl)}=\lambda _m+k(\lambda _{i(1)}+\mathrm{}+\lambda _{i(l)}).`$
Note that $`\lambda _{i(1)}+\mathrm{}\lambda _{i(l)}1`$ (since $`i(1),\mathrm{},i(l)`$ are all distinct).
As in the proof of the previous Lemma, write
$$\rho _\pi ^{}(\lambda )=\underset{i=1}{\overset{N}{}}\frac{1}{l_{i1}(\lambda )l_{i2}(\lambda )\mathrm{}l_{im}(\lambda )},$$
whence
$$\frac{\rho _\pi ^{}(\mathrm{\Lambda }^{(kl)}))}{\rho _\pi ^{}(\lambda )}\underset{i}{\mathrm{min}}\frac{l_{i1}(\lambda )l_{i2}(\lambda )\mathrm{}l_{im}(\lambda )}{l_{i1}(\mathrm{\Lambda }^{(kl)})l_{i2}(\mathrm{\Lambda }^{(kl)})\mathrm{}l_{im}(\mathrm{\Lambda }^{(kl)})}.$$
(17)
For any linear form $`l(\lambda )=a_1\lambda _1+\mathrm{}+a_m\lambda _m`$, $`a_i0`$, we have
$$\frac{l(\mathrm{\Lambda }^{(kl)})}{l(\lambda )}\frac{\lambda _m}{\lambda _m+k(\lambda _{i(1)}+\mathrm{}+\lambda _{i(l)}}\frac{\lambda _m}{\lambda _m+k},$$
and the Lemma follows.
## 7 An estimate on the probability of stopping.
###### Lemma 9
For any $`\gamma >0`$, there exists $`c(\gamma )>0`$ such that if $`\lambda _{i(N)}>\gamma `$, then
$$\frac{๐ฉ_N(\lambda ,\pi )}{_{n=N+1}^{\mathrm{}}๐ฉ_n(\lambda ,\pi )}c(\gamma )$$
From Lemma 8 we immediately have the following Corollary.
###### Corollary 6
For any $`\gamma >0`$, there exists $`c(\gamma )>0`$ such that the following is true.
Assume $`(\lambda ,\pi )\mathrm{\Delta }^+`$, $`\lambda _{i(N)}>\gamma `$, $`\lambda _{\pi ^1m}>\gamma `$. Then
$$๐ฉ_N\frac{c(\gamma )}{N^m}.$$
Similarly, if $`(\lambda ,\pi )\mathrm{\Delta }^{}`$, $`\lambda _{i(N)}>\gamma `$, $`\lambda _m>\gamma `$, then
$$๐ฉ_N\frac{c(\gamma )}{N^m}.$$
If $`(\lambda ,\pi )\mathrm{\Delta }^+`$, then, by the definition of $`๐ฉ_n(\lambda ,\pi )`$ and by Lemma 6, we have
$$๐ฉ_N(\lambda ,\pi )=\frac{\rho _{a^N\pi }^{}(T_{a^N}^{(\pi )}\lambda )}{\rho _\pi ^+(\lambda )},$$
$$\underset{n=N+1}{\overset{\mathrm{}}{}}๐ฉ_n(\lambda ,\pi )=\frac{\rho _{a^N\pi }^+(T_{a^N}^{(\pi )}(\lambda ))}{\rho _\pi ^+(\lambda )},$$
and, therefore,
$$\frac{๐ฉ_N(\lambda ,\pi )}{_{n=N+1}^{\mathrm{}}๐ฉ_n(\lambda ,\pi )}=\frac{\rho _{a^N\pi }^{}(T_{a^N}^{(\pi )}\lambda )}{\rho _{a^N\pi }^+(T_{a^N}^{(\pi )}(\lambda ))}.$$
Lemma 9 follows now from the following
###### Lemma 10
For any $`\gamma >0`$ there exists a constant $`c(\gamma )>0`$ such that the following is true. Let $`(\lambda ,\pi )\mathrm{\Delta }()`$. If $`\lambda _{\pi ^1m+1}>\gamma `$, then
$$\frac{\rho _\pi ^{}(\lambda )}{\rho _\pi ^+(\lambda )}c(\gamma ).$$
If $`\lambda _{\pi ^1(\pi (m)+1)}>\gamma `$, then
$$\frac{\rho _\pi ^+(\lambda )}{\rho _\pi ^{}(\lambda )}c(\gamma ).$$
The proof of Lemma 10 will take the remainder of this section.
First, we modify Veechโs coordinates on the space of zippered rectangles. Take a zippered rectangle $`(\lambda ,h,a,\pi )\mathrm{\Delta }()`$, and introduce the vector $`\delta =(\delta _1,\mathrm{},\delta _m)^m`$ by the formula
$$\delta _i=a_{i1}a_i,i=1,\mathrm{},m$$
(here we assume, as always, $`a_0=a_{m+1}=0`$).
###### Proposition 7
The data $`(\lambda ,\pi ,\delta )`$ determine the zippered rectangle $`(\lambda ,h,a,\pi )`$ uniquely.
Remark. The coordinates $`(\lambda ,\pi ,\delta )`$ on the space of zippered rectangles have a natural interpretation in terms of the cohomological coordinates of Hubbard and Masur : namely, the $`\lambda _i`$ are the real parts of the corresponding cycles, and the $`\delta _i`$ are (minus) the imaginary parts.
Proof of Proposition 7. For any $`i=1,\mathrm{},m`$, we have
$$a_i=\delta _1\mathrm{}\delta _i,$$
(18)
so the vector $`a`$ is uniquely defined by $`\delta `$. It remains to show that the vector $`h`$ is uniquely defined by $`\delta `$, and, to do this, we shall express the $`h`$ through the $`a`$. First note that
$$h_{\pi ^1m}=a_{\pi ^1m}a_m.$$
Now, if $`i\pi ^1m`$, then $`i=\pi ^1(k1)`$ for some $`k\{1,\mathrm{},m\}`$. The equation
$$h_ia_i=h_{\pi ^1(\pi (i)+1)}a_{\pi ^1(\pi (i)+1)1}.$$
(19)
then takes the form
$$h_{\pi ^1(k1)}a_{\pi ^1(k1)}=h_{\pi ^1(k)}a_{\pi ^1(k)1},$$
or, equivalently,
$$h_{\pi ^1(k)}=a_{\pi ^1(k)1}+h_{\pi ^1(k1)}a_{\pi ^1(k1)}.$$
Since
$$h_{\pi ^11}=a_{\pi ^111},$$
by induction, we obtain
$$h_{\pi ^1k}=a_{\pi ^1k1}+\underset{l=1}{\overset{k1}{}}(a_{\pi ^1l1}a_{\pi ^1l})$$
for any $`k=1,\mathrm{},m`$, and the Lemma is proved.
The above computations give us the following expression for $`h`$ in terms of $`\delta `$:
$$h_{\pi ^1k}=\underset{i=1}{\overset{\pi ^1k1}{}}\delta _i+\underset{l=1}{\overset{k1}{}}\delta _{\pi ^1(l)}$$
(20)
or, equivalently,
$$h_r=\underset{i=1}{\overset{r1}{}}\delta _i+\underset{l=1}{\overset{\pi (r)1}{}}\delta _{\pi ^1l}.$$
(21)
Rewriting the inequalities defining the zippered rectangle in terms of $`\delta `$, we obtain by a straightforward computation the following system:
$$\delta _1+\mathrm{}+\delta _i0,i=1,\mathrm{},m1.$$
$$\delta _{\pi ^11}+\mathrm{}+\delta _{\pi ^1i}0,i=1,\mathrm{},m1.$$
The parameter $`a_m=(\delta _1+\mathrm{}+\delta _m)`$ can be both positive and negative.
Introduce the following cones in $`^m`$:
$$K_\pi =\{\delta =(\delta _1,\mathrm{},\delta _m):\delta _1+\mathrm{}+\delta _i0,\delta _{\pi ^11}+\mathrm{}+\delta _{\pi ^1i}0,i=1,\mathrm{},m1\},$$
$$K_\pi ^+=K_\pi \{\delta :\underset{i=1}{\overset{m}{}}\delta _i0\},K_\pi ^{}=K_\pi \{\delta :\underset{i=1}{\overset{m}{}}\delta _i0\}.$$
We have established the following
###### Proposition 8
For $`(\lambda ,\pi )\mathrm{\Delta }()`$ and an arbitrary $`\delta K_\pi `$ there exists a unique zippered rectangle $`(\lambda ,h,a,\pi )`$ corresponding to the parameters $`(\lambda ,\pi ,\delta )`$.
In what follows, we shall simply refer to the zippered rectangle $`(\lambda ,\pi ,\delta )`$.
Remark. It would be interesting to write down explicitly the genrating vectors for the cones $`K_\pi `$, $`K_\pi ^+`$, $`K_\pi ^{}`$; in particular, that would allow to give an explicit expression for the invariant densities of Veech and Zorich .
Denote by $`Area(\lambda ,\pi ,\delta )`$ the area of the zippered rectangle $`(\lambda ,\pi ,\delta )`$. We have:
$$Area(\lambda ,\pi ,\delta )=\underset{r=1}{\overset{m}{}}\lambda _rh_r=\underset{r=1}{\overset{m}{}}\lambda _r(\underset{i=1}{\overset{r1}{}}\delta _i+\underset{l=1}{\overset{\pi (r)1}{}}\delta _{\pi ^1l})=$$
$$\underset{i=1}{\overset{m}{}}\delta _i(\underset{r=i+1}{\overset{m}{}}\lambda _r+\underset{r=\pi (i)+1}{\overset{m}{}}\lambda _{\pi ^1r})=1.$$
(22)
A straightforward computation shows that in the coordinates $`(\lambda ,\pi ,\delta )`$ the Rauzy induction map is written as follows:
$$๐ฏ(\lambda ,\pi ,\delta )=\{\begin{array}{cc}(\frac{A(\pi ,b)^1\lambda }{|A(\pi ,b)^1\lambda |},b\pi ,A(\pi ,b)^1\delta |A(\pi ,b)^1\lambda |),\hfill & \text{if }\lambda \mathrm{\Delta }_\pi ^+\text{;}\hfill \\ (\frac{A(\pi ,a)^1\lambda }{|A(\pi ,a)^1\lambda |},a\pi ,A(\pi ,a)^1\delta |A(\pi ,a)^1\lambda |),\hfill & \text{if }\lambda \mathrm{\Delta }_\pi ^{}\text{.}\hfill \end{array}$$
For $`\lambda _+^m`$, denote
$$K(\lambda ,\pi )=K_\pi \{\delta :Area(\lambda ,\pi ,\delta )1\},$$
$$K^+(\lambda ,\pi )=K_\pi ^+\{\delta :Area(\lambda ,\pi ,\delta )1\},$$
$$K(\lambda ,\pi )=K_\pi ^{}\{\delta :Area(\lambda ,\pi ,\delta )1\}.$$
Denote by $`vol_m`$ the Lebesgue measure in $`^m`$.
Set
$$๐ซ(\lambda ,\pi )=vol_m(K(\lambda ,\pi )),๐ซ^+(\lambda ,\pi )=vol_m(K^+(\lambda ,\pi )),๐ซ^{}(\lambda ,\pi )=vol_m(K^{}(\lambda ,\pi )).$$
By definition, the functions $`๐ซ,๐ซ^+,๐ซ^{}`$ are positive rational functions, homogeneous of degree $`m`$.
###### Lemma 11
1. $`๐ซ^{}(\lambda ,\pi )=๐ซ(T_{b^1}(\lambda ,\pi ))`$.
2. $`๐ซ^+(\lambda ,\pi )=๐ซ(T_{a^1}(\lambda ,\pi ))`$.
3. $`๐ซ(\lambda ,\pi )=๐ซ(T_{a^1}(\lambda ,\pi ))+๐ซ(T_{b^1}(\lambda ,\pi ))`$.
Proof. If
$$\delta =(\delta _1,\mathrm{},\delta _m)K^{}(\lambda ,\pi ),$$
then
$$\stackrel{~}{\delta }=(\delta _1,\mathrm{},\delta _{m1},\delta _m+\delta _{\pi ^1m})K(T_{b^1}(\lambda ,\pi )),$$
and vice versa. This gives a volume-preserving bijection between $`K^{}(\lambda ,\pi )`$ and $`K(T_{b^1}(\lambda ,\pi ))`$, whence $`๐ซ^{}(\lambda ,\pi )=๐ซ(T_{b^1}(\lambda ,\pi ))`$. The second assertion is proved in the same way, and the third follows from the first two.
###### Corollary 7
$$๐ซ^+(\lambda ,\pi )=\underset{n=1}{\overset{\mathrm{}}{}}๐ซ^{}(T_{a^n}(\lambda ,\pi )).$$
$$๐ซ^{}(\lambda ,\pi )=\underset{n=1}{\overset{\mathrm{}}{}}๐ซ^+(T_{b^n}(\lambda ,\pi )).$$
We only prove the first assertion. We have
$$๐ซ^+(\lambda ,\pi )=๐ซ(T_{a^1}(\lambda ,\pi ))=๐ซ^+(T_{a^1}(\lambda ,\pi )+๐ซ^{}(T_{a^1}(\lambda ,\pi )=๐ซ(T_{a^2}(\lambda ,\pi ))+๐ซ^{}(T_{a^1}(\lambda ,\pi )).$$
Proceeding by induction,
$$๐ซ^+(\lambda ,\pi )=\underset{n=1}{\overset{N}{}}๐ซ^{}(T_{a^n}(\lambda ,\pi ))+๐ซ(T_{a^{N1}}(\lambda ,\pi )).$$
Since
$$T_{a^{N1}}(\lambda ,\pi )=(T_{a^{N1}}^{(\pi )}(\lambda ),a^{N1}\pi ),$$
and $`|T_{a^{N1}}^{(\pi )}(\lambda )|\mathrm{}`$ as $`N\mathrm{}`$, we obtain $`๐ซ(T_{a^{N1}}(\lambda ,\pi ))0`$ as $`N\mathrm{}`$, and the Corollary is proved.
Since the functions $`๐ซ,๐ซ^+,๐ซ^{}`$ are positive, rational and homogeneous of degree $`m`$, Corollary 7 implies that, for some positive constant $`C()`$, depending only on the Rauzy class $``$, we have
$$\rho ^+(\lambda ,\pi )=C(\pi )๐ซ^+(\lambda ,\pi ),\rho ^{}(\lambda ,\pi )=C(\pi )๐ซ^{}(\lambda ,\pi ).$$
By construction, for any $`\lambda _+^m`$ we have
$$๐ซ_\pi ^+(\lambda _1,\mathrm{}\lambda _m)=๐ซ_{\pi ^1}^{}(\lambda _{\pi (1)},\mathrm{}\lambda _{\pi (m)}).$$
In view of this observation, it suffices to prove only the first assertion of the Lemma 10, as the second one follows automatically.
Take $`\delta =(\delta _1,\mathrm{},\delta _m)^m`$, and, for $`\theta >0`$, define
$$J_\theta ^{(m)}\delta =(\delta _1,\mathrm{},\delta _m+\theta ),J_\theta ^{(\pi ^1m)}\delta =(\delta _1,\mathrm{},\delta _{\pi ^1m}\theta ,\mathrm{},\delta _m).$$
###### Proposition 9
Let $`\theta >0`$. If $`\delta K_\pi `$, then $`J_\theta ^{(m)}\delta K_\pi ,J_\theta ^{(\pi ^1m)}\delta K_\pi `$. If $`\delta K_\pi ^{}`$, then $`J_\theta ^{(m)}\delta K_\pi ^{}`$. If $`\delta K_\pi ^+`$, then $`J_\theta ^{(\pi ^1m)}\delta K_\pi ^+`$.
This follows directly from the definition of the cones $`K_\pi ,K_\pi ^{},K_\pi ^+`$. From (22) we obtain
$$Area(\lambda ,\pi ,J_\theta ^{(m)}\delta )=Area(\lambda ,\pi ,\delta )+\theta (\underset{r=\pi (m)+1}{\overset{m}{}}\lambda _{\pi ^1r}),$$
$$Area(\lambda ,\pi ,J_\theta ^{(\pi ^1m)}\delta )=Area(\lambda ,\pi ,\delta )+\theta (\underset{r=\pi ^1(m)+1}{\overset{m}{}}\lambda _r),$$
which implies
###### Proposition 10
$$Area(\lambda ,\pi ,\delta )Area(\lambda ,\pi ,J_\theta ^{(m)}\delta )Area(\lambda ,\pi ,\delta )+\theta |\lambda |.$$
$$Area(\lambda ,\pi ,\delta )Area(\lambda ,\pi ,J_\theta ^{(\pi ^1m)}\delta )Area(\lambda ,\pi ,\delta )+\theta |\lambda |.$$
For $`s`$ and a hyperplane of the form $`\delta +\mathrm{}+\delta _m=s`$, let $`vol_{m1}`$ stand for the induced $`(m1)`$-dimensional volume form on the hyperplane.
Denote
$$K_{s,\pi }=K_\pi \{\delta :\underset{i=1}{\overset{m}{}}\delta _i=s\},$$
$$K_s(\lambda ,\pi )=K(\lambda ,\pi )K_{s,\pi },$$
$$V_s(\lambda ,\pi )=vol_{m1}(K_s(\lambda ,\pi )).$$
Denote by $`๐_{\mathrm{๐ฆ๐๐ฑ}}^{}`$ the maximal possible value of $`\delta _1+\mathrm{}+\delta _m=a_m`$ in $`K(\lambda ,\pi )`$.
###### Proposition 11
Assume $`0s๐_{\mathrm{๐ฆ๐๐ฑ}}^{}`$. Then
$$V_s(\lambda ,\pi )(1+s)^{m1}V_0(\lambda ,\pi ).$$
Proof: Indeed, if $`(\lambda ,\pi ,\delta )V_s(\lambda ,\pi ),`$ then Proposition 10 implies
$$(\lambda ,\pi ,\frac{J_s^{(\pi ^1m)}\delta }{1+s})V_0(\lambda ,\pi ),$$
and the assertion follows.
###### Proposition 12
Assume $`s`$, $`0s1`$ is such that $`\frac{s}{1s}๐_{\mathrm{๐ฆ๐๐ฑ}}^{}`$. Then
$$(\frac{1}{1s})^{m1}V_s(\lambda ,\pi )V_0(\lambda ,\pi ).$$
Denote $`\theta =\frac{s}{1s}`$, then $`s=\frac{\theta }{1+\theta }`$. If $`(\lambda ,\pi ,\delta )V_0(\lambda ,\pi ),`$ then
$$(\lambda ,\pi ,\frac{J_\theta ^{(m)}\delta }{1+\theta })V_s(\lambda ,\pi ),$$
and, again, the assertion follows.
Propositions 11, 12 imply
###### Lemma 12
For any $`C_1>0`$ there exists $`C_2>0`$ such that the following is true.
Let $`๐_{\mathrm{๐ฆ๐๐ฑ}}^{}(\lambda ,\pi )<C_1`$. Then
$$๐ซ^{}(\lambda ,\pi )<C_2V_{m1}^0(\lambda ,\pi ).$$
Note that there exists $`ฯต>0`$, depending only on $``$ and such that for any $`(\lambda ,\pi )\mathrm{\Delta }()`$, we have $`๐_{\mathrm{๐ฆ๐๐ฑ}}^{}>ฯต`$. In conjunction with Propositions 11, 12, this implies
###### Lemma 13
There exists a constant $`C_3`$ such that for any $`(\lambda ,\pi )\mathrm{\Delta }()`$, we have
$$\rho ^{}(\lambda ,\pi )C_3V_0(\lambda ,\pi ).$$
Since $`a_mh_{\pi ^1m+1}`$, we have
$$๐_{\mathrm{๐ฆ๐๐ฑ}}^{}(\lambda ,\pi )\frac{1}{\lambda _{\pi ^1m+1}},$$
which implies the following
###### Corollary 8
For any $`C_4>0`$ there exists $`C_5>0`$ such that the following is true.
Assume $`\lambda _{\pi ^1m+1}>C_4`$. Then
$$\frac{๐ซ^{}(\lambda ,\pi )}{๐ซ^+(\lambda ,\pi )}<C_5,$$
which implies Lemma 10.
## 8 Kerckhoff names
In the following two sections, we shall use Kerckhoffโs convention of numbering the subintervals of an interval exchange ; to avoid confusion, we shall speak of Kerckhoff names of subintervals.
Take an interval exchange $`(\lambda ,\pi )`$. A Kerckhoff naming on the subintervals of $`(\lambda ,\pi )`$ is defined by an arbitrary permutation $`i_1,\mathrm{},i_m`$ of the symbols $`\{1,\mathrm{},m\}`$. Once such a permutation is given, we asign names $`I_{i_1},\mathrm{},I_{i_m}`$ to the subintervals of $`(\lambda ,\pi )`$, from the left to the right (i.e., the subinterval $`[0,\lambda _1)`$ is named $`I_{i_1}`$, the subinterval $`[\lambda _1,\lambda _1+\lambda _2)`$ is named $`I_{i_2}`$ and so forth).
A Kerckhoff naming of the subintervals of $`(\lambda ,\pi )`$ induces a naming on the subintervals of $`๐ฏ(\lambda ,\pi )`$ in the following way. Assume $`\lambda _m<\lambda _{\pi ^1m}`$ and the Rauzy operation $`a`$ was applied to $`(\lambda ,\pi )`$ in order to obtain $`๐ฏ(\lambda ,\pi )`$. Then the subintervals of $`๐ฏ(\lambda ,\pi )`$ are named, from the left to the right, by $`I_{i_1},\mathrm{},I_{i_{\pi ^1m}},I_m`$, $`I_{\pi ^1m+1},\mathrm{},I_{m1}`$. If $`\lambda _m>\lambda _{\pi ^1m}`$ and the Rauzy operation $`b`$ was applied, then the subintervals of $`๐ฏ(\lambda ,\pi )`$ are just named, as before, by $`I_{i_1},\mathrm{},I_{i_m}`$, from the left to the right. Proceeding inductively, we obtain a naming for any $`๐ข^n(\lambda ,\pi )`$. Conversely, if we have a Kerckhoff naming of subintervals of $`(\lambda ,\pi )`$, then, for any word $`w๐ฒ_{๐,B}`$ compatible with $`(\lambda ,\pi )`$, we automatically obtain a Kerckhoff naming on the subintervalsof $`t_w(\lambda ,\pi )`$ and $`T_w(\lambda ,\pi )`$.
Let $`(\lambda ,\pi )`$ be an interval exchange with a Kerckhoff naming $`I_{i_1},\mathrm{},I_{i_m}`$. If $`(\lambda ,\pi )\mathrm{\Delta }^+`$, then we say that $`I_{i_{\pi ^1m}}`$ is the subinterval in the critical position (we shall also sometimes say โin the $`a`$-critical positionโ). If $`(\lambda ,\pi )\mathrm{\Delta }^{}`$, then we say that $`I_{i_m}`$ is the subinterval in the critical position (we shall also sometimes say โin the $`b`$-critical positionโ).
## 9 Exponential growth.
Let $`๐ฑ\overline{\mathrm{\Delta }}`$, that is, $`๐ฑ=(\mathrm{},(\lambda (n),\pi (n),\mathrm{},(\lambda ,\pi ))`$, where, as usual, $`๐ข(\lambda (n),\pi (n))=(\lambda (1n),\pi (1n))`$. Define the words $`w(n)`$ by the relation $`(\lambda (n),\pi (n))=t_{w(n)}(\lambda ,\pi )`$. Set $`(\mathrm{\Lambda }(n),\pi (n))=T_{w(n)}(\lambda ,\pi )`$.
###### Lemma 14
There exists $`N`$ such that the following is true. For any $`๐ฑ\overline{\mathrm{\Delta }}()`$, there exist $`i_1,i_2\{1,\mathrm{},m\}`$ such that
$$\mathrm{\Lambda }(N)_{i_1}+\mathrm{\Lambda }(N)_{i_2}2(\lambda (0)_{i_1}+\lambda (0)_{i_2})$$
Proof:
Take a point $`x\overline{\mathrm{\Delta }}`$,
$$x=(\mathrm{},(\lambda (n),\pi (n)),\mathrm{},(\lambda ,\pi )).$$
Give Kerckhoff names $`I_1,\mathrm{},I_m`$ to the subintervals of the exchange $`(\lambda ,\pi )`$ from the left to the right, so that the length of $`I_i`$ is $`\lambda _i`$. We thus automatically obtain a Kerckhoff naming for the subintervals of $`((\lambda (n),\pi (n))`$ for any $`n`$.
Let $`I_{j_n}`$ be the critical subinterval for $`(\lambda (n),\pi (n))`$.
Consider the infinite sequence
$$I_{j_1}\mathrm{}I_{j_n}\mathrm{}.$$
(23)
Note that $`j_nj_{n+1}`$. A subword $`I_{j_k}\mathrm{}I_{j_{k+l}}`$ will be called a simple cycle if $`I_{j_k}=I_{j_{k+l}}`$ whereas $`I_{j_k},\mathrm{}I_{j_{k+l1}}`$ are all distinct. Naturally, $`lm`$. There are finitely many possible simple cycles, therefore there exists $`N`$, depending only on $`m`$, such that for any word of length $`N`$ in the alphabet $`\{I_1,\mathrm{},I_m\}`$, some simple cycle occurs at least $`m`$ times. Now take the word
$$I_{j_1}\mathrm{}I_{j_N},$$
(24)
the beginning of the sequence (23), and take a simple cycle which occurs $`m`$ times, say
$$I_{l_1}\mathrm{}I_{l_r},$$
(25)
Here, of course, $`rm`$. Now estimate the non-renormalized length of the subintervals $`I_{l_1},\mathrm{},I_{l_r}`$ ($`rm)`$. In the beginning, these are $`\lambda _{l_1},\mathrm{},\lambda _{l_r}`$. The key observation is, as usual, that the interval in critical position at a given inverse Zorich step was, at the previous step, added to the previous critical interval. After the first occurrence of the cycle (25), therefore, the (non-normalized) length of $`I_{l_1}`$ is at least $`\lambda _{l_1}+\lambda _{l_2}`$, that of $`I_{l_2}`$ is at least $`\lambda _{l_2}+\lambda _{l_3}`$ and so forth. After the second occurrence of (25), the length of $`I_{l_1}`$ is at least $`\lambda _{l_1}+\lambda _{l_2}+\lambda _{l_3}`$, that of $`I_{l_2}`$ is at least $`\lambda _{l_2}+\lambda _{l_3}+\lambda _{l_4}`$, and so forth. Finally, after the $`r`$-th occurrence of (25), the length of $`I_{l_1}`$ is not less than $`\lambda _{l_1}+\lambda _{l_2}+\mathrm{}+\lambda _{l_r}`$, that is, not less than $`2\lambda _{l_1}`$, since $`\lambda _{l_1}=\lambda _{l_r}`$. The Lemma is proven.
## 10 Proof of the Lemma 5
An informal sketch of the proof of Lemma 5. One divides the subintervals into โbigโ ones and โsmallโ ones: the aim is to obtain one more โbigโ interval. For this, one must first put a small subinterval into critical position. This is achieved by Lemma 15. In the previous ection, we have seen that the total length of the (non-renormalized) interval grows exponentially with the number of Zorich steps (with an exponent depending on $`ฯต`$). When the total length of the interval doubles, we obtain a new โbigโ subinterval.
### 10.1 Putting a small interval into critical position
Take an interval exchange $`(\lambda ,\pi )`$ and name the subintervals $`I_1,\mathrm{},I_m`$, from the right to the left.
###### Proposition 13
Any interval can be put both in the $`a`$-critical and in the $`b`$-critical position.
Proof: First note that if an interval can be put in the $`a`$-critical position, then it can also be put into the $`b`$-critical position just by performing the entire $`a`$-cycle of the corresponding permutation. Since the permutation is irreducible, it suffices to prove that, if $`I_i`$ can be put into critical position, then also all $`I_j`$ for $`j>i`$. To prove this, take the shortest word $`w`$ that puts $`I_i`$ into the $`a`$-critical position. Then, in the preimage, all $`I_j`$, $`j>i`$, still stand to the right of $`I_i`$,though perhaps in a different order (because an inversion of order between $`I_i`$ and $`I_j`$ can only happen once $`I_i`$ reaches the critical position). Therefore, we can immediately place any of the $`I_j`$, $`j>i`$, into the $`b`$-critical position, but then also into the $`a`$-critical position.
More precisely, pick a positive integer $`km`$ and a real $`\gamma >0`$. We say that we have a $`(k,\gamma )`$-big-small decomposition if the intervals of the exchange are divided into two groups: $`I_{i_1},\mathrm{},I_{i_k}`$, each of length at least $`\gamma `$, and the remaining ones (nothing is said about the length of the remaining ones).
Under the Kerckhoff convention, a big-small decomposition of $`(\lambda ,\pi )`$ is inherited by all $`t_w(\lambda ,\pi )`$ (one just takes the intervals with the same names).
###### Lemma 15
For any $`\gamma >0`$, there exist constants $`p(\gamma ),L(\gamma )`$ such that the following is true. Let $`(\lambda ,\pi )\mathrm{\Delta }_{k,\gamma }`$ with a fixed big-small decomposition. Then there exists $`w๐ฒ_{๐,B}`$ such that
1. $`(w|\lambda ,\pi )p(\gamma )`$.
2. $`|T_w(\lambda ,\pi )|<L(\gamma )`$.
3. the exchange $`t_w(\lambda ,\pi )`$ has a small interval in critical position.
Proof: Take the shortest word (in terms of the number of Zorich operations) that puts a small interval into critical position. Among all such words, pick the one that involves the smallest number of Rauzy operations. The length of this word, as well as the number of Rauzy operations involved, only depends on the Rauzy class. At each intermediate Rauzy step, all subintervals following the critical one either in the preimage or in the image must be big, otherwise there would exist a shorter word placing a small interval into critical position. Therefore, by Lemma 9 and the Corollary 6, the probability of each Zorich operation involved is bounded from below by a constant that only depends on $`\gamma `$. The Lemma is proved.
### 10.2 Completion of the proof.
Proof: Take any $`x\overline{\mathrm{\Delta }}`$. Take the first $`n`$ such that $`|\mathrm{\Lambda }(n)|>2`$. By Lemma 14, $`n<K|\mathrm{log}ฯต|`$. By Lemma 7, with positive probability depending only on $`M`$, we can assume $`|\mathrm{\Lambda }(n)|<2M`$. Consider two cases:
1. at all steps from $`1`$ to $`n`$, only small intervals were added between themselves.
2. at some step a large interval was added to a small one.
Note, that since we start with a small interval in critical position, either one or the other case holds (for, in order that a small interval be added to a big interval, a big interval must first be placed into critical position, and for that it must first be added to a small one).
In the first case, the lengths of all large intervals remain the same, and after renormalization at step $`n`$, each large interval has length at least $`\gamma /2M`$. However, since $`|\mathrm{\Lambda }(n)|>2`$, there must be another interval of length at least $`1/2mM`$, and the Lemma is proved.
In the second case, let $`n_1`$ be the first moment, at which a big interval is added to a small one. Then $`|\mathrm{\Lambda }(n_1)|<2`$, and, since at previous moments only small intervals were added between themselves, we have $`k+1`$ intervals of length at least $`\gamma /2`$, and the Lemma is proved completely.
## 11 Return times for the Teichmรผller flow.
We have in fact proven a stronger statement, namely, the following Lemma.
###### Lemma 16
For any word $`๐ช๐ฒ_{๐,B}`$ such that all entries of the matrix $`A(๐ช)`$ are positive, there exist constants $`K_0(๐ช),p(๐ช)`$, depending only on $`๐ช`$ and such that the following is true. For any $`KK_0`$ and any $`(\lambda ,\pi )\mathrm{\Delta }()`$,
$$(n:(\lambda (n),\pi (n))\mathrm{\Delta }_๐ช,|\mathrm{\Lambda }(n)|<K)|(\lambda ,\pi ))p(๐ช)$$
This statement has the following Corollary for the Teichmรผller flow on the space of zippered rectangles.
Take an arbitrary word $`๐ช=q_1\mathrm{}q_{2l+1}๐ฒ_{๐,B}`$ such that all entries of the matrix $`A(q_1\mathrm{}q_l)`$ are positive and all entries of the matrix $`A(๐ช)`$ are positive. As usually, set
$$\mathrm{\Delta }_๐ช=\{(\lambda ,\pi ):\mathrm{\Phi }(\lambda ,\pi )=\omega _1\mathrm{}\omega _n\mathrm{},\omega _1=q_1,\mathrm{},\omega _{2l+1}=q_{2l+1}\}.$$
Consider also the cylinder
$$\overline{\mathrm{\Delta }}_๐ช=\{\omega \mathrm{\Omega }_{๐,B}^{},\omega _l=q_1,\mathrm{},\omega _l=q_{2l+1}\}.$$
Consider the flow $`P^t`$ as a special flow over $`\overline{\mathrm{\Delta }}_๐ช`$. Denote the roof function of the flow by $`\tau _๐ช`$.
We shall now see that Lemma 16 implies
###### Corollary 9
There exists $`ฯต>0`$ such that
$$_{\overline{\mathrm{\Delta }}_๐ช}\mathrm{exp}(ฯต\tau _๐ช(\omega ))๐(\omega )<+\mathrm{}.$$
Take $`\omega \mathrm{\Omega }_{๐,B}^{}`$. As usually, set
$$(\lambda (n),\pi (n))=\mathrm{\Phi }^1(\omega _n\mathrm{}\omega _0\omega _1\mathrm{}),(\mathrm{\Lambda }(n),\pi (n))=T_{\omega _n\mathrm{}\omega _1\omega _0}(\lambda (0),\pi (0)).$$
Set $`n_๐ช(\omega )`$ to be the smallest $`n`$ such that
$$\omega _n=q_1,\mathrm{},\omega _{n+2l}=q_{2l+1}.$$
Finally, set $`L_๐ช(\omega )=\mathrm{log}|\mathrm{\Lambda }(n_๐ช(\omega ))|`$. Informally, $`L_๐ช(\omega )`$ is the โTeichmรผller flow timeโ it takes $`\omega `$ to reach $`\overline{\mathrm{\Delta }}_๐ช`$.
To establish the Corollary 9, it suffices to prove
###### Proposition 14
There exists $`ฯต>0`$ such that
$$_{\mathrm{\Omega }_{๐,B}^{}}\mathrm{exp}(ฯตL_๐ช(\omega ))๐(\omega )<+\mathrm{}.$$
Proof of Proposition 14. Our main tool will be Lemma 16. Take a $`K>K_0`$ such that $`1p(๐ช)+\frac{1}{K}<1`$. Define a random time $`k_1(\omega )`$ to be the first moment $`n`$ such that $`|\mathrm{\Lambda }(n)(\omega )|>K`$. Note that the map
$$\stackrel{~}{\sigma }(\omega )\sigma ^{k_1(\omega )}(\omega )$$
is invertible (here, as always, $`\sigma `$ is the shift on $`\mathrm{\Omega }_{๐,B}`$).
Introduce a function $`\eta :\mathrm{\Omega }_{๐,B}^{}`$ by the formula
$$\eta (\omega )=[\frac{\mathrm{log}|\mathrm{\Lambda }(k_1(\omega ))|}{\mathrm{log}K}].$$
In other words, $`\eta (\omega )=n`$ if
$$K^n|\mathrm{\Lambda }(k_1(\omega ))|K^{n+1}.$$
###### Proposition 15
There exists a constant $`C`$ such that the following is true for any $`K>K_0`$.
For any $`c_1\mathrm{}c_n\mathrm{}\mathrm{\Omega }_{๐,B}^+`$,
$$(\{\omega :\eta (\omega )=n|\omega _1=c_1,\mathrm{}\omega _n=c_n\mathrm{})\}\frac{C}{K^n}.$$
This immediately follows from Lemma 7.
###### Proposition 16
$$(\{\omega :\eta (\omega )=1,\omega _{k_1(\omega )}\mathrm{}\omega _0\mathrm{does}\mathrm{not}\mathrm{contain}\mathrm{the}\mathrm{word}๐ช\})1p(๐ช)$$
Finally, take a large $`N`$ and let
$$n_N(\omega )=\mathrm{min}n:k_1(\omega )+\mathrm{}k_1(\stackrel{~}{\sigma }^n(\omega ))N.$$
Note that, by definition,
$$K^N|\mathrm{\Lambda }(n_1)(\omega )|K^{2N}.$$
Now consider the set
$$\stackrel{~}{\mathrm{\Omega }}(N)=\{\omega :\omega _{n_N(\omega )}\mathrm{}\omega _0\mathrm{does}\mathrm{not}\mathrm{contain}\mathrm{the}\mathrm{word}๐ช\}.$$
Note that
$$\{\omega :L_๐ช(\omega )>2N\}\stackrel{~}{\mathrm{\Omega }}(N).$$
It suffices, therefore, to prove that there exists $`r<1`$ such that
$$(\stackrel{~}{\mathrm{\Omega }}(N))r^N.$$
But by the previous two propositions, we immediately have
$$(\stackrel{~}{\mathrm{\Omega }})C(1p(๐ช)+\frac{1}{K})^N,$$
and, since $`1p(๐ช)+\frac{1}{K}<1`$, the Proposition follows.
This Proposition admits an equivalent formulation in terms of the norms of renormalization matrices on the space of of interval exchange transformations.
More precisely, for $`(\lambda ,\pi )\mathrm{\Delta }_๐ช`$, $`\mathrm{\Phi }(\lambda ,\pi )=\omega _1\mathrm{}\omega _n\mathrm{}`$, we let $`n^๐ช(\lambda ,\pi )`$ to be the smallest $`n>0`$ such that $`๐ข^n(\lambda ,\pi )\mathrm{\Delta }_๐ช`$, and we set
$$๐ฉ(\lambda ,\pi )=A(\omega _1\mathrm{}\omega _{n^๐ช(\omega )}).$$
###### Corollary 10
There exists $`ฯต>0`$ such that
$$_{\mathrm{\Delta }_๐ช}N(\lambda ,\pi )^ฯต๐<+\mathrm{}.$$
Remark. First results on exponential decay for the probabilities of return times were obtained by Jayadev Athreya. In his approach, Athreya used the dynamics of $`SL(2,)`$-action, which allowed him to obtain optimal exponents. The argument above is an attempt to recover some of Athreyaโs theorems using the language of interval exchange transformations; the argument above does not, however, give an optimal exponent.
Avila, Gouรซzel, and Yoccoz have recently announced exponential decay of correlations for the Teichmรผller flow. One of the steps in their proof is, again, an exponential estimate for return times, which they have obtained independently (Avila \[oral communication\]). Their exponent is optimal.
## 12 Estimate of the measure.
###### Lemma 17
There exists a constant $`C()`$ depending only on the Rauzy class $``$ such that
$$\nu (\mathrm{\Delta }()\mathrm{\Delta }_ฯต())<Cฯต$$
The proof repeats that of Proposition 13.2 in Veech .
Lemma 4 and Corollary 5 therefore imply the following
###### Corollary 11
Let $`๐ชW_{๐,B}`$, $`๐ช=q_1\mathrm{}q_l`$ be such that all entries of the matrix $`A(๐ช)`$ are positive.
There exist $`C>0,\alpha >0`$ such that the following is true for any $`n`$.
$$((\lambda ,\pi ):๐ข^{2k}(\lambda ,\pi )\mathrm{\Delta }(๐ช)\mathrm{for}\mathrm{all}k,1kn)C\mathrm{exp}(\alpha \sqrt{n}).$$
Proof: Let $`n=r^2`$ and denote
$$X(n,๐ช)=\{(\lambda ,\pi ):๐ข^{2k}(\lambda ,\pi )\mathrm{\Delta }(๐ช)\mathrm{for}\mathrm{all}k,1kn)\}.$$
Take
$$B(n)=\{(\lambda ,\pi ):๐ข^{2k}(\lambda ,\pi )\mathrm{\Delta }_{\mathrm{exp}(r)}\mathrm{for}\mathrm{some}k,1kn)\}$$
Then, by the previous Lemma, $`\nu (B(n))Cr^2\mathrm{exp}(r)`$, whereas, by Corollary 5,
$$\nu (X(n,๐ช)B(n))(1p(๐ช))^r,$$
and Corollary 11 is proven.
Remark. This result allows to use the tower method of L.-S. Young and to obtain the decay rate $`\mathrm{exp}(\alpha \sqrt{n})`$ for correlations of bounded Hรถlder functions. For bounded Lipschitz functions, one can also use the method of V. Maume-Deschamps and obtain the uniform rate of decay at the rate $`\mathrm{exp}(\alpha n^{1/2ฯต})`$. It is not clear to me, however, how to use either of these methods in the invertible case.
## 13 Inequalities
Let
$$W_{๐,B}^+=\{wW_{๐,B}:|w|\mathrm{is}\mathrm{even},\mathrm{\Delta }(w)\mathrm{\Delta }^+\}.$$
###### Lemma 18
For any $`C_1,C_2>0`$ there exists $`C_3>0`$ such that the following is true.
Suppose $`row(A)<C_1`$ and $`\lambda \mathrm{\Delta }_{C_2}`$.
Then
$$\frac{1}{C_3}\frac{|A\lambda |^m}{\mathrm{\Pi }_{j=1}^m_{i=1}^mA_{ij}}C_3$$
Proof:
Denote $`A_j=_{i=1}^mA_{ij}`$, so that $`|A|=_{j=1}^mA_j`$.
Then
$$\frac{A_j}{A_k}row(A),$$
whence
$$\frac{A_j}{|A|}\frac{1}{mrow(A)}.$$
Finally, if $`\lambda \mathrm{\Delta }_{C_2}`$, then
$$|A\lambda |C_2|A|,$$
which completes the proof.
###### Corollary 12
For any $`C_4>0`$, $`C_5>0`$ there exists $`C_6>0`$ such that the following is true. Suppose $`(\lambda ,\pi )\mathrm{\Delta }_{C_4}`$. Suppose $`w๐ฒ_{๐,B}`$ is compatible with $`(\lambda ,\pi )`$ and such that $`row(A(w))<C_5`$. Then
$$\frac{1}{C_6}\frac{๐ฆ(C(w))}{(w|(\lambda ,\pi ))}C_6$$
###### Corollary 13
For any $`C_7>0`$, $`C_8>0`$ $`C_9>0`$, there exists $`C_{10}>0`$ such that the following is true.
Suppose $`(\lambda ,\pi )\mathrm{\Delta }_{C_7}`$.
Suppose $`w๐ฒ_{๐,B}`$ is compatible with $`(\lambda ,\pi )`$ and furthermore satisfies
$$row(A(w))<C_8,\mathrm{\Delta }(w)\mathrm{\Delta }_{C_9}$$
Then
$$\frac{1}{C_{10}}\frac{(C(w))}{(w|(\lambda ,\pi ))}C_{10}$$
###### Corollary 14
Let $`M`$ be such that for any $`n>M`$ any two vertices in the Rauzy graph can be joined in $`n`$ steps.
Then for any $`C_{17}>0`$, $`C_{18}>0`$ $`C_{19}>0`$, there exists $`C_{20}>0`$ such that the following is true.
Suppose $`(\lambda ,\pi )\mathrm{\Delta }^+\mathrm{\Delta }_{C_{17}}`$.
Suppose $`wW_{๐,B}^+`$ satisfies
$$row(A(w))<C_{18},\mathrm{\Delta }(w)\mathrm{\Delta }^+\mathrm{\Delta }_{C_{19}}$$
Then for any $`nM`$, we have
$$\frac{1}{C_{20}}\frac{(C(w))}{^{(2n)}(w|(\lambda ,\pi ))}C_{20}$$
From the definition (4) of the Hilbert metric it easily follows that for any $`\lambda ,\lambda ^{}\mathrm{\Delta }_{m1}`$ we have
$$e^{d(\lambda ,\lambda ^{})}\lambda _i^{}\lambda _ie^{d(\lambda ,\lambda ^{})}\lambda _i^{}.$$
(26)
###### Proposition 17
Assume $`\lambda ,\lambda ^{}\mathrm{\Delta }_\pi ^+`$. Then
$$\mathrm{exp}(md(\lambda ,\lambda ^{}))\frac{\rho (\lambda ,\pi )}{\rho (\lambda ^{},\pi )}\mathrm{exp}(md(\lambda ,\lambda ^{}))$$
Proof. Indeed, there exist linear forms
$$l_i^{(j)}(\lambda )=\underset{k=1}{\overset{m}{}}a_{ik}^{(j)}\lambda _k,$$
where $`a_{ik}^{(j)}`$ are nonnegative integers (in fact, either $`0`$ or $`1`$, but we do not need this here),
such that
$$\rho (\lambda ,\pi )=\underset{j=1}{\overset{s}{}}\frac{1}{l_1^{(j)}(\lambda )l_2^{(j)}(\lambda )\mathrm{}l_m^{(j)}(\lambda )}.$$
Clearly, if for all $`i=1,\mathrm{},m`$ and some $`\alpha >0`$, we have $`\alpha ^1\lambda _i\lambda _i^{}\alpha \lambda _i`$, then
$$\alpha ^m\frac{\rho (\lambda ,\pi )}{\rho (\lambda ^{},\pi )}\alpha ^m,$$
and the Proposition is proved.
For similar reasons we have
###### Proposition 18
Assume $`\lambda ,\lambda ^{}\mathrm{\Delta }_\pi ^+`$ and let $`A`$ be an arbitrary matrix with nonnegative integer entries. Then
$$\mathrm{exp}(md(\lambda ,\lambda ^{}))\frac{\rho (A\lambda ,\pi )}{\rho (A\lambda ^{},\pi )}\mathrm{exp}(md(\lambda ,\lambda ^{}))$$
From these propositions and the formula 11 we obtain
###### Corollary 15
Let $`c๐`$ be compatible with $`\pi `$. Then for any $`\lambda ,\lambda ^{}\mathrm{\Delta }_\pi ^+`$ we have
$$\mathrm{exp}(2md(\lambda ,\lambda ^{}))\frac{(c|(\lambda ,\pi ))}{(c|(\lambda ^{},\pi ))}\mathrm{exp}(2md(\lambda ,\lambda ^{}))$$
This Corollary implies the following
###### Lemma 19
Let $`wW_{๐,B}^+`$ be such that the cylinder $`C(w)`$ has finite Hilbert diameter.
Then for any $`c`$ compatible with $`w`$ and any $`(\lambda _0,\pi )C(w)`$ we have
$$\mathrm{exp}(2mdiamC(w))\frac{(c|(\lambda _0,\pi ))}{(\omega _0=c|\omega |_{[1,|w|]}=w)}\mathrm{exp}(2mdiamC(w))$$
Proof: We have
$$\nu (C(cw))=_{C(w)}(c|(\lambda ,\pi ))๐\nu (\lambda ,\pi )$$
Let $`d=diamC(w)`$. For any $`(\lambda _,\pi ),(\lambda ^{},\pi )C(w)`$, we have, by Corollary 15,
$$\mathrm{exp}(2md)\frac{(c|(\lambda ,\pi ))}{(c|(\lambda ^{},\pi ))}\mathrm{exp}(2md).$$
Fix an arbitrary $`(\lambda _0,\pi )\mathrm{\Delta }_w`$.
Then, from the above,
$$\nu (C(w))P(c|(\lambda _0,\pi ))\mathrm{exp}(2md)_{C(w)}P(c|(\lambda ,\pi ))๐\nu (\lambda ,\pi )$$
$$\nu (C(w))P(c|(\lambda _0,\pi ))\mathrm{exp}(2md),$$
and, since, by definition, we have
$$(\omega _0=c|\omega |_{[1,|w|]}=w)=\frac{(cw)}{(w)},$$
the Lemma is proved.
For $`N`$ and $`A\mathrm{\Delta }()`$, we denote $`^{(N)}(A|(\lambda ,\pi ))=((\lambda (N),\pi (N))A|(\lambda (0),\pi (0))=(\lambda ,\pi ))`$; for $`w๐ฒ_{๐,B}`$, we write $`^{(N)}(w|(\lambda ,\pi ))=^{(N)}(\mathrm{\Delta }(w)|(\lambda ,\pi ))`$.
###### Lemma 20
Let $`M`$ be a number such that for any $`NM`$ any two vertices of the Rauzy graph can be connected in $`N`$ steps. For any $`\gamma >0`$, $`NM`$ there exists a constant $`C_0`$ depending only on $`\gamma `$ and $`N`$ such that for any word $`w๐ฒ_{๐,B}^+`$ and any $`(\lambda ,\pi )\mathrm{\Delta }_\gamma `$
$$^{(2N)}(w|(\lambda ,\pi ))\frac{C_0}{|A(w)\lambda |^m}$$
Proof:
Let $`w=w_1\mathrm{}w_{2n}`$, and let $`w_{2n}=(a,m_1,\pi _1)`$.
Let $`\pi _1^{}\pi _2^{}\mathrm{}\pi _{2N}^{}`$ a path of length $`2N`$ between $`\pi `$ and $`\pi _1`$ (here $`\pi _1^{}=\pi `$, $`\pi _{2n}^{}=\pi _1`$, $`\pi _{2k+1}=a\pi _{2k}`$, $`\pi _{2k+2}=b\pi _{2k+1}`$.
Denote $`w_{n+2i+1}=(a,1,\pi _{2i+1})`$, $`w_{n+2i}=(b,1,\pi _{2i})`$. In other words, the word $`=w_{2n+1}\mathrm{}w_{2n+2N}๐ฒ_{๐,B}`$ is the word correspoding to the path $`\pi _1^{}\pi _2^{}\mathrm{}\pi _{2N}^{}`$ in the Rauzy graph. Then $`w^{}=w_1\mathrm{}w_{2n+2N}`$ is a word compatible with $`(\lambda ,\pi )`$. Besides,
$$|A(w_{2n+1}c_{n+2}\mathrm{}w_{2n+2N})|<(2N)^{(2N)}.$$
We have
$$P^{(2n)}(w|(\lambda ,\pi ))P(w^{}|(\lambda ,\pi ))=\frac{\rho (T_w^{}(\lambda ),w^{}\pi )}{|A(w^{})\lambda |^m\rho (\lambda ,\pi )},$$
There exists a universal constant $`C_1`$ such that $`\rho (\lambda ^{},\pi ^{})>C_1`$ for any $`(\lambda ^{},\pi ^{})\mathrm{\Delta }^+`$ (the density of the invariant measure is bounded from below).
Then, $`|A(w^{})\lambda |^m|A(w^{})|^m(2N)^{2mN}|A(w)|^m.`$
Finally, there exists a $`C_2`$ depending on $`c`$ only such that if $`\lambda _i>c`$ for all $`i`$ then $`\rho (\lambda ,\pi )>C_2`$.
Combining all of the above, we obtain the result of the Lemma.
## 14 Markov approximation and the Doeblin condition
### 14.1 Good cylinders
Let $`๐ช=q_1\mathrm{}q_l`$ be a word such that all entries of the matrix $`A(๐ช)`$ are positive. Fix $`ฯต>0`$ and let et $`k_0`$ be such that
$$(\mathrm{\Delta }(๐ช)๐ข^{2n}\mathrm{\Delta }(๐ช))ฯต\mathrm{for}n>k_0.$$
(27)
Note that, due to mixing, Corollary 5 implies the following
###### Proposition 19
Let $`๐ชW_{๐,B}`$, $`๐ช=q_1\mathrm{}q_l`$ be such that all entries of the matrix $`A(๐ช)`$ are positive and that $`\mathrm{\Delta }(๐ช)\mathrm{\Delta }^+`$. Then there exist positive constants $`K(๐ช),p(๐ช)`$ such that the following is true for any $`ฯต>0`$. Suppose $`(\lambda ,\pi )\mathrm{\Delta }_ฯต\mathrm{\Delta }^+`$ and set $`n`$ to be the integer part of $`K(๐ช)|\mathrm{log}ฯต|`$. Then
$$\{(\lambda (2n),\pi (2n))\mathrm{\Delta }(๐ช)|(\lambda (0),\pi (0))=(\lambda ,\pi ))\}p(๐ช).$$
Take $`kk_0`$. Let $`r=2(K+1)k+2M`$, where $`K`$ is the constant from the Lemma 4 and $`M`$ is the connecting constant of the Rauzy graph from Lemma 20.
Let $`\theta `$, $`0<\theta <1`$ be arbitrary. A word $`w=w_1\mathrm{}w_k`$ is called good if
1. $`\mathrm{\Delta }(w)\mathrm{\Delta }_{\mathrm{exp}(k)}`$.
2. the word $`๐ช`$ appears at least $`\frac{k^\theta }{l}`$ times in $`w`$ (we only count disjoint appearances).
A word $`w_1\mathrm{}w_r`$ is called good if $`w_1\mathrm{}w_k`$ is good, a word $`w_1\mathrm{}w_{Nr}`$ is called good if all words $`w_1\mathrm{}w_r`$, $`w_{r+1}\mathrm{}w_{2r}`$, โฆ$`w_{(N1)r+1}\mathrm{}w_{Nr}`$ are good, and a word $`w_1\mathrm{}w_{Nr+L}`$, $`L<r`$, is good if $`w_1\mathrm{}w_{Nr}`$ is good and either $`L<k`$ or $`w_{Nr+1}\mathrm{}w_{Nr+k}`$ is good.
We denote by $`๐(N)`$ the set of all good words of length $`N`$.
Let
$$\mathrm{\Delta }(๐(N))=_{w๐(N)}\mathrm{\Delta }(w),$$
and
$$\mathrm{\Delta }(B(N))=\mathrm{\Delta }^+\mathrm{\Delta }(๐(N))$$
By Corollary 11, there exist constants $`C_{31},C_{32}`$ such that for all $`r`$ we have
$$(\mathrm{\Delta }(B(N))C_{31}N\mathrm{exp}(C_{32}r^{(1\theta )/2}).$$
(28)
and, for any $`(\lambda ,\pi )\mathrm{\Delta }(๐ช)`$, also
$$((\lambda (1),\pi (1))\mathrm{\Delta }(B(N))|(\lambda (0),\pi (0)=(\lambda ,\pi ))C_{31}N\mathrm{exp}(C_{32}r^{(1\theta )/2}).$$
(29)
### 14.2 Preliminary estimates for the Doeblin condition.
From Corollary 15 we deduce that there exists a constant $`C_{33}`$ such that for any $`(\lambda ,\pi ),(\lambda ^{},\pi )\mathrm{\Delta }(๐ช)`$, and any word $`w`$ compatible with $`๐ช`$, we have
$$\frac{1}{C}_{33}\frac{(w|(\lambda ,\pi ))}{(w|(\lambda ^{},\pi ))}C_{33}.$$
Finally, by Lemma 20, there exists a constant $`C_{34}`$ such that for any $`w๐ฒ_{๐,B}`$ and for any $`N>M`$ we have
$$\frac{1}{C}_{34}\frac{^{(2N)}(w|(\lambda ,\pi ))}{^{(2N)}(w|(\lambda ^{},\pi ))}C_{34}.$$
Take an arbitrary point $`(\lambda ,\pi )\mathrm{\Delta }_๐ช`$. Define a new measure $`\phi `$ on $`\mathrm{\Delta }^+`$. Namely, for a set $`A\mathrm{\Delta }^+`$ put
$$\phi (A)=(\lambda (2M),\pi (2M))A|\lambda (0),\pi (0)=(\lambda ,\pi ))$$
(30)
###### Lemma 21
There exists a constant $`\alpha >0`$ such that the following is true for any $`r`$. Let $`๐_1,๐_2๐(r)`$.
Then
$$(\omega |_{[1,r]}=๐_1,\omega |_{[r+1,2r]}๐(r)|\omega |_{[2r+1,3r]}=๐_2)\alpha \phi (๐_1)$$
Indeed, we have the following propositions:
###### Proposition 20
There exist a constant $`p_1`$ such that the following is true for all $`r`$ and all $`nr`$.
Let $`C_2๐(r)`$, $`(\lambda ,\pi )๐_2`$. Then
$$((\lambda (2n),\pi (2n))\mathrm{\Delta }(๐ช)|(\lambda (0),\pi (0))=(\lambda ,\pi ))p_1.$$
This follows from the definition of a good cylinder and Corollary 5.
###### Proposition 21
There exists a constant $`p_2`$ such that the following is true for all $`k`$.
$$(\omega |_{[1,r]}๐(r),\omega |_{[2M+1,l+2M+1]}=๐ช|\omega |_{[r+1,r+l+1]}=๐ช)p_2$$
This follows from the estimates (28),(29) on the measure of bad cylinbders and from Proposition 19.
###### Proposition 22
There exists a constant $`p_3`$ such that the following is true for all $`r`$. Let $`c_1\mathrm{}c_n\mathrm{}\mathrm{\Delta }(๐ช)`$.
$$(\omega |_{[1,r]}=๐_1|\omega _{r+2M+1}=c_1,\omega _{r+2M+2}=c_2,\mathrm{})p_3\phi (C_1)$$
This follows directly from Lemma 20.
The three Propositions imply Lemma 21.
### 14.3 Approximation by a Markov measure
We define a new measure $`๐ฉ_{r,\theta }`$ on the set $`๐(r^2)`$ of good cylinders of length $`r^2`$.
Let $`๐=c_1\mathrm{}c_{r^2}`$ be a $`(r,\theta )`$-good cylinder. Set $`๐_i=c_{ir+1}\mathrm{}c_{(i+1)r}`$.
Define
$$๐ฉ_{r,\theta }(๐)=(\omega |_{[1,r]}=๐_1|\omega |_{[r+1,2r]}=๐_2)(\omega |_{[r+1,2r]}=๐_2|\omega |_{[2r+1,3r]}=๐_3)\mathrm{}(\omega |_{[r^2r+1,r^2]}=๐_r).$$
If $`D`$ is not a good cylinder, then $`๐ฉ_{r,\theta }(D)=0`$.
Normalize to get a probability measure:
$$๐_{r,\theta }(๐)=\frac{๐ฉ_{r,\theta }(๐)}{_{๐๐(r^2)}๐ฉ_{r,\theta }(๐)}.$$
$`๐_{r,\theta }`$ is a Markov measure of memory $`r`$ (in general, non-homogeneous), as is shown by the following well-known Lemma .
###### Lemma 22
For any $`k`$, $`0<k<r`$, we have
$$๐_{r,\theta }(\omega |_{[kr+1,(k+1)r]}=๐_k|\omega |_{[(k+1)r+1,r^2]})=๐_{k+1}\mathrm{}๐_r)=$$
$$๐_{r,\theta }(\omega |_{[kr+1,(k+1)r]}=๐_k|\omega |_{[(k+1)r+1,(k+2)r]})=๐_{k+1}).$$
From the Hรถlder property for the transition probability, we have
###### Proposition 23
There exist constants $`C_{41},C_{42}`$ such that the following is true for any $`r`$.
Let $`c_1\mathrm{}c_n\mathrm{}\mathrm{\Omega }_{๐,B}`$ and assume $`c_{n+1}\mathrm{}c_{n+r}๐(r)`$. Then
$$\mathrm{exp}(C_{41}\mathrm{exp}(C_{42}k^\theta ))$$
$$\frac{P(\omega _1=c_1,\mathrm{},\omega _n=c_n|\omega _{n+1}=c_{n+1},\mathrm{},\omega _{n+r}=c_{n+r})}{(\omega _1=c_1,\mathrm{},\omega _n=c_n|\omega _{n+1}=c_{n+1},\mathrm{},\omega _{n+i}=c_{n+i},\mathrm{})}$$
$$\mathrm{exp}(C_{41}\mathrm{exp}(C_{42}k^\theta ))$$
###### Corollary 16
There exist constants $`C_{43},C_{44}`$ such that the following is true for any $`r`$. Let $`A_n`$, let $`c_{n+1}\mathrm{}c_{n+i}\mathrm{}\mathrm{\Omega }_๐`$, and assume $`c_{n+1}\mathrm{}c_{n+r}๐(r)`$. Then
$$\mathrm{exp}(C_{43}\mathrm{exp}(C_{44}k^\theta ))\frac{(A|\omega _{n+1}=c_{n+1},\mathrm{},\omega _{n+r}=c_{n+r})}{(A|\omega _{n+1}=c_{n+1},\mathrm{},\omega _{n+i}=c_{n+i},\mathrm{})}\mathrm{exp}(C_{43}\mathrm{exp}(C_{44}k^\theta ))$$
Applying $`l`$ times, we obtain
###### Lemma 23
There exist constants $`C_{45},C_{46},C_{47},C_{48}`$ such that the following is true for any $`r`$. Let $`c_1\mathrm{}c_{r^2}๐(r^2)`$. Then for any $`l`$, $`1lr`$, we have
$$\mathrm{exp}(C_{45}l\mathrm{exp}(C_{46}k^\theta ))$$
$$\frac{(\omega _1=c_1,\mathrm{},\omega _{lr}=c_{lr}|\omega _{lr+1}=c_{lr+1},\mathrm{},\omega _{r^2}=c_{r^2})}{๐ฉ_{r,\theta }(\omega _1=c_1,\mathrm{},\omega _{lr}=c_{lr}|\omega _{lr+1}=c_{lr+1},\mathrm{},\omega _{r^2}=c_{r^2})}$$
$$\mathrm{exp}(C_{45}l\mathrm{exp}(C_{46}k^\theta ))$$
and
$$\mathrm{exp}(C_{47}l\mathrm{exp}(C_{48}k^\theta ))\frac{(\omega _1=c_1,\mathrm{},\omega _{lr}=c_{lr})}{๐ฉ_{r,\theta }(\omega _1=c_1,\mathrm{},\omega _{lr}=c_{lr})}\mathrm{exp}(C_{47}l\mathrm{exp}(C_{48}k^\theta ))$$
Summing over cylinders of length $`lr`$, we obtain
###### Corollary 17
There exist constants $`C_{49},C_{50}`$ such that the following is true for any $`r`$. Let $`c_1\mathrm{}c_{r^2}๐(r^2)`$. Then for any $`l`$, $`1lr`$, and any $`A_{lr}`$, we have
$$\mathrm{exp}(C_{49}l\mathrm{exp}(C_{50}k^\theta ))\frac{(A๐(lr)|\omega _{lr+1}=c_{lr+1},\mathrm{},\omega _{r^2}=c_{r^2})}{๐ฉ_{r,\theta }(A|\omega _{lr+1}=c_{lr+1},\mathrm{},\omega _{r^2}=c_{r^2})}\mathrm{exp}(C_{49}l\mathrm{exp}(C_{50}k^\theta ))$$
and
$$\mathrm{exp}(C_{49}l\mathrm{exp}(C_{50}k^\theta ))\frac{(A๐(lr))}{๐ฉ_{r,\theta }(A)}\mathrm{exp}(C_{49}l\mathrm{exp}(C_{50}k^\theta ))$$
Using (28), we can estimate the total mass of the measure $`๐ฉ_{r,\theta }`$.
###### Corollary 18
There exist constants $`C_{51},C_{52}`$ such that for any $`r`$ we have
$$๐ฉ_{r,\theta }(๐(r^2))\mathrm{exp}(C_{51}r\mathrm{exp}(C_{52}k^{(1\theta )/2}))$$
We now have normalized versions of previous statements.
###### Corollary 19
There exist constants $`C_{53},C_{54},C_{55},C_{56}`$ such that the following is true for any $`r`$. Let $`c_1\mathrm{}c_{r^2}๐(r^2)`$. Then for any $`l`$, $`1lr`$, and any $`A_{lr}`$, we have
$$\mathrm{exp}(C_{53}l\mathrm{exp}(C_{54}k^\theta )C_{55}r\mathrm{exp}(C_{56}k^{(1\theta )/2}))$$
$$\frac{(A๐(lr)|\omega _{lr+1}=c_{lr+1},\mathrm{},\omega _{r^2}=c_{r^2})}{๐_{r,\theta }(A|\omega _{lr+1}=c_{lr+1},\mathrm{},\omega _{r^2}=c_{r^2})}$$
$$\mathrm{exp}(C_{53}l\mathrm{exp}(C_{54}k^\theta )+C_{55}r\mathrm{exp}(C_{56}k^{(1\theta )/2})$$
and
$$\mathrm{exp}(C_{53}l\mathrm{exp}(C_{54}k^\theta )C_{55}r\mathrm{exp}(C_{56}k^{(1\theta )/2}))$$
$$\frac{(A๐(lr))}{๐_{r,\theta }(A)}$$
$$\mathrm{exp}(C_{53}l\mathrm{exp}(C_{54}k^\theta )+C_{55}r\mathrm{exp}(C_{56}k^{(1\theta )/2}).$$
Using the Markov approximation, we can estimate conditional measure of good cylinders for the measure $``$:
###### Corollary 20
There exist constants $`C_{57},C_{58},C_{59},C_{60}`$ such that the following is true for any $`r`$. Let $`c_1\mathrm{}c_{r^2}๐(r^2)`$. Then for any $`l`$, $`1lr`$, we have
$$((\omega _1\mathrm{}\omega _{lr})๐(lr)|\omega _{lr+1}=c_{lr+1},\mathrm{},\omega _{r^2}=c_{r^2})\mathrm{exp}(C_{57}l\mathrm{exp}(C_{58}k^\theta )C_{59}r\mathrm{exp}(C_{60}k^{(1\theta )/2}))$$
Proof: Indeed,
$$๐_{r,\theta }((\omega _1\mathrm{}\omega _{lr})๐(lr)|\omega _{lr+1}=c_{lr+1},\mathrm{},\omega _{r^2}=c_{r^2})=1.$$
### 14.4 Doeblin Condition
###### Proposition 24
There exists $`C_{61}`$ such that the following holds for any $`r`$. For any $`๐_1\mathrm{\Delta }(๐ช)`$, $`C_2\mathrm{\Delta }_๐ช`$, and any $`๐_3๐(r)`$, we have either
$$\frac{1}{C_{61}}\frac{๐ฉ_{r,\theta }(C_3|C_2)}{๐ฉ_{r,\theta }(C_3|C_1)}C_{61},$$
or $`๐ฉ_{r,\theta }(C_3|C_2)=๐ฉ_{r,\theta }(C_3|C_1)=0`$.
Considering $`n`$-step transition probabilities, we obtain
###### Proposition 25
There exists a constant $`C_{62}`$ such that the following holds for any $`r`$. For any $`๐_1\mathrm{\Delta }(๐ช)`$, $`C_2\mathrm{\Delta }_๐ช`$ any $`๐_3๐(r)`$, and any $`nM`$, we have
$$\frac{1}{C_{62}}\frac{๐ฉ_{r,\theta }(\omega |_{[1,r]}=C_1|\omega _{[2n+r,2n+2r]}=C_2)}{๐ฉ_{r,\theta }(\omega |_{[1,r]}=C_1|\omega _{[2n+r,2n+2r]}=C_3)}C_{62}.$$
Now, mixing, Proposition 19 and Proposition 20, and the definition oif a good cylinder imply that
###### Proposition 26
There exists a constant $`C_{63}`$ such that the following holds for any $`r`$. For any $`๐_1,๐_2,๐_3๐(r)`$ we have
$$\frac{1}{C_{63}}\frac{๐ฉ_{r,\theta }(\omega |_{[1,r]}=C_1|\omega _{[2r,3r]}=C_2)}{๐ฉ_{r,\theta }(\omega |_{[1,r]}=C_1|\omega _{[2r,3r]}=C_3)}C_{63}.$$
Now let $`c_1\mathrm{}c_{r^2}๐(r^2)`$. Denote $`๐_i=c_{ir+1}\mathrm{}c_{(i+1)r}`$. Lemma 21, together with the above estimates, implies the following
###### Corollary 21
There exist constants $`C_{71},C_{72}`$ such that the following is true. For any $`l`$, $`1lr`$, we have
$$(\omega |_{[1,lr]}๐(lr),\omega |_{[lr+1,(l+1)r]}=๐_l,\omega |_{[(l+1)r+1,(l+2)r]}๐(r)|\omega _{(l+2)r+1,(l+3)r])}=๐_3)C_{71}\times \phi (๐_l)$$
and
$$๐_{r,\theta }(\omega |_{[1,lr]}๐(lr),\omega |_{[lr+1,(l+1)r]}=๐_l,\omega |_{[(l+1)r+1,(l+2)r]}๐(r)|\omega _{(l+2)r+1,(l+3)r])}=๐_3)C_{72}\times \phi (๐_l)$$
This is the Doeblin Condition for the measure $`๐_{r,\theta }`$ (see , , ). The Doeblin Condition implies that there exist constants $`C_{73},C_{74}`$ such that for any $`๐_1,๐_2๐(r)`$, we have
$$\mathrm{exp}(C_{73}\mathrm{exp}(C_{74}r))\frac{๐_{r,\theta }(\omega |_{[1,r]}=๐_1|\omega |_{[r^2,r^2+r]}=๐_2)}{๐_{r,\theta }(๐_1)}\mathrm{exp}(C_{73}\mathrm{exp}(C_{74}r)),$$
whence we obtain
###### Proposition 27
There exist constants $`C_{81},C_{82},C_{83},C_{84}`$ such that the following is true for any $`r`$.
$$\mathrm{exp}(C_{81}(\mathrm{exp}(C_{82}r)+\mathrm{exp}(C_{83}r^\theta )+\mathrm{exp}(C_{84}r^{(1\theta )/2})))$$
$$\frac{(\omega |_{[1,r]}=๐_1|\omega |_{[r+1,r^2]}๐(r^2r),\omega |_{[r^2,r^2+r]}=๐_2)}{(๐_1)}$$
$$\mathrm{exp}(C_{81}\mathrm{exp}(C_{82}r)+\mathrm{exp}(C_{83}r^\theta )+\mathrm{exp}(C_{84}r^{(1\theta )/2)}))).$$
Moreover, in view of mixing, Proposition 19, and Proposition 20, the same estimate, upto a constant, takes place for any $`nr^2`$.
###### Proposition 28
There exist constants $`C_{85},C_{86},C_{87},C_{88}`$ such that the following is true for all $`r`$ and all $`nr^2`$.
$$\mathrm{exp}(C_{85}(\mathrm{exp}(C_{86}r)+\mathrm{exp}(C_{87}r^\theta )+\mathrm{exp}(C_{88}r^{(1\theta )/2}))))$$
$$\frac{(\omega |_{[1,r]}=๐_1|\omega |_{[r+1,n]}๐(nr),\omega |_{[n,n+r]}=๐_2)}{(๐_1)}$$
$$\mathrm{exp}(C_{85}(\mathrm{exp}(C_{86}r)+\mathrm{exp}(C_{87}r^\theta )+\mathrm{exp}(C_{88}r^{(1\theta )/2)})))).$$
## 15 Approximation of Hรถlder Functions and Completion of the Proof of Theorems 4, 7, 8.
We shall prove the decay of correlations for a slightly more general class of functions on $`\mathrm{\Delta }()`$ than Hรถlder functions. (we shall need this slightly more general class in the proof of the Central Limit Theorem).
Namely, we shall only require that a function be Hรถlder in restriction to cylinders of some given length and we shall also allow a moderate growth of the Hรถlder constant at infinity.
Formally, say that a function $`\varphi :\mathrm{\Delta }()`$ is weakly $`l,\alpha `$-Hรถlder if the following holds. Let $`k`$ be a positive integer, and let $`w๐ฒ_{๐,B}`$, $`|w|l`$ be such that $`\mathrm{\Delta }(w)\mathrm{\Delta }_{\mathrm{exp}(k)}`$. Then there exists a constant $`C(\varphi )`$ such that for any $`(\lambda ,\pi ),(\lambda ^{},\pi )\mathrm{\Delta }(w)`$, we have
$$|\varphi (\lambda ,\pi )\varphi (\lambda ^{},\pi )|Ckd(\lambda ,\lambda ^{})^\alpha .$$
The smallest such $`C`$ for a given $`\varphi `$ will be denoted $`C_{l,\alpha }^{weak}(\varphi )`$. Clearly, if $`\varphi `$ is Hรถlder with exponent $`\alpha `$, then it is also weakly $`l,\alpha `$-Hรถlder for any $`l`$ and $`C_{l,\alpha }^{weak}(\varphi )C_\alpha (\varphi )`$.
Recall that $`_n`$ is the $`\sigma `$-algebra of sets of the form $`๐ข^n(A)`$, $`A\mathrm{\Delta }()`$.
To prove the decay of correlations, it suffices to estimate the $`L_2`$-norm of $`E(\varphi |_{2n})`$ for a given weakly $`l`$-$`\alpha `$-Hรถlder $`\varphi `$.
It will be convenient to assume that $`\varphi 1`$ (by linearity, it suffices to consider that case).
###### Proposition 29
Let $`\theta `$, $`0<\theta <1`$. Let $`p>2`$ and $`\alpha >0`$. There exist constants $`C_{91},C_{92}`$, $`C_{93}`$ such that the following is true for any $`r`$ and any $`nr^2`$.
Let $`lr`$. Let $`\varphi L_p(\mathrm{\Delta }()^+,\nu )`$ be weakly $`l,\alpha `$-Hรถlder and satisfy $`\varphi 1`$.
Then $`\varphi =\varphi _1+\varphi _2+\varphi _3`$ where
1. $`\varphi _11`$ on $`๐(n)`$ and $`\varphi _1=\varphi _2=0`$ on $`\mathrm{\Delta }(B(n))`$.
2. for any $`(\lambda ,\pi )๐(n)`$, we have
$$|\frac{E(\varphi _1|_n)(\lambda ,\pi )}{E(\varphi _1)}1|\mathrm{exp}(C_{91}(r^{(1\theta )/2}+r^\theta ).$$
3. for $`(\lambda ,\pi )G(n)`$, we have $`|\varphi _2|C_{l,\alpha }^{weak}(\varphi )\mathrm{exp}(C_{92}r^\theta ).`$
4. $`\varphi _3_{L_2}\mathrm{exp}(C_{93}r^{(1\theta )/2})\varphi _{L_p}.`$
Proof: For any good word $`w=w_1\mathrm{}w_{n+r}`$, consider its beginning $`w_1\mathrm{}w_r`$ and choose a point $`x_{w_1\mathrm{}w_r}\mathrm{\Delta }(w_1\mathrm{}w_r)`$.
Denote by $`\chi _{\mathrm{\Delta }(w)}`$ the characteristic function of $`\mathrm{\Delta }(w)`$ and set
$$\varphi _1=\underset{w๐(n+r)}{}\varphi (x_{w_1\mathrm{}w_r})\chi _{\mathrm{\Delta }(w)}.$$
Proposition 28 yields the required properties of $`\varphi _1`$ (note that we sum over all good words of length $`n+r`$ in order to be able to apply the Proposition).
We set $`\varphi _2=(\varphi \varphi _1)\chi _{G(n+r)}`$ and $`\varphi _3=\varphi \chi _{\mathrm{\Delta }(B(n+r))}`$. The estimate for $`\varphi _2`$ is satisfied by the definition of a Hรถlder function.
Finally, we have
$$\varphi _3_{L_2}^2=E(|\varphi \chi _{\mathrm{\Delta }(B(n))}|^2),$$
whence, by Hรถlderโs inequality, using the estimate (28), we obtain the desired estimate for $`\varphi _3`$, and the Proposition is proved completely.
Proposition 29 with $`\theta =1/3`$ yields Theorem 4.
We now complete the proof of Theorem 7.
For a word $`wW_{๐,B}`$, $`|w|=2n+1`$, $`w=w_1\mathrm{}w_{2n+1}`$, denote $`C^{[n,n]}(w)=\{\omega \mathrm{\Omega }_{๐,B}^{}:\{\omega _n=w_1,\mathrm{},\omega _n=w_{2n+1}\}`$ and set $`\overline{\mathrm{\Delta }}(w)=\overline{\mathrm{\Phi }}^1C^{[n,n]}(w)`$. Denote by $`_{[n,n]}`$ the sigma-algebra generated by $`\overline{\mathrm{\Delta }}(w)`$ for all $`wW_{๐,B}`$.
Also, for $`ฯต>0`$, denote
$$\overline{\mathrm{\Delta }}_ฯต=\{(\lambda ,h,a,\pi )\overline{\mathrm{\Delta }}():\lambda \mathrm{\Delta }_ฯต.$$
Again, we shall prove the Theorem for a slightly larger class of functions.
Say that a function $`\varphi :\overline{\mathrm{\Delta }}()`$ is weakly $`l,\alpha `$-Hรถlder if the following holds. Let $`k`$ be a positive integer, and let $`w๐ฒ_{๐,B}`$, $`|w|2l+1`$ be such that $`\overline{\mathrm{\Delta }}(w)\mathrm{\Delta }_{\mathrm{exp}(k)}`$. Then there exists a constant $`C(\varphi )`$ such that for any $`(\lambda ,h,a,\pi ),(\lambda ^{},h^{},a^{},\pi )\overline{\mathrm{\Delta }}(w)`$, we have
$$|\varphi (\lambda ,h,a,\pi )\varphi (\lambda ^{},,h^{},a^{},\pi )|Ckd((\lambda ,h,a,\pi ),(\lambda ^{},,h^{},a^{},\pi ))^\alpha .$$
The smallest such $`C`$ for a given $`\varphi `$ will be denoted $`C_{l,\alpha }^{weak}(\varphi )`$. Clearly, if $`\varphi `$ is Hรถlder with exponent $`\alpha `$, then it is also weakly $`l,\alpha `$-Hรถlder for any $`l`$ and $`C_{l,\alpha }^{weak}(\varphi )C_\alpha (\varphi )`$.
Denote by $`\overline{G}(2n+1)`$ the union of all $`\overline{\mathrm{\Delta }}(w)`$ for good $`w`$, by $`\overline{B}(2n+1)`$ the complement of $`\overline{G}(2n+1)`$.
###### Proposition 30
Let $`\theta `$, $`0<\theta <1`$. Let $`p>2`$ and $`\alpha >0`$. There exist constants $`C_{101},C_{102}`$, such that the following is true for any $`r`$ and any $`nr^2`$.
Let $`lr`$. Let $`\varphi L_p(\overline{\mathrm{\Delta }}()^,\overline{\nu })`$ be weakly $`l,\alpha `$-Hรถlder and satisfy $`\varphi 1`$. Then there exist functions $`\varphi _1`$, $`\varphi _2`$, $`\varphi _3`$ such that
1. $`\varphi =\varphi _1+\varphi _2+\varphi _3`$.
2. $`\varphi _1`$ is $`_{[n,n]}`$-measurable and supported on $`\overline{G}(2n+1)`$.
3. $`|\varphi _2|C_{101}C_\alpha (\varphi )\mathrm{exp}(r^{(1\theta )/2}+r^\theta )`$.
4. $`|\varphi _3|_{L_2}C_{102}\mathrm{exp}(r^{(1\theta )/2)}||\varphi ||_{L_p}`$.
For any good $`w`$, $`|w|=2n+1`$, take an arbitrary point $`x_w`$ in $`\overline{\mathrm{\Delta }}(w)`$. Set
$$\varphi _1=\underset{wG(2n+1)}{}\varphi (x_w)\chi _{\overline{\mathrm{\Delta }}(w)},$$
$$\varphi _2=(\varphi \varphi _1)\chi _{G(2n+1)},$$
$$\varphi _3=\varphi \chi _{B(2n+1)},$$
and the Proposition is proved.
Proposition 30 with $`\theta =1/3`$ yields Theorem 7.
It remains to establish the Central Limit Theorem for the flow $`P^t`$. Consider the special function $`\stackrel{~}{\tau }`$ of the flow $`P^t`$ over the transformation $``$. Note that $`\stackrel{~}{\tau }(\lambda ,h,a,\pi )`$ only depends on $`(\lambda ,\pi )`$. Consider the restriction of $`\stackrel{~}{\tau }`$ on a cylinder of the form $`\mathrm{\Delta }(w_1)`$, $`w_1๐`$. Then there exist distinct $`j(1),\mathrm{},j(l)\{1,\mathrm{},m\}`$ such that
$$\stackrel{~}{\tau }(\lambda ,\pi )=\mathrm{log}(\lambda _{j(1)}+\lambda _{(j(2)}+\mathrm{}+\lambda _{j(l)}),$$
which shows that the function $`\stackrel{~}{\tau }`$, restricted to an arbitrary $`\mathrm{\Delta }(w_1)`$ is Lipshitz with respect to the Hilbert metric on $`\mathrm{\Delta }()`$.
Now for a Hรถlder $`\varphi `$ consider the function
$$\stackrel{~}{\varphi }(x)=_0^{\stackrel{~}{\tau }(x)}\varphi (P^tx).$$
For any $`k>1`$, if $`(\lambda ,\pi )\mathrm{\Delta }_{\mathrm{exp}(k)}`$, then, by definition, $`\stackrel{~}{\tau }(\lambda ,\pi )k`$. Therefore, if $`\varphi `$ is Hรถlder of exponent $`\alpha `$, then $`\stackrel{~}{\varphi }`$ is weakly $`1,\alpha `$-Hรถlder.
It is easy to see that $`\stackrel{~}{\tau }(\lambda ,\pi )L_r(\mathrm{\Delta }(),\nu )`$ for any $`r>1`$, whence, if $`\varphi L_p(\mathrm{\Omega }_0(),\mu _{})`$ for some $`p>2`$, then there exists $`p^{}>2`$ such that the function
$$\stackrel{~}{\varphi }(x)=_0^{\stackrel{~}{\tau }(x)}\varphi (P^tx)$$
satisfies $`\stackrel{~}{\varphi }L_p^{}(๐ด^\pm ,\overline{\nu })`$.
Therefore, the Theorem of Melbourne and Tรถrรถk implies Theorem 8, the Central Limit Theorem for the flow $`P^t`$.
Acknowledgements.
I am deeply grateful to Yakov G. Sinai, who introduced me to interval exchange transformations, explained to me the method of Markov approximations, and encouraged me in every possible way as the work progressed (more importantly, when it did not).
I am deeply grateful to Alexander Eskin, who suggested to me the problem of the decay of correlations for the induction map.
I am deeply grateful to Giovanni Forni who introduced me to Teichmรผller theory.
I am deeply grateful to Corinna Ulcigrai and Pavel Batchourine for their suggestions that have been of invaluable help to me.
I am deeply grateful to Jayadev Athreya, Valdo Durrleman, Charles L. Fefferman, Boris M. Gurevich, Carlangelo Liverani, Michael Ludkovski, Ian Melbourne, Klaus Schmidt, Andrei Tรถrรถk and Anton V. Zorich for useful discussions.
Part of this work was done at The Erwin Schrรถdinger Institute in Vienna, at The Institute of Mathematics โGuido Castelnuovoโ of the University of Rome โLa Sapienzaโ, and at the CIRM-IML in Marseille. I am deeply grateful to these institutions for their hospitality.
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# Morita Duality and Noncommutative Wilson Loops in Two Dimensions
## 1 Introduction
The calculation of Wilson loop obervables is an ongoing area of activity in the study of noncommutative gauge theories (see for reviews). Early calculations performed in the corresponding dual supergravity theories found that, in even spacetime dimension and for maximal rank noncommutativity, large area Wilson loop correlators behave exactly as their commutative counterparts up to a rescaling of the YangโMills coupling constant, while noncommutative effects dominate small area loops. In contrast, numerical studies based on twisted reduced models in two dimensions have revealed that small noncommutative Wilson loops follow an area law behaviour while large loops become complex-valued with a phase that rises linearly with the area. Numerical results in four dimensions and for non-maximal rank noncommutativity are qualitatively similar .
In this paper we will analyse some nonperturbative properties of Wilson loops in two-dimensional noncommutative gauge theory. The partition function and open Wilson line observables in this theory have been computed nonperturbatively in terms of instanton expansions which are manifestly invariant under gauge Morita equivalence and area-preserving diffeomorphisms of the spacetime. In marked contrast, closed Wilson line observables have thus far only been amenable to a variety of perturbative studies . Foremost among the interesting effects that have been unveiled is the loss of invariance under area-preserving diffeomorphisms of the two-dimensional spacetime . Commutative Wilson loop correlators in two dimensions are well-known to be independent of the shape of the contour on which they are defined and depend only on the area enclosed by the loop. However, in noncommutative gauge theory on $`^2`$ the loop correlators depend on the path shape . Thus, for example, one obtains different correlation functions associated to a circular loop and a square loop which encircle the same area. The simplest way to understand the violation of this invariance is through the noncommutative loop equation , which relates an infinitesimal variation in the loop geometry of a closed Wilson line correlator to a non-vanishing correlation function of open Wilson lines. This symmetry breaking may be related to the fact that, unlike its commutative version, the lattice regularization of the noncommutative gauge theory is not invariant under subdivision of plaquettes which have long-ranged interactions with one another. The standard GrossโWitten reduction breaks down due to UV/IR mixing in this case . On the other hand, it is expected that at least an $`SL(2,)`$ subgroup of area-preserving diffeomorphisms remains a symmetry of the quantum averages in perturbation theory.
In the following we will study the shape dependence of loop correlators from a nonperturbative perspective. Our fundamental point of view will be to look at Morita equivalent formulations of the gauge theory on a two-dimensional noncommutative torus. Since rational noncommutative YangโMills theory is equivalent to ordinary YangโMills theory on a torus, one would naively expect that in this case Wilson loop correlators are shape-independent. By continuity one could then try to extrapolate this result to the irrational noncommutative torus and by decompactification even to the noncommutative plane. The reason this argument breaks down is that closed Wilson lines, unlike the open ones, have a very intricate transformation property under Morita equivalence. The Morita dual of a closed simple curve can be a very complicated loop with many self-intersections and windings around itself. We describe these transformations in detail, and show that Morita equivalence maps a simple noncommutative Wilson loop on the torus into a non-planar graph realizing a triangulation of the dual torus. The problem of computing the loop correlator is in this way mapped onto a combinatorial problem. Loops which enclose the same area but have a different shape can yield topologically inequivalent graphs and hence different correlation functions. This is in fact also true of loops which differ only in their relative orientation, a feature which distinguishes observables on the torus from those on the plane which are rotationally invariant . The loss of invariance under area-preserving diffeomorphisms from this perspective is then attributed to the different graph combinatorics induced by contours of varying shape. The spacetime transformations which leave a given loop correlator invariant are determined by the automorphism group of the non-planar graph induced by the contour under Morita equivalence.
The organisation of this paper is as follows. In Section 2 we review some aspects of Morita equivalence and spell out in detail how it acts on closed Wilson lines. In Section 3 we present various explicit constructions and calculations in rational noncommutative gauge theory on the torus. In this case the Morita dual gauge theory can be taken to be ordinary YangโMills theory, in which we can perform calculations of self-intersecting loop correlators using combinatorial techniques. Our explicit nonperturbative expressions indeed do suggest the claimed shape dependence, and we present various supporting arguments for this claim. In Section 4 we make some remarks concerning irrational noncommutative gauge theories. Although we cannot make progress with analytical determinations of loop correlators in this case, we can give a heuristic picture of irrational noncommutative Wilson loops as infinitely wound and self-intersecting contours in some dual gauge theory. In Section 5 we then give an explicit realization of this infinite winding property, and derive a nonperturbative expression for the Wilson loop correlator on the noncommutative plane in this case which coincides with the result of resumming commutative planar diagrams in perturbation theory. Finally, in Section 6 we summarize our findings and make some further remarks about the relation between our nonperturbative approach and existing perturbative calculations.
## 2 Morita Equivalence of Wilson Loops
In this section we will recall some basic features of two dimensional noncommutative YangโMills theory. When this theory is defined on a noncommutative torus, there exists a powerful tool to perform explicit computations called Morita equivalence. This is a duality that relates observables in the noncommutative gauge theory to observables in a dual gauge theory. When the noncommutativity parameter is a rational number, the equivalence can be arranged so that the dual theory is a commutative gauge theory and the standard techniques of ordinary YangโMills theory in two dimensions can be applied to compute quantum correlation functions.
Consider $`U(1)`$ YangโMills theory defined on a square noncommutative torus $`๐_\theta ^2`$ with noncommutativity parameter $`\theta `$, so that $`[x^1,x^2]_{}=\mathrm{i}\theta `$ with $`๐=(x^1,x^2)`$ local coordinates on the torus. The radius of $`๐_\theta ^2`$ is $`r^{}`$ so that one has the identifications
$$x^\mu x^\mu +2\pi r^{},\mu =1,2.$$
(1)
While the main features below will hold for general $`\theta `$, we will mostly refer to the gauge theory with rational-valued dimensionless noncommutativity parameter of the form $`\mathrm{\Theta }=\frac{\theta }{2\pi r_{}^{}{}_{}{}^{2}}=\frac{c}{N}`$ with $`c,N`$ relatively prime positive integers. The YangโMills action is given by
$$S_{\mathrm{NCYM}}[๐]=\frac{1}{2g_{}^{}{}_{}{}^{2}}_{๐_\theta ^2}\mathrm{d}^2๐\left(+\mathrm{\Phi }\right)^2,$$
(2)
where the YangโMills field strength
$$=_1๐_2_2๐_1\mathrm{i}\left(๐_1๐_2๐_2๐_1\right)$$
(3)
with $`_\mu =/x^\mu `$ is defined in terms of the abelian noncommutative gauge field $`๐_\mu `$ which has a Fourier series expansion
$$๐_\mu (๐)=\underset{๐^2}{}a_{๐;\mu }\mathrm{e}^{\mathrm{i}๐๐/r^{}},a_{๐;\mu }.$$
(4)
Hermiticity of the gauge field requires $`a_{๐;\mu }=\overline{a_{๐;\mu }}`$. The star-product of fields is defined as
$$(fg)(๐)=f(๐)\mathrm{exp}\left(\frac{\mathrm{i}\theta }{2}ฯต^{\mu \nu }\stackrel{}{_\mu }\stackrel{}{_\nu }\right)g(๐),$$
(5)
and we have introduced a constant abelian background flux $`\mathrm{\Phi }`$.
Observables of noncommutative gauge theories are given by closed and open Wilson lines . In this paper we will focus only on closed paths, whose corresponding Wilson lines are defined as
$$๐ช_{}\left(๐\right)=_{๐_\theta ^2}\mathrm{d}^2๐๐ฐ(x;๐)$$
(6)
where $`๐`$ is a closed contour on $`๐_\theta ^2`$ with embedding $`๐=(\xi ^1,\xi ^2):[0,1]๐_\theta ^2`$, $`\xi ^\mu (0)=\xi ^\mu (1)`$ and
$`๐ฐ(x;๐)`$ $`=`$ $`\mathrm{P}_{}\mathrm{exp}\left(\mathrm{i}{\displaystyle _๐}d\xi ^\mu ๐_\mu \left(๐+๐\right)\right)`$ (7)
$`=`$ $`1+{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}\mathrm{i}^n{\displaystyle _0^1}ds_1{\displaystyle _0^{s_1}}ds_2\mathrm{}{\displaystyle _0^{s_{n1}}}ds_n\dot{\xi }^{\mu _1}(s_1)\dot{\xi }^{\mu _2}(s_2)\mathrm{}\dot{\xi }^{\mu _n}(s_n)`$
$`\times ๐_{\mu _1}\left(๐+๐(s_1)\right)๐_{\mu _2}\left(๐+๐(s_2)\right)\mathrm{}๐_{\mu _n}\left(๐+๐(s_n)\right)`$
with $`\dot{\xi }^\mu (s)=\mathrm{d}\xi ^\mu (s)/\mathrm{d}s`$ is the noncommutative holonomy. The technique we will employ to compute noncommutative Wilson loop correlators is to implement Morita equivalence at the level of these observables in the case of a rational noncommutativity parameter where the target dual gauge theory is commutative, and use the known techniques to compute the correlators in ordinary YangโMills theory. As we will see in the following, this procedure, though naively well-defined, is full of subtleties that need to be dealt with.
Generally, gauge Morita equivalence is a map between the noncommutative gauge theory with action given by (2) and a $`U(N)`$ noncommutative gauge theory on another torus $`๐_{\stackrel{~}{\theta }}^2`$ with $`m`$ units of background magnetic flux whose parameters are related to those of the original theory by the action of an $`SL(2,)`$ duality group . Explicitly, the parameters of the two gauge theories are related as
$`\left(\begin{array}{c}m\\ N\end{array}\right)`$ $`=`$ $`\left(\begin{array}{cc}a& b\\ c& d\end{array}\right)\left(\begin{array}{c}0\\ 1\end{array}\right),`$ (14)
$`\stackrel{~}{\mathrm{\Theta }}`$ $`=`$ $`{\displaystyle \frac{c+d\mathrm{\Theta }}{a+b\mathrm{\Theta }}},`$
$`\stackrel{~}{r}`$ $`=`$ $`|a+b\mathrm{\Theta }|r,`$
$`\stackrel{~}{g}^2`$ $`=`$ $`(a+b\mathrm{\Theta })g^2,`$
$`\stackrel{~}{\mathrm{\Phi }}`$ $`=`$ $`(a+b\mathrm{\Theta })^2\mathrm{\Phi }{\displaystyle \frac{b(a+b\mathrm{\Theta })}{2\pi r^2}},`$ (15)
where $`a,b,c,d`$ satisfy the Diophantine relation $`adbc=1`$. These relations guarantee invariance of the YangโMills action under the Morita transformation.
For the particular case in which the original $`U(1)`$ gauge theory has a rational-valued noncommutativity parameter $`\mathrm{\Theta }=\frac{c}{N}`$, the $`SL(2,)`$ element above (having $`b=m,d=N`$) yields a vanishing $`\stackrel{~}{\mathrm{\Theta }}`$ and the dual theory is a commutative $`U(N)`$ YangโMills theory with coupling constant
$$g^2=\frac{g_{}^{}{}_{}{}^{2}}{N}$$
(16)
defined on a torus of radius
$$r=\frac{r^{}}{N}$$
(17)
with a non-trivial magnetic flux. This relation will play an important role in the following, as it implies that the target torus is smaller than the original one since its area shrinks by a factor $`N^2`$. When considering Wilson loops on the torus, we will have to deal with this shrinking. We will be primarily concerned with this form of the duality, since some explicit nonperturbative computations can be done on the commutative torus. Moreover, since any irrational number is the limit of an infinite sequence of rational numbers, we expect that the results obtained in this way hold at general values of $`\mathrm{\Theta }`$, or equivalently that they are continuous functions of the noncommutativity parameter. In the following we will conventionally refer to the $`U(1)`$ noncommutative gauge theory with primed variables and to its commutative Morita dual gauge theory with unprimed variables.
We will make use of Morita duality to compute correlators in the noncommutative gauge theory by computing their commutative counterparts. This idea has been exploited in and it enables one to perform calculations very explicitly. To complete this program, we have to exhibit the transformation law of the observables under Morita duality. This problem has been solved in . For example, to the operator (6) we associate the commutative Wilson loop
$$๐ช\left(๐\right)=_{๐^2}\mathrm{d}^2๐\mathrm{Tr}_N^{}\mathrm{P}\mathrm{exp}\left(\mathrm{i}_๐d\xi ^\mu A_\mu \left(๐+๐\right)\right)$$
(18)
where $`\mathrm{Tr}_N^{}`$ is the trace in the fundamental representation of the $`U(N)`$ gauge group, and the commutative path-ordering operator P is defined as the analog of (7) with star-products replaced by ordinary matrix products and the noncommutative gauge fields $`๐_\mu `$ by their commutative counterparts $`A_\mu `$. The identification of the observables is then given by
$$๐ช_{}\left(๐\right)=N๐ช\left(๐\right).$$
(19)
As this equation is of fundamental importance to us, let us briefly review its derivation following .
Consider commutative pure $`U(N)`$ gauge theory on $`๐^2`$ with $`m`$ units of background magnetic flux. A non-trivial flux implies that the gauge fields obey twisted boundary conditions on the torus. They are solved by the Fourier expansions
$$A_\mu (๐)=\underset{๐^2}{}a_{๐;\mu }Q^{cq^1}P^{q^2}\mathrm{e}^{\pi \mathrm{i}cq^1q^2/N}\mathrm{e}^{\mathrm{i}๐๐/Nr}$$
(20)
where $`P`$ and $`Q`$ are the usual shift and clock matrices of rank $`N`$ which obey the commutation relation $`PQ=\mathrm{e}^{2\pi \mathrm{i}/N}QP`$. The $`n^{\mathrm{th}}`$ term in the expansion of the Wilson loop observable $`N๐ช(๐)`$ given by (18) then takes the form
$`\mathrm{i}^nN{\displaystyle _0^{2\pi r}}dx^1{\displaystyle _0^{2\pi r}}dx^2{\displaystyle \underset{i=1}{\overset{n}{}}}{\displaystyle _0^{s_{i1}}}ds_i{\displaystyle \underset{๐_i^2}{}}\dot{\xi }^{\mu _i}(s_i)a_{๐_i;\mu _i}`$
$`\times \mathrm{Tr}_N^{}\left(Q^{cq_1^1}P^{q_1^2}\mathrm{}Q^{cq_n^1}P^{q_n^2}\right){\displaystyle \underset{i=1}{\overset{n}{}}}\mathrm{e}^{\pi \mathrm{i}cq_i^1q_i^2/N}\mathrm{e}^{\mathrm{i}๐_i(๐+๐(s_i))/Nr}`$ (21)
where we have defined $`s_0=1`$. By using the commutation properties of the clock and shift matrices it follows that
$$\mathrm{Tr}_N^{}\left(Q^{cq_1^1}P^{q_1^2}\mathrm{}Q^{cq_n^1}P^{q_n^2}\right)=\mathrm{Tr}_N^{}\left(Q^{c(q_1^1+\mathrm{}+q_n^1)}P^{q_1^2+\mathrm{}+q_n^2}\right)\mathrm{e}^{\frac{2\pi \mathrm{i}c}{N}\underset{i>j}{}q_i^1q_j^2}.$$
(22)
The trace on the right-hand side of this equation vanishes unless $`q_1^\mu +\mathrm{}+q_n^\mu =q^\mu N`$ for some integers $`q^\mu `$. If these conditions are satisfied, then since $`P^N=Q^N=11_N`$ the trace is equal to $`N`$. Finally, due to these momentum constraints the integrals over $`๐^2`$ give Kronecker delta-functions, and we can thereby formally rewrite (2) as
$`\mathrm{i}^n{\displaystyle _0^{2\pi r^{}}}dx^1{\displaystyle _0^{2\pi r^{}}}dx^2{\displaystyle \underset{i=1}{\overset{n}{}}}{\displaystyle _0^{s_{i1}}}ds_i{\displaystyle \underset{๐_i^2}{}}\dot{\xi }^{\mu _i}(s_i)a_{๐_i;\mu _i}\mathrm{e}^{\mathrm{i}๐_i๐(s_i)/r^{}}\mathrm{e}^{\pi \mathrm{i}cq_i^1q_i^2/N}`$
$`\times \mathrm{e}^{\frac{2\pi \mathrm{i}c}{N}\underset{i>j}{}q_i^1q_j^2}\mathrm{e}^{\mathrm{i}(q_1^1+\mathrm{}+q_n^1)x^1/r^{}}\mathrm{e}^{\mathrm{i}(q_1^2+\mathrm{}+q_n^2)x^2/r^{}}.`$ (23)
By repeatedly using the properties of the star-product, it is straightforward to recast (2) into the form
$`\mathrm{i}^n{\displaystyle _0^{2\pi r^{}}}dx^1{\displaystyle _0^{2\pi r^{}}}dx^2{\displaystyle \underset{i=1}{\overset{n}{}}}{\displaystyle _0^{s_{i1}}}ds_i{\displaystyle \underset{๐_i^2}{}}\dot{\xi }^{\mu _i}(s_i)a_{๐_i;\mu _i}`$
$`\times \mathrm{e}^{\mathrm{i}๐_1(๐+๐(s_1))/r^{}}\mathrm{}\mathrm{e}^{\mathrm{i}๐_n(๐+๐(s_n))/r^{}}`$ (24)
which matches the $`n^{\mathrm{th}}`$ term in the expansion of (6).
It follows that the Morita correspondence between closed Wilson lines, unlike the case of open Wilson lines , does not involve any transformation of the quantum numbers associated with the loop. Thus if we take a closed Wilson line in the noncommutative gauge theory which encloses an area $`\rho `$, then it maps into a closed Wilson line in the Morita equivalent commutative gauge theory with the same shape and the same area $`\rho `$, because in the steps that led to (19) the parametrization $`๐(s)`$ of the loop never played any role. But, because of the relation (17), while the area of the loop remains fixed, the area of the target torus is smaller. Thus the path can start to wind in a non-trivial way and in general self-intersections of the loop may appear in the dual gauge theory. While in the above analysis we have focused only on the particular case when the noncommutativity parameter is rational-valued, our conclusions concerning loop areas also hold in the more general case when $`\mathrm{\Theta }`$ is an irrational number .
## 3 Dual Loop Correlators: Rational Case
In this section we will explicitly compute some noncommutative Wilson loop correlators. After choosing the closed path, we will apply Morita duality to the observable (6) thus mapping it into a Wilson loop on a new torus. Since the target torus is smaller, the dual Wilson loop can wind around the torus and self-intersect in a very non-trivial way. To keep the discussion as simple as possible and to provide concrete examples, we will restrict ourselves to the case where the noncommutativity parameter is a rational number and thus the target torus is a commutative space. We will describe the case of irrational $`\mathrm{\Theta }`$ in the next section.
As we will find, the actual computation of the Wilson loop average depends heavily on the geometrical shape and orientation of the closed path in the original torus. As the area of the target torus is smaller than that of the original torus, under Morita equivalence a simple closed curve can transform into a rather intricate (self-intersecting) loop. In particular, it can happen that two contractible loops, of different shape but equal area, transform into two topologically inequivalent loops. The same is true of contours which have the same area and shape but different orientations on the torus. Paths which exhibit such behaviour can in principle have very different quantum averages. In this section we shall argue that this is indeed the case. We will focus on the behaviour of a loop of fixed area on a shrinking torus. A complete characterization of this phenomenon would take us deep into non-planar graph theory, which is beyond the scope of this paper. Instead, we will develop a working knowledge of this behaviour by discussing a necessary criterion for a path to become self-intersecting under Morita equivalence.
Consider a closed contractible path $`๐`$ with no self-intersections on a square torus of radius $`r`$. Let $`๐=(\xi ^1,\xi ^2):[0,1]๐^2`$ be a parametrization of $`๐`$. Introduce the two characteristic lengths
$$\mathrm{}^\mu \left(๐\right)=\underset{s,s^{}[0,1]}{sup}\left|\xi ^\mu (s)\xi ^\mu (s^{})\right|,\mu =1,2$$
(25)
which measure the width and the height of the loop. Given $`a,b`$ and $`\mathrm{\Theta }`$, consider the behaviour of the path $`๐`$ as the torus shrinks to a torus of radius
$$r_\mathrm{c}^{a,b}=|a+b\mathrm{\Theta }|r.$$
(26)
If $`\mathrm{}^1(๐)`$ and $`\mathrm{}^2(๐)`$ are both smaller than $`r_\mathrm{c}^{a,b}`$, then the path will not self-intersect on the dual torus. We thereby arrive at a necessary condition that the loop $`๐`$ should satisfy in order to self-intersect on the dual torus given by
$$\mathrm{}^\mu \left(๐\right)r_\mathrm{c}^{a,b}$$
(27)
for $`\mu =1`$ or $`\mu =2`$. We stress that the bound (27) is not a sufficient condition. It is not difficult to draw loops that do indeed satisfy the bound (27) but do not self-intersect on the dual torus. In fact, a little practice with drawing loops on the torus shows how involved the task of providing necessary and sufficient conditions for self-intersections is.
Through Morita equivalence, we can compute quantum averages of Wilson loops on a noncommutative torus by mapping the observable to a smaller but commutative torus and then resorting to the known techniques of commutative YangโMills theory. But, according to (27), given a loop on the original torus, there exists a critical radius $`r_\mathrm{c}^{a,b}`$ such that the loop can become a self-intersecting closed contour on the target torus. In the Morita transformation to commutative gauge theory, the critical radius is $`r_\mathrm{c}^{a,b}=r/N`$. We will now explore some of the physical consequences of this statement.
### 3.1 General Construction
In the previous section we have reduced the problem of evaluating a noncommutative Wilson loop correlator to the computation of its Morita dual correlator. We will now describe how this is done in practice. According to , the partition function of two-dimensional YangโMills theory on a torus $`๐^2`$ (and more generally on any Riemann surface) can be conveniently evaluated through a combinatorial approach wherein one covers the surface with a set of simplices (plaquettes) and works with the lattice regularization of the original gauge theory . The continuum limit is recovered in the limit as the triangulation becomes finer. The partition function is invariant under subdivision of the lattice, and thus the lattice regularization provides a concrete definition of two dimensional quantum YangโMills theory. In the lattice gauge theory, the partition function is a sum over local factors associated to all of the plaquettes, which each have the topology of a disk. It is natural to associate to each plaquette $`D_\lambda `$ the holonomies $`U_\sigma `$ of a gauge connection $`A`$ along its links $`L_\sigma `$. Gauge invariance requires that the local factor corresponding to each plaquette be a class function of the holonomies.
The local factors $`\mathrm{\Gamma }(๐ฐ_\lambda ;D_\lambda )`$ associated to each simplex $`D_\lambda `$ of area $`\rho _\lambda `$ are given by
$$\mathrm{\Gamma }(๐ฐ_\lambda ;D_\lambda )=\underset{R_\lambda }{}dimR_\lambda \mathrm{e}^{\frac{g^2\rho _\lambda }{2}C_2(R_\lambda )}\chi _{R_\lambda }^{}\left(๐ฐ_\lambda \right)$$
(28)
where the sum runs through all isomorphism classes of $`U(N)`$ representations, $`C_2(R_\lambda )`$ is the second Casimir invariant of the representation $`R_\lambda `$, and $`\chi _{R_\lambda }^{}(๐ฐ_\lambda )=\mathrm{Tr}_{R_\lambda }^{}๐ฐ_\lambda `$ are the characters of the representation $`R_\lambda `$ evaluated on the holonomy $`๐ฐ_\lambda =_\sigma U_\sigma `$ along the perimeter of the simplex $`D_\lambda `$ with respect to a fixed orientation of its edges. The factors appearing in the formula (28) can be understood as follows. The representation dimension $`dimR_\lambda `$ is a normalization factor which ensures that the holonomy around a loop of area $`\rho _\lambda `$ approaches $`1`$ as $`\rho _\lambda 0`$. The characters appear since they form a basis for the vector space of class functions. Finally, the exponential factor is essentially the exponential of the YangโMills hamiltonian in the representation basis with $`g^2`$ the YangโMills coupling constant. For a more detailed account see .
Let us now consider the vacuum expectation value of a Wilson loop in ordinary YangโMills theory. It is defined by the functional integral
$$W_{๐;R}(\rho _๐^{})=\mathrm{D}A\mathrm{e}^{S_{\mathrm{YM}}[A]}\mathrm{Tr}_R^{}\mathrm{P}\mathrm{exp}\left(\mathrm{i}_๐A\right),$$
(29)
where $`\rho _๐^{}`$ is the area enclosed by the path $`๐`$, $`S_{\mathrm{YM}}[A]=\frac{1}{2g^2}_{๐^2}\mathrm{d}^2๐\mathrm{Tr}_N^{}F^2`$ is the YangโMills action functional, and we have explicitly indicated the dependence on the representation $`R`$ of the character used to compute the holonomy of the connection $`A`$ around the path $`๐๐^2`$. In our case, the Wilson loop will always be taken to lie in the fundamental representation $`R=N`$ of the $`U(N)`$ gauge group.
The Wilson loop provides a natural division of the torus into plaquettes $`D_\lambda `$ bounded by line segments $`L_\sigma `$ (links in the lattice formulation) in which the loop is divided by its self-intersections. Each plaquette $`D_\lambda `$ has area $`\rho _\lambda `$ such that $`_\lambda \rho _\lambda =(2\pi r/N)^2`$ is the area of the Morita dual torus. In this way we can write (29) as
$$W_{๐;R}(\rho _๐^{})=\underset{\sigma }{}_{U(N)}\left[\mathrm{d}U_\sigma \right]\underset{\lambda }{}\mathrm{\Gamma }(๐ฐ_\lambda ;D_\lambda )\chi _R^{}\left(๐ฐ^1\right)$$
(30)
where $`๐ฐ=_\sigma U_\sigma `$ is the holonomy along the corresponding edges $`L_\sigma `$ and $`\left[\mathrm{d}U_\sigma \right]`$ is the invariant Haar measure on the $`U(N)`$ gauge group. In (30) we have implicitly assumed that each plaquette has the topology of a disk. If this is not the case, then it suffices to consider a finer triangulation of the torus. For a simplex of different topology, the local factor (28) becomes
$$\mathrm{\Gamma }(๐ฐ_\lambda ;D_\lambda )=\underset{R_\lambda }{}\left(dimR_\lambda \right)^{22h_\lambda b_\lambda }\mathrm{e}^{\frac{g^2\rho _\lambda }{2}C_2(R_\lambda )}\chi _{R_\lambda }^{}\left(๐ฐ_\lambda \right)$$
(31)
when the simplex $`D_\lambda `$ has $`h_\lambda `$ handles and $`b_\lambda `$ boundaries.
We can recast (30) in a simpler form by noticing that each group element $`U_\sigma `$ representing the holonomy along the edge $`L_\sigma `$ appears three times in the integral (once in the Wilson line insertion and once for each of the two simplices that has $`L_\sigma `$ as part of its boundary). Thus if we denote by $`R_\alpha (U)_b^a`$, $`a,b=1,\mathrm{},dimR_\alpha `$ the matrix representing the group element $`U`$ in the representation $`R_\alpha `$, then from the identity
$$\chi _{R_\alpha }^{}\left(UU^{}\right)=R_\alpha \left(U\right)_b^aR_\alpha \left(U^{}\right)_a^b$$
(32)
it follows that the computation of (30) reduces to the evaluation of integrals of the form $`_{U(N)}[\mathrm{d}U]R_\alpha (U)_b^aR_\beta (U)_d^cR_\gamma (U)_f^e`$. Such group integrals give information about the fusion numbers $`\mathrm{N}_{R_\alpha R_\beta }^{R_\gamma }`$ which count the multiplicity of the irreducible representation $`R_\gamma `$ in the ClebschโGordan decomposition $`R_\alpha R_\beta =_{R_\gamma }\mathrm{N}_{R_\alpha R_\beta }^{R_\gamma }R_\gamma `$. We can collect these coefficients into factors associated with each vertex of the triangulation which combine into a local object. We may thereby write a final compact expression for the Wilson loop average as
$$W_{๐;R}(\rho _๐^{})=\underset{R_\lambda }{}\underset{\epsilon _\sigma }{}\underset{\lambda }{}(dimR_\lambda )^{22h_\lambda b_\lambda }\mathrm{e}^{\frac{g^2\rho _\lambda }{2}C_2(R_\lambda )}\underset{\delta }{}\mathrm{G}_\delta (R,R_\lambda ;\epsilon _\sigma )$$
(33)
where the index $`\delta `$ runs over all vertices of the lattice, while $`\epsilon _\sigma `$ runs over a basis for the vector space of intertwiners between the representations $`R_\gamma `$ and $`R_\alpha R_\beta `$. In the particular case that the vertex $`\delta `$ is four-valent, the local factor $`\mathrm{G}_\delta `$ is a $`6j`$-symbol . We will see explicitly how this works in some concrete examples below.
When the commutative gauge theory is related to noncommutative YangโMills theory on $`๐_\theta ^2`$ by Morita duality, one uses the global group isomorphism $`U(N)=U(1)\times SU(N)/_N`$ to cancel the $`U(1)`$ contribution to the partition function by the background abelian gauge field generated in the Morita transformation (15. One is then left with an $`SU(N)/_N`$ gauge theory in a certain discrete theta-vacuum of โt Hooft flux $`k=0,1,\mathrm{},N`$ which labels the isomorphism classes of principal $`SU(N)/_N`$ bundles over the torus. The $`U(1)`$ phases only contribute non-trivially when one sums over the topological sectors. For trivial bundles ($`k=0`$), all formulas above hold using $`SU(N)`$ representations in place of $`U(N)`$ representations. For non-trivial bundles ($`k0`$), one incorporates the background flux as follows. It contributes a factor $`\mathrm{exp}(\mathrm{i}_๐\alpha )`$ to the Wilson loop average , where $`\alpha `$ is any abelian gauge potential that gives rise to the constant background flux $`\mathrm{\Phi }=\mathrm{d}\alpha `$. Then the dependence of the correlator (29) on the $`k`$ units of magnetic flux follows from
$$\frac{1}{2\pi }_๐\alpha =\frac{1}{2\pi }_\mathrm{\Sigma }\mathrm{\Phi }=\frac{k}{N},$$
(34)
where $`\mathrm{\Sigma }=๐`$ is any surface spanned by the loop $`๐`$. When $`๐`$ contains self-intersections, one has to give a precise meaning to this integration. The path $`๐`$ admits a unique decomposition into simple closed paths as $`๐=_i๐_i`$. To obtain the appropriate flux factors one then splits the holonomy line integral over $`๐`$ into line integrals along the individual paths $`๐_i`$ and repeatedly applies (34). This modifies the characters in the above formulas by products of the characters $`\chi _{R_\lambda }^{}(\mathrm{e}^{2\pi \mathrm{i}k/N})`$ evaluated on elements in the center of the $`SU(N)`$ gauge group.
### 3.2 Simple Loops
We will now perform several explicit calculations in the rational noncommutative gauge theory. To illustrate the ideas in a somewhat general setting, we begin by comparing the Wilson loop correlators associated to two paths which have the generic forms depicted in Figs. 2 and 2 (The torus $`๐^2`$ is throughout represented as a square of sides $`r`$ with opposite edges identified). The paths enclose the same area $`\rho _1`$, but the second one satisfies the inequality (27).
Using (30) and (31) the first loop correlator can be associated with the formal expression
$`W_{๐_1;R}^k(\rho _1)`$ $`=`$ $`{\displaystyle \underset{R_1,R_2}{}}{\displaystyle \frac{dimR_1}{dimR_2}}\mathrm{e}^{\frac{g^2\rho _1}{2}C_2(R_1)\frac{g^2\rho _2}{2}C_2(R_2)}\chi _{R_1}^{}\left(\mathrm{e}^{2\pi \mathrm{i}k/N}\right)`$ (35)
$`\times {\displaystyle _{SU(N)}}[\mathrm{d}U_1]\chi _{R_1}^{}(U_1)\chi _{R_2}^{}\left(U_1^1\right)\chi _R^{}(U_1)`$
$`=`$ $`{\displaystyle \underset{R_1,R_2}{}}{\displaystyle \frac{dimR_1}{dimR_2}}\mathrm{N}_{R_1R}^{R_2}\mathrm{e}^{\frac{g^2\rho _1}{2}C_2(R_1)\frac{g^2\rho _2}{2}C_2(R_2)}\chi _{R_1}^{}\left(\mathrm{e}^{2\pi \mathrm{i}k/N}\right)`$
with $`\rho _1+\rho _2=(2\pi r/N)^2`$, where we have performed the group integrals to obtain the fusion coefficients of the three representations. The irreducible representations $`R`$ of $`SU(N)`$ can be labelled by decreasing sets $`๐^R=(n_1^R,\mathrm{},n_N^R)`$ of $`N`$ integers, $`+\mathrm{}>n_1^R>n_2^R>\mathrm{}>n_N^R>\mathrm{}`$, which satisfy the linear Casimir constraint $`_{a=1}^Nn_a^R=0`$. They determine the lengths of the rows of the corresponding Young tableaux. In particular, the integer $`_{a=1}^{N1}n_a^R`$ is the total number of boxes in the Young diagram describing $`R`$. In terms of these integers, the second Casimir invariant of $`R`$ can be written as
$$C_2(R)=C_2\left(๐^R\right)=\underset{a=1}{\overset{N}{}}\left(n_a^R\frac{N1}{2}\right)^2\frac{N}{12}\left(N^21\right)+\frac{\left(n_N^R\right)^2}{N},$$
(36)
while the dimension of $`R`$ can be expressed as the Vandermonde determinant
$$dimR=\mathrm{\Delta }\left(๐^R\right)=\underset{a<b}{}\left(n_a^Rn_b^R\right).$$
(37)
To compute the fusion numbers, we use the Weyl formula for the $`SU(N)`$ characters
$$\chi _R^{}(U)=\chi _{๐^R}^{}\left(\mathrm{e}^{2\pi \mathrm{i}๐}\right)=\frac{\underset{1a,bN}{det}\left[\mathrm{e}^{2\pi \mathrm{i}n_a^R\lambda _b}\right]}{\mathrm{\Delta }\left(\mathrm{e}^{2\pi \mathrm{i}๐}\right)}$$
(38)
where $`\mathrm{e}^{2\pi \mathrm{i}๐}=(\mathrm{e}^{2\pi \mathrm{i}\lambda _1},\mathrm{},\mathrm{e}^{2\pi \mathrm{i}\lambda _N})`$, $`\lambda _a[0,1]`$, $`a=1,\mathrm{},N`$ are the eigenvalues of the unitary matrix $`U`$ with $`_{a=1}^N\lambda _a=0\mathrm{mod}`$. Then the integration over the group variables $`U`$ can be transformed into an integration over the eigenvalues at the price of introducing a jacobian $`\mathrm{\Delta }(\mathrm{e}^{2\pi \mathrm{i}๐})^2`$. With these identifications, we can finally write (35) for $`R=N`$ the fundamental representation as
$`W_{๐_1;N}^k(\rho _1)`$ $`=`$ $`{\displaystyle \underset{๐^{R_1},๐^{R_2}}{}}{\displaystyle \frac{\mathrm{\Delta }\left(๐^{R_1}\right)}{\mathrm{\Delta }\left(๐^{R_2}\right)}}\mathrm{e}^{\frac{g^2\rho _1}{2}C_2(๐^{R_1})\frac{g^2\rho _2}{2}C_2(๐^{R_2})}\mathrm{e}^{\frac{2\pi \mathrm{i}k}{N}\underset{a=1}{\overset{N1}{}}n_a^{R_1}}`$ (39)
$`\times {\displaystyle \underset{a=1}{\overset{N}{}}}{\displaystyle _0^1}\mathrm{d}\lambda _a\delta \left(\underset{a=1}{\overset{N}{}}\lambda _a\right){\displaystyle \underset{c=1}{\overset{N}{}}}\mathrm{e}^{2\pi \mathrm{i}\lambda _c}`$
$`\times \underset{1a,bN}{det}\left[\mathrm{e}^{2\pi \mathrm{i}n_a^{R_1}\lambda _b}\right]\underset{1a,bN}{det}\left[\mathrm{e}^{2\pi \mathrm{i}n_a^{R_2}\lambda _b}\right].`$
We will return to the evaluation of this expression in Section 5.
The calculation is much different for the second path. Let us associate to the domains depicted in Fig. 3 the local factors (31) given by
$`\mathrm{\Gamma }(๐ฐ_1;D_1)={\displaystyle \underset{R_1}{}}dimR_1\mathrm{e}^{\frac{g^2\rho _1^{}}{2}C_2(R_1)}\chi _{R_1}^{}(U_2U_4),`$
$`\mathrm{\Gamma }(๐ฐ_2;D_2)={\displaystyle \underset{R_2}{}}dimR_2\mathrm{e}^{\frac{g^2\rho _2^{}}{2}C_2(R_2)}\chi _{R_2}^{}\left(U_1U_4^1U_3U_2^1\right),`$
$`\mathrm{\Gamma }(๐ฐ_3;D_3)={\displaystyle \underset{R_3}{}}{\displaystyle \frac{1}{dimR_3}}\mathrm{e}^{\frac{g^2\rho _3^{}}{2}C_2(R_3)}\chi _{R_3}^{}\left(U_1^1\right)\chi _{R_1}^{}\left(U_3^1\right),`$ (40)
where the dual area parameters obey $`\rho _1^{}+\rho _2^{}+\rho _3^{}=(2\pi r/N)^2`$ and $`2\rho _1^{}+\rho _2^{}=\rho _1`$. The last factor can be understood by regarding the contribution from the third simplex as a cylinder amplitude whose initial and final states are parametrized by the holonomies $`U_1`$ and $`U_3`$.
Then the general formula (30) for the path $`๐_2`$ becomes
$`W_{๐_2;R}^k(\rho _1)`$ $`=`$ $`{\displaystyle \underset{R_1,R_2,R_3}{}}{\displaystyle \frac{dimR_1dimR_2}{dimR_3}}\mathrm{e}^{\frac{g^2\rho _1^{}}{2}C_2(R_1)\frac{g^2\rho _2^{}}{2}C_2(R_2)\frac{g^2\rho _3^{}}{2}C_2(R_3)}`$ (41)
$`\times {\displaystyle \underset{\sigma =1}{\overset{4}{}}}{\displaystyle _{SU(N)}}[\mathrm{d}U_\sigma ]\chi _{R_1}^{}(U_2U_4)\chi _{R_2}^{}\left(U_1U_4^1U_3U_2^1\right)\chi _{R_3}^{}\left(U_1^1\right)\chi _{R_3}^{}\left(U_3^1\right)`$
$`\times \chi _R^{}(U_1U_2U_3U_4)\chi _{R_1}^{}\left(\mathrm{e}^{2\pi \mathrm{i}k/N}\right)^2\chi _{R_2}^{}\left(\mathrm{e}^{2\pi \mathrm{i}k/N}\right),`$
where the โt Hooft flux factors arise from the decomposition of the holonomy integral
$$_{๐_2}\alpha =\underset{\sigma =1}{\overset{4}{}}_{L_\sigma }\alpha =\left(_{L_1L_4^1L_3L_2^1}+2_{L_4L_2}\right)\alpha $$
(42)
and the line segment $`L_\sigma `$ refers to the path labelled by the holonomy $`U_\sigma `$ in Fig. 3. Thus one of the central characters squares in (41). Employing the same $`SU(N)`$ representation machinery used to arrive at (39), we can rewrite (41) as
$`W_{๐_2;R}^k(\rho _1)`$ $`=`$ $`{\displaystyle \underset{๐^{R_1},๐^{R_2},๐^{R_3}}{}}{\displaystyle \frac{\mathrm{\Delta }\left(๐^{R_1}\right)\mathrm{\Delta }\left(๐^{R_2}\right)}{\mathrm{\Delta }\left(๐^{R_3}\right)}}\mathrm{e}^{\frac{g^2\rho _1^{}}{2}C_2(๐^{R_1})\frac{g^2\rho _2^{}}{2}C_2(๐^{R_2})\frac{g^2\rho _3^{}}{2}C_2(๐^{R_3})}`$ (43)
$`\times {\displaystyle \underset{\sigma =1}{\overset{4}{}}}{\displaystyle _{SU(N)}}[\mathrm{d}U_\sigma ]\chi _{R_1}^{}(U_2U_4)\chi _{R_2}^{}\left(U_1U_4^1U_3U_2^1\right)\chi _{R_3}^{}\left(U_1^1\right)\chi _{R_3}^{}\left(U_3^1\right)`$
$`\times \chi _R^{}(U_1U_2U_3U_4)\mathrm{e}^{\frac{2\pi \mathrm{i}k}{N}\underset{a=1}{\overset{N1}{}}(2n_a^{R_1}+n_a^{R_2})}.`$
This is as far as we can proceed with general expressions for the Wilson loop correlators (39) and (43). Superficially, these two analytic expressions look quite different. For example, while (39) depends only on the loop area $`\rho _1`$, the function (43) effectively depends on two independent areas, say $`\rho _1`$ and $`\rho _1^{}`$. If the dependence on $`\rho _1^{}`$ is non-trivial, then evidently the two loop correlators are distinct, even though in the original theory they enclosed the same area $`\rho _1`$. To perform a more direct comparison of these correlation functions and get an idea of the nature of this extra area dependence, let us simplify matters enormously by turning to the special example of $`SU(2)`$ gauge theory. In this case we may appeal to various well-known angular momentum identities from the representation theory of the group $`SU(2)`$ .
Irreducible representations $`R_j`$ of $`SU(2)`$ are labelled by an angular momentum quantum number $`j\frac{1}{2}_0`$. The dimension of $`R_j`$ is given by $`dimR_j=2j+1`$, the quadratic Casimir invariant is $`C_2(R_j)=j(j+1)`$, and the total number of boxes in the Young diagram representing $`R_j`$ is $`2j`$. The integrations over the group variables in (35) give the fusion numbers $`\mathrm{N}_{j_1j}^{j_2}`$ which count the multiplicity of the irreducible representations $`R_{j_2}`$ in the ClebschโGordan decomposition of $`R_{j_1}R_j`$. These coefficients are equal to $`1`$ if $`|j_1j|j_2j_1+j`$ and $`0`$ otherwise. Thus we can write the quantum average (35) for $`SU(2)`$ gauge group in the explicit form
$$W_{๐_1;j}^k(\rho _1)=\underset{j_1\frac{1}{2}_0}{}\underset{j_2=|j_1j|}{\overset{j_1+j}{}}(1)^{2j_1k}\frac{2j_1+1}{2j_2+1}\mathrm{e}^{\frac{g^2\rho _1}{2}j_1(j_1+1)\frac{g^2\rho _2}{2}j_2(j_2+1)}.$$
(44)
Now let us rewrite the expression (41) by using the explicit form of the characters for $`SU(2)`$ representations. In this case we can introduce as representation matrices the Wigner functions of angular momentum $`j`$, so that (32) becomes
$$\chi _{R_j}^{}\left(UU^{}\right)=๐ฃ_{mm^{}}^j\left(U\right)๐ฃ_{m^{}m}^j\left(U^{}\right)$$
(45)
with $`jm,m^{}j`$, where throughout we implicitly assume that repeated indices represented by lower case Latin letters are summed over. In this way we can better organize the integration over $`SU(2)`$ group variables and write (41) as
$`W_{๐_2;j}^k(\rho _1)`$ $`=`$ $`{\displaystyle \underset{j_1,j_2,j_3\frac{1}{2}_0}{}}(1)^{2j_2k}{\displaystyle \frac{(2j_1+1)(2j_2+1)}{2j_3+1}}`$ (46)
$`\times \mathrm{e}^{\frac{g^2\rho _1^{}}{2}j_1(j_1+1)\frac{g^2\rho _2^{}}{2}j_2(j_2+1)\frac{g^2\rho _3^{}}{2}j_3(j_3+1)}`$
$`\times {\displaystyle _{SU(2)}}[\mathrm{d}U_1]๐ฃ_{a_2b_2}^{j_2}(U_1)๐ฃ_{a_3a_3}^{j_3}\left(U_1^1\right)๐ฃ_{ab}^j(U_1)`$
$`\times {\displaystyle _{SU(2)}}[\mathrm{d}U_2]๐ฃ_{a_1b_1}^{j_1}(U_2)๐ฃ_{d_2a_2}^{j_2}\left(U_2^1\right)๐ฃ_{bc}^j(U_2)`$
$`\times {\displaystyle _{SU(2)}}[\mathrm{d}U_3]๐ฃ_{c_2d_2}^{j_2}(U_3)๐ฃ_{b_3b_3}^{j_3}\left(U_3^1\right)๐ฃ_{cd}^j(U_3)`$
$`\times {\displaystyle _{SU(2)}}[\mathrm{d}U_4]๐ฃ_{b_1a_1}^{j_1}(U_4)๐ฃ_{b_2c_2}^{j_2}\left(U_4^1\right)๐ฃ_{da}^j(U_4).`$
If we regard the path drawn in Fig. 3 as a triangulation of the torus as before, then each Wigner function $`๐ฃ_{mm^{}}^j(U)`$ is associated with an oriented edge of the triangulation, with the first index $`m`$ representing the origin of the line segment and the second index $`m^{}`$ representing its endpoint. The reason for this identification is that it is more convenient to understand the quantum average (46) as a product of contributions arising from the vertices of the triangulation, rather than as integrals over edge variables. This procedure implements the general construction of Section 3.1 and provides an explicit realization of the expression (33).
To this end we use the formula
$$[{}_{m_1}{}^{j_1}{}_{m_2}{}^{j_2}{}_{m_3}{}^{j_3}][{}_{m_1^{}}{}^{j_1}{}_{m_2^{}}{}^{j_2}{}_{m_3^{}}{}^{j_3}]=\frac{2j_3+1}{8\pi ^2}_{SU(2)}[\mathrm{d}U]๐ฃ_{m_1m_1^{}}^{j_1}(U)๐ฃ_{m_2m_2^{}}^{j_2}(U)\overline{๐ฃ_{m_3m_3^{}}^{j_3}(U)}$$
(47)
relating the integral over edge variables to a product of two ClebschโGordan coefficients, each one associated with an endpoint of the given edge. We can now perform the integration over group variables in (46) and collect together the ClebschโGordan coefficients associated to each vertex, which in the present case are all of valence $`4`$. The crucial identity is
$`{\displaystyle \underset{m_1,m_2,m_3,m_{12},m_{23}}{}}[{}_{m_{12}}{}^{j_{12}}{}_{m_3}{}^{j_3}{}_{m}{}^{j}][{}_{m_1}{}^{j_1}{}_{m_2}{}^{j_2}{}_{m_{12}}{}^{j_{12}}][{}_{m_1}{}^{j_1}{}_{m_{23}}{}^{j_{23}}{}_{m^{}}{}^{j^{}}][{}_{m_2}{}^{j_2}{}_{m_3}{}^{j_3}{}_{m_{23}}{}^{j_{23}}]`$
$`=\delta _{jj^{}}\delta _{mm^{}}(1)^{j_1+j_2+j_3+j}\sqrt{(2j_{12}+1)(2j_{23}+1)}\{{}_{j_3}{}^{j_1}{}_{j}{}^{j_2}{}_{j_{23}}{}^{j_{12}}\}.`$ (48)
We have introduced the classical Wigner $`6j`$-symbol whose explicit form is provided by the Racah formula. We will not require this detailed expression here, except for noting that the square of the $`6j`$-symbol in (3.2) is proportional to a product of completely symmetric combinatorial factors $`\mathrm{}(j_1,j_2,j_{12})\mathrm{}(j_1,j,j_{23})\mathrm{}(j_3,j_2,j_{23})\mathrm{}(j_3,j,j_{12})`$ which are each non-vanishing only if the triangle inequality
$$\mathrm{}(j_1,j_2,j_3):j_1j_2+j_3,j_2j_1+j_3,j_3j_1+j_2,j_1+j_2+j_3_0+\frac{1}{2}$$
(49)
is obeyed by the corresponding angular momenta. In computing (46) the triangle inequalities imply that the average is non-zero only for half-integer spin $`j`$, in which case the loop correlator is given explicitly by
$`W_{๐_2;j}^k(\rho _1)`$ $`=`$ $`{\displaystyle \underset{2j_2j+\frac{1}{2}}{}}{\displaystyle \underset{\genfrac{}{}{0pt}{}{\alpha <\beta }{\alpha 2}}{}}{\displaystyle \underset{\genfrac{}{}{0pt}{}{j_\alpha =|j_\beta j|}{jj_\alpha +j_\beta _0}}{\overset{j_\beta +j}{}}}(1)^{2j_2k}\frac{(2j_1+1)(2j_2+1)}{2j_3+1}\{{}_{j}{}^{j}{}_{j_2}{}^{j_1}{}_{j_3}{}^{j_2}\}^2`$ (50)
$`\times \mathrm{e}^{\frac{g^2\rho _1^{}}{2}j_1(j_1+1)\frac{g^2\rho _2^{}}{2}j_2(j_2+1)\frac{g^2\rho _3^{}}{2}j_3(j_3+1)}`$
up to an irrelevant overall numerical factor.
The key point now is that the expressions (44) and (50), while bearing certain similarities, are very different. In particular, the correlator (50) appears to be a non-trivial function of the extra area $`\rho _1^{}`$. After writing the areas of the second and third simplices as $`\rho _2^{}=\rho _12\rho _1^{}`$ and $`\rho _3^{}=\rho _2+\rho _1^{}`$, we can differentiate (50) with respect to $`\rho _1^{}`$ to get
$`{\displaystyle \frac{W_{๐_2;j}^k(\rho _1)}{\rho _1^{}}}`$ $`=`$ $`{\displaystyle \frac{g^2}{2}}{\displaystyle \underset{2j_2j+\frac{1}{2}}{}}{\displaystyle \underset{\genfrac{}{}{0pt}{}{\alpha <\beta }{\alpha 2}}{}}{\displaystyle \underset{\genfrac{}{}{0pt}{}{j_\alpha =|j_\beta j|}{jj_\alpha +j_\beta _0}}{\overset{j_\beta +j}{}}}(1)^{2j_2k}\{{}_{j}{}^{j}{}_{j_2}{}^{j_1}{}_{j_3}{}^{j_2}\}^2`$ (51)
$`\times \mathrm{e}^{\frac{g^2\rho _1}{2}j_2(j_2+1)\frac{g^2\rho _2}{2}j_3(j_3+1)}\mathrm{e}^{\frac{g^2\rho _1^{}}{2}[j_1(j_1+1)2j_2(j_2+1)+j_3(j_3+1)]}`$
$`\times \frac{(2j_1+1)(2j_2+1)\left[j_1(j_1+1)2j_2(j_2+1)+j_3(j_3+1)\right]}{2j_3+1}.`$
We have not been able to rigorously prove that this quantity is non-vanishing. But we have also not been able to find any angular momentum identities implying that it is $`0`$, and we strongly doubt the existence of any such identity. Asymptotically, while the $`6j`$-symbol has an exponential decay for certain configurations of large angular momenta , generically it has only a trigonometric behaviour for large $`j`$โs. When $`\rho _1^{}0`$, one can construct an area-preserving diffeomorphism on the noncommutative torus which changes $`\rho _1^{}`$ and therefore likely gives a different correlator on the commutative torus. Thus the convergent series (51) does not appear to produce a vanishing result. This heuristic argument is strong evidence in favour of the non-vanishing of the expression (51). We therefore propose that for the class of simple loops considered here, the corresponding Wilson averages are strongly dependent on the shapes and even the orientations of the contours on $`๐_\theta ^2`$.
### 3.3 Circular Loops
To explore further the shape and orientation dependence of rational noncommutative Wilson loops, let us now consider a more specific smooth path with the circular geometry of Fig. 5. Under Morita equivalence it is mapped to the complicated self-intersecting path of Fig. 5.
Using the general formula (30), we must associate to each simplex $`D_\lambda `$, which in this case all have the topology of a disk, the local factor (28). We label the simplices and the corresponding edges, each with their proper orientation, as shown in Fig. 5. We take the Wilson loop in the representation $`R`$. With this notation, the circular Wilson loop correlator then reads
$`W_{;R}^k(\rho _1)`$ $`=`$ $`{\displaystyle \underset{R_1,\mathrm{},R_8}{}}{\displaystyle \underset{\lambda =1}{\overset{8}{}}}dimR_\lambda \mathrm{e}^{\frac{g^2\rho _\lambda ^{}}{2}C_2(R_\lambda )}\chi _{R_1}^{}\left(\mathrm{e}^{2\pi \mathrm{i}k/N}\right)^3\chi _{R_2}^{}\left(\mathrm{e}^{2\pi \mathrm{i}k/N}\right)^2`$ (52)
$`\times \chi _{R_3}^{}\left(\mathrm{e}^{2\pi \mathrm{i}k/N}\right)^2\chi _{R_4}^{}\left(\mathrm{e}^{2\pi \mathrm{i}k/N}\right)^2\chi _{R_5}^{}\left(\mathrm{e}^{2\pi \mathrm{i}k/N}\right)^2\chi _{R_6}^{}\left(\mathrm{e}^{2\pi \mathrm{i}k/N}\right)`$
$`\times \chi _{R_7}^{}\left(\mathrm{e}^{2\pi \mathrm{i}k/N}\right){\displaystyle \underset{\sigma =1}{\overset{16}{}}}{\displaystyle _{SU(N)}}[\mathrm{d}U_\sigma ]\chi _{R_1}^{}\left(U_1U_5U_9U_{14}\right)\chi _{R_2}^{}\left(U_2U_8U_5^1\right)`$
$`\times \chi _{R_3}^{}\left(U_6U_{12}U_9^1\right)\chi _{R_4}^{}\left(U_{16}U_{13}^1U_{10}\right)\chi _{R_5}^{}\left(U_{14}U_4U_1^1\right)`$
$`\times \chi _{R_6}^{}\left(U_{15}U_{10}^1U_{12}^1U_7U_2^1U_4^1\right)\chi _{R_7}^{}\left(U_8^1U_3U_{14}^1U_{16}^1U_{11}U_6^1\right)`$
$`\times \chi _{R_8}^{}\left(U_7^1U_{11}^1U_{15}^1U_3^1\right)\chi _R^{}\left(\underset{\sigma =1}{\overset{16}{}}U_\sigma \right)`$
where we have used
$$_{}\alpha =\underset{\sigma =1}{\overset{16}{}}_{L_\sigma }\alpha =\left(3_{D_1}+2_{D_2}+2_{D_3}+2_{D_4}+2_{D_5}+_{D_6}+_{D_7}\right)\alpha $$
(53)
and the dual areas $`\rho _\lambda ^{}`$ obey
$$\underset{\lambda =1}{\overset{8}{}}\rho _\lambda ^{}=\left(\frac{2\pi r}{N}\right)^2,4\rho _1^{}+3\rho _2^{}+3\rho _3^{}+3\rho _4^{}+3\rho _5^{}+2\rho _6^{}+2\rho _7^{}+\rho _8^{}=\rho _1.$$
(54)
As above, in order to be as explicit as possible we will limit the analysis to the case of an $`SU(2)`$ gauge group. We write the characters in terms of Wigner functions and integrate over each group variable individually using (47). Then we collect together all ClebschโGordan coefficients relative to each vertex, which are again all of valence $`4`$, to get
$`W_{;j}^k(\rho _1)={\displaystyle \underset{j_1,\mathrm{},j_8\frac{1}{2}_0}{}}\frac{(1)^{2(j_1+j_6+j_7)k}}{\left[(2j_1+1)(2j_6+1)(2j_7+1)(2j_8+1)\right]^3}\mathrm{e}^{\underset{\lambda =1}{\overset{8}{}}\frac{g^2\rho _\lambda ^{}}{2}j_\lambda (j_\lambda +1)}`$
$`\times \left([{}_{b_1}{}^{j_1}{}_{b}{}^{j}{}_{c_5}{}^{j_5}][{}_{b_1}{}^{j_1}{}_{e}{}^{j}{}_{a_2}{}^{j_2}][{}_{c_5}{}^{j_5}{}_{e}{}^{j}{}_{f_6}{}^{j_6}][{}_{a_2}{}^{j_2}{}_{b}{}^{j}{}_{f_6}{}^{j_6}]\right)\left([{}_{c_1}{}^{j_1}{}_{f}{}^{j}{}_{c_2}{}^{j_2}][{}_{c_2}{}^{j_2}{}_{i}{}^{j}{}_{a_2}{}^{j_7}][{}_{a_3}{}^{j_3}{}_{f}{}^{j}{}_{a_7}{}^{j_7}][{}_{c_1}{}^{j_1}{}_{i}{}^{j}{}_{a_3}{}^{j_3}]\right)`$
$`\times \left([{}_{d_1}{}^{j_1}{}_{q}{}^{j}{}_{c_3}{}^{j_3}][{}_{c_3}{}^{j_3}{}_{m}{}^{j}{}_{c_6}{}^{j_6}][{}_{c_4}{}^{j_4}{}_{q}{}^{j}{}_{c_6}{}^{j_6}][{}_{d_1}{}^{j_1}{}_{m}{}^{j}{}_{f_4}{}^{j_4}]\right)\left([{}_{a_1}{}^{j_1}{}_{n}{}^{j}{}_{b_4}{}^{j_4}][{}_{b_4}{}^{j_4}{}_{a}{}^{j}{}_{d_7}{}^{j_7}][{}_{a_5}{}^{j_5}{}_{n}{}^{j}{}_{d_7}{}^{j_7}][{}_{a_1}{}^{j_1}{}_{a}{}^{j}{}_{a_5}{}^{j_5}]\right)`$
$`\times \left([{}_{b_2}{}^{j_2}{}_{c}{}^{j}{}_{e_6}{}^{j_6}][{}_{e_6}{}^{j_6}{}_{h}{}^{j}{}_{a_8}{}^{j_8}][{}_{b_7}{}^{j_7}{}_{c}{}^{j}{}_{a_8}{}^{j_8}][{}_{a_2}{}^{j_2}{}_{h}{}^{j}{}_{b_7}{}^{j_7}]\right)\left([{}_{b_3}{}^{j_3}{}_{g}{}^{j}{}_{f_7}{}^{j_7}][{}_{f_7}{}^{j_7}{}_{l}{}^{j}{}_{b_8}{}^{j_8}][{}_{d_6}{}^{j_6}{}_{g}{}^{j}{}_{b_8}{}^{j_8}][{}_{b_3}{}^{j_3}{}_{l}{}^{j}{}_{d_6}{}^{j_6}]\right)`$
$`\times \left([{}_{a_4}{}^{j_4}{}_{r}{}^{j}{}_{b_6}{}^{j_6}][{}_{b_6}{}^{j_6}{}_{p}{}^{j}{}_{c_8}{}^{j_8}][{}_{a_4}{}^{j_4}{}_{p}{}^{j}{}_{e_7}{}^{j_7}][{}_{e_7}{}^{j_7}{}_{k}{}^{j}{}_{c_8}{}^{j_8}]\right)\left([{}_{b_5}{}^{j_5}{}_{o}{}^{j}{}_{c_7}{}^{j_7}][{}_{c_7}{}^{j_7}{}_{d}{}^{j}{}_{b_8}{}^{j_8}][{}_{a_6}{}^{j_6}{}_{o}{}^{j}{}_{d_8}{}^{j_8}][{}_{b_5}{}^{j_5}{}_{d}{}^{j}{}_{a_6}{}^{j_6}]\right)`$
where except for the $`j`$โs all Latin indices are implicitly summed over. Each term in parentheses is the local representation of a self-intersection on the Morita dual circle. It can be written in a more compact way by repeatedly applying the formula (3.2), as well as various symmetry properties of the ClebschโGordan coefficients that can be found in , to convert four ClebschโGordan coefficients into a $`6j`$-symbol. Up to an overall numerical factor one finally finds
$`W_{;j}^k(\rho _1)`$ $`=`$ $`{\displaystyle \underset{j_1,\mathrm{},j_8D_{}^j}{}}(1)^{2(j_1+j_6+j_7)k}{\displaystyle \underset{\lambda =1}{\overset{8}{}}}(2j_\lambda +1)\mathrm{e}^{\frac{g^2\rho _\lambda ^{}}{2}j_\lambda (j_\lambda +1)}`$ (56)
$`\times \{{}_{j}{}^{j}{}_{j_6}{}^{j_1}{}_{j_5}{}^{j_2}\}\{{}_{j}{}^{j}{}_{j_7}{}^{j_1}{}_{j_2}{}^{j_3}\}\{{}_{j}{}^{j}{}_{j_6}{}^{j_1}{}_{j_3}{}^{j_4}\}\{{}_{j}{}^{j}{}_{j_7}{}^{j_1}{}_{j_4}{}^{j_5}\}`$
$`\times \{{}_{j}{}^{j}{}_{j_8}{}^{j_2}{}_{j_6}{}^{j_7}\}\{{}_{j}{}^{j}{}_{j_8}{}^{j_3}{}_{j_7}{}^{j_6}\}\{{}_{j}{}^{j}{}_{j_8}{}^{j_4}{}_{j_7}{}^{j_6}\}\{{}_{j}{}^{j}{}_{j_8}{}^{j_5}{}_{j_7}{}^{j_6}\},`$
where by the triangle inequalities the sum over spins is restricted to the range
$$D_{}^j=\underset{\genfrac{}{}{0pt}{}{\alpha =2,3,4,5}{\beta =1,6,7}}{}D_{j_\alpha j_\beta }^j\underset{\alpha =6,7}{}D_{j_8j_\alpha }^j$$
(57)
with
$$D_{j_\alpha j_\beta }^j=\{|j_\beta j|j_\alpha j_\beta +j,j_\alpha +j_\beta j,j+j_\alpha +j_\beta _0+\frac{1}{2}\}.$$
(58)
In contrast to the intersecting Wilson loop average over the contour $`๐_2`$ of Section 3.2, the correlator (56) is generically non-vanishing for all angular momenta $`j\frac{1}{2}_0`$.
### 3.4 Square Loops
For our final explicit example, we will consider the case of the polygonal contour with the geometry of the square Wilson loop of Fig. 7.
After a Morita transformation this loop is mapped into the loop of Fig. 7. This dual path is much more complicated than the previously considered dual circle, because the edges which bound the inner square of Fig. 7 are covered twice in computing the Wilson loop holonomy using the combinatorial construction of Section 3.1. Thus the group elements associated with these particular edges will appear four times in (30), twice because of the Wilson loop holonomy and once for each of the two faces that are bounded by this edge due to (28). We would then need the generalization of the group integral (47) involving four group elements, but these generalizations are difficult to handle. Because of this technical difficulty, instead of computing the square Wilson loop of Fig. 7, we will perform an area-preserving deformation of the square contour as illustrated in Fig. 9. After a Morita transformation, this path is mapped to the loop of Fig. 9. Each simplex $`D_\lambda `$ in this case has the topology of a disk.
The condition that the loop of Fig. 9 encloses the same area as the loop of Fig. 7 implies for the areas $`\rho _\lambda ^{}`$ of the simplices $`D_\lambda `$ of Fig. 9 that
$$\rho _6^{}+\rho _7^{}=\rho _5^{}+\rho _8^{}.$$
(59)
The deformed square Wilson loop correlator (30) in the commutative dual gauge theory thereby reads
$`W_{\mathrm{};R}^k(\rho _1)`$ $`=`$ $`{\displaystyle \underset{R_1,\mathrm{},R_8}{}}{\displaystyle \underset{\lambda =1}{\overset{8}{}}}dimR_\lambda \mathrm{e}^{\frac{g^2\rho _\lambda ^{}}{2}C_2(R_\lambda )}\chi _{R_1}^{}\left(\mathrm{e}^{2\pi \mathrm{i}k/N}\right)^3\chi _{R_2}^{}\left(\mathrm{e}^{2\pi \mathrm{i}k/N}\right)`$ (60)
$`\times \chi _{R_3}^{}\left(\mathrm{e}^{2\pi \mathrm{i}k/N}\right)\chi _{R_5}^{}\left(\mathrm{e}^{2\pi \mathrm{i}k/N}\right)^2\chi _{R_6}^{}\left(\mathrm{e}^{2\pi \mathrm{i}k/N}\right)^2\chi _{R_7}^{}\left(\mathrm{e}^{2\pi \mathrm{i}k/N}\right)^2`$
$`\times \chi _{R_8}^{}\left(\mathrm{e}^{2\pi \mathrm{i}k/N}\right)^2{\displaystyle \underset{\sigma =1}{\overset{12}{}}}{\displaystyle _{SU(N)}}[\mathrm{d}U_\sigma ]\chi _{R_1}^{}\left(U_3U_4U_7U_{12}\right)\chi _{R_2}^{}\left(U_{11}U_9^1U_5U_1^1\right)`$
$`\times \chi _{R_3}^{}\left(U_2U_{10}^1U_8U_6^1\right)\chi _{R_4}^{}\left(U_8^1U_{11}^1U_2^1U_5^1\right)\chi _{R_5}^{}\left(U_1U_3^1\right)`$
$`\times \chi _{R_6}^{}\left(U_6U_4^1\right)\chi _{R_7}^{}\left(U_9U_7^1\right)\chi _{R_8}^{}\left(U_{10}U_{12}^1\right)\chi _R^{}\left(\underset{\sigma =1}{\overset{12}{}}U_\sigma \right)`$
where we have used
$$_{\mathrm{}}\alpha =\underset{\sigma =1}{\overset{12}{}}_{L_\sigma }\alpha =\left(3_{D_1}+2_{D_5}+2_{D_6}+2_{D_7}+2_{D_8}+_{D_2}+_{D_3}\right)\alpha $$
(61)
and the dual areas $`\rho _\lambda ^{}`$ obey, in addition to (59), the constraints
$$\underset{\lambda =1}{\overset{8}{}}\rho _\lambda ^{}=\left(\frac{2\pi r}{N}\right)^2,3\rho _1^{}+2\underset{\lambda =2}{\overset{8}{}}\rho _\lambda ^{}=\rho _1.$$
(62)
Again we take the gauge group to be $`SU(2)`$ and follow our combinatorial procedure. Each group integration is performed by using the formula (47). Each edge contributes to (60) with a ClebschโGordan coefficient for each one of its endpoints. The sum over edges (holonomies) is converted into a sum over vertices (collections of ClebschโGordan coefficients). However, now each vertex of the triangulation depicted in Fig. 9 is of valence $`6`$ and so will generally have associated to it a more complicated object than a $`6j`$-symbol. By collecting the ClebschโGordan coefficients for each of the four vertices, the quantum average (60) becomes
$`W_{\mathrm{};j}^k(\rho _1)`$ $`=`$ $`{\displaystyle \underset{j_1,\mathrm{},j_8\frac{1}{2}_0}{}}\frac{(1)^{2(j_1+j_2+j_3)k}(2j_1+1)}{(2j_2+1)(2j_3+1)(2j_4+1)^3}\mathrm{e}^{\underset{\lambda =1}{\overset{8}{}}\frac{g^2\rho _\lambda ^{}}{2}j_\lambda (j_\lambda +1)}`$ (63)
$`\times \left([{}_{a_5}{}^{j_5}{}_{a}{}^{j}{}_{a_2}{}^{j_2}][{}_{a_1}{}^{j_1}{}_{c}{}^{j}{}_{a_5}{}^{j_5}][{}_{a_1}{}^{j_1}{}_{a}{}^{j}{}_{b_8}{}^{j_8}][{}_{b_8}{}^{j_8}{}_{b}{}^{j}{}_{b_3}{}^{j_3}][{}_{b_3}{}^{j_3}{}_{c}{}^{j}{}_{c_4}{}^{j_4}][{}_{a_2}{}^{j_2}{}_{k}{}^{j}{}_{c_4}{}^{j_4}]\right)_\mathrm{A}`$
$`\times \left([{}_{b_5}{}^{j_5}{}_{b}{}^{j}{}_{d_2}{}^{j_2}][{}_{d_2}{}^{j_2}{}_{f}{}^{j}{}_{d_4}{}^{j_4}][{}_{a_3}{}^{j_3}{}_{b}{}^{j}{}_{d_4}{}^{j_4}][{}_{a_6}{}^{j_6}{}_{f}{}^{j}{}_{a_3}{}^{j_3}][{}_{b_1}{}^{j_1}{}_{d}{}^{j}{}_{a_6}{}^{j_6}][{}_{b_1}{}^{j_1}{}_{d}{}^{j}{}_{b_5}{}^{j_5}]\right)_\mathrm{B}`$
$`\times \left([{}_{c_1}{}^{j_1}{}_{g}{}^{j}{}_{a_7}{}^{j_7}][{}_{c_1}{}^{j_1}{}_{e}{}^{j}{}_{b_5}{}^{j_6}][{}_{b_6}{}^{j_6}{}_{g}{}^{j}{}_{d_3}{}^{j_3}][{}_{d_3}{}^{j_3}{}_{i}{}^{j}{}_{a_4}{}^{j_4}][{}_{c_2}{}^{j_2}{}_{e}{}^{j}{}_{c_4}{}^{j_4}][{}_{c_1}{}^{j_1}{}_{g}{}^{j}{}_{c_2}{}^{j_2}]\right)_\mathrm{C}`$
$`\times \left([{}_{d_1}{}^{j_1}{}_{h}{}^{j}{}_{b_7}{}^{j_7}][{}_{b_7}{}^{j_7}{}_{q}{}^{j}{}_{b_2}{}^{j_2}][{}_{b_2}{}^{j_2}{}_{l}{}^{j}{}_{b_4}{}^{j_4}][{}_{c_3}{}^{j_3}{}_{h}{}^{j}{}_{b_4}{}^{j_4}][{}_{a_8}{}^{j_8}{}_{q}{}^{j}{}_{c_3}{}^{j_3}][{}_{d_1}{}^{j_1}{}_{l}{}^{j}{}_{a_8}{}^{j_8}]\right)_\mathrm{D}`$
where for later reference we have labelled each vertex contribution with an upper case Latin letter.
Let us now consider in more detail the individual vertex contributions in (63). Their computation relies on a number of angular momentum identities which can all be found in . We begin with the vertex labelled โAโ. The first three ClebschโGordan coefficients can be summed by using the formula
$$\underset{\alpha ,\beta ,\delta }{}[{}_{\alpha }{}^{a}{}_{\beta }{}^{b}{}_{\gamma }{}^{c}][{}_{\delta }{}^{d}{}_{\beta }{}^{b}{}_{ฯต}{}^{e}][{}_{\alpha }{}^{a}{}_{\varphi }{}^{f}{}_{\delta }{}^{d}]=(1)^{b+c+d+f}\sqrt{(2c+1)(2d+1)}[{}_{\gamma }{}^{c}{}_{\varphi }{}^{f}{}_{ฯต}{}^{e}]\{{}_{e}{}^{a}{}_{f}{}^{b}{}_{d}{}^{c}\},$$
(64)
while the last three coefficients can be summed in a similar way thanks to the identity
$$\underset{\alpha ,\beta ,\delta }{}[{}_{\beta }{}^{b}{}_{\gamma }{}^{c}{}_{\alpha }{}^{a}][{}_{\beta }{}^{b}{}_{ฯต}{}^{e}{}_{\delta }{}^{d}][{}_{\alpha }{}^{a}{}_{\varphi }{}^{f}{}_{\delta }{}^{d}]=(1)^{a+b+e+f}\sqrt{\frac{2a+1}{2e+1}}(2d+1)[{}_{\gamma }{}^{c}{}_{\varphi }{}^{f}{}_{ฯต}{}^{e}]\{{}_{e}{}^{a}{}_{f}{}^{b}{}_{d}{}^{c}\}.$$
(65)
The first three ClebschโGordan coefficients of the vertex labelled โBโ can be summed similarly by again applying (65), while the remaining ClebschโGordan contributions sum to Kronecker delta-functions according to the orthogonality relations
$`{\displaystyle \underset{\alpha ,\beta }{}}[{}_{\alpha }{}^{a}{}_{\beta }{}^{b}{}_{\gamma }{}^{c}][{}_{\alpha }{}^{a}{}_{\beta }{}^{b}{}_{\gamma ^{}}{}^{c^{}}]`$ $`=`$ $`\delta _{cc^{}}\delta _{\gamma \gamma ^{}},`$
$`{\displaystyle \underset{\alpha ,\gamma }{}}[{}_{\alpha }{}^{a}{}_{\beta }{}^{b}{}_{\gamma }{}^{c}][{}_{\alpha }{}^{a}{}_{\beta ^{}}{}^{b^{}}{}_{\gamma }{}^{c}]`$ $`=`$ $`\frac{2c+1}{2b+1}\delta _{bb^{}}\delta _{\beta \beta ^{}}.`$ (66)
The vertex C has the same structure as vertex A.
The final vertex D has a completely different structure. By means of the reflection identity
$$[{}_{\alpha }{}^{a}{}_{\beta }{}^{b}{}_{\gamma }{}^{c}]=(1)^{a\alpha }\sqrt{\frac{2c+1}{2b+1}}[{}_{\gamma }{}^{c}{}_{\alpha }{}^{a}{}_{\beta }{}^{b}]$$
(67)
its contribution can be rewritten as
$`(1)^{j_7j_2+j_8j_3}\frac{\sqrt{(2j_7+1)(2j_8+1)}}{2j+1}`$
$`\times [{}_{q}{}^{j}{}_{b_7}{}^{j_7}{}_{b_2}{}^{j_2}][{}_{a_8}{}^{j_8}{}_{d_1}{}^{j_1}{}_{l}{}^{j}][{}_{b_2}{}^{j_2}{}_{l}{}^{j}{}_{b_4}{}^{j_4}][{}_{q}{}^{j}{}_{a_8}{}^{j_8}{}_{c_3}{}^{j_3}][{}_{b_7}{}^{j_7}{}_{d_1}{}^{j_1}{}_{h}{}^{j}][{}_{c_3}{}^{j_3}{}_{h}{}^{j}{}_{b_4}{}^{j_4}].`$ (68)
We can then apply the identity
$`{\displaystyle \underset{\alpha ,\beta ,\mathrm{},\nu }{}}[{}_{\alpha }{}^{a}{}_{\beta }{}^{b}{}_{\gamma }{}^{c}][{}_{\delta }{}^{d}{}_{ฯต}{}^{e}{}_{\varphi }{}^{f}][{}_{\gamma }{}^{c}{}_{\varphi }{}^{f}{}_{\nu }{}^{q}][{}_{\alpha }{}^{a}{}_{\delta }{}^{d}{}_{\eta }{}^{g}][{}_{\beta }{}^{b}{}_{ฯต}{}^{e}{}_{\mu }{}^{h}][{}_{\eta }{}^{g}{}_{\mu }{}^{h}{}_{\nu }{}^{q}]`$
$`=\sqrt{(2c+1)(2f+1)(2g+1)(2h+1)}(2q+1)\left\{\begin{array}{ccc}a& d& g\\ b& e& h\\ c& f& q\end{array}\right\}`$ (69)
to obtain a final expression for the contribution from vertex D in terms of $`9j`$-symbols of the second kind.
By grouping together all of these contributions, the deformed square Wilson loop correlator (63) thus becomes
$`W_{\mathrm{};j}^k(\rho _1)`$ $`=`$ $`{\displaystyle \underset{j_1,\mathrm{},j_8\frac{1}{2}_0}{}}(1)^{2(j_1+j_2+j_3)k}\delta _{j_5j_6}\delta _{j_6j_5}\delta _{a_5b_6}\delta _{b_6a_5}\mathrm{e}^{\underset{\lambda =1}{\overset{8}{}}\frac{g^2\rho _\lambda ^{}}{2}j_\lambda (j_\lambda +1)}`$ (70)
$`\times (2j_1+1)(2j_4+1)(2j_7+1)(2j_8+1)\sqrt{\frac{2j_5+1}{2j_6+1}}`$
$`\times [{}_{b_8}{}^{j_8}{}_{c}{}^{j}{}_{a_2}{}^{j_2}][{}_{b_8}{}^{j_8}{}_{c}{}^{j}{}_{a_2}{}^{j_2}][{}_{a_7}{}^{j_7}{}_{e}{}^{j}{}_{d_3}{}^{j_3}][{}_{a_7}{}^{j_7}{}_{e}{}^{j}{}_{d_3}{}^{j_3}]\{{}_{j_3}{}^{\stackrel{j}{j_8}}{}_{j}{}^{\stackrel{j_7}{j_1}}{}_{j_4}{}^{\stackrel{j_2}{j}}\}`$
$`\times \{{}_{j_2}{}^{j_1}{}_{j}{}^{j}{}_{j_5}{}^{j_8}\}\{{}_{j_2}{}^{j_3}{}_{j}{}^{j}{}_{j_4}{}^{j_8}\}\{{}_{j_3}{}^{j_2}{}_{j}{}^{j}{}_{j_4}{}^{j_5}\}\{{}_{j_3}{}^{j_2}{}_{j}{}^{j}{}_{j_4}{}^{j_7}\}\{{}_{j_3}{}^{j_1}{}_{j}{}^{j}{}_{j_6}{}^{j_7}\}.`$
We can rewrite this expression in a manner which resembles more closely the circular Wilson loop correlator (56) by expressing the $`9j`$-symbol in terms of $`6j`$-symbols, at the price of having to introduce an additional angular momentum sum. This is accomplished via the identity
$$\{{}_{j_3}{}^{\stackrel{j}{j_8}}{}_{j}{}^{\stackrel{j_7}{j_1}}{}_{j_4}{}^{\stackrel{j_2}{j}}\}=\underset{j_9\frac{1}{2}_0}{}(1)^{3j+j_1+j_2+j_3+j_4+j_7+j_8+2j_9}(2j_9+1)\{{}_{j}{}^{j}{}_{j_7}{}^{j_3}{}_{j_8}{}^{j_9}\}\{{}_{j_4}{}^{j_7}{}_{j_2}{}^{j}{}_{j_1}{}^{j_9}\}\{{}_{j}{}^{j_2}{}_{j_3}{}^{j_4}{}_{j}{}^{j_9}\}.$$
(71)
Doing the implicit sums left over in (70) then gives the final form
$`W_{\mathrm{};j}^k(\rho _1)`$ $`=`$ $`{\displaystyle \underset{2j_4=0}{\overset{4j1}{}}}{\displaystyle \underset{j_1,\mathrm{},j_9D_{\mathrm{}}^j}{}}(1)^{(j_1+j_2+j_3)(2k+1)+3j+j_4+j_7+j_8+2j_9}{\displaystyle \underset{\genfrac{}{}{0pt}{}{\alpha =1}{\alpha 6}}{\overset{9}{}}}(2j_\alpha +1)`$ (72)
$`\times \mathrm{e}^{\frac{g^2(\rho _5^{}+\rho _6^{})}{2}j_5(j_5+1)}{\displaystyle \underset{\genfrac{}{}{0pt}{}{\lambda =1}{\lambda 5}}{\overset{8}{}}}\mathrm{e}^{\frac{g^2\rho _\lambda ^{}}{2}j_\lambda (j_\lambda +1)}`$
$`\times \{{}_{j_2}{}^{j_1}{}_{j}{}^{j}{}_{j_5}{}^{j_8}\}\{{}_{j_2}{}^{j_3}{}_{j}{}^{j}{}_{j_4}{}^{j_8}\}\{{}_{j_3}{}^{j_2}{}_{j}{}^{j}{}_{j_4}{}^{j_5}\}\{{}_{j_3}{}^{j_2}{}_{j}{}^{j}{}_{j_4}{}^{j_7}\}`$
$`\times \{{}_{j_3}{}^{j_1}{}_{j}{}^{j}{}_{j_6}{}^{j_7}\}\{{}_{j}{}^{j}{}_{j_7}{}^{j_3}{}_{j_8}{}^{j_9}\}\{{}_{j_4}{}^{j_7}{}_{j_2}{}^{j}{}_{j_1}{}^{j_9}\}\{{}_{j}{}^{j_2}{}_{j_3}{}^{j_4}{}_{j}{}^{j_9}\}`$
where by the triangle inequalities the sum over spins is restricted to the range
$$D_{\mathrm{}}^j=D_{j_2j_7}^{j_1}D_{j_4j_9}^{j_2}D_{j_1j_6}^jD_{j_2j_3}^j\underset{\genfrac{}{}{0pt}{}{\alpha =4,5,7,8}{\beta =1,2,3}}{}D_{j_\alpha j_\beta }^j\underset{\alpha =3,8}{}D_{j_9j_\alpha }^j.$$
(73)
We can now compare (72) with the circular Wilson loop correlator (56) that encloses the same area $`\rho _1`$ as the square on the original torus. Again, while bearing some similarities, the two formulas have a very different angular momentum structure and a different functional dependence on the areas involved. It is thus very likely that they are different. Of course this is not a rigorous proof that the two expressions obtained are really not equal, and to accomplish this one should perform the sum over all angular momenta. Unfortunately it is very difficult to handle these sums analytically.
These calculations can be straightforwardly generalized to more complicated polygonal contours on $`๐_\theta ^2`$. The differences will lie in the nature of the corresponding triangulation of the dual torus. The generic contribution from a local vertex will involve a $`3nj`$-symbol of the second kind, which can be represented as a sum over products of $`n`$ $`6j`$-symbols . The higher the valencies of these vertices the more angular momentum sums that are introduced, yielding apparently distinct expressions for the corresponding loop correlators. This is evident even in the additional area dependences that the self-intersecting contours contain. While in principle the dual areas $`\rho _\lambda ^{}`$ depend on the original area $`\rho _1`$ and the rank $`N`$ of the Morita dual commutative gauge theory, an infinitesimal total area-preserving variation of the parameters $`\rho _\lambda ^{}`$ generally produces a non-vanishing result and accounts for the distinct correlation functions obtained.
The claimed shape dependence of Wilson loops on the noncommutative torus is much more drastic than on the noncommutative plane . For example, it is clear that a circular contour and an ellipsoidal contour can produce distinct loop correlators, even though the two loops can be mapped into one another by a unimodular linear transformation. This can be understood from the fact that the global $`U(\mathrm{})`$ group of area-preserving diffeomorphisms on $`๐^2`$ is different from that on $`^2`$ . Because of the smaller invariance group on $`๐^2`$, rotational symmetry is lost. Thus the loop correlators depend crucially on the orientation in the torus and other geometrical factors in addition to the shape of the contour. A similar feature has been observed numerically in the lattice regularization of the noncommutative gauge theory . Within the present combinatorial approach, the shape dependence of closed Wilson line correlators is understood through an intricate graph theoretic problem. Note that, conversely, an intricate self-intersecting Wilson loop described by a graph in commutative non-abelian gauge theory can be mapped to a simple Wilson loop in $`U(1)`$ noncommutative gauge theory. The self-intersections can be thought of as being absorbed into the noncommutativity of spacetime, in much the same way that the rank $`N`$ can.
## 4 Dual Loop Correlators: Irrational Case
Let us now examine the general form of Morita equivalent loop correlators that arises when the noncommutativity parameter $`\mathrm{\Theta }`$ is an irrational number. In this case, the target theory is necessarily another noncommutative gauge theory. This dual gauge theory is once again defined on a torus whose size $`\stackrel{~}{r}`$ depends on the noncommutativity parameter as prescribed in (15). This means that as we go from our original noncommutative gauge theory to its Morita dual, the size of the torus may change drastically. More precisely, we recall that to every Morita equivalence parameterized by $`SL(2,)`$ integers $`a,b,c,d`$, there exists a critical radius $`r_\mathrm{c}^{a,b}`$ given by (26) which is associated to each path such that if the radius of the target torus $`\stackrel{~}{r}`$ is smaller than $`r_\mathrm{c}^{a,b}`$, then the path can self-intersect in the dual gauge theory. Whether or not self-intersections actually occur depends on the shape of the path itself, as well as on its width, length and orientation. The key point is that, if the noncommutativity parameter is irrational-valued, then the critical radius $`r_\mathrm{c}^{a,b}`$ can be made vanishingly small. This is a consequence of the well-known number theoretic property that, given $`\mathrm{\Theta }\backslash `$, the subset $`+\mathrm{\Theta }=\{a+b\mathrm{\Theta }|a,b\}`$ is dense on the real line $``$. In particular, given any $`\epsilon >0`$, we can always find $`a,b`$ such that $`|a+b\mathrm{\Theta }|<\epsilon `$ and hence $`r_\mathrm{c}^{a,b}<\epsilon r`$. In other words, since the area of a closed Wilson line does not change under a Morita transformation, any Wilson loop in irrational noncommutative YangโMills theory is dual to a Wilson loop with arbitrarily many self-intersections and windings around the torus. This gives a combinatorial picture of irrational Wilson loops as densely wound and interesecting loops on arbitrarily small tori.
On more heuristic grounds, we can rephrase our argument as follows. Let us take $`\mathrm{\Theta }\backslash `$ and approximate it by a sequence of rational numbers as
$$\mathrm{\Theta }=\underset{n\mathrm{}}{lim}\frac{c_n}{N_n},$$
(74)
where both sequences of integers $`c_n`$ and $`N_n`$ tend to infinity such that their ratio is held fixed in the limit. For every fixed $`n`$, we can choose a Morita transformation such that $`r_\mathrm{c}^{a_n,b_n}=r/N_n`$. In the limit $`n\mathrm{}`$, one has $`r_\mathrm{c}^{a_n,b_n}0`$. With $`\mathrm{}^\mu (๐)`$ the characteristic lengths associated with the path $`๐`$ which we introduced in (25), it follows that it is always possible to find a bound of the form (27). In other words, whenever $`\mathrm{\Theta }`$ is an irrational number, there is a target torus on which the given path self-intersects and winds an (uncountably) infinite number of times. Recalling the analysis of the previous section, we see that the apparent violation of invariance under area-preserving diffeomorphisms is in fact due to the self-intersecting nature of dual Wilson loops. Differently shaped loops can have drastically different self-intersection and winding images under the same Morita transformation.
This self-intersecting property presents a serious technical obstruction to obtaining exact nonperturbative expressions for correlation functions of irrational noncommutative Wilson loops. In particular, the loop functional is not a smooth function of $`\theta `$, and the geometrical path parameters display a drastic change under Morita equivalence. It is thus not clear what a Morita duality-invariant expression for closed Wilson line correlators should look like. However, there is a natural and obvious regime in which exact results can be obtained. If one considers a certain double scaling limit in which the area enclosed by the Wilson loop vanishes faster than the area of the target torus, then the dual Wilson loop will be of the same (non-intersecting) type. This particular limit is the topic of the next section.
## 5 Loop Correlators in the Double Scaling Limit
In this section we will compute a particular class of noncommutative loop correlators that can be consistently obtained through the use of Morita equivalence. Consider a loop winding $`n`$ times around itself and encircling an area $`\rho _1`$. As in Section 3.2, the area outside the loop is denoted $`\rho _2`$ so that the total area of the torus is $`(2\pi r^{})^2=\rho _1+\rho _2`$. Having in mind the picture of noncommutative Wilson loops drawn out in the previous section, we will take the limit $`n\mathrm{}`$ with the product $`n^2\rho _1=\lambda `$ held fixed. Because the loop area can be taken arbitrarily small in the $`\stackrel{~}{r}0`$ limit required to induce gauge theory on the noncommutative plane , the final result should be consistent with the known expression obtained by resumming the small loop area perturbation series on $`^2`$ .
Our starting point is the general expression (35) for the Wilson loop correlator in the $`k^{\mathrm{th}}`$ topological sector of the dual $`SU(N)`$ gauge theory on $`๐^2`$. For the representation $`R`$ we take $`R=N^n`$ which, since $`\chi _{N^n}^{}(U)=\chi _N^{}(U^n)`$, describes a Wilson loop in the fundamental representation with $`n`$ windings. We will compute the corresponding normalized correlation function
$$๐ฒ_n^k(\rho _1)=\frac{W_{๐_1;N^n}^k(\rho _1)}{NZ_k}$$
(75)
where
$$Z_k=\underset{R}{}\mathrm{e}^{\frac{g^2(\rho _1+\rho _2)}{2}C_2(R)}\chi _R^{}\left(\mathrm{e}^{2\pi \mathrm{i}k/N}\right)$$
(76)
is the partition function of YangโMills theory on the torus in the $`k^{\mathrm{th}}`$ โt Hooft sector. As in (39), the required ClebschโGordan coefficients can be computed from the explicit expression
$`\mathrm{N}_{R_1N^n}^{R_2}`$ $`=`$ $`{\displaystyle \underset{a=1}{\overset{N}{}}}{\displaystyle _0^1}d\lambda _a\delta \left(\underset{a=1}{\overset{N}{}}\lambda _a\right){\displaystyle \underset{c=1}{\overset{N}{}}}\mathrm{e}^{2\pi \mathrm{i}n\lambda _c}`$ (77)
$`\times \underset{1a,bN}{det}\left[\mathrm{e}^{2\pi \mathrm{i}n_a^{R_1}\lambda _b}\right]\underset{1a,bN}{det}\left[\mathrm{e}^{2\pi \mathrm{i}n_a^{R_2}\lambda _b}\right].`$
It is convenient to introduce integers $`l_N^{R_i}`$, $`i=1,2`$ through the identities
$$1=\frac{1}{\sqrt{\pi }}_0^1d\alpha _i\underset{l_N^{R_i}=\mathrm{}}{\overset{\mathrm{}}{}}\mathrm{e}^{(2\pi )^2\left(\alpha _i\frac{1}{N}\underset{a=1}{\overset{N1}{}}n_a^{R_i}l_N^{R_i}\right)^2},$$
(78)
and to change summation variables from Young tableau boxes to integers $`l_a^{R_i}`$, $`a=1,\mathrm{},N1`$ and $`l^{R_i}`$ defined by
$$l_a^{R_i}=n_a^{R_i}+l_N^{R_i}a+N,l^{R_i}=\underset{a=1}{\overset{N}{}}l_a^{R_i}.$$
(79)
In terms of these new integers, the quadratic Casimir invariant and dimension of the representation $`R_i`$ are given by
$$C_2\left(R_i\right)=C_2\left(๐^{R_i}\right)=\underset{a=1}{\overset{N}{}}\left(l_a^{R_i}\frac{l^{R_i}}{N}\right)^2\frac{N}{12}\left(N^21\right),dimR_i=\mathrm{\Delta }\left(๐^{R_i}\right).$$
(80)
By exploiting the complete symmetry of the correlator in the summation integers we thereby arrive at the expression
$`๐ฒ_n^k(\rho _1)`$ $`=`$ $`{\displaystyle \frac{1}{NZ_k}}{\displaystyle \frac{1}{(2\pi )^N\pi (N!)^2}}{\displaystyle \underset{a=1}{\overset{N}{}}}{\displaystyle _0^1}d\lambda _a\delta \left(\underset{a=1}{\overset{N}{}}\lambda _a\right){\displaystyle \underset{c=1}{\overset{N}{}}}\mathrm{e}^{2\pi \mathrm{i}n\lambda _c}`$ (81)
$`\times {\displaystyle \underset{๐^{R_1},๐^{R_2}}{}}{\displaystyle \frac{\mathrm{\Delta }\left(๐^{R_1}\right)}{\mathrm{\Delta }\left(๐^{R_2}\right)}}\mathrm{e}^{\frac{g^2\rho _1}{2}C_2(๐^{R_1})\frac{g^2\rho _2}{2}C_2(๐^{R_2})}\mathrm{e}^{2\pi \mathrm{i}kl^{R_1}/N}{\displaystyle \underset{a=1}{\overset{N}{}}}\mathrm{e}^{2\pi \mathrm{i}\left(l_a^{R_2}l_a^{R_1}\right)\lambda _a}`$
$`\times {\displaystyle _0^1}\mathrm{d}\alpha _1\mathrm{e}^{(2\pi )^2\left(\alpha _1\frac{l^{R_1}}{N}\right)^2}{\displaystyle _0^1}\mathrm{d}\alpha _2\mathrm{e}^{(2\pi )^2\left(\alpha _2\frac{l^{R_2}}{N}\right)^2}.`$
Expressing the delta-function as a Fourier series and integrating over $`\alpha _1`$, we convert this expression into the form
$`๐ฒ_n^k(\rho _1)`$ $`=`$ $`{\displaystyle \frac{1}{NZ_k}}{\displaystyle \frac{1}{(2\pi )^N\sqrt{\pi }(N!)^2}}{\displaystyle \underset{a=1}{\overset{N}{}}}{\displaystyle _0^1}d\lambda _a{\displaystyle \underset{c=1}{\overset{N}{}}}\mathrm{e}^{2\pi \mathrm{i}n\lambda _c}`$ (82)
$`\times {\displaystyle \underset{๐^{R_1},๐^{R_2}}{}}{\displaystyle \frac{\mathrm{\Delta }\left(๐^{R_1}\right)}{\mathrm{\Delta }\left(๐^{R_2}\right)}}\mathrm{e}^{\frac{g^2\rho _1}{2}C_2(๐^{R_1})\frac{g^2\rho _2}{2}C_2(๐^{R_2})}\mathrm{e}^{2\pi \mathrm{i}kl^{R_1}/N}{\displaystyle \underset{a=1}{\overset{N}{}}}\mathrm{e}^{2\pi \mathrm{i}\left(l_a^{R_2}l_a^{R_1}\right)\lambda _a}`$
$`\times {\displaystyle _0^1}\mathrm{d}\alpha \mathrm{e}^{(2\pi )^2\left(\alpha \frac{l^{R_2}}{N}\right)^2}.`$
We use the complete symmetry again to now fix the summation index $`c=1`$, which produces an additional factor of $`N`$. The $`\lambda _a`$ integrals can now be performed explicitly giving the constraints $`l_a^{R_1}=l_a^{R_2}`$, $`a1`$ and $`l_1^{R_1}=l_1^{R_2}+n`$. This leads to our final explicit result
$`๐ฒ_n^k(\rho _1)`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{\pi }N!Z_k}}\mathrm{e}^{\frac{g^2n^2\rho _1}{2}\left(1\frac{1}{N}\right)}{\displaystyle \underset{๐}{}}{\displaystyle \frac{\mathrm{\Delta }(n_1+n,n_2,\mathrm{},n_N)}{\mathrm{\Delta }(๐)}}\mathrm{e}^{\frac{g^2(\rho _1+\rho _2)}{2}C_2(๐)}`$ (83)
$`\times \mathrm{e}^{\frac{g^2\rho _1}{2}\left(2nn_1\frac{2n}{N}\underset{a=1}{\overset{N}{}}n_a\right)}\mathrm{e}^{\frac{2\pi \mathrm{i}k}{N}\underset{a=1}{\overset{N}{}}n_a}{\displaystyle _0^1}d\alpha \mathrm{e}^{(2\pi )^2\left(\alpha \frac{1}{N}\underset{a=1}{\overset{N}{}}n_a\right)^2}.`$
We can check the normalization here by observing that these same steps can be used to write the partition function (76) as
$$Z_k=\frac{1}{\sqrt{\pi }N!}\underset{๐}{}\mathrm{e}^{\frac{g^2(\rho _1+\rho _2)}{2}C_2(๐)}\mathrm{e}^{\frac{2\pi \mathrm{i}k}{N}\underset{a=1}{\overset{N}{}}n_a}_0^1d\alpha \mathrm{e}^{(2\pi )^2\left(\alpha \frac{1}{N}\underset{a=1}{\overset{N}{}}n_a\right)^2}.$$
(84)
Thus our conventions imply the normalization condition $`๐ฒ_0^k(\rho _1)=1`$.
Let us now take the $`n\mathrm{}`$ limit. For this, we insert the explicit expression for the Vandermonde determinant (37) to recast (83) as
$`๐ฒ_n^k(\rho _1)`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{\pi }N!Z_k}}{\displaystyle \underset{๐}{}}\mathrm{e}^{\frac{g^2(\rho _1+\rho _2)}{2}C_2(๐)}\mathrm{e}^{\frac{2\pi \mathrm{i}k}{N}\underset{a=1}{\overset{N}{}}n_a}{\displaystyle _0^1}d\alpha \mathrm{e}^{(2\pi )^2\left(\alpha \frac{1}{N}\underset{a=1}{\overset{N}{}}n_a\right)^2}`$ (85)
$`\times \mathrm{e}^{\frac{g^2\rho _1}{2}\left(n^2(1\frac{1}{N})+2nn_1\frac{2n}{N}\underset{a=1}{\overset{N}{}}n_a\right)}{\displaystyle \underset{m=0}{\overset{N1}{}}}{\displaystyle \frac{(N1)!}{(N1m)!}}{\displaystyle \underset{j=2}{\overset{m+1}{}}}{\displaystyle \frac{n^m}{n_1n_j}}.`$
We thus obtain a Laurent series in $`\frac{1}{n}`$ from expanding the exponential term to get
$`๐ฒ_n^k(\rho _1)`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{\pi }N!Z_k}}{\displaystyle \underset{๐}{}}\mathrm{e}^{\frac{g^2(\rho _1+\rho _2)}{2}C_2(๐)}\mathrm{e}^{\frac{2\pi \mathrm{i}k}{N}\underset{a=1}{\overset{N}{}}n_a}{\displaystyle _0^1}d\alpha \mathrm{e}^{(2\pi )^2\left(\alpha \frac{1}{N}\underset{a=1}{\overset{N}{}}n_a\right)^2}`$ (86)
$`\times {\displaystyle \underset{l=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{m=0}{\overset{N1}{}}}{\displaystyle \underset{p=0}{\overset{l}{}}}{\displaystyle \frac{(N1)!}{(N1m)!}}n^{ml}{\displaystyle \frac{\left(g^2\lambda \right)^l}{p!(lp)!}}`$
$`\times \left({\displaystyle \frac{1}{N}}{\displaystyle \underset{a=1}{\overset{N}{}}}n_a^2\right)^{lp}{\displaystyle \underset{j=2}{\overset{m+1}{}}}{\displaystyle \frac{n_1^p}{n_1n_j}}\mathrm{e}^{\frac{g^2\lambda }{2}\left(1\frac{1}{N}\right)}`$
where $`\lambda =n^2\rho _1`$. The very same structure appears in the computation of $`n`$-winding Wilson loops on the sphere and it is clear that this result generalizes to arbitrary genus Riemann surfaces. Let us thus proceed as in .
First of all, we observe that (86) is actually an expansion in $`\frac{1}{n^2}`$. This point can be understood by changing $`๐๐`$, which produces an overall factor $`(1)^{ml}`$ weighting the sum over $`๐`$. This implies that $`ml`$ must be an even integer in order to contribute a non-vanishing result. It is useful to now rewrite the sum over $`๐`$ for fixed $`l,m,p`$ as
$$\frac{1}{(m+1)!}\underset{๐^N}{}\mathrm{e}^{\frac{g^2(\rho _1+\rho _2)}{2}\underset{a=1}{\overset{N}{}}n_a^2}\left(\frac{1}{N}\underset{a=1}{\overset{N}{}}n_a^2\right)^{lp}\underset{\pi S_{m+1}}{}\underset{j=2}{\overset{m+1}{}}\frac{n_{\pi (1)}^p}{n_{\pi (1)}n_{\pi (j)}}.$$
(87)
Let us evaluate the zeroth-order contribution to (86). For $`l=p=m`$ one can write the sum over permutations as
$$\underset{\pi S_{m+1}}{}\underset{j=2}{\overset{m+1}{}}\frac{n_{\pi (1)}^m}{n_{\pi (1)}n_{\pi (j)}}=\frac{1}{\mathrm{\Delta }(n_1,\mathrm{},n_{m+1})}\underset{\pi S_{m+1}}{}f_\pi (n_1,\mathrm{},n_{m+1}),$$
(88)
where the Vandermonde determinant arises from the common denominator and the quantity $`_{\pi S_{m+1}}f_\pi (n_1,\mathrm{},n_{m+1})`$ is a polynomial of degree $`\frac{1}{2}m(m+1)`$ in $`m+1`$ variables. Since the Vandermonde determinant $`\mathrm{\Delta }(n_1,\mathrm{},n_{m+1})`$ is completely antisymmetric in its arguments, the non-vanishing contribution to (87) comes from the completely antisymmetric part of $`_{\pi S_{m+1}}f_\pi (n_1,\mathrm{},n_{m+1})`$ implying that
$$\underset{\pi S_{m+1}}{}f_\pi (n_1,\mathrm{},n_{m+1})=C\mathrm{\Delta }(n_1,\mathrm{},n_{m+1}).$$
(89)
The proportionality constant is easily found by inspection to be $`C=1`$. It is not difficult to prove that the potentially divergent contributions in the limit $`n\mathrm{}`$, coming from the terms with $`l<m`$ in (86), vanish. For this, we again use (87) to notice that we can still factorize a Vandermonde determinant in the denominator as in (88), but now $`_{\pi S_{m+1}}f_\pi (n_1,\mathrm{},n_{m+1})`$ is a polynomial of degree less than $`\frac{1}{2}m(m+1)`$ in $`m+1`$ variables because $`p<m`$. Its completely antisymmetric part thus vanishes.
In this way we arrive finally at
$$๐ฒ_{\mathrm{}}\left(g^2\lambda \right)=\underset{n\mathrm{}}{lim}๐ฒ_n^k\left(\frac{\lambda }{n^2}\right)=\mathrm{e}^{\frac{g^2\lambda }{2}\left(1\frac{1}{N}\right)}\underset{m=0}{\overset{N1}{}}\frac{(N1)!}{(N1m)!}\frac{\left(g^2\lambda \right)^m}{m!(m+1)!}.$$
(90)
Corrections to this formula are of order $`\frac{1}{n^2}`$. Note that the partition function cancels in this limit. The Wilson loop average can be expressed in terms of a generalized Laguerre polynomial as
$$๐ฒ_{\mathrm{}}\left(g^2\lambda \right)=\frac{1}{N}\mathrm{e}^{\frac{g^2\lambda }{2}\left(1\frac{1}{N}\right)}L_{N1}^1\left(g^2\lambda \right).$$
(91)
Rather remarkably, this result coincides with the analogous result for YangโMills theory on the sphere. In fact, it is completely independent of the genus of the original Riemann surface. Differences would appear only at sub-leading order in $`\frac{1}{n^2}`$.
At this point we can take the large $`N`$ limit to reach gauge theory on the noncommutative plane. The noncommutative YangโMills coupling constant in this case is defined through $`g^2=\mathrm{\Theta }\widehat{g}^2=\widehat{g}^2/N`$, and the Wilson loop correlator can be expressed in terms of a Bessel function as
$$\widehat{๐ฒ}_{\mathrm{}}\left(\widehat{g}^2\lambda \right)=\underset{N\mathrm{}}{lim}๐ฒ_{\mathrm{}}\left(\frac{\widehat{g}^2\lambda }{N}\right)=\frac{J_1\left(2\sqrt{\widehat{g}^2\lambda }\right)}{\sqrt{\widehat{g}^2\lambda }}.$$
(92)
This expression coincides exactly with the result obtained, at this order, by resumming the perturbation series . The coincidence of the correlator of the noncommutative Wilson loop in the present limit with that of the commutative Wilson loop obtained by resumming planar diagrams confirms the general expectation that noncommutativity modifies only large area Wilson loops. For small loops the usual commutative behaviour at large $`N`$ is recovered, while large area loops become complex-valued . The double scaling limit we have considered in this section effectively singles out small area loops. The fact that noncommutativity only modifies the long wavelength behaviour of Wilson loops is indicative of some nonperturbative form of UV/IR mixing. This mixing only affects the closed Wilson line observables of the noncommutative gauge theory and is another manifestation of the loss of invariance under area-preserving diffeomorphisms. Note that the present double scaling limit โzooms inโ on only a very small portion of the torus, so that the final correlator in the limit is completely independent of any global properties of the two-dimensional spacetime. This small area limit is equivalent to the limit $`\theta =\mathrm{}`$, as one might have naively expected, and therefore eliminates all higher order traces of the $`\frac{1}{\theta }`$-expansion.
## 6 Summary and Discussion
In this paper we have explored new aspects of the shape dependence of Wilson loop correlators on a two-dimensional noncommutative torus. Because of the non-trivial topology and the compactness of the spacetime, correlation functions associated to loops $`๐`$ of the same area apparently depend not only on their shape but also on their characteristic lengths $`\mathrm{}^\mu (๐)`$ defined in (25) (i.e. their heights and widths) and on their orientation in the torus. We illustrated this dependence through several explicit calculations using Morita equivalence along with a combinatorial approach. From our perspective the observed breaking of invariance of Wilson loop correlators under area-preserving diffeomorphisms of the two-dimensional spacetime may be attributed to the wrapping and self-intersecting nature of Morita dual Wilson loops. Only those contours whose images under Morita equivalence lead to isomorphic non-planar graphs will give rise to identical correlation functions. In irrational noncommutative gauge theory, there always exist dual loops which wind infinitely many times around the torus. Motivated by this picture, we have also explicitly computed an infinitely wound Wilson loop correlator. Since the limit of infinite winding considered corresponds to small loop area, our results agree with those obtained by resumming planar diagrams in commutative gauge theory, the planarity arising essentially as a combined large $`N`$ and large $`\theta `$ effect. This limit also eliminates the non-perturbative topological degrees of freedom which are expected to restore the usual large $`N`$ GrossโWitten area law behaviour for small noncommutative Wilson loops.
It is interesting to examine in the present context the perturbative anomaly that comes from the contribution of the non-planar diagram of order $`\widehat{g}^4`$ to the average of the noncommutative Wilson loop on $`^2`$ . The leading term in the $`\frac{1}{\theta }`$ expansion of the correlator is proportional to $`\widehat{g}^4\rho _1^2`$ in our notation, and thus it survives the limit $`\theta \mathrm{}`$ due to the singular infrared behaviour of the gauge propagator in two dimensions. This term appears to be in conflict with both general arguments of noncommutative perturbation theory and with the representation of noncommutative gauge theories via large $`N`$ twisted reduced models . However, the anomalous term vanishes in the double scaling limit considered in Section 5 and so does not show up in our calculations which capture the entire small loop area perturbation series.
There may be yet another way to eliminate this anomalous behaviour. We offer the following argument only as a somewhat speculative conjecture at this stage. The perturbative calculations which unveil this anomalous term are performed in the axial gauge where the self-interactions of the gauge field disappear. As is well-known, the axial gauge is forbidden on the torus (or on any spacetime of non-trivial topology) due to the existence of topologically non-trivial field configurations (transforming under large gauge transformations) which yield non-trivial Polyakov loops along the axial direction. In commutative gauge theory on $`^2`$ only topologically trivial gauge fields exist (transforming under gauge transformations connected to the identity) and there is no problem with the axial gauge choice. However, this is not the case for gauge theory on the noncommutative plane. In contrast to the commutative situation the gauge theory now contains topologically non-trivial backgrounds called fluxons owing to the fact that $`_\theta ^2`$ has, like the torus, a non-trivial K-theory group. The fluxons can be regarded as the surviving degrees of freedom left over from the usual instanton configurations in the limit where one decompactifies the noncommutative torus onto the noncommutative plane. An $`L`$-fluxon solution is labelled by a set of moduli $`๐_1,\mathrm{},๐_L^2`$, which specify the locations of the vortices on the plane, and by a collection of magnetic charges $`m_1,\mathrm{},m_L`$. One can compute the semi-classical average of an open Wilson line operator along a straight infinite contour pointing in a direction $`\widehat{๐}`$ of $`^2`$ in the fluxon background with the result $`W_{\mathrm{open}}(\widehat{๐})=_{a=1}^L\mathrm{e}^{\mathrm{i}\widehat{๐}๐_a}`$ (independently of the vortex charges). Generically, this expectation value cannot be trivialized by any noncommutative gauge transformation and the correlator thus presents an obstruction to choosing the axial gauge. Axial gauge choices are also forbidden in the lattice regularization of noncommutative YangโMills theory due to UV/IR mixing . This fact suggests that the observed anomalous behaviour of noncommutative Wilson loops could be due to the choice of a wrong vacuum, and that the correct perturbative calculation should instead expand about the background of a fluxon. It would be interesting to investigate this point further.
###### Acknowledgments.
We thank A. Bassetto, Y. Makeenko and F. Vian for helpful discussions. The work of M.C. and R.J.S. was supported in part by PPARC Grant PPA/G/S/2002/00478 and by the EU-RTN Network Grant MRTN-CT-2004-005104. The work of R.J.S. was supported in part by a PPARC Advanced Fellowship.
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# Giant resonances in exotic spherical nuclei within the RPA approach with the Gogny force
## 1 Introduction
Giant multipole resonances (GR) are collective excitations of nuclei that lie at excitation energies above the nucleon separation energy (8-10 MeV), have different multipolarities and carry different spin-isospin quantum numbers. They have been observed for stable nuclei throughout the mass table with large cross sections, close to the maximum allowed by sum rule arguments, implying that a large number of nucleons participate in a very collective nuclear motion . It is a challenge both to experimentalists and theorists to study the properties of these states for nuclei far from the valley of stability. Not too much has been done from the experimental side yet: let us just mention the two measurements of the electric dipole GR (GDR) made in neutron-rich oxygen isotopes . Beside GR, there are also low-lying collective excitations, in particular quadrupole and octupole states, which reflect much more than the GR the detail of shell structure. More experimental data are available for such states in the case of unstable nuclei, giving us information on the modifications of the shell structure far from stability.
From the theoretical side, more and more calculations of GR and low-lying states are performed nowadays in the framework of microscopic HF+RPA or HFB+QRPA approaches. The effective nucleon-nucleon interactions used are taken as non-relativistic effective two-body potentials or relativistic Lagrangians for meson exchange . Such microscopic approaches, although less accurate than more phenomenological ones, usually describe reasonably well the properties of these states in stable nuclei.
Among the effective forces used in the non-relativistic approaches, the Gogny force is one of those which has been extensively employed for the description of GR and low-lying states in doubly closed shell nuclei with the RPA method . Recently, this force has been used for the first time in full Quasi-Particle RPA (QRPA) calculations. Chains of isotopes in the oxygen, nickel and tin regions have been studied in order to derive the properties of low-lying states .
The purpose of this paper is to present the results of calculations performed in three spherical exotic nuclei: <sup>78</sup>Ni, <sup>100</sup>Sn and <sup>132</sup>Sn, and to compare them with those obtained in stable nuclei. More precisely, GR and low-lying states in these nuclei will be analyzed and comparisons will be made with systematics and with analogous quantities in the well-known <sup>208</sup>Pb. The latter nucleus will serve as a reference and, for this reason, results for <sup>208</sup>Pb will be displayed along with those of the three exotic nuclei in most Tables and Figures. Let us point out that the results presented here for <sup>208</sup>Pb are new. They have been derived with the D1S parameterization of the Gogny force which is the one currently used now. They slightly differ from those of Ref. where the older parameterization D1 was employed.
A point we pay special attention to in the present work is the effect of the full consistency of the residual particle-hole (p-h) interaction with the mean field produced by the same force, as allowed by the use of consistently combined HF and RPA approaches. In order to analyze this effect, we present results where different components of the residual p-h interaction such as those generated by the spin-orbit or the Coulomb force are switched off. As will be seen, the influence of these often omitted components are far from being negligible.
In the following Section details concerning the parameters of the two-body force, the numerical methods used for solving the RPA equations are briefly recalled along with a few useful formulas. Results are presented and discussed in Section 3. The main conclusions of this work are summarized in Section 4. Let us mention that a preliminary account of the present results has appeared in the workshop Proceedings of Ref. .
## 2 The HF+RPA approach with the Gogny force
The RPA approach employed here is described in Refs. . The effective force D1S proposed by Gogny is used. This finite-range density-dependent interaction describes the mean field of the nucleus, and the residual interaction in the RPA calculations is obtained via the functional second derivative of the mean field with respect to the one-body density matrix. We want to stress that all the terms of the effective force are considered in the HF mean-field and in the residual p-h interaction, including the spinโspin component, the Coulomb force and the terms produced by the two-body spin-orbit interaction. Only the two-body terms coming from the two-body center of mass correction are not included in the RPA matrix elements. Therefore, they have been also left out from the mean field calculations. In order to get equivalent binding energies and radii, the coefficient of the spin-orbit component of D1S has been reduced from 130 MeV to 115 MeV. Such a procedure was previously employed in calculations with the D1 force, as explained in Ref. . The Gogny force D1S including this change of the spin-orbit strength will be called D1Sโ.
In the results presented here, spherical symmetry is imposed. Consequently nuclear states can be characterized by their angular momentum J and their parity $`\pi `$. The individual Hartree-Fock wave functions are expanded on finite sets of spherical harmonic oscillator (HO) wave-functions containing 15 major shells for all nuclei. For each nucleus, the value of the parameter $`\mathrm{}\omega `$ of the HO basis is taken as the one minimizing the HF total nuclear energy.
The RPA equations are solved in matrix form in the p-h representation. RPA energies do not appear very sensitive to the value adopted for the HO parameter of the basis. For instance, by changing the optimal HF value $`\mathrm{}\omega =8.7`$ MeV in <sup>208</sup>Pb by $`10\%`$, the variation of the ISGMR energy ($`13.46`$ MeV) is less than $`.5\%`$ and the energy of the first $`2^+`$ at $`4.609`$ MeV is changed by less than $`5`$ keV.
Electric transition operators are defined according to:
$$\widehat{Q}_{JM}=\frac{e}{2}\underset{i}{\overset{A}{}}\left(1\tau _z(i)\right)j_J(qr_i)Y_{JM}(\theta _i,\varphi _i),$$
(1)
where $`j_J`$ is a spherical Bessel function of order $`J`$, $`q`$ a transferred momentum, $`\tau _z`$ the third component of the nucleon isospin and $`Y_{JM}`$ the usual spherical harmonics.
The degree of collectivity of the excited states is measured from their contribution to the Energy Weighted Sum Rule (EWSR)
$$M_1(\widehat{Q}_{JM})=\underset{N}{}(E_NE_0)|N|\widehat{Q}_{JM}|0|^2$$
(2)
where $`|0`$ and $`|N`$ are the RPA correlated ground state and excited states, respectively and $`E_NE_0`$ their excitation energies. Eq.(2) can also be expressed as the average in the HF ground state $`|HF`$ of a double commutator :
$$M_1(\widehat{Q}_{JM})=\frac{1}{2}HF|[\widehat{Q}_{JM},[\widehat{H},\widehat{Q}_{JM}]]|HF.$$
(3)
Therefore, exact values of $`M_1(\widehat{Q}_{JM})`$ can be computed from expression (3) whereas smaller values will be obtained from (2), reflecting the finiteness of the particle-hole space used in the RPA calculations.
A comparison between the values calculated from (2) and (3) is shown in Figure 1 for <sup>78</sup>Ni as an example. As can be seen, with the 15 major shell basis employed, RPA calculations are able to describe with a reasonable accuracy the nuclear response for $`J^\pi =0^+`$, $`2^+`$, $`3^{}`$, $`4^+`$ and $`5^{}`$ up to transferred momenta $`q`$=1.5 fm<sup>-1</sup>.
## 3 Results
First, we will discuss the validity of the doubly-magic nature of these exotic nuclei. The single-particle neutron spectra obtained in <sup>78</sup>Ni, <sup>100</sup>Sn and <sup>132</sup>Sn are shown in Figure 2. The N=50 gap in <sup>78</sup>Ni and <sup>100</sup>Sn and the N=82 one in <sup>132</sup>Sn are of the order of 5 MeV, which is less than 20$`\%`$ smaller than the gaps obtained for stable spherical nuclei with same neutron numbers. The same is true for the proton gaps at Z=28 in <sup>78</sup>Ni and at Z=50 in tin isotopes. That is, no significant reduction of the magic gaps are observed in these nuclei. Therefore, the three exotic nuclei are still doubly magic ones and the HF+RPA method is applicable to them.
In what follows, results for states with multipolarities $`0^+`$, $`2^+`$, $`1^{}`$ and $`3^{}`$ are presented for four nuclei <sup>78</sup>Ni, <sup>100</sup>Sn <sup>132</sup>Sn and <sup>208</sup>Pb, the latter nucleus being included as a reference.
The strengths shown in the Figures are given in percentage of the EWSR calculated in the long wavelength limit $`q0`$. The relevant formulas to be used in this limit for the different values of $`J`$ are given in the appendix of Ref. .
In the present calculations the continuum spectrum of the HF Hamiltonian is approximated by a discrete one. As a consequence, the RPA strength functions appear in the form of discrete peaks. In order to make comparisons with experiments more meaningful, energy centroids will be defined in terms of the moments
$$M_k\left(\widehat{Q}_{JM}\right)=\underset{N}{}(E_NE_0)^k|N|\widehat{Q}_{JM}|0|^2.$$
(4)
of the strength function. Two of these centroids will be used in the following: the mean value of the energy $`M_1/M_0`$, and the so-called โhydrodynamicโ energy $`\sqrt{M_1/M_1}`$ for isoscalar monopole resonances.
As experimental data on GR energies is scarce in exotic nuclei, comparisons will often be made with the systematic $`A^{1/3}`$ empirical laws approximately verified in stable nuclei . Values from these systematics as well as available experimental data are given in the Tables.
### 3.1 Monopole states
Figure 3 and Table 2 display the results obtained for the Isoscalar Giant Monopole Resonance (ISGMR).
As is well known, the excitation energies of this resonance strongly depends on the compression modulus $`K_{nm}`$ calculated in infinite nuclear matter . One observes in Table 2 that the theoretical energies in <sup>208</sup>Pb, although in good agreement with the empirical $`80A^{1/3}`$ law, are 5% lower than the experimental value of Ref. . This difference is consistent with the compression modulus found in infinite nuclear matter with D1Sโ, $`K_{nm}`$=209 MeV, which is slightly outside the interval 220-235 MeV that explains the bulk of experimental data within non-relativistic approaches .
Concerning the three exotic nuclei, we note that resonance energies significantly differ from the empirical law only in <sup>78</sup>Ni. It must be noted that, of all three nuclei, <sup>78</sup>Ni is the one where the squared neutron-proton asymetry $`\left(\left(NZ\right)/A\right)^2`$ most differs from the one of the stable isotope: $`\left(\left(NZ\right)/A\right)^2\left(\left(NZ\right)/A\right)_{stable}^2`$=0.78, 0.36 and -0.23 in <sup>78</sup>Ni, <sup>132</sup>Sn and <sup>100</sup>Sn, respectively. It is therefore tempting to correlate the $``$ 1.5 MeV lowering of the ISGMR found in <sup>78</sup>Ni with this large neutron excess, the contribution of the symmetry term $`K_{sym}`$ to the finite nucleus incompressibility $`K_A`$ being negative .
The strengths displayed in Figure 3 show that the major part of the EWSR is concentrated in a single peak in all four nuclei. This feature explains why the two sets of theoretical energies listed in Table 2 are very close to each other. One notes that the fragmentation of the strength is almost zero in the $`N`$=$`Z`$ nucleus <sup>100</sup>Sn, whereas it is slightly bigger in the other three nuclei which have neutron-proton asymmetry $`(NZ)/A`$ in the range .21โ.28.
In Table 2, we show the values of the mean monopole energies $`M_1/M_0`$ obtained when different terms of the residual particle-hole (p-h) interaction are left out of the RPA calculation. Columns $`(1)`$, $`(2)`$ and $`(3)`$ refer to the mean energies calculated by leaving out the spin-orbit and the Coulomb terms, the Coulomb term and the spin-orbit term, respectively.
One observes that the spin-orbit part of the residual interaction gives a contribution to ISGMR energies ranging from 8% in <sup>78</sup>Ni to 5% in <sup>208</sup>Pb. In contrast, the Coulomb contribution is larger in Pb (3%) and almost negligible in Ni. These results are consistent with those discussed in Ref. where <sup>40</sup>Ca, <sup>90</sup>Zr and <sup>208</sup>Pb were analyzed with the SLy4 interaction. In the latter work, the inclusion in the constrained HF (CHF) of the Coulomb force and of the spin-orbit component of the Skyrme interaction was proved to be essential in order to reconcile the value of $`K_{nm}`$ obtained with the Skyrme and Gogny forces.
### 3.2 Quadrupole states
Figure 4 and Tables 5, 5 and 5 display the results obtained for isoscalar quadrupole states. Figure 4 shows that in all four nuclei the quadrupole strength is divided essentially between two states: the isoscalar Giant Quadrupole Resonance (ISGQR) exhausting $``$ 80% of the EWSR with an energy in the range 12-16 MeV and a lower-lying state at $``$ 3-5 MeV carrying $``$ 10%-15% of the quadrupole strength. We will label the latter $`2_1^+`$.
The theoretical ISGQR energies are calculated using $`M_1/M_0`$ excluding the $`2_1^+`$ state. The results shown in Table 5 are seen to be higher than the $`A^{1/3}`$ systematics by 1.0โ1.5 MeV. As the latter agrees well with the experimental value in <sup>208</sup>Pb, it is difficult to draw definite conclusions concerning the behaviour of our results in the three exotic nuclei. Let us mention that such large ISGQR energies can be understood from a too large spreading of the particle-hole spectrum in the $`2^+`$ channel at high energies. Such spreading is a consequence of the value of the effective mass of the D1Sโ interaction (m/m = 0.7) which is the one giving correct single-particle properties in mean-field calculations. As is well known, taking into account the coupling of RPA configurations to 2-particleโ2-hole (2p-2h) states would reduce this disagreement . Clearly, such a coupling should be introduced in the present calculations before reliable predictions for the ISGQR in exotic nuclei can be made . Let us mention that the same is true for the other giant resonances, with some dependence on the mode quantum numbers . Nevertheless, few results have been obtained up to now with such a coupling and it is difficult to foresee the magnitude of energy shifts, except for quadrupole and dipole states.
Our theoretical results for low-lying $`2_1^+`$ states are presented in Table 5. For these states, experimental data exist both for <sup>208</sup>Pb and <sup>132</sup>Sn . As can be seen, a fair agreement between experiment and theory is found in <sup>208</sup>Pb and an even better one in <sup>132</sup>Sn, with B(E2) values being of the same order of magnitude as experimental ones. Let us point out that QRPA calculations applied to quadrupole states have been made recently with the D1S interaction for a series of tin isotopes including <sup>132</sup>Sn . In these calculations, the spin-orbit part and the coulomb part of the residual interaction were omitted for simplicity reasons. The $`2^+`$ energies were found larger than the experimental ones by 400 keV in <sup>102</sup>Sn and 1 Mev in <sup>132</sup>Sn. The corresponding theoretical B(E2) values were lower than experimental ones by at least a factor of two.
These results are consistent with those shown in Table 5 where the same quantities as those of Table 5 are displayed. They have been calculated by leaving out from the D1Sโ p-h interaction the spin-orbit and the Coulomb terms, the Coulomb term, the spin-orbit term and no term, respectively. One observes that, as previously for monopole vibrations, taking into account the spin-orbit part of the residual interaction is essential to get results consistent with experimental data.
Going back to Table 5, $`2_1^+`$ energies are similar in <sup>100</sup>Sn and <sup>132</sup>Sn, whereas a comparatively low value is predicted in <sup>78</sup>Ni. Let us note that the $`2_1^+`$ state in <sup>78</sup>Ni is still higher than the one in <sup>56</sup>Ni, the other doubly magic Ni isotope, where the experimental value of the $`2_1^+`$state is 2.7 MeV and the RPA calculated one is 2.42 MeV with D1Sโ.
The collectivity of this $`2^+`$ state appears larger in <sup>100</sup>Sn than in <sup>132</sup>Sn and rather weak in <sup>78</sup>Ni. Figure 5 displays the transition density $`\rho _{TR}`$ of this first $`2_1^+`$ state in <sup>78</sup>Ni. The definition of the transition density is the same as the one given in appendix of Ref. . One observes that the two transition densities are in phase and that the neutron transition density is higher than the proton one and displaced to a larger radius. This mode can therefore be interpreted as an isoscalar surface mode dominated by neutron excitation.
### 3.3 Dipole states
Results for the isovector dipole resonance (IVGDR) are presented in Figure 6 and Table 9. <sup>100</sup>Sn is the nucleus where the giant dipole mode is the least fragmented with 70% of the strength concentrated into two peaks. The dipole responses of <sup>208</sup>Pb and <sup>132</sup>Sn and to a lesser extent of <sup>78</sup>Ni also appear concentrated into two main energy regions. It is expected, that the fragmentation is somewhat reduced by the coupling of the RPA modes to 2pโ2h states, producing smoother strength functions, as in Refs. where Skyrme forces were used.
In <sup>100</sup>Sn the mean value $`M_1/M_0=`$ 19.98 MeV is 3 MeV larger than the systematic 79$`A^{1/3}`$ law (17.02 MeV). The EWSR value given in Thomas-Reiche-Kuhn (TRK) unit is 1.59, which is large compared to typical experimental values . The IVGDR in <sup>132</sup>Sn is more fragmented than in <sup>100</sup>Sn. As in <sup>100</sup>Sn the mean energy value, 18.33 MeV, is much larger than systematics (79$`A^{1/3}=`$ 15.52 MeV) and the EWSR value is 1.58. In the case of <sup>78</sup>Ni, the IVGDR is quite fragmented with one major peak and smaller ones at higher energy. The mean energy value, 20.31 MeV, remains higher than systematics (79$`A^{1/3}=`$ 18.49 MeV) and the EWSR in TRK unit is 1.57.
It must be said that IVGDR excitation energies calculated with the Gogny force usually overestimate experimental data. In the case of <sup>208</sup>Pb, the calculated mean value is 16.50 MeV, which is quite large compared to experiment (13.43 MeV ), but smaller than the result of Ref. . Let us note that, ignoring the higher part of the IVGDR response by keeping only the strength around the main lower energy peak, considerably improves the agreement with systematic estimations : mean energy values become 19.28 MeV, 18.16 MeV, 16.81 MeV and 14.99 MeV in <sup>78</sup>Ni, <sup>100</sup>Sn, <sup>132</sup>Sn and <sup>208</sup>Pb, respectively.
In fact, calculated IVGDR energies and EWSR appear quite sensitive to the energy interval considered and also to the components of the effective interaction included in the p-h residual interaction. This is shown in Table 9 where mean IVGDR energies and EWSR in <sup>208</sup>Pb are listed for three energy integration intervals and for RPA calculations where Coulomb and/or spin-orbit terms are not included in the RPA matrix elements. One can see that the overestimation obtained with the Gogny force decreases by $``$ 700 keV when the Coulomb and the spin-orbit forces are ignored, which is usually done in RPA calculations employing Skyrme forces, see however Ref. . By taking all the terms of the Gogny force and considering the largest energy interval, the calculated EWSR given is 1.59 in TRK units. This value is higher than the experimental one obtained for a 10-20 MeV energy interval (1.37) but lower than the one obtained for a energy interval going up to 140 MeV (1.78) . In this case, however, another mechanism, the โquasideuteron effectโ, is expected to play a major role in the photon absorption .
It is of great interest, beyond nuclear physics itself, to study the amount of excited low-lying dipole strength, that is the often called โpygmyโ resonances. In terms of EWSR, we obtain much less than 1% strength below 10 MeV in Ni and Sn nuclei, and about that amount in <sup>208</sup>Pb. The result for Pb is in agreement with the data of Ref. . The absence of collective states in the low-lying region is at variance with the results of relativistic RPA calculations , but agrees with the arguing in Ref. . There, it is pointed out that the soft dipole strength should decrease in nuclei displaying a neutron skin, compared to that in light halo nuclei because of a more efficient coupling to the IVGDR. On the other hand, the coupling to 2pโ2h can significantly increase the amount of low-lying strength .
By introducing a very small renormalization factor (1.01-1.03) of the residual interaction the isoscalar spurious mode can be made to appear at zero frequency. This factor is introduced only in the $`J^\pi =1^{}`$ subspace. In Table 9, the values of the energy of this state are shown as calculated with or without different parts of the D1Sโ p-h interaction. For each nucleus the same renormalisation factor is used in the four cases. The symbol $`\mathrm{}`$ means that the RPA eigenvalue is imaginary. These results show, as expected, that the consistency between the HF field and the residual interaction is important for the treatment of the spurious states.
### 3.4 Octupole states
As shown in Figure 7, the $`J^\pi =3^{}`$ states belong to two well-separated energy regions. Only the component at energies larger than $``$15 MeV can be considered as a genuine giant resonance, the High Energy Octupole Resonance (HEOR). Keeping only high energy regions (19-35 MeV for <sup>100</sup>Sn, 22-31 MeV for <sup>132</sup>Sn, 22-44 MeV for <sup>78</sup>Ni and 13-28 MeV for <sup>208</sup>Pb), the mean calculated HEOR energies are 28.16 MeV, 26.06 MeV, 29.51 MeV and 23.20 MeV, respectively. These values give systematics $`E_0A^{1/3}`$, with $`E_0=`$ 130, 132, 126, and 137 in the four nuclei, to be compared with the usual estimate $`110A^{1/3}`$ . Previous studies in stable nuclei gave values between $`130A^{1/3}`$ and $`140A^{1/3}`$ for heavy nuclei and around $`120A^{1/3}`$ in lighter ones. We therefore do not observe a strongly different behaviour of HEOR energies in exotic nuclei compared to the one previously obtained along the valley of stability.
The characteristics of the low energy $`3^{}`$ states are reported in Table 9. The influence of the different components of the D1Sโ force included in the p-h interaction is also shown. The effect of the spin-orbit term appears to be smaller than for the quadrupole states in Table 5, especially for <sup>78</sup>Ni.
### 3.5 Isovector strength
In Figures 810, the fractions of the isovector EWSR carried by the $`J^\pi =`$ $`0^+`$, $`2^+`$, $`3^{}`$ states is drawn. In this case, systematics for stable nuclei are not yet well known and is not reported. Note that only the transition operator is changed compared to the isoscalar case in Figures 3, 4 and 7. From the comparison between the two sets of figures, a much larger fragmentation of the strength is found in the isovector case, and a mixed (isoscalar-isovector) character of several states appears, as expected, in particular in <sup>78</sup>Ni.
## 4 Conclusion
To summarize, we have presented the results obtained for different giant resonances in three doubly magic exotic nuclei, using the HF+RPA approach and the Gogny force. The largest difference with usual doubly magic nuclei inside the valley of stability occurs in <sup>78</sup>Ni where the ISGMR appears significantly lower than systematics. This seems to be due to the large proton-neutron asymmetry of this nucleus.
The fragmentation of the isovector dipole strength has to be explored further in order to see the correlation or the no-correlation with proton-neutron radius differences. In particular, the nature of the double-peaks obtained in tin isotopes remains to be determined.
Results obtained in the three exotic nuclei for the ISGQR and HEOR resonances are similar to those of <sup>208</sup>Pb, but more exotic systems have to be studied to confirm such a trend.
Low energy states and B(E2) values appear to be well reproduced within the present approach, in particular the first $`2^+`$ in <sup>132</sup>Sn.
From a more general point of view, we have found that the spin-orbit component of the p-h residual interaction plays a very important role in the structure of the low-lying quadrupole and octupole states, as it strongly influences both excitation energies and transition probabilities. Similarly, our results show that including the Coulomb force in the RPA p-h matrix elements significantly affects IVGDR energies and EWSR.
## 5 Acknowledgments
The authors want to thank D. Gogny for his interest in this work and useful comments. P.F.B. acknowledges the Service de Physique Nuclรฉaire, CEA/DAMโIleโdeโFrance at BruyรจresโleโChรขtel for financial support and warm hospitality during the periods in which parts of this work were performed.
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# Bulk Viscosity Effects on the Early Universe Stability
## I Introduction
The very early stages of the (homogeneous and isotropic) Universe evolution, say at temperatures greater than $`๐ช(10^{16}GeV)`$, are characterized by a thermal history which can not be regarded as settled down into the equilibrium. Indeed at sufficiently high temperatures the cross sections of the microphysical processes, responsible for the thermal equilibrium, decay like $`๐ช(1/T^2)`$ and they are no longer able to restore the equilibrium while the Universe expands. Thus, going backward, we meet stages where the expansion has an increasing rate and induces non-equilibrium phenomena in the primordial bath. The average effect of having a microphysics unable to follow the Universe expansion by equilibrium stages, results into dissipative processes appropriately described by the presence of bulk (or second) viscosity.
As shown in bk76 ; bk77 ; bnk79 the thermodynamical consequences which bulk viscosity induces on the very early Universe history, can profoundly modify its dynamics and configurations in which matter is created near the singularity become possible. Such scenarios arise because the bulk viscosity coefficient $`\zeta `$ is a function of the Universe energy density $`\rho `$ and can be modeled in terms of a power-law as $`\zeta =\zeta _0\rho ^\nu `$ where $`\zeta _0,\nu =const`$ (for a discussion of this constitution equation see m73 and see also m95 ; dhi01 ).
Aim of the present work is to investigate the effects that bulk viscosity has on the stability of the isotropic Universe, i.e. the dynamics of cosmological perturbations is analyzed when viscous phenomena affect the zero and first order evolution of the system. We consider a background corresponding to a Robertson-Walker model filled with ultrarelativistic viscous matter, whose coefficient $`\zeta `$ corresponds to the choice $`\nu =1/2`$ and then we develop a perturbation theory which generalizes the Lifshitz works l46 ; lk63 to the presence of bulk viscosity. Though the analysis is performed for the case of a flat model, nevertheless it holds in general as soon as the perturbations scale remains much smaller than the Universe radius of curvature. In this respect we deal with perturbations such that $`\eta q1`$, $`2\pi /q`$ being the size of the coordinate scale and $`\eta `$ the conformal time variable. Since the dynamics we consider holds near the singularity for $`\eta 1`$, then we make allowance for arbitrarily large values of $`q`$ and therefore the condition for the general validity $`q2\pi |๐ฆ|^{1/2}`$ ($`๐ฆ`$ being the Robertson-Walker curvature parameter) can be always fulfilled.
As issue of our analysis we find that two different dynamical regimes appear when viscosity is taken into account and the transition from one regime to the other one takes place when the parameter $`\zeta _0`$ overcomes a given threshold value. However in both these stages of evolution the Universe results to be stable as it expands; the effect of increasing viscosity is that the density contrast begins to decrease with increasing $`\eta `$ when $`\zeta _0`$ is over the threshold. It follows that a real new feature arises with respect to the Lifshitz analysis when the collapsing point of view is addressed. In fact, as far as $`\zeta _0`$ remains below the threshold value, the isotropic Universe approaches the initial *Big-Bang* with vanishing density contrast and its stability is preserved in close analogy to the non-viscous behavior. But if $`\zeta _0`$ overcomes its critical value, then the density contrast explodes asymptotically ($`\eta 0`$) and the isotropic Universe results unstable approaching the initial singularity. In the Lifshitz non-viscous analysis this same backward in time instability takes place only when tensor perturbations (gravitational waves) are taken into account, since their amplitude increases backward as the inverse of the cosmic scale factor.
The new feature induced by the bulk viscosity consists of having instability simply in correspondence to scalar perturbations induced by fluctuations in the matter filling the Universe. The cosmological interest in such instability of the primordial Universe (toward scalar perturbations) comes out reversing the picture from collapse to expansion and taking into account the time reversibility of the Einstein equations. In fact if the early Universe does not emerge from the Planck era peaked around the Robertson-Walker geometry (indeed a good degree of generality in its structure is predicted either by classical and quantum argumentation bm04 ) then it can not reach (according to our analysis) an homogeneous and isotropic stage of evolution before the viscous effect become sufficiently small.
Since a reliable estimation k fixes the appearance of thermal bath into the equilibrium below temperatures $`๐ช(10^{16}GeV)`$ and this limit corresponds to the pre-inflationary age, our result supports the idea that an isotropic universe outcomes only after a vacuum phase transition settled down. In Sec. II we present the paradigm underlying perturbations theory to the Robertson-Walker Universe and provide the key expressions to determine the inhomogeneous corrections corresponding to different states of the matter filling the cosmological space. In Sec. III a brief review of the Lifshitz analysis is presented in order to compare our subsequent results with those ones outcoming from a non-viscous model. In Sec. IV we provide the basic equations governing the perturbations dynamics in a viscous isotropic Universe. In particular we generalize the energy-momentum tensor of the cosmological fluid in view of including the viscous effects and then derive the new zero and first order dynamics; the Friedmann equations corresponding to the viscous case are solved for the choice $`\nu =1/2`$. In Sec. V we solve the perturbations equations for the same value of the parameter $`\nu `$; the behavior of the 3-metric perturbations and of the density contrast is provided in the asymptotic limit to the *Big-Bang* and the outcoming issues are discussed. In Sec. VI brief concluding remarks follow and possible upgradings for the description of the bulk viscosity cosmological effects are discussed.
## II Perturbations theory to the Einstein equations
In order to describe the temporal evolution of the energy density small fluctuation, we develop a perturbations theory on the Einstein equations. We limit our work to the study of space regions having small dimensions compared with the scale factor of the Universe $`a`$ ll . According to this approximation, we can consider a 3-dimensional Euclidean (time dependent) metric as spatial component of the background line element
$$ds^2=dt^2a^2(dx^2+dy^2+dz^2).$$
(1)
In the linear approximation, perturbed Einstein equations are
$$\delta R_\mu ^\nu \frac{1}{2}\delta _\mu ^\nu \delta R=8\pi G\delta T_\mu ^\nu ,$$
(2)
where the term $`\delta T_\mu ^\nu `$ represents the perturbation of the energy-momentum tensor which describes the properties of the matter involved in the cosmological collapse. The perturbations of the Ricci tensor $`\delta R_\mu ^\nu `$ can be derived from the metric perturbations $`h_\mu ^\nu =\delta g_\mu ^\nu `$, since the general expression for the perturbed curvature tensor is
$$\delta R_{\mu \nu \rho }^\sigma =\frac{1}{2}(h_{\mu ;\rho ;\nu }^\sigma +h_{\rho ;\mu ;\nu }^\sigma h_{\mu \rho ;\nu }^{;\sigma }h_{\mu ;\nu ;\rho }^\sigma h_{\nu ;\mu ;\rho }^\sigma +h_{\mu \nu ;\rho }^{;\sigma }).$$
(3)
For convenience let us introduce a new temporal variable $`\eta `$, given by the relation $`dt=ad\eta `$, and use the symbol $`(^{})`$ for its derivatives; we moreover impose, without loss of generality, that the synchronous reference system is still preserved under perturbations, so that
$$h_{00}=h_{0\alpha }=0.$$
(4)
If we consider the background metric (1) the perturbations of the mist components of the Ricci tensor and of the curvature scalar read:
$`\delta R_0^0={\displaystyle \frac{1}{2a^2}}h^{\prime \prime }{\displaystyle \frac{a^{}}{2a^3}}h^{},\delta R_0^\alpha ={\displaystyle \frac{1}{2a^2}}\left(h^{,\alpha ^{}}h_\beta ^{\alpha ,\beta ^{}}\right),`$ (5a)
$`\delta R_\alpha ^\beta ={\displaystyle \frac{1}{2a^2}}\left(h_{\alpha ,\gamma }^{\gamma ,\beta }+h_{\gamma ,\alpha }^{\beta ,\gamma }h_{\alpha ,\gamma }^{\beta ,\gamma }h_{,\alpha }^{,\beta }\right){\displaystyle \frac{1}{2a^2}}h_\alpha ^{\beta ^{\prime \prime }}{\displaystyle \frac{a^{}}{a^3}}h_\alpha ^\beta ^{}{\displaystyle \frac{a^{}}{2a^3}}h^{}\delta _\alpha ^\beta ,`$ (5b)
$`\delta R={\displaystyle \frac{1}{a^2}}\left(h_{\alpha ,\gamma }^{\gamma ,\alpha }h_{,\gamma }^{,\gamma }\right){\displaystyle \frac{1}{a^2}}h^{\prime \prime }{\displaystyle \frac{3a^{}}{a^3}}h^{}.`$ (5c)
By using these expressions we are able to rewrite the left-hand side of Einstein equations through the metric perturbations $`h_\beta ^\alpha `$ in order to develop the perturbations theory after describing the matter properties via an appropriate energy-momentum tensor.
## III Primordial Universe in absence of viscosity
In this Section we assume that the Universe, in its primordial expansion, behaves like a perfect fluid. This hypothesis was made by Lifshitz in his works l46 ; lk63 and can be expressed writing the energy-momentum tensor in the form:
$$T_{\mu }^{}{}_{}{}^{\nu }=(\rho +p)u_\mu u^\nu pg_{\mu }^{}{}_{}{}^{\nu },$$
(6)
where $`p`$, $`\rho `$ are respectively the pressure a the energy density of the fluid, and $`u_\mu `$ is its 4-velocity expressed in the comoving system we consider, i.e.
$$u^0=1/au^\alpha =0.$$
(7)
Using the synchronous character of the perturbed metric, we are now able to write the perturbations of the above energy-momentum tensor:
$$\delta T_0^0=\delta \rho ,\delta T_0^\alpha =a(p+\rho )\delta u^\alpha ,\delta T_\alpha ^\beta =\delta _\alpha ^\beta v_s^2\delta \rho ,$$
(8)
$`v_s=\sqrt{\delta p/\delta \rho }`$ being the sound speed of the fluid.
Since we use an Euclidean background metric, we can expand the perturbations in plane waves of the form $`e^{i\text{q}\text{r}}`$, where q is the adimensional comoving wave vector being the physical one $`\text{k}=\text{q}/a`$. Here we investigate the gravitational stability which is described by the behavior of the energy density perturbation expressible only by a scalar function; thus we have to choose the scalar representation of the metric perturbations in order to involve a change, not only in the gravitational field, but also in the velocity and in the energy density ll ; lk63 . Such a development is made by the scalar harmonics $`Q=e^{i\text{q}\text{r}}`$, from which the following tensor
$$Q_\alpha ^\beta =\frac{1}{3}\delta _\alpha ^\beta Q,P_\alpha ^\beta =\left[\frac{1}{3}\delta _\alpha ^\beta \frac{q_\alpha q^\beta }{q^2}\right]Q.$$
(9)
can be constructed. We can now express the time dependence of the gravitational perturbations through two functions $`\lambda (\eta )`$, $`\mu (\eta )`$ and write the tensor $`h_\alpha ^\beta `$ in the form
$$h_\alpha ^\beta =\lambda (\eta )P_\alpha ^\beta +\mu (\eta )Q_\alpha ^\beta ,h=\mu (\eta )Q.$$
(10)
Let us now consider the primordial stages of the Universe expansion, i.e. $`\eta 1`$, when the radiation-like density dominates the matter one. The equation of state is $`p=\rho /3`$, from which the relations (for a flat Universe $`๐ฆ=0`$) arise
$$\rho =Ca^4,a=a_1\eta ,v_s^2=1/3,$$
(11)
where $`C`$ is an integration constant and $`a_1=\sqrt{8\pi GC/3}`$. In this approximation we can obtain the basic equations which describe the temporal evolution of the perturbations. Expressing equations (5) through the representation (10) and using expressions (8) in the form $`\delta T_\alpha ^\beta =\delta _\alpha ^\beta v_s^2\delta T_0^0`$, the perturbed Einstein equations give, respectively for $`\alpha \beta `$ and for contraction over these indexes, two equations for the metric perturbations
$$\lambda ^{\prime \prime }+\frac{2}{\eta }\lambda ^{}\frac{q^2}{3}(\lambda +\mu )=0,\mu ^{\prime \prime }+\frac{3}{\eta }\mu ^{}+\frac{2q^2}{3}(\lambda +\mu )=0.$$
(12)
Furthermore, taking the 0-0 components of (2), we can express the energy density directly from the adopted functions $`\lambda `$ and $`\mu `$ in the form
$$\delta \rho =\frac{Q}{24\pi Ga^2}\left[q^2(\lambda +\mu )+\frac{3a^{}}{a}\mu ^{}\right].$$
(13)
Among the solutions there are some which can be removed by a simple transformation of the reference system (compatible with its synchronous character), and therefore they do not represent any real physical change in the metric. The corresponding expression for the metric perturbations can be established, *a priori*, through a coordinates transformation ll taking into account the constraint (4):
$$\stackrel{~}{h}_\alpha ^\beta =f_{0,\alpha }^{,\beta }\frac{d\eta }{a}+\frac{a^{}}{a^2}f_0\delta _\alpha ^\beta +\left(f_\alpha ^{,\beta }+f_{,\alpha }^\beta \right),$$
(14)
where $`f_0`$, $`f_\alpha `$ are arbitrary (small) functions of the coordinates.
In the assumption $`\eta q1`$ the equations (12) admit, in the leading order, the solutions
$$\lambda =\frac{3C_1}{\eta }+C_2,\mu =\frac{2q^2}{3}C_1\eta +C_2,$$
(15)
where the fictitious solutions (14), which in our ultrarelativistic approach assume the form $`\lambda \mu =const`$ ($`f_0=0`$, $`f_\alpha =P_\alpha `$) and $`\lambda +\mu 1/\eta ^2`$ ($`f_0=Q`$, $`f_\alpha =0`$), are excluded. The final expressions for the gravitational perturbations and for the density contrast $`\delta \rho /\rho `$ can be obtained substituting this solutions in (10) and (13)
$$h_\alpha ^\beta =\frac{3C_1}{\eta }P_\alpha ^\beta +C_2(Q_\alpha ^\beta +P_\alpha ^\beta )$$
(16)
$$\frac{\delta \rho }{\rho }=\frac{q^2}{9}(C_1\eta +C_2\eta ^2)Q.$$
(17)
Here the constants $`C_1`$, $`C_2`$ must satisfy the conditions expressing the smallness of the perturbations at the moment $`\eta _0`$ when they arise; assuming that harmonics $`Q`$ are of the unity order magnitude, the inequalities $`\lambda 1`$, $`\mu 1`$ give the constraints $`C_1\eta _01`$ and $`C_21`$.
The expression of the cosmological perturbation (17) contains terms which increase, in an expanding Universe, proportionally to positive powers of the scale-factor $`a=a_1\eta `$. This expansion canโt, nevertheless, imply the gravitational instability: if we consider the magnitude order $`\eta 1/q`$, the conditions satisfied by the constants $`C_1`$, $`C_2`$ imply that the density perturbation remains small even in the higher order of approximation. This behavior of the cosmological fluctuation yields the gravitational stability of the primordial Universe; the only stability we can found in a non-viscous Universe lk63 is provided by the tensor perturbations $`h_\beta ^\alpha `$and takes place approaching backward the *Big-Bang*.
## IV Dynamical description of the viscous Universe
In the last Section we assumed that the Universe could be described by the energy-momentum tensor of an ultrarelativistic perfect fluid. The immediate generalization is to consider the presence of dissipative processes within the fluid dynamics as expected at temperatures above $`๐ช(10^{16}GeV)`$; this correction is represented by an additional term in the expression of the energy-momentum tensor (6) and it can be derived from thermodynamical properties of the fluid (in particular from the law of increasing entropy ll-fluid ; e40 ). Using the conservation law $`T_{\mu ;\nu }^\nu =0`$ we arrive clw92 ; mont01 to express the energy-momentum tensor of an ultrarelativistic viscous fluid in the form
$$T_{\mu \nu }=(\stackrel{~}{p}+\rho )u_\mu u_\nu \stackrel{~}{p}g_{\mu \nu },\stackrel{~}{p}=p\zeta u_{;\rho }^\rho ,$$
(18)
where $`p`$ denotes the usual thermostatic pressure and $`\zeta `$ is the *bulk viscosity* coefficient. In this work we neglect the so called shear viscosity (first viscosity) since in the case of isotropic cosmological evolution there is no displacement of the matter layers with respect to each other and this kind of viscosity represents the energy dissipation due to this effect.
The coefficient $`\zeta `$ is not constant and we have to express its dependence on the state parameters of the fluid. In the homogeneous models this quantity depends only on time, and therefore we may consider it as a function of the Universe energy density $`\rho `$. According to literature developments bk76 ; bk77 ; bnk79 ; m95 we assume that $`\zeta `$ depends on $`\rho `$ via a power-law of the form
$$\zeta =\zeta _0\rho ^\nu ,$$
(19)
where $`\zeta _0`$ is a constant and $`\nu `$ is an adimensional parameter. The behavior of this parameter is derived by V.A. Belinskii et al. bk76 for asymptotic values of the density energy yielding the constraint $`0\nu 1/2`$ in the region of large $`\rho `$.
Let us now perturb the viscous energy-momentum tensor in our synchronous reference system obtaining the expressions
$$\delta T_0^0=\delta \rho ,\delta T_0^\alpha =a(\stackrel{~}{p}+\rho )\delta u^\alpha ,\delta T_\alpha ^\beta =\delta _\alpha ^\beta \left[c_s^2\delta \rho +\zeta \left(\delta u_{,\gamma }^\gamma +\frac{h^{}}{2a^2}\right)\right],$$
(20)
and the relation between the mist components of the tensor
$$\delta T_\alpha ^\beta =\delta _\alpha ^\beta \left[c_s^2\delta T_0^0+\zeta \frac{h^{}}{2a}+\frac{\zeta \delta T_{0,\gamma }^\gamma }{a(\stackrel{~}{p}+\rho )}\right];$$
(21)
using the background metric (1) and the expressions (7) of the 4-velocity $`u_\mu `$ in a comoving system, the viscous pressure $`\stackrel{~}{p}`$ becomes
$$\stackrel{~}{p}=p3\zeta \frac{a^{}}{a^2},$$
(22)
and the quantity $`c_s`$ is given by
$$c_s^2v_s^23\zeta _0\frac{a^{}}{a^2}\nu \rho ^{\nu 1}.$$
(23)
Of course the presence of viscosity does not influence the expression of the Ricci tensor and its perturbations, thus we can still keep the expressions (5) and use the formula (21) to build up the equations which describe the dynamics of $`h_\alpha ^\beta `$ and $`\delta \rho `$. It is convenient to choose, as final equations, the one obtained from the Einstein ones for $`\alpha \beta `$ and for contraction over $`\alpha `$ and $`\beta `$, which read respectively
$$\left(h_{\alpha ,\gamma }^{\gamma ,\beta }+h_{\gamma ,\alpha }^{\beta ,\gamma }h_{\alpha ,\gamma }^{\beta ,\gamma }h_{,\alpha }^{,\beta }\right)+h_\alpha ^{\beta ^{\prime \prime }}+\frac{2a^{}}{a}h_\alpha ^\beta ^{}=0,\alpha \beta ,$$
(24)
$`{\displaystyle \frac{1}{2}}\left(h_{\alpha ,\gamma }^{\gamma ,\alpha }h_{,\gamma }^{,\gamma }\right)`$ $`\left(1+3c_s^2\right)+h^{\prime \prime }+`$ (25)
$`+{\displaystyle \frac{a^{}}{a}}\left(2+3c_s^212\pi G{\displaystyle \frac{a}{a^{}}}\zeta \right)h^{}+`$
$`{\displaystyle \frac{3\zeta }{2a(\stackrel{~}{p}+\rho )}}\left(h_{,\alpha }^{,\alpha ^{}}h_{\alpha ,\gamma }^{\gamma ,\alpha ^{}}\right)=0.`$
Furthermore the fictitious solutions (14) stand also in presence of dissipative processes because they are founded by a transformation of synchronous reference system.
Since we want to describe the gravitational instability in presence of viscosity, we consider the scalar representation (10) which, once substituted in the last expressions, yields the equations describing the perturbations temporal evolution. Let us now express the time dependence of the model variables; as in the last Section, we consider the earlier stages of a flat Universe corresponding to $`\eta 1`$ and with the equation of state $`p=\rho /3`$. The Universe zero-order dynamics is described by the equation of energy conservation and the Friedmann one, which are respectively
$$\rho ^{}+3\frac{a^{}}{a}(\rho +\stackrel{~}{p})=0,\frac{a^{}}{a^2}=\sqrt{\frac{8}{3}\pi G\rho }.$$
(26)
In order to integrate these equations we assume, according to the large energy density of the primordial expansion, $`\nu =1/2`$ dhi01 ; bhm02 . Substituting (22) into the above equations we obtain, for $`\nu =1/2`$:
$$\rho =Ca^{(2+2\omega )},a=a_1\eta ^{1/\omega },\omega =1\sqrt{54\pi G}\zeta _0,$$
(27)
being $`C`$ an integration constant and $`a_1=(8\omega ^2\pi CG/3)^{1/2\omega }`$. Since we consider an expanding Universe, the factor $`a`$ must increase with positive power of the temporal variable thus we obtain the constraint $`0<\omega 1`$.
Using these explicit dependences we get two equations for the $`\lambda `$, $`\mu `$ time functions
$$\lambda ^{\prime \prime }+\frac{2}{\omega \eta }\lambda ^{}\frac{q^2}{3}\left(\lambda +\mu \right)=0,$$
(28)
$`\mu ^{\prime \prime }+\left({\displaystyle \frac{2+3c_s^2}{\omega \eta }}\right)\mu ^{}({\displaystyle \frac{12\pi \sqrt{C}G\zeta _0}{a_1^{1+\omega }\eta ^{1+1/\omega }}}`$ $`)\mu ^{}+{\displaystyle \frac{q^2}{3}}(\lambda +\mu \left)\right(1+3c_s^2)+`$ (29)
$`+{\displaystyle \frac{q^2\zeta _0}{4\sqrt{C}/3a_1^\omega 3\zeta _0/\omega }}\eta \left(\mu ^{}+\lambda ^{}\right)=0,`$
whose solutions describe the evolution of the gravitational perturbations and the energy density fluctuation.
## V The behavior of density perturbation
In the model developed in this paper, we study the gravitational collapse dynamics of the primordial Universe near the initial *Big-Bang* for $`\eta 1`$. As in Lifshitzโs works, we now analyze the case of perturbations scale sufficiently large to use the approximation $`\eta q1`$. In the non-viscous model the cosmological stability of the isotropic Universe is guarantied by the positive power-law exponents of the density contrast evolution and by the constraint for the constants $`C_1`$, $`C_2`$. In our scheme equations (28) and (29) admit asymptotic analytic solutions for the functions $`\lambda `$ and $`\mu `$; in the leading order $`\lambda `$ takes the form
$$\lambda =\frac{C_1}{\eta ^{2/\omega 1}}+C_2,$$
(30)
where $`C_1`$, $`C_2`$ are two integration constants. Substituting this expression in (29) we get, in the same order of approximation, the behavior of the function $`\mu `$ as
$$\mu =\frac{\stackrel{~}{C}_1}{\eta ^{1/\omega 3}}+C_2,$$
(31)
where we have excluded the non-physical solutions (14) as written in the form $`\lambda \mu =const`$. The constant $`\stackrel{~}{C}_1`$ is given by the expression $`\stackrel{~}{C}_1=A/B(31/\omega )`$, $`A`$ and $`B`$ being constants involved in the equation for $`\mu `$ having the form
$$A=\frac{C_1q^2}{3}\left(1+3c_s^2\right)+\frac{C_1(12/\omega )q^2\zeta _0}{4\sqrt{C}/3a_1^\omega 3\zeta _0/\omega },B=\frac{12\pi \sqrt{C}G\zeta _0}{a_1^{1+\omega }}.$$
(32)
Let us now write the final form of the perturbations pointing out their temporal dependence in the viscous Universe. The gravitational perturbations (10) become
$$h_\alpha ^\beta =\frac{C_1}{\eta ^{2/\omega 1}}P_\alpha ^\beta +\frac{\stackrel{~}{C}_1}{\eta ^{1/\omega 3}}Q_\alpha ^\beta +C_2\left(Q_\alpha ^\beta +P_\alpha ^\beta \right),$$
(33)
and the density contrast reads
$$\frac{\delta \rho }{\rho }=F_\omega \left[C_1\eta ^{32/\omega }+C_2\eta ^2+C_3\eta ^{31/\omega }+\stackrel{~}{C}_1\eta ^{51/\omega }\right],$$
(34)
where $`C_3=3A/q^2\omega B`$, and $`F_\omega =\omega ^2Qq^2/9`$. As in the non-viscous case, we now impose the conditions expressing the smallness of perturbations at the initial time $`\eta _0`$. The inequalities $`h_\alpha ^\beta 1`$ and $`\delta \rho /\rho 1`$ yield only two fundamental constraints for the integration constants: $`C_1\eta _0^{2/\omega 1}`$ and $`C_21`$ for any $`\omega `$-value within the interval $`(0,\mathrm{\hspace{0.17em}1}]`$. Furthermore we find an additional condition which involves the wave number $`q`$ and the integration constant $`C`$; in particular a rough estimate for $`\omega <1/3`$ of the inequalities $`\stackrel{~}{C}_1\eta _0^{1/\omega 3}`$ and $`C_3\eta _0^{1/\omega 3}`$ yields the condition $`q(GC\eta _0)^{1/2\omega }`$ which ensures the smallness of the cosmological perturbations.
Using the hypothesis $`\eta 1`$ we can get the asymptotic form the corrections to the cosmological background. The exponents of the variable $`\eta `$ can be positive or negative according to the value of the viscous parameter $`\omega `$, which is always positive but less than unity. This behavior produces two different regimes of the density contrast evolution: in the first case $`2/3<\omega 1`$, the perturbation increases forward in time which corresponds qualitatively to the same picture of the non-viscous Universe where the density fluctuation increases but remains small; while in the other case $`0<\omega <2/3`$ the density contrast is suppressed behaving like a negative power of $`\eta `$. When the density contrast results to be increasing, the presence of viscosity induces a *damping* of the perturbation evolution in the direction of the expanding Universe, so the cosmological stability is fortified since the leading $`\eta `$ powers are smaller than the non-viscous ones (17).
For $`0<\omega <2/3`$, i.e. for a large coefficient $`\zeta _0`$, the density contrast evolution is deeply modify by the presence of viscosity. We get the leading order expression
$$\frac{\delta \rho }{\rho }=\frac{C_1F_\omega }{\eta ^\kappa },\kappa =2/\omega 3>0.$$
(35)
This behavior is very different from the non-viscous one since it yields a strong *damping* of the density contrast. In this regime the density fluctuation decreases forward in time but the most interesting result is the instability which the isotropic and homogeneous Universe acquires in the direction of the collapse toward the *Big-Bang*. For $`\omega <2/3`$ the density contrast diverges when approaching the cosmological singularity, i.e. for $`\eta 0`$. In a non-viscous Universe the only perturbations which are able to generate the same kind of asymptotic instability are the tensor fluctuations (gravitational waves) whereas now the scalar perturbations destroy asymptotically the primordial Universe symmetry. The dynamical implication of this issue is that an isotropic and homogeneous stage of the Universe can not be generated, from generic initial conditions, as far as the viscosity becomes smaller than the critical value associated to the condition $`\omega =2/3`$.
## VI Concluding remarks
The main issue of our investigation is to have shown that the isotropic Universe acquires, backward in time, a regime of instability corresponding to sufficiently high values of the viscous parameter $`\zeta _0`$. Such a window of instability implies that, if the Universe was born sufficiently far from the homogeneous and isotropic stage, than the bulk viscosity (i.e. the absence of a stable thermal equilibrium) works against isotropization mechanisms and the inflation becomes the scenario from which a Robertson-Walker geometry arises (at least on a given scale). The explanation of this result is in the real physical meaning of the bulk viscosity: such viscous effects come out from the difficulty that microphysics finds to restore the thermal equilibrium against the rapid Universe expansion. As a natural consequence of this physical context, bulk viscosity makes unfavored the establishment of an homogeneous stage from a general cosmological dynamics. On the other hand in a Robertson-Walker Universe, already settled down, we expect that, as we find, the viscous effects depress the density contrast because the particles inside the inhomogeneous fluctuations undergo dissipative processes which frozen the growth of the structures. Despite of the reliable feature of our result, the present investigation, as well as the whole previous literature on this subject, relies on a phenomenological ground; in fact the description of the viscous effect is based on the constitutive equation relating the viscosity coefficient to a power-law of the system energy density. This statement appears well-grounded, but nevertheless it requires to be carefully considered in a precise derivation of the viscosity coefficient from a real kinetic theory of matter ll-kinetic . We will address for such a point in a further investigation, which will be aimed to yield an upgrading of the present cosmological issue.
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# Suppression of Shot-Noise in Quantum Cavities: Chaos vs. Disorder
## Abstract
We investigate the behavior of the shot-noise power through quantum mechanical cavities in the semiclassical limit of small electronic wavelength. In the absence of impurity scattering, the Fano factor $`F`$, giving the noise to current ratio, was previously found to disappear as more and more classical, hence deterministic and noiseless transmission channels open up. We investigate the behavior of $`F`$ as diffractive impurities are added inside the cavity. We find that $`F`$ recovers its universal value provided (i) impurities cover the full cavity so that only a set of zero measure of classical trajectories may avoid them, and (ii) the impurity scattering rate exceeds the inverse dwell time through the cavity. If condition (i) is not satisfied, $`F`$ saturates below its universal value, even in the limit of strong scattering. Our results corroborate the validity of the two-phase fluid model according to which the electronic flow splits into two well separated components, a classical deterministic fluid and a stochastic quantum-mechanical fluid. Only the latter carries shot-noise.
###### Keywords:
mesoscopic transport, noise, quantum chaos, random matrix theory
Time-resolved transport measurements through quantum mechanical systems invariably observe current fluctuations, even in the (experimentally unrealistic) situation of a noiseless measurement apparatus and at zero temperature. This intrinsically quantal noise is usually referred to as shot-noise. It results from the quantization of charge together with the statistical nature of quantum mechanical transport Blan00 . As but one of the consequences of the quantum-classical correspondence at large quantum numbers, it has been predicted that, shot-noise through a chaotic ballistic cavity disappears as the system becomes more and more classical, i.e. when the ratio of the electronic Fermi wavelength to the linear cavity size vanishes $`\lambda _\mathrm{F}/L0`$ Ben91 . The purpose of this article is to discuss when and how shot-noise starts to be reduced by an emergent classical, deterministic behavior.
Recent technological advances have made it possible to make electronic systems small and clean enough that the resulting electronic mean free path is larger than the size of the confining potential defining the device Revdot1 . The electronic motion in these quantum dots is thus ballistic, and provided that their wavelength is short enough, the electrons have a dynamics strongly related to the dynamics that a classical particle would have. When this classical dynamics is chaotic, that is, when the shape of the dot differs significantly from a circle or an ellipse, the transport properties are usually universal and well-captured by the Random Matrix Theory (RMT) of transport Ben97 . The starting point of RMT is the scattering approach scatg , which relates transport properties to the systemโs scattering matrix
$`๐ฎ=\left(\begin{array}{cc}๐ซ\hfill & ๐ญ^{}\hfill \\ ๐ญ\hfill & ๐ซ^{}\hfill \end{array}\right).`$ (3)
Here we consider a symmetric two terminal geometry (the cavity is connected to two external leads with equal number $`N`$ of propagating channels) for which $`๐ฎ`$ is a 2-block by 2-block matrix, written in terms of $`N\times N`$ transmission ($`๐ญ`$ and $`๐ญ^{}`$) and reflection ($`๐ซ`$ and $`๐ซ^{}`$) matrices. From $`๐ฎ`$, the systemโs conductance is given by $`g=\mathrm{Tr}(๐ญ^{}๐ญ)=_nT_n`$ ($`g`$ is expressed in units of $`e^2/h`$ and the $`T_n`$โs are the $`N`$ eigenvalues of $`๐=๐ญ^{}๐ญ`$). RMT provides a statistical theory of transport where $`๐ฎ`$ is assumed to be uniformly distributed over one of Dysonโs circular ensemble of random matrices Meh91 . Transport properties can be calculated from this sole assumption. For instance, within RMT, and in the limit $`N1`$ the transmission eigenvalues have a probability distribution Ben97
$$P_{\mathrm{RMT}}(T)=\frac{1}{\pi }\frac{1}{\sqrt{T(1T)}}$$
(4)
for any $`T[0,1]`$. Note that classical particles would be either deterministically transmitted, $`T=1`$ or reflected $`T=0`$.
The distribution of transmission eigenvalues is all one needs to get fluctuations and higher moments of the current at low frequency. The zero-frequency shot-noise power in particular is given by $`S=2eV_nT_n(1T_n)`$ Blan00 . According to (4), $`S`$ is suppressed below its Poissonian value of $`S_\mathrm{p}=2eI`$ ($`V`$ is the applied voltage and $`I`$ the time-averaged current) by the Fano factor which reads
$$F=\frac{_nT_n(1T_n)}{_nT_n};F_{\mathrm{RMT}}=\frac{1}{4}.$$
(5)
The RMT predictions (4) and (5) have been confirmed in various transport experiments and numerical simulations on open chaotic cavities Blan00 ; Ben97 .
In closed chaotic systems, the semiclassical limit $`\lambda _\mathrm{F}/L\mathrm{}_{\mathrm{eff}}0`$ usually results in a better and better agreement with the Hamiltonian RMT of spectral fluctuations Haake . One may thus expect that the same applies to transport in open systems. That is not so, as illustrated in Fig. 1. The numerical data presented there (see also Ref. Jac04 ) show that instead of the RMT prediction of Eq. (4), the transmission eigenvalues appear to be distributed according to
$$P_\alpha (T)=\alpha P_{\mathrm{RMT}}(T)+\frac{1\alpha }{2}\left[\delta (T)+\delta (1T)\right],$$
(6)
with an increasingly deterministic behavior $`\alpha 0`$ as $`\mathrm{}_{\mathrm{eff}}0`$ (all classical parameters being fixed). The presence of $`\delta `$-peaks at $`T=0`$ and $`T=1`$ in $`P_\alpha (T)`$ becomes evident once the integrated distribution $`I(T)=_0^TP(T^{})๐T^{}`$ is plotted. One has
$$I_\alpha (T)=\frac{2\alpha }{\pi }\mathrm{sin}^1\sqrt{T}+\frac{1\alpha }{2}(1+\delta _{1,T}),$$
(7)
so that $`I_\alpha (0)=(1\alpha )/2`$ vanishes only for $`\alpha =1`$. It turns out that the parameter $`\alpha `$ is well approximated by $`\alpha \mathrm{exp}(\tau _e/\tau _d)`$, in term of the new time scale $`\tau _e=\lambda ^1\mathrm{ln}[\mathrm{}_{\mathrm{eff}}\tau _d^2]`$ and the average dwell time $`\tau _d`$ through the cavity Jac04 . In short, for a classically fixed configuration (i.e. considering an ensemble of systems with fixed $`\lambda `$ and $`\tau _d`$), the fraction $`(1\alpha )/2`$ of deterministic transmission eigenvalues $`T=0,1`$ increases as one goes deeper and deeper into the semiclassical limit, $`\mathrm{}_{\mathrm{eff}}0`$. The rate of the crossover is set by a partially quantum, partially classical time scale, the Ehrenfest time $`\tau _e`$ Zas81 . Compared to closed systems, the emergence of a finite $`\tau _e`$ has more profound an impact on transport properties once it becomes comparable to the dwell time $`\tau _d(N\mathrm{}_{\mathrm{eff}})^1`$.
Inserting (6) into (5) with the numerically extracted value $`\alpha \mathrm{exp}(\tau _e/\tau _d)`$, one directly recovers $`F\mathrm{exp}[\tau _e/\tau _d]`$, in agreement with the analytical prediction of Refs. Agam00 ; Sil03 , the experimental results of Ref. Ober02 and the numerical data of Ref. Two03 . They can be qualitatively understood by first realizing that complex quantum systems split into the two classes of quantum chaotic and quantum disordered systems Ale96 . Ballistic cavities in particular belong to the first class, for which electronic wavepackets are carried along very few classical paths until the time $`\tau _e`$ after which they have a finite probability to be found on trajectories that their center of mass would not follow classically. This dynamical diffraction process restores quantum mechanical stochasticity for larger times, however it does not affect short trajectories with $`\tau <\tau _e`$. Thus, in quantum chaotic systems, two classes of classical trajectories emerge, depending on their dwell time through the cavity. Short trajectories with $`\tau <\tau _e`$ are able to carry an electronic wavepacket deterministically through the cavity (i.e. with transmission probability $`T=0`$ or 1). If the electronic wavepacket sits on longer trajectories with $`\tau >\tau _e`$ on the other hand, diffraction splits it into pieces before its exit, and quantum mechanical stochasticity (with $`T]0,1[`$) prevails. The fraction of scattering trajectories in the stochastic subset is obtained via the dwell time distribution $`\rho (\tau )`$ by
$$\alpha _{\tau _e}^{\mathrm{}}\rho (\tau )d\tau .$$
(8)
In a chaotic system one has $`\rho (\tau )=\tau _d^1\mathrm{exp}[\tau /\tau _d]`$, hence the fraction of stochastically transmitted channels gives $`\alpha \mathrm{exp}[\tau _e/\tau _d]`$ for the weight $`\alpha `$ of Eq. (6), in agreement with the numerics shown in Fig. 1. This is the essence of the two-phase-fluid model Jac04 , originally proposed in Ref. Sil03 and given a microscopic foundation in Ref. wj2004
Fig. 1 provides us with a direct evidence for the validity of the two-phase fluid. Other evidences of this kind have been found in investigations of the excitation spectrum of Andreev billiards, i.e. ballistic cavities in contact with a superconductor Goo05 . In a previous work wj2004 , we used an approach based on a semiclassical expansion for the Greenโs function in term of a sum over classical trajectories. Together with the construction of a quantum-mechanical phase-space basis, this allowed us to import classical concepts such as Liouville conservation and determinism into quantum mechanics, thereby providing with a microscopic foundation for the two-phase fluid model comment . Our purpose in the reminder of this article is to check numerically the model by investigating its quantum chaos โ quantum disorder crossover.
To this end, we use the open kicked rotator model and follow the procedure described e.g. in Ref. Two03 . So far, the model has been implemented only in its quantum chaotic version. Here we add a diffraction term to it so that its Floquet (time-evolution) operator has matrix elements
$`U_{m,m^{}}`$ $`=`$ $`\mathrm{}_{\mathrm{eff}}^{1/2}e^{(iK/4\pi \mathrm{}_{\mathrm{eff}})[\mathrm{cos}(2\pi m\mathrm{}_{\mathrm{eff}})+\mathrm{cos}(2\pi m^{}\mathrm{}_{\mathrm{eff}})]}e^{(i/4\pi \mathrm{}_{\mathrm{eff}})[\eta (m)+\eta (m^{})]}`$ (9)
$`\times {\displaystyle \underset{l}{}}e^{2\pi il(mm^{})\mathrm{}_{\mathrm{eff}}}e^{(\pi i\mathrm{}_{\mathrm{eff}}/2)l^2},`$
with a randomly distributed function $`\eta (m)\eta (m^{})=u^2\delta _{m,m^{}}\mathrm{\Theta }(|mm_0|\xi /2)`$. Such a point-like diffractive impurity potential corresponds to the extreme quantum disorder limit. The length scale $`\xi `$ allows to cover all or only a part of the system, and $`u`$ sets the strength of the impurity scattering. We extracted the scattering rate $`\mathrm{\Gamma }_Q`$ from the width of the Local Spectral Density of States (LDOS) induced by $`\eta (m)`$ for the closed kicked rotator. As is commonly the case Jac95 , we found that the LDOS has a Lorentzian shape with a width $`\mathrm{\Gamma }_Q0.016u^2\sqrt{\xi }/\mathrm{}_{\mathrm{eff}}`$ in a wide range of parameter. One expects that, as $`\mathrm{\Gamma }_Q`$ becomes comparable to the Ehrenfest time, quantum diffraction effects start to eat away the determinism of short trajectories.
We first show in Fig. 2 the behavior of the Fano factor as the diffractive scattering rate is cranked up, all other parameter being fixed. For $`\mathrm{\Gamma }_Q=0`$, $`F<0.25`$ lies below its universal value, indicating a finite $`\tau _e`$, thus a finite fraction of non-diffractive scattering orbits. This fraction gets reduced once $`\mathrm{\Gamma }_Q`$ increases, and accordingly $`F/\mathrm{\Gamma }_Q>0`$. The most striking feature of Fig. 2, however, is that as long as the diffractive disorder does not cover the whole system volume, i.e. for $`\xi <\mathrm{}_{\mathrm{eff}}^1(1\tau _d^1)`$, $`F`$ saturates below its universal value, even in the limit $`\mathrm{\Gamma }_Q\tau _e\mathrm{}`$. This reflects the fact that some short trajectories are able to avoid the diffractive potential and thus remain deterministic. The existence of two separated phase-space fluids allows only that part of the deterministic fluid which directly scatters off the impurities to become stochastic. It is not clear from our numerics if, for $`\xi >\mathrm{}_{\mathrm{eff}}^1(1\tau _d^1)`$, one has an exponential or a power-law behavior of the Fano factor as predicted in Ref. Ober02 ; Bul05 . More detailed investigations are necessary to draw definite conclusions.
We finally show on Fig. 3 the behavior of the distribution of transmission eigenvalues in the saturated regime $`\mathrm{\Gamma }_Q\tau _e1`$. Clearly, the distribution and integrated distribution follow Eqs. (6) and (7). In contrast to the quantum chaotic case, we qualitatively found a dependence $`\alpha \mathrm{exp}[\tau _e/\tau _d]\xi \mathrm{}_{\mathrm{eff}}/(1\tau _d^1)`$, again reflecting a reduction of that part of the deterministic component which directly touches the diffractive potential.
All these findings support the two-phase fluid hypothesis, that is, the splitting of the cavity into a stochastic and a deterministic cavity. The latter being noiseless, shot-noise is suppressed by a factor reflecting its phase-space measure relative to the total phase-space. This measure is reduced by the presence of diffractive disorder, however, only homogeneously spread impurities are able to diffract all trajectories, thus only in this case does one recover universality at large diffractive scattering rate. This complement recent investigations of shot-noise with homogeneously spread diffractive disorder in regular cavities Aig05 .
This work has been supported by the Swiss National Science Foundation.
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# Nuclear fusion in muonic deuterium-helium complex
## I Introduction
The formation of muonic molecules of hydrogen isotopes and their nuclear reactions have been the subject of many experimental and theoretical studies Marshall et al. (2001); Petitjean (2001, 1992); Cohen (1990/91); Ponomarev (1990, 2001); Nagamine (2001); Bogdanova and Markushin (1990). As to the studies of formation of charge-asymmetrical muonic molecules like $`h\mu Z`$ ($`h=p,d,t`$, $`Z`$ are nuclei with a charge $`Z>1`$) and their respective nuclear fusion, the situation slightly different. What gave an impetus to study such systems was the theoretical prediction and experimental observation of the molecular mechanism for charge exchange (MMCE) of $`p\mu `$ atoms on He nuclei Aristov et al. (1981); Bystritsky et al. (1983). Essentially, the mechanism is reduced to the following. Colliding with a He atom in a $`H`$โHe mixture ($`H=\mathrm{H}_2,\mathrm{D}_2,\mathrm{T}_2`$ and $`\mathrm{He}={}_{}{}^{3}\mathrm{He},{}_{}{}^{4}\mathrm{He}`$), the muonic hydrogen atom forms a muonic complex $`h\mu \mathrm{He}`$ in the excited $`2p\sigma `$ state. In the case of a deuteriumโhelium mixture, the complex may then decays from this state (see Fig. 1) via one of three channels
$`d\mu +\mathrm{He}\stackrel{\lambda _{d\mathrm{He}}}{}`$ $`[(d\mu \mathrm{He})^{}e^{}]^+`$ $`+e^{}`$
$``$
$`\stackrel{\lambda _\gamma }{}`$ $`[(d\mu \mathrm{He})^+e^{}]+\gamma `$ (1a)
$`\stackrel{\lambda _p}{}`$ $`[(\mu \mathrm{He})_{1s}^+e^{}]+d`$ (1b)
$`\stackrel{\lambda _e}{}`$ $`(\mu \mathrm{He})_{1s}^++d+e^{}.`$ (1c)
If $`\mathrm{He}={}_{}{}^{3}\mathrm{He}`$, fusion reactions may occur
$`d\mu {}_{}{}^{3}\mathrm{He}`$ $`\stackrel{\stackrel{~}{\lambda }_f}{}`$ $`\alpha +\mu +\mathrm{p}(14.64\mathrm{MeV})`$ (2a)
$`\stackrel{\stackrel{~}{\lambda }_{f\mathrm{\Gamma }}}{}`$ $`\mu ^5\mathrm{Li}+\gamma (16.4\mathrm{MeV}).`$ (2b)
Thus, the fusion proceeds by the formation of a $`d\mu `$ atom, which, when incident on a $`{}_{}{}^{3}\mathrm{He}`$ atom, forms the $`d\mu {}_{}{}^{3}\mathrm{He}`$ molecular system. This molecule has two primary spin states, $`J=1`$ and $`J=0`$ <sup>1</sup><sup>1</sup>1$`J`$ denotes the total angular momentum of the three particles.; formation favors the former, fusion the latter Bogdanova et al. (1998). In Eqs. (Iโc), $`\lambda _\gamma `$ is the $`(d\mu \mathrm{He})^{}`$ molecular decay channel for the 6.85 keV $`\gamma `$โray emission, $`\lambda _e`$ for the Auger decay, and $`\lambda _p`$ for the breakโup process. The $`d\mu \mathrm{He}`$ molecule is formed with a rate $`\lambda _{d\mathrm{He}}`$. The main fusion process, Eq. (2aa), occurs with the rate $`\stackrel{~}{\lambda }_f`$, whereas the reaction (2ab), with the associated rate $`\stackrel{~}{\lambda }_{f\mathrm{\Gamma }}`$ has a branching ratio on the order of $`10^{(4,5)}`$ Cecil et al. (1985).
Interests in further study of charge-asymmetrical systems was caused by first getting information on characteristics of the strong interaction in the region of ultralow energies. Secondly, it allows us to test the problem of three bodies interacting via the Coulomb law. More precisely, these studies may allow us to
* check fundamental symmetries and to measure the main characteristics of the strong interaction in the region of astrophysical particle collision energies ($``$keV) in the entrance channel. It should be mentioned that nuclear fusion reactions in charge-asymmetrical muonic molecules are characterized by the same astrophysical range of energies Friar et al. (1991).
* test the calculation algorithm for rates of nuclear fusion reactions in $`\mu `$-molecular complexes as well as for partial rates of decay of these asymmetrical complexes via various channels.
* solve some existing astrophysical problems.
By now the experimental discovery of the MMCE has been confirmed in a number of experiments on study of muon transfer from $`h\mu `$ to the He isotopes.
Formation rates of the charge-asymmetrical $`d\mu \mathrm{He}`$, and $`p\mu \mathrm{He}`$ systems were measured Balin et al. (1985); von Arb et al. (1989); Bystritsky et al. (1993a, 1990/91); Bystritsky et al. (1990a); Bystritsky et al. (1990b); Tresch et al. (1998a); Gartner et al. (2000); Bystritsky et al. (1995, 2003) and calculated Ivanov et al. (1986); Kravtsov et al. (1986); Kino and Kamimura (1993); Gershtein and Gusev (1993); Korobov et al. (1993); Czapliลski et al. (1997a); Belyaev et al. (1995); Czapliลski et al. (1996a); Belyaev et al. (1997) with quite a good accuracy, and partial decay rates of such complexes were found.
In the past five years interest in studying charge-asymmetrical complexes and in particular fusion in the $`d\mu {}_{}{}^{3}\mathrm{He}`$ system has revived. Table 1 presents the calculated fusion rates of deuterium and $`{}_{}{}^{3}\mathrm{He}`$ nuclei in the $`d\mu {}_{}{}^{3}\mathrm{He}`$ complex in its states with the orbital momenta $`J=0`$ and $`J=1`$ and the experimental upper limits of the effective fusion rate, $`\stackrel{~}{\lambda }_f`$, in the molecule averaged over the populations of fine-structure states of the $`d\mu {}_{}{}^{3}\mathrm{He}`$ complex.
The experimental study of nuclear fusion in the $`d\mu {}_{}{}^{3}\mathrm{He}`$ molecule is quite justified as far as detection of the process is concerned because there might exist an intermediate resonant compound state $`{}_{}{}^{5}\mathrm{Li}_{}^{}`$ leading to the expected high fusion rate which results from quite a large value of the $`S`$โfactor for the $`d{}_{}{}^{3}\mathrm{He}`$ reaction Gaughlan and Fowler (1988). However, as follows from the calculations presented in Table 1, the theoretical predictions of the fusion rate in this molecule show a wide spread in value from $`10^5\text{s}^1`$ to $`10^{11}\text{s}^1`$.
The nuclear fusion rate in muonic molecules is usually calculated on the basis of Jacksonโs idea Jackson (1957) which allows the factorization of nuclear and molecular coordinates. In this case the nuclear fusion rate $`\lambda _{nf}`$ is given by
$$\lambda _{nf}=\frac{S}{(\pi MZ_1Z_2)}\times \left|\mathrm{\Psi }_{sc}(0)\right|^2,$$
(3)
which is defined by the astrophysical $`S`$โfactor, the reduced mass of the system $`M`$, the charges of nuclei in the muonic molecule $`Z_1`$ and $`Z_2`$, and the threeโbody system wave function $`\mathrm{\Psi }_{sc}(0)`$ averaged over the muon degrees of freedom and taken at distances comparable with the size of the nuclei, i.e., for $`r0`$ because of the short-range nature of the nuclear forces.
It should be mentioned that, strictly speaking, asymmetrical muonic molecules ($`Z_1Z_2`$) do not form bound states but correspond to resonant states of the continuous spectrum. In this case an analogue of Eq. (3) is given in Ref. Penโkov (1997) as
$$\lambda _{nf}=\frac{S}{(\pi MZ_1Z_2)}\times \frac{1}{2l+1}\frac{Mk_0}{4\pi }\mathrm{\Gamma }\left|\mathrm{\Psi }_{sc}(0)\right|^2,$$
(4)
where $`l`$ is the orbital quantum number of the resonant state, $`k_0`$ is the relative momentum corresponding to the resonant energy, $`\mathrm{\Gamma }`$ is the width of the molecular state and $`\mathrm{\Psi }_{sc}(0)`$ is the wave function for the state of scattering at resonant energy. In the limit of a very narrow resonance when $`\mathrm{\Gamma }0`$ Eqs. (3) and (4) coincide. However, one should take into account the asymptotic part of the wave function responsible for an inโflight fusion, including the possible interferences between the resonant and nonresonant channels.
Let us briefly discuss the calculated nuclear fusion rates in the $`d\mu {}_{}{}^{3}\mathrm{He}`$ reaction presented in Table 1. The value given in Refs. Nagamine et al. (1989); Kino and Kamimura (1993) were given with some references to a calculation by Kamimura but without any references to the calculation method. In Ref. Penโkov (1997) the author used a small variation basis and the experimental value of the astrophysical factor $`S6.32\text{MeV}\times \text{b}`$ and found the nuclear fusion rate in the $`d\mu {}_{}{}^{3}\mathrm{He}`$ molecule in the $`J=0`$ state to be $`3.8\times 10^6\mathrm{s}^1`$.
In Refs. Czapliลski et al. (1996b, 1998) the nuclear fusion rate in the $`d\mu {}_{}{}^{3}\mathrm{He}`$ complex from the $`J=0`$ state was calculated by various methods. Since the nuclear fusion rate in the $`1s\sigma `$ states of the $`d\mu {}_{}{}^{3}\mathrm{He}`$ molecule is much higher than the fusion rate form the $`2p\sigma `$ state (because of a far smaller potential barrier), the under-barrier $`2p\sigma 1s\sigma `$ transition was calculated with finding the transition point in the complex $`r`$-plane. This procedure is not quite unambiguous and therefore the nuclear fusion rate in the $`d\mu {}_{}{}^{3}\mathrm{He}`$ molecule was calculated in an alternative way by reducing it to the $`S`$โfactor and using experimental data on low-energy scattering in $`{}_{}{}^{3}\mathrm{He}(dp){}_{}{}^{4}\mathrm{He}`$ reactions from Ref. Czapliลski et al. (1996b). However, the procedure of an approximation of the experimental data for the ultralow energy region leads to some ambiguity of the results. The results of the calculation by the above two methods may differ by a factor of five for the $`t\mu {}_{}{}^{3}\mathrm{He}`$ molecule and by a factor of three for the $`d\mu {}_{}{}^{3}\mathrm{He}`$ molecule in question Czapliลski et al. (1998).
The highest nuclear fusion rate was obtained in Ref. Harley et al. (1989). Unlike the case in Ref. Czapliลski et al. (1998), where the barrier penetration factor in the $`2p\sigma 1s\sigma `$ transition was evaluated, in Ref. Harley et al. (1989) the contribution from the $`1s\sigma `$ state to the total wave function for the at small internuclear distances $`r`$ was determined. The determination of the contribution from this state to the total mesomolecule wave function at small distance requires the solution of a multichannel system of differential equations, which is a complicated problem because of the singularity of the expansion coefficients at small distances $`r0`$. As to the results of Ref. Bogdanova et al. (1999) given in the last column of Table 1, it is difficult to judge the calculation method used because the method for calculation the wave function at small distances was not presented in the paper.
Different results of calculations of the fusion rate in the $`d\mu {}_{}{}^{3}\mathrm{He}`$ molecule reflect different approximations of the solution to the Schrรถdinger equation for three particles with Coulomb interaction. The main uncertainty is associated with the results at small distances and hence follows the spread of the calculated values for the nuclear fusion rate in the $`d\mu {}_{}{}^{3}\mathrm{He}`$ molecule given in Table 1. When the adiabatic expansion is used, the important problem of convergence of this expansion at small distances is usually ignored. Such problems vanish if the direct solution of the Faddeev equations in the configuration space is performed in Refs. Kostrykin et al. (1989); Hu et al. (1992); Hu and Kvitsinsky (1992). For this reason the calculation of the fusion rate in the $`d\mu {}_{}{}^{3}\mathrm{He}`$ molecule using Faddeev equations in order to adjudge discrepancies between different theoretical results becomes very actual problem.
Much less has been done to study the nuclear fusion reaction in the $`d\mu {}_{}{}^{3}\mathrm{He}`$ experimentally. The estimations of the lower limit for the fusion reaction (2aa) rate, has been done by a Gatchina โ PSI collaboration using an ionization chamber Balin et al. (1992, 1998); Maev et al. (1999). Their results (see Table 1) differ by several orders of magnitude. Another experiment aimed to measure the effective rate, $`\stackrel{~}{\lambda }_{f,p}`$, of nuclear fusion reaction (2aa) was performed by our team Del Rosso et al. (1999). A preliminary result, also as estimation of lower limit, is shown in Table 1.
The purpose of this work was to measure the effective rate, $`\stackrel{~}{\lambda }_f`$, of nuclear fusion reaction (2aa) in the $`d\mu {}_{}{}^{3}\mathrm{He}`$ complex with the formation of a 14.64 MeV proton at two $`\mathrm{D}_2+{}_{}{}^{3}\mathrm{He}`$ mixture density values.
## II Measurement method
Figure 1 shows a slightly simplified version of the kinetics to be considered, when negative muons stop in the $`\mathrm{D}_2+{}_{}{}^{3}\mathrm{He}`$ mixture. The information on the fusion reaction (2aa) rate in the $`d\mu {}_{}{}^{3}\mathrm{He}`$ complex can be gained by measuring the time distribution, $`dN_p/dt`$, and the total yield, $`N_p`$, of 14.64 MeV protons. These quantities are derived from the differential equations governing the evolution of the $`J=1,0`$ states of the $`d\mu {}_{}{}^{3}\mathrm{He}`$ molecules
Establishing the time dependence of the number of $`d\mu {}_{}{}^{3}\mathrm{He}`$ molecules, $`N_{d\mu {}_{}{}^{3}\mathrm{He}}^J(t)`$, for the two possible states $`J`$ is sufficient to predict the time spectrum of the fusion products. In the following, we will include the effective transition rate $`\stackrel{~}{\lambda }_{10}`$ of the $`d\mu {}_{}{}^{3}\mathrm{He}`$ complex between the states $`J=1`$ and $`J=0`$. The $`\stackrel{~}{\lambda }_{10}`$ transition is important if the $`\stackrel{~}{\lambda }_f^1`$ and $`\stackrel{~}{\lambda }_f^0`$ rates differ strongly from one to another, and an appropriate value of $`\stackrel{~}{\lambda }_{10}`$ permits the two rates to be measured. This possibility can be checked by measuring the fusion using different concentrations and densities which should also help clear up the questions surrounding the mechanism of the $`\stackrel{~}{\lambda }_{10}`$ transition Bystritsky et al. (1999a), which is predicted to scale nonlinearly with the density.
There is a direct transfer rate from ground state $`d\mu `$โs to $`{}_{}{}^{3}\mathrm{He}`$โs but that rate is about 200 times smaller than the $`\lambda _{d{}_{}{}^{3}\mathrm{He}}`$ rate and will be ignored Matveenko and Ponomarev (1972). No hyperfine dependence on the $`\lambda _{d{}_{}{}^{3}\mathrm{He}}`$ formation rate is expected since the molecular formation involves an Auger electron and bound state energies of many tens of electron volts Aristov et al. (1981). Using the expectation that the $`d\mu {}_{}{}^{3}\mathrm{He}`$ is formed almost exclusively in the $`J=1`$ state, the solution for the fusion products from the $`J=0`$ and $`J=1`$ states is relatively straightforward given the $`d\mu `$ population. The recycling of the muon after $`d\mu {}_{}{}^{3}\mathrm{He}`$ fusion will be ignored due to the extremely small probability of the fusion itself, and thus the system of equations decouples into the $`d\mu {}_{}{}^{3}\mathrm{He}`$ sector, and the $`dd`$โfusion sector (where cycling will be considered). Since there is no expectation of a $`J=0`$ to $`J=1`$ transition, i.e., $`\lambda _{01}`$, the $`d\mu {}_{}{}^{3}\mathrm{He}`$ sector is easily solved.
Formation of $`d\mu d`$ molecules from a $`d\mu `$ in hyperfine state $`F=3/2`$ and $`F=1/2`$ is given by the effective rate $`\stackrel{~}{\lambda }_F`$, whereas the branching ratio $`\beta _F`$ and sticking probability $`\omega _d`$ model the number of muons lost from the cycle by sticking. In both the initial condition on the number of $`d\mu `$ atoms, and in the cycling efficiency after $`dd`$ fusion, $`q_{1s}`$ represents the probability for a $`d\mu `$ atom formed in an excited state to reach the ground state Bystritsky et al. (1990a). Finally, $`W_d`$, represents the probability that the muon will be captured by a deuterium atom given that there are both $`\mathrm{D}_2`$ and $`{}_{}{}^{3}\mathrm{He}`$ in the mixture:
$$W_d=\frac{\mathrm{c}_d}{\mathrm{c}_d+A\mathrm{c}_{{}_{}{}^{3}\mathrm{He}}}=\frac{X_{\mathrm{D}_2}}{X_{\mathrm{D}_2}+A^{}X_{{}_{}{}^{3}\mathrm{He}}}$$
(5)
where $`\mathrm{c}_d`$ and $`\mathrm{c}_{{}_{}{}^{3}\mathrm{He}}`$ are the deuterium and helium atomic concentration. $`A`$ is the relative muon atomic capture probability by a $`{}_{}{}^{3}\mathrm{He}`$ atom compared to deuterium atom, and $`A^{}`$ is the same ratio measured with respect to gas fraction concentrations ($`X`$). An previous experimental measure exists for $`\mathrm{D}_2+{}_{}{}^{3}\mathrm{He}`$ $`(A=1.7\pm 0.2)`$ Bystritsky (1993); Balin et al. (1992); Bystritsky et al. (1993b), and theoretical calculations for $`A^{}`$ have been made by J. S. Cohen Cohen (1999): for $`\mathrm{D}_2+{}_{}{}^{3}\mathrm{He}:A^{}=0.78`$ and for $`\mathrm{HD}+{}_{}{}^{3}\mathrm{He}:A^{}=0.68`$. Our gas mixtures have $`\mathrm{c}_{{}_{}{}^{3}\mathrm{He}}=0.0496(10)`$ and thus $`X_{{}_{}{}^{3}\mathrm{He}}=0.0946(20)`$. By atomic concentration, and using the experimental value, we get $`W_d=0.92(2)`$. Using theory and the gas fraction the result is the same, $`W_d=0.92`$. Using our own experiment Bystritsky et al. (1994), $`A=1.67_{0.33}^{+0.35}`$, to determine $`W_d`$ leads also to the exact same value.
The differential equations governing the evolution of the $`J=1,0`$ spin states of the $`d\mu {}_{}{}^{3}\mathrm{He}`$ molecules are (see Fig. 1):
$`{\displaystyle \frac{dN_{d\mu {}_{}{}^{3}\mathrm{He}}^1}{dt}}`$ $`=`$ $`+\phi \mathrm{c}_{{}_{}{}^{3}\mathrm{He}}\lambda _{d{}_{}{}^{3}\mathrm{He}}N_{d\mu }\lambda _\mathrm{\Sigma }^1N_{d\mu {}_{}{}^{3}\mathrm{He}}^1`$ (6)
$`{\displaystyle \frac{dN_{d\mu {}_{}{}^{3}\mathrm{He}}^0}{dt}}`$ $`=`$ $`+\stackrel{~}{\lambda }_{10}N_{d\mu {}_{}{}^{3}\mathrm{He}}^1\lambda _\mathrm{\Sigma }^0N_{d\mu {}_{}{}^{3}\mathrm{He}}^0`$ (7)
where $`N_{d\mu }`$ is the number of $`d\mu `$ atoms and with the definition
$`\lambda _\mathrm{\Sigma }^1`$ $`=`$ $`\left(\lambda _0+\lambda _p^{J=1}+\lambda _\gamma ^{J=1}+\lambda _e^{J=1}+\lambda _f^{J=1}\right)`$ (8)
$`\lambda _\mathrm{\Sigma }^0`$ $`=`$ $`\left(\lambda _0+\lambda _p^{J=0}+\lambda _\gamma ^{J=0}+\lambda _e^{J=0}+\lambda _f^{J=0}\right),`$ (9)
and
$`\lambda _{d\mu }`$ $`=`$ $`\lambda _0+\phi \mathrm{c}_{{}_{}{}^{3}\mathrm{He}}\lambda _{d{}_{}{}^{3}\mathrm{He}}`$ (10)
$`+`$ $`\phi \mathrm{c}_d\stackrel{~}{\lambda }_F\left[1W_dq_{1s}(1\beta _F\omega _d)\right].`$
The yield for protons between two given times after the muon arrival, $`t_1`$ and $`t_2`$, is:
$`Y_p(t_1,t_2)`$ $`=`$ $`Y_p^1(t_1,t_2)+Y_p^0(t_1,t_2)`$ (11)
$`=`$ $`N_\mu ^{\mathrm{D}/\mathrm{He}}{\displaystyle \frac{\stackrel{~}{\lambda }_f}{\lambda _\mathrm{\Sigma }}}{\displaystyle \frac{\phi \mathrm{c}_{{}_{}{}^{3}\mathrm{He}}\lambda _{d{}_{}{}^{3}\mathrm{He}}W_dq_{1s}\epsilon _Y\epsilon _p}{\lambda _{d\mu }}},`$
where the difference in time exponents has been defining as the yield efficiency:
$$\epsilon _Y=\left(e^{\lambda _{d\mu }t_1}e^{\lambda _{d\mu }t_2}\right).$$
(12)
and with the effective fusion rate defined as
$`\stackrel{~}{\lambda }_f`$ $`=`$ $`\left(\lambda _f^{J=1}{\displaystyle \frac{\lambda _\mathrm{\Sigma }^0}{\stackrel{~}{\lambda }_{10}+\lambda _\mathrm{\Sigma }^0}}+\lambda _f^{J=0}{\displaystyle \frac{\stackrel{~}{\lambda }_{10}}{\stackrel{~}{\lambda }_{10}+\lambda _\mathrm{\Sigma }^0}}\right)`$ (13)
$`\lambda _\mathrm{\Sigma }`$ $`=`$ $`\lambda _\mathrm{\Sigma }^0\left({\displaystyle \frac{\stackrel{~}{\lambda }_{10}+\lambda _\mathrm{\Sigma }^1}{\stackrel{~}{\lambda }_{10}+\lambda _\mathrm{\Sigma }^0}}\right).`$ (14)
In the above equations, $`N_\mu ^{\mathrm{D}/\mathrm{He}}`$ is the number of muons stopped in the $`\mathrm{D}_2+{}_{}{}^{3}\mathrm{He}`$ mixture and $`\phi `$ is the mixture atomic density relative to the liquid hydrogen density (LHD, $`N_0=4.25\times 10^{22}\text{cm}^3`$).
When protons are detected in coincidences with muon decay electrons, later on called the del-$`e`$ criterion, the fusion rate from Eq. (11) takes the form
$$\stackrel{~}{\lambda }_f=\frac{Y_p(t_1,t_2)\lambda _{d\mu }\lambda _\mathrm{\Sigma }}{N_\mu ^{\mathrm{D}/\mathrm{He}}W_dq_{1s}\phi \mathrm{c}_{{}_{}{}^{3}\mathrm{He}}\lambda _{d{}_{}{}^{3}\mathrm{He}}\epsilon _p\epsilon _e\epsilon _t\epsilon _Y},$$
(15)
where $`\epsilon _e`$ is the detection efficiency for muon decay electrons and $`\epsilon _t`$ defined as
$$\epsilon _t=e^{\lambda _0t_{ini}}e^{\lambda _0t_{fin}}$$
(16)
is the time efficiency depending on the interval during which we accept the muon decay electrons. Note that Eqs. (1115) are valid when the proton detection times are $`t1/\lambda _\mathrm{\Sigma }`$. The values $`\epsilon _p`$ and $`\lambda _\mathrm{\Sigma }`$ are found through calculation. Note an important feature of this experimental setup: $`\stackrel{~}{\lambda }_f`$ is found by using the experimental values of $`\lambda _{d\mu }`$, $`\epsilon _e`$, $`W_d`$ , $`\lambda _{d{}_{}{}^{3}\mathrm{He}}`$, and $`q_{1s}`$.
The information on these quantities corresponds to the conditions of a particular experiment and is extracted by the analysis of yields and time distributions of the 6.85 keV $`\gamma `$ rays from reaction (I), prompt and delayed x rays of $`\mu {}_{}{}^{3}\mathrm{He}`$ atoms in the $`\mathrm{D}_2+{}_{}{}^{3}\mathrm{He}`$ mixture and muon decay electrons. The quantity $`\lambda _{d{}_{}{}^{3}\mathrm{He}}`$ is determined from Eq. (10) where $`\beta _F=0.58`$, $`\omega _d=0.122(3)`$ are taken from Refs. Balin et al. (1984). $`\stackrel{~}{\lambda }_F=0.05\times 10^6\text{s}^1`$ is taken from Ref. Petitjean et al. (1996). The rate $`\lambda _{d\mu }`$ is the slope of the time distribution of $`\gamma `$ ray from reaction (I).
The procedure of measuring $`q_{1s}`$, $`\lambda _{d{}_{}{}^{3}\mathrm{He}}`$, $`W_d`$ , $`\epsilon _e`$, $`A`$ and $`\lambda _\gamma `$ (the partial probability for the radiative $`d\mu {}_{}{}^{3}\mathrm{He}`$ complex decay channel) as well as our results are described in detail in our previous work Bystritsky et al. (2004).
## III Experimental setup
The experimental layout (see Fig. 2) was described in details in Refs. Boreiko et al. (1998); Bystritsky et al. (2004, 2003). The experimental facility was located at the $`\mu `$E4 beam line of the PSI meson factory (Switzerland) with the muon beam intensity around $`2\times 10^4\text{s}^1`$. After passing through a thin plastic entrance monitoring counter muons hit the target and stopped there initiating a sequence of processes shown in Fig. 1. The electronics are protected from muon pileup within a $`\pm 10\mu `$s time gate so pileup causes a 30% reduction in the effective muon beam. Thus, we have a number of โgood muonsโ, called $`N_\mu `$, stopping in our target.
Three pairs of Si($`dEE`$) telescopes were installed directly behind 135 $`\mu `$m thick kapton windows and a 0.17 cm<sup>3</sup> germanium detector behind a 55 $`\mu `$m thick kapton window to detect the 14.64 MeV protons from reaction (2aa) and the 6.85 keV $`\gamma `$ rays from reaction (I), respectively. The Si telescopes with a 42 mm diameter were made of a 4 mm thick Si($`E`$) detector and a thin, 360 $`\mu `$m thick, Si($`dE`$) detector, respectively. An assembly of Si detectors like that give a good identification of protons, deuterons, and electrons based on different energy losses of the above particles in those detectors. Muon decay electrons were detected by four pairs of scintillators, E<sub>UP</sub>, E<sub>DO</sub>, E<sub>RI</sub> and E<sub>LE</sub>, placed around the vacuum housing of the target. The total solid angle of the electron detectors was $`17`$%. The cryogenic target was located inside the vacuum housing. The design of the target is described in detail in Refs. Stolupin et al. (1999); Boreiko et al. (1998).
The analysis of the 6.85 keV $`\gamma `$โray time distributions allows us to determine the disappearance rate, $`\lambda _{d\mu }`$, for the $`d\mu `$ atoms in the $`\mathrm{D}_2+{}_{}{}^{3}\mathrm{He}`$ mixture. Note that the presence of a signal from the electron detectors during a certain time interval (the del-$`e`$ criterion) whose beginning corresponds to the instant of time when the $`K\alpha `$, $`K\beta `$, and $`K\gamma `$ lines of $`\mu \mathrm{He}`$ atoms is detected makes it possible to determine uniquely the detection efficiency for muon decay electrons. When the del-$`e`$ criterion is used in the analysis of events detected by the Si($`dEE`$) telescopes one obtains a suppression factor of $`300400`$ of the background, which is quite enough to meet the requirements of the experiment on the study of nuclear fusion in the $`d\mu {}_{}{}^{3}\mathrm{He}`$ complex.
Our experiment included two runs with the $`\mathrm{D}_2+{}_{}{}^{3}\mathrm{He}`$ mixture. The experimental conditions are listed in Table 2. In addition, we performed different measurements with pure $`\mathrm{D}_2`$, $`{}_{}{}^{3}\mathrm{He}`$, and $`{}_{}{}^{4}\mathrm{He}`$ at different pressures and temperature. Details are given in Refs. Bystritsky et al. (2003).
The germanium detector was calibrated using <sup>55</sup>Fe and <sup>57</sup>Co sources. The Si($`dEE`$) detectors were calibrated using a radioactive <sup>222</sup>Rn source. Before the cryogenic target was assembled, a surface saturation of the Si($`dE`$) and Si($`E`$) detectors by radon was carried out. The <sup>222</sup>Rn decay with the emission of alpha-particles of energies 5.3, 5.5, 6.0, and 7.7 MeV were directly detected by each of the Si detectors. The linearity of the spectrometric channels of the Si detectors in the region of detection of protons with energies $`815`$ MeV was checked using exact-amplitude pulse generators.
## IV Analysis of the experimental data
### IV.1 Determination of the $`d\mu {}_{}{}^{3}\mathrm{He}`$complex formation rate
By way of example Fig. 3 shows energy spectra of events detected by the germanium detector in run I without and with the del-$`e`$ criterion. The rather wide left peak corresponds to the $`\gamma `$ rays with an average energy of 6.85 keV and the three right peaks correspond to the $`K\alpha `$, $`K\beta `$, $`K\gamma `$ lines of $`\mu \mathrm{He}`$ atoms with energies 8.17, 9.68, and 10.2 keV, respectively. As seen in Fig. 3, the suppression factor for the background detected by the germanium detector with the del-$`e`$ criterion is of the order of $`10^3`$.
Figure 4 shows time distributions of 6.85 keV $`\gamma `$ rays resulting from radiative de-excitation of the $`d\mu {}_{}{}^{3}\mathrm{He}`$ complex in runs I and II. The distributions were measured in coincidences with delayed muon decay electrons. The experimental time distributions of $`\gamma `$ rays shown in Fig. 4 were approximated by the following expression
$$\frac{dN_\gamma }{t}=B^\gamma e^{\lambda _{d\mu }t}+C^\gamma e^{\lambda _0t}+D^\gamma ,$$
(17)
where $`B^\gamma `$, $`C^\gamma `$, and $`D^\gamma `$ are the normalization constants. The second and third terms in Eq. (17) describe the contribution from the background. The analysis of the time distributions of the 6.85 keV $`\gamma `$ rays yielded values of $`\lambda _{d\mu }`$ and thus the formation rates $`\lambda _{d{}_{}{}^{3}\mathrm{He}}`$. Results are given in Table 3.
The systematic error is larger than the uncertainty of the result caused by various possible model of the background, including the case where it is equal to zero (e.g., when time structure of the background is inaccurately known). We describe the procedure of determining $`\lambda _{d{}_{}{}^{3}\mathrm{He}}`$ in more detail in Ref. Czapliลski et al. (1996b).
### IV.2 Number of muon stops in the $`\mathrm{D}_2+{}_{}{}^{3}\mathrm{He}`$ mixture
The number of muon stops in the $`\mathrm{D}_2+{}_{}{}^{3}\mathrm{He}`$ mixture was determined by analyzing time distributions of events detected by the four electron counters, We detailed this matter in Refs. Czapliลski et al. (1996a); Boreiko et al. (1998); Bystritsky et al. (2003). Here it is pertinent to dwell upon some particular points in determination of this value.
By way of example Fig. 5 shows the time distribution of muon decay electrons measured in run I. To determine the number of muon stops in the $`d\mu {}_{}{}^{3}\mathrm{He}`$ the time distribution of the detected electrons, $`dN_e/dt`$, is approximated by an expression which is superposition of four exponents and a background of accidental coincidences
$$\frac{dN_e}{dt}=A_{\mathrm{Al}}^ee^{\lambda _{\mathrm{Al}}t}+A_{\mathrm{Au}}^ee^{\lambda _{\mathrm{Au}}t}+A_{\mathrm{He}}^ee^{\lambda _{\mathrm{He}}t}+A_\mathrm{D}^ee^{\lambda _0t}+B^e,$$
(18)
where $`A_{\mathrm{Al}}^e`$, $`A_{\mathrm{Au}}^e`$, $`A_{\mathrm{He}}^e`$ and $`A_\mathrm{D}^e`$, are the normalized amplitudes with
$$A_i^e=N_\mu ^iQ_i\lambda _0\epsilon _ei=\mathrm{Al},\mathrm{Au},\mathrm{He},\mathrm{D},$$
(19)
and
$`\lambda _{\mathrm{Al}}`$ $`=`$ $`Q_{\mathrm{Al}}\lambda _0+\lambda _{\mathrm{cap}}^{\mathrm{Al}},`$
$`\lambda _{\mathrm{Au}}`$ $`=`$ $`Q_{\mathrm{Au}}\lambda _0+\lambda _{\mathrm{cap}}^{\mathrm{Au}},`$ (20)
$`\lambda _{\mathrm{He}}`$ $`=`$ $`\lambda _0+\lambda _{\mathrm{cap}}^{\mathrm{He}},`$
are the muon disappearance rates in the different elements (the rates are the inverse of the muon lifetimes in the target wall materials). In reality, Eq. (18) is an approximation of a more complex equation, which can be found in Ref. Knowles (1999). The different rates are $`\lambda _0=0.455\times 10^6\text{s}^1`$ and $`\lambda _{cap}^{\mathrm{He}}=2216(70)\text{s}^1`$ Maev et al. (1996). The nuclear capture rates in aluminum and gold, $`\lambda _{\mathrm{cap}}^{\mathrm{Al}}=0.7054(13)\times 10^6\text{s}^1`$ and $`\lambda _{\mathrm{cap}}^{\mathrm{Au}}=13.07(28)\times 10^6\text{s}^1`$, are taken from Ref. Suzuki et al. (1987). $`Q_{\mathrm{Al}}`$ and $`Q_{\mathrm{Au}}`$ are the Huff factors, which take into account that muons are bound in the $`1s`$ state of the respective nuclei when they decay. This factor is negligible for helium but necessary for aluminum $`Q_{\mathrm{Al}}=0.993`$ and important for gold $`Q_{\mathrm{Au}}=0.850`$ Suzuki et al. (1987). The constant $`B^e`$ characterizes the random coincidence background.
We denote $`N_\mu `$ as the total number of muon stops in the target, $`N_\mu ^{\mathrm{Al}}`$, $`N_\mu ^{\mathrm{Au}}`$, and $`N_\mu ^{\mathrm{D}/\mathrm{He}}`$ as the numbers of muon stops in Al, Au, and the gaseous $`\mathrm{D}_2+{}_{}{}^{3}\mathrm{He}`$ mixture, respectively. Thus, we have the relation
$$N_\mu =N_\mu ^{\mathrm{Al}}+N_\mu ^{\mathrm{Au}}+N_\mu ^{\mathrm{D}/\mathrm{He}}.$$
(22)
Since the muon decay with emission electrons in the $`\mathrm{D}_2+{}_{}{}^{3}\mathrm{He}`$ mixture take place from the $`1s`$ state of the $`d\mu `$ or $`\mu {}_{}{}^{3}\mathrm{He}`$ atom, the third and fourth terms in Eq(18) will differ only by the values of the amplitudes $`A_{\mathrm{He}}^e`$ and $`A_\mathrm{D}^e`$ because the slopes of both exponents are practically identical ($`\lambda _{\mathrm{He}}=0.457\mu \text{s}^1`$, $`\lambda _0=0.455\mu \mathrm{s}^1`$). In this connection the following simplified expression was used to approximate experimental time distributions of
$$\frac{dN_e}{dt}=A_{\mathrm{Al}}^ee^{\lambda _{\mathrm{Al}}t}+A_{\mathrm{Au}}^ee^{\lambda _{\mathrm{Au}}t}+A_{\mathrm{D}/\mathrm{He}}^ee^{\stackrel{~}{\lambda }_{\mathrm{D}/\mathrm{He}}t}+B^e,$$
(23)
Under our experimental conditions of runs I and II, we obtained the effective rates $`\stackrel{~}{\lambda }_{D/He}=0.4563\mu \text{s}^1`$ and 0.4567 $`\mu \text{s}^1`$, respectively. With these effective muon decay rates, the uncertainty in the calculated number of muon stops in the gaseous $`\mathrm{D}_2+{}_{}{}^{3}\mathrm{He}`$ mixture is negligibly small as compared with the more rigorous calculation of this value by Eq. (18).
The amplitudes in Eq. (19) are expressed in terms of the factors $`a_{Al}`$, $`a_{Au}`$, and $`a_{D/He}`$, defined as the partial muon stopping in Al, Au, and $`\mathrm{D}_2+{}_{}{}^{3}\mathrm{He}`$ mixture,
$$a_i=\frac{N_\mu ^i}{N_\mu },\underset{i}{}a_i=1i=\mathrm{Al},\mathrm{Au},\mathrm{D}/\mathrm{He};$$
(24)
take the new form
$$A_i^e=N_\mu \lambda _0Q_i\epsilon _ea_i.$$
(25)
The electron detection efficiency, $`\epsilon _e`$, of the detectors E<sub>UP</sub>, E<sub>DO</sub>, E<sub>RI</sub> and E<sub>LE</sub> was determined as a ratio between the number of events belonging to the $`K`$โlines of the $`\mu {}_{}{}^{3}\mathrm{He}`$ atoms, found from the analysis of the data with and without the del-$`e`$ criterion,
$$\epsilon _e=\frac{N_{xe}}{N_x},$$
(26)
where $`N_{xe}`$ and $`N_x`$ are the numbers of events belonging to $`K`$โlines of the $`\mu {}_{}{}^{3}\mathrm{He}`$ atom and detected by the germanium detector with and without coincidence with the electron detectors. The thus measured experimental value is electron detection efficiency averaged over the target volume. Table 4 presents the results.
The electron detection efficiency of the detector E<sub>LE</sub> is considerably lower than that of each of the other three electron detectors. This is because the material (Al, Fe) layer which the muon decay electron has to pass through in the direction of the detector E<sub>LE</sub> is thicker than material layers in the direction of the other electron detectors.
Table 5 lists the values of the fraction of muons stopped in the $`\mathrm{D}_2+{}_{}{}^{3}\mathrm{He}`$ mixture, $`a_{D/He}`$, found from the analysis of the time distributions of the events detected by the four electron detectors in runs I and II. Note that when the $`a_{D/He}`$ fraction, was calculated by Eqs. (24) and (25) it was assumed that the electron detection efficiency by each of the detectors E<sub>UP</sub>, E<sub>DO</sub>, E<sub>RI</sub> and E<sub>LE</sub> did not depend on the coordinates of the muon stop point in the target (be it in the target walls or in the $`\mathrm{D}_2+{}_{}{}^{3}\mathrm{He}`$ mixture).
The systematic errors were determined as one half of the maximum spread between the $`a_{D/He}`$ values found from analysis of the time distributions of the electrons detected by each of the electron detectors E<sub>UP</sub>, E<sub>DO</sub>, E<sub>RI</sub> and E<sub>LE</sub>. Note that the fraction of muons stopped in gas, $`a_{D/He}`$, is a result of simultaneously fitting all time distributions obtained with each of the electron detectors (and not a result of averaging all four distributions corresponding to each of the four detectors).
### IV.3 Determination of the detection efficiency for 14.64 MeV protons
To determine the proton detection efficiency, $`\epsilon _p`$, of the three Si($`dEE`$) telescopes, one should know the distribution of muon stops over the target volume in runs I and II. The average muon beam momentum $`\overline{\mathrm{P}}_\mu `$ corresponding to the maximum fraction $`a_{D/He}`$ of muons stopped in the $`\mathrm{D}_2+{}_{}{}^{3}\mathrm{He}`$ mixture in runs I and II was found by varying the muon beam momentum $`\overline{\mathrm{P}}_\mu `$ and analyzing the time distributions of the detected electrons by Eq. (23). Next, knowing the average momentum $`\overline{\mathrm{P}}_\mu `$ and the beam momentum spread, we simulated the real distribution of muon stops in runs I and II by the Monte Carlo (MC) method Jacot-Guillarmod (1997). The results of the simulation were used in another Monte Carlo program to calculate the detection efficiency of each pair of Si($`dEE`$) detectors for protons from reaction (2aa) Woลบniak et al. (1996). The algorithm of the calculation program included simulation of the muon stop points in the $`\mathrm{D}_2+{}_{}{}^{3}\mathrm{He}`$ mixture and the $`d\mu `$ and $`\mu {}_{}{}^{3}\mathrm{He}`$ atom formation points, the consideration of the entire chain of processes occurring in the $`\mathrm{D}_2+{}_{}{}^{3}\mathrm{He}`$ mixture from the instant when the muon hits the target to the instant of possible production of 14.64 MeV protons in the fusion reaction in the $`d\mu {}_{}{}^{3}\mathrm{He}`$ complex. The calculation program took into account the proton energy loss in the gas target, kapton windows and Si($`dEE`$) detectors themselves (in the thin Si($`dE`$) and thick Si($`E`$) detectors). The proton detection efficiency $`\epsilon _p`$ was calculated at the $`q_{1s}`$, $`W_d`$, and $`\lambda _{d\mu }`$ values (see Table 3) corresponding to our experimental conditions. The scattering cross sections of $`d\mu `$ atoms form $`\mathrm{D}_2`$ molecules were taken from Refs. Chiccoli et al. (1992); Adamczak et al. (1996); Melezhik and Woลบniak (1992).
We ceased tracing the muon stopped in the target when
* the muon decays ($`\mu ^{}e^{}+\nu _\mu +\stackrel{~}{\nu }_e`$)
* the muon is transferred from the deuteron to the $`{}_{}{}^{3}\mathrm{He}`$ nucleus with the formation of a $`{}_{}{}^{3}\mathrm{He}`$$`\mu `$ atom
* nuclear fusion occurs in the $`d\mu {}_{}{}^{3}\mathrm{He}`$ complex
* a $`dd\mu p+t+\mu `$ reaction proceeds in the $`dd\mu `$ molecule.
Note that the algorithm of the program also involved the consideration of the background process resulting from successive occurrence of the reactions
$`d\mu +ddd\mu `$ $`{}_{}{}^{3}\mathrm{He}`$ $`(0.8\mathrm{MeV})+n`$ (27)
$`+`$
$`d`$ $`\alpha +p(14.64\mathrm{MeV}).`$
This reaction (27) is called $`d{}_{}{}^{3}\mathrm{He}`$ โfusion in flightโ.
In our calculations we used the dependence of the cross section for reaction (27) on the $`{}_{}{}^{3}\mathrm{He}`$ deuteron collision energy, averaged over the data of Refs. White et al. (1997); Kunz (1955); Kljuchaiev et al. (1956); Allred et al. (1952); Argo et al. (1952); Freier and Holmgren (1954). Figure 6 displays the cross section dependence on the $`{}_{}{}^{3}\mathrm{He}`$ deuteron collision energy. The program also took into account the energy loss of $`{}_{}{}^{3}\mathrm{He}`$ nuclei in the $`\mathrm{D}_2+{}_{}{}^{3}\mathrm{He}`$ mixture caused by ionization of $`{}_{}{}^{3}\mathrm{He}`$ atoms and deuterium molecules. The time distributions of protons from reactions (2aa) and (27) under the same experimental conditions have completely different shapes in accordance with the kinetics of processes in the $`\mathrm{D}_2+{}_{}{}^{3}\mathrm{He}`$ mixture.
Figures. 7 and 8 show the calculated time dependencies of the expected yields of protons from reactions (2aa) and (27) under the conditions of runs I and II. Thus, there arises a possibility of selecting a time interval of detection of events by the Si($`dEE`$) detectors where the ratio of the reaction (2aa) and (27) yields is the largest. This, in turn, makes it possible to suppress the detected background from reaction (27) to a level low enough to meet the requirement of the experiment on the study of nuclear fusion in the $`d\mu {}_{}{}^{3}\mathrm{He}`$ complex. Table 6 presents the calculated values of some quantities describing kinetics of muonic processes in the $`\mathrm{D}_2+{}_{}{}^{3}\mathrm{He}`$ mixture and the process of detecting protons from reactions (2aa) and (27).
$`W_{{}_{}{}^{3}\mathrm{He}}`$ is the total probability for the $`{}_{}{}^{3}\mathrm{He}`$ formation ($`E_{{}_{}{}^{3}\mathrm{He}}=0.8`$ MeV) in the $`\mathrm{D}_2+{}_{}{}^{3}\mathrm{He}`$ mixture, as a result of the fusion reaction in the $`dd\mu `$ molecule. $`W_{d\mu {}_{}{}^{3}\mathrm{He}}`$ is the $`d\mu {}_{}{}^{3}\mathrm{He}`$ complex formation probability and $`W_{d{}_{}{}^{3}\mathrm{He}}`$ is the probability for $`d{}_{}{}^{3}\mathrm{He}`$ fusion in flight, following reaction (27), and $`W_{\mu e}`$ is the branching ratio of the muon decay via the $`\mu ^{}e^{}+\nu _\mu +\overline{\nu }_e`$ channel. $`\epsilon _p`$ and $`\epsilon _p^{ff}`$ are the detection efficiencies of one Si($`dEE`$) telescope for protons from reactions (2aa) and (27), respectively. $`\eta _p`$ and $`\eta _p^{ff}`$ are the yields of protons from reactions (2aa) and (27) detected by the Si($`dEE`$) telescope per muon stop in the gaseous $`\mathrm{D}_2+{}_{}{}^{3}\mathrm{He}`$ mixture (the value of $`d\mu {}_{}{}^{3}\mathrm{He}`$ fusion rate $`\lambda _f=10^6\text{s}^1`$ was used for calculation of $`\eta _p`$). There are some noteworthy intermediate results in the calculation of the detection efficiencies for protons from reactions (2aa) and (27). Table 7 presents average energy losses of protons on their passage through various material in the direction of the Si($`dEE`$) detectors.
Figures 9 show the two-dimensional distributions of events detected by the Si($`dEE`$) detectors without coincidences with muon decay electrons in runs I and II. The xโaxis represents the energy losses in the thin Si($`dE`$) counters and the yโaxis shows the total energies losses by the particle in both the Si($`dE`$) and Si($`E`$) detectors connected in coincidence. The distributions of events in Figs. 9 correspond to the detection of protons arising both from reactions (2aa) and (27) and from the background reactions such as
$`\mu +{}_{}{}^{3}\mathrm{He}`$ $``$ $`p+2n+\nu _\mu `$
$`\mu +\mathrm{Al}`$ $`\mathrm{Na}^{}`$ $`+p+n+\nu _\mu `$ (28)
$`+p+\nu _\mu `$
$`+p+2n+\nu _\mu `$
$`\mu +\mathrm{Fe}`$ $`\mathrm{Cr}^{}`$ $`+p+n+\nu _\mu `$
$`+p+\nu _\mu `$
$`+p+2n+\nu _\mu .`$
In addition, the background which is not correlated with muon stops in the target (background of accidental coincidences) contributes to these distributions.
Figures 10 show the two-dimensional Si($`dE(E+dE)`$) distributions obtained in coincidences with muon decay electrons. As seen, the use of the del-$`e`$ criterion leads to an appreciably reduction of the background, which in turn makes it possible to identify a rather weak effect against the intensive background signal. To suppress muon decay electrons in the Si($`dEE`$) telescope, provision was made in the electronic logic of the experiment to connect each of the electron detectors in anti-coincidence with the corresponding Si($`dEE`$) telescope. The choice of optimum criteria in the analysis of the data from the Si($`dEE`$) telescopes was reduced to the determination of the boundaries and widths of the time and energy intervals where the background is substantially suppressed in absolute value and the effect-to-background ratio is the best. To this end the two-dimensional Si($`dE(dE+E)`$) distributions corresponding to the detection of protons were simulated by MC method for runs I and II. On the basis of these distributions boundaries were determined for the energy interval of protons from reaction (2aa) where the loss of the โusefulโ event statistics collected by the Si telescope would be insignificant.
Figures 11 and 12 show the two-dimensional Si($`dE(dE+E)`$) distributions corresponding to the proton detection which were simulated by the MC method for runs I and II. Based on these distributions, we chose some particular proton energy intervals named $`\mathrm{\Delta }E_\mathrm{\Sigma }`$ when considering the total energy deposited and $`\delta E`$ when looking only at the Si($`dE`$) detector (see Table 8) for further analysis. The regions of events corresponding to the intervals $`\delta E`$ and $`\mathrm{\Delta }E_\mathrm{\Sigma }`$ are shown in the form of rectangles on the two- dimensional distributions presented on the Fig. 9 and 10.
It is noteworthy that the proton detection efficiencies given in Table 6 correspond to these chosen proton energy intervals for runs I and II. The next step in the data analysis was to choose a particular times interval of detection of events by the Si($`dEE`$) telescope. Figures 11 and 12 show the simulated time distributions of protons corresponding to the chosen energy loss intervals $`\delta E`$, for the energy loss in the Si($`dE`$) detector and $`\mathrm{\Delta }E_\mathrm{\Sigma }=E+\delta E`$ the energy loss in both silicon detector. For the chosen proton energy intervals Table 8 presents the statistics suppression factors corresponding to different initial time,$`t_{thr}`$ (with respect to the instant of the muon stop in the target) of the time intervals of detection of proton events. These factors correspond to the $`\epsilon _Y`$ value in Eq. (11). The data in Table 8 are derived from time dependencies of the yields of protons from reactions (2aa) and (27) (see Figs. 7 and 8).
According to the data given in Table 8, we took the following time intervals $`\mathrm{\Delta }t_{\mathrm{Si}}`$ (with $`t_{\mathrm{Si}}`$ the time for the Si signal to appear) for analyzing the events
$`\mathrm{\Delta }t_{\mathrm{Si}}\text{(run I)}:0.7`$ $`t_{\mathrm{Si}}`$ $`2.2\mu \text{s}`$
$`\mathrm{\Delta }t_{\mathrm{Si}}\text{(run II)}:0.4`$ $`t_{\mathrm{Si}}`$ $`1.2\mu \text{s}.`$ (29)
Figures 13 display the two-dimensional distributions of Si($`dEE`$) events obtained in coincidences with muon decay electrons in runs I and II with this time criteria imposed. With these time intervals $`\mathrm{\Delta }t_{\mathrm{Si}}`$ and the proton energy loss $`\mathrm{\Delta }E_\mathrm{\Sigma }`$ and $`\delta E`$ intervals, the statistics collection suppression factors for events from reactions (2aa) and (27) are
$`k_{d\mu {}_{}{}^{3}\mathrm{He}}=2.9,`$ $`k_{d{}_{}{}^{3}\mathrm{He}}`$ $`=11.2,\text{Run I,}`$
$`k_{d\mu {}_{}{}^{3}\mathrm{He}}=3.2,`$ $`k_{d{}_{}{}^{3}\mathrm{He}}`$ $`=12.1,\text{Run II.}`$
Another stage of the data analysis was the determination of the number of events detected by the Si($`dEE`$) telescopes in runs I and II under the following criteria
* the coincidence of signals from the Si telescopes and electron detectors in the time interval $`0.2<(t_et_{\mathrm{Si}})<5.5\mu `$s ($`t_e`$ is the time when the E detector signal appear). Such a requirement add the efficiency factor $`\epsilon _t=0.83`$ when determining the rates.
* the total energy release in the Si($`dE`$) detector is $`\delta E`$ as given in Table 8. This particular $`\delta E`$ interval will be called $`\delta E`$. For the thin and thick Si detector together, we choose the smallest interval, namely $`\mathrm{\Delta }E_\mathrm{\Sigma }=[11.714.2]`$ MeV for run I and $`\mathrm{\Delta }E_\mathrm{\Sigma }=[8.013.4]`$ MeV for run II.
* the time when the signal from the Si telescope appears falls in the $`\mathrm{\Delta }t_{\mathrm{Si}}`$ intervals.
Table 9 presents the numbers of events $`N_p`$ detected in runs I and II under the above mentioned criteria.
The contribution of the background events, $`N_p^{ff}`$, given in Table 9 from the reaction (27) is found in the following way. The expected number of detected protons from reaction (27) in runs I and II is calculated by
$$N_p^{ff}=\frac{N_\mu a_{D/He}W_{{}_{}{}^{3}\mathrm{He}}W_{d{}_{}{}^{3}\mathrm{He}}\epsilon _p^{ff}N_{\mathrm{Si}}\epsilon _e\epsilon _t}{k_{d{}_{}{}^{3}\mathrm{He}}}.$$
(31)
$`N_{\mathrm{Si}}`$ is the number of Si($`dEE`$) telescopes and $`1/k_{d{}_{}{}^{3}\mathrm{He}}`$ is the factor of background suppression by imposing the criteria (ii) and (iii). Using the values of $`a_{D/He}`$ and $`N_\mu `$ measured in runs I and II, the calculated values of $`W_{{}_{}{}^{3}He}`$, $`W_{d^3He}`$, $`\epsilon _p^f`$, $`N_{\mathrm{Si}}`$, $`k_{d{}_{}{}^{3}\mathrm{He}}`$, $`\epsilon _t`$, and Eq. (31), we obtained $`N_p^{ff}`$, which is given in Table 9. Errors of the calculated $`N_p^{ff}`$ arose from the inaccurate dependence of the cross sections $`\sigma _{d{}_{}{}^{3}\mathrm{He}}`$ for the $`d{}_{}{}^{3}\mathrm{He}`$ reaction in flight on the $`{}_{}{}^{3}\mathrm{He}`$ deuteron collision energy and from the errors in the calculations of the detection efficiency of the Si telescopes for protons from reaction (27). These errors were found by substituting various experimental $`\sigma _{d^3He}(E_{d^3He})`$ dependencies White et al. (1997); Kunz (1955); Kljuchaiev et al. (1956); Allred et al. (1952); Argo et al. (1952); Freier and Holmgren (1954). into the program for Monte Carlo calculation of the in-flight $`d{}_{}{}^{3}\mathrm{He}`$ fusion probability $`W_{d{}_{}{}^{3}\mathrm{He}}`$.
Now it is necessary to find the level of the accidental coincidence background by analyzing the experimental data from runs I and II. To this end the two-dimensional distribution of events detected by the Si($`dEE`$) telescopes was divided into three regions which did not include the separated region of events belonging to the process (2aa). Considering the boundaries of the intervals $`\delta E`$ and $`\mathrm{\Delta }E_\mathrm{\Sigma }`$ of energy losses of the protons from reaction (2aa) we used three regions, A,B, and C, of the two-dimensional $`\delta E\mathrm{\Delta }E_\mathrm{\Sigma }`$ distributions for determining the background level. The regions are given in Table 10.
The level $`N_p^{acc}`$ of the background of the accidental coincidences of signals from the Si($`dEE`$) telescopes and the electron detectors for the given three region of the two-dimensional $`\delta E\mathrm{\Delta }E_\mathrm{\Sigma }`$ distributions and the corresponding suppression factor of the accidental background in the Si telescopes, $`\eta _{SiE}`$, are defined as
$$N_p^{acc}=N_{\mathrm{Si}}^f\eta _{SiE},$$
(32)
$$\eta _{SiE}=\frac{\underset{i}{}N_{SiE}^i}{\underset{i}{}N_{Si}^i},$$
(33)
where $`N_{\mathrm{Si}}^f`$ is the number of events detected by the three Si($`dEE`$) telescopes and belonging to the selected $`(\delta E\mathrm{\Delta }E_\mathrm{\Sigma })`$ region of detection of protons from reaction (2aa). $`N_{SiE}^i`$ and $`N_{Si}^i`$ are the numbers of events detected by the $`i`$th Si($`dEE`$) telescope with and without del-$`e`$ coincidences and belonging to the other $`(\delta E\mathrm{\Delta }E_\mathrm{\Sigma })`$ intervals. Note that the degree of suppression of the accidental coincidence background was determined not only by averaging the data obtained with the $`\mathrm{D}_2+{}_{}{}^{3}\mathrm{He}`$ mixture but also in additional experiments with the targets filled with pure $`{}_{}{}^{4}\mathrm{He}`$, $`\mathrm{D}_2`$, and $`{}_{}{}^{3}\mathrm{He}`$ whose densities were $`\phi 0.17`$, $`\phi 0.09`$, and $`\phi 0.035`$,respectively. This guaranteed an identical ratio of stops in the target walls and in the gas in the experiments with $`{}_{}{}^{4}\mathrm{He}`$, $`\mathrm{D}_2`$, and the $`\mathrm{D}_2+{}_{}{}^{3}\mathrm{He}`$ mixture ($`\phi =0.168`$) and in the experiments with the $`\mathrm{D}_2+{}_{}{}^{3}\mathrm{He}`$ mixture ($`\phi =0.0585`$) and $`{}_{}{}^{3}\mathrm{He}`$. Figures 1415, and 16 display the two-dimensional distributions of background events detected by the Si($`dEE`$) telescopes in the experiments with $`{}_{}{}^{4}\mathrm{He}`$, $`\mathrm{D}_2`$, and $`{}_{}{}^{3}\mathrm{He}`$.
The values of $`\eta _{SiE}`$ and $`N_p^{acc}`$ are given in Table 9 for runs I and II. The total numbers of detected background events,$`N_p^{bckg}`$, which belongs to the analyzed region of energies $`(\delta E\mathrm{\Delta }E_\mathrm{\Sigma })`$ of protons from reaction (2aa) and met the criteria (i)โ(iii) were defined as
$$N_p^{bckg}=N_p^{ff}+N_p^{acc}$$
(34)
and are also given in Table 9. The uncertainties of $`N_p^{bckg}`$ include both statistical and systematical errors.
Based on the measured values $`N_p`$ and the calculated values $`N_p^{bckg}`$ and following Refs. Helene (1984, 1983); Feldman and Cousins (1998), we found the yields of detected protons ,$`Y_p`$, from reaction (2aa) in runs I and II.
$`Y_p`$ $`=`$ $`7.7_{3.4}^{+4.4}\text{run I}`$
$`Y_p`$ $`=`$ $`7.5_{3.2}^{+3.8}\text{run II}`$ (35)
The errors of $`Y_p`$ are found in accordance with Refs. Helene (1984, 1983); Feldman and Cousins (1998) dealing with analysis of small statistical samples. In view of Eq. (11) and the measured values $`Y_p`$, the effective rate of nuclear fusion in the $`d\mu {}_{}{}^{3}\mathrm{He}`$ complex is obtained from Eq. (15). It can be written as
$$\stackrel{~}{\lambda }_f=\frac{\lambda _{d\mu }\lambda _\mathrm{\Sigma }}{N_\mu a_{D/He}W_dq_{1s}\phi \mathrm{c}_{{}_{}{}^{3}\mathrm{He}}\lambda _{d{}_{}{}^{3}\mathrm{He}}}\frac{Y_p}{\epsilon _p\epsilon _e\epsilon _t\epsilon _Y},$$
(36)
The values of $`\stackrel{~}{\lambda }_f`$ and $`\lambda _\mathrm{\Sigma }`$ corresponding to the conditions of runs I and II are given in Table 11. Using Eq. (13) and the measured effective rates of nuclear fusion and assuming that $`\lambda _f^1\lambda _f^0`$ Bogdanova et al. (1999), one can get hypothetical estimates of the partial fusion rate in the $`d\mu {}_{}{}^{3}\mathrm{He}`$ complex in its states with the total orbital momentum $`J=0`$
$$\lambda _f^{J=0}=\frac{\stackrel{~}{\lambda }_f(\stackrel{~}{\lambda }_{10}+\lambda _\mathrm{\Sigma }^0)}{\stackrel{~}{\lambda }_{10}}.$$
(37)
Table 11 also presents the values for $`\lambda _f^{J=0}`$ found in runs I and II.
The averages $`\lambda _\mathrm{\Sigma }^0=6\times 10^{11}\text{s}^1`$ and $`\lambda _\mathrm{\Sigma }^1=7\times 10^{11}\text{s}^1`$ (averaging over the data Kino and Kamimura (1993); Gershtein and Gusev (1993); Czapliลski et al. (1997b); Belyaev et al. (1995); Czapliลski et al. (1996a); Belyaev et al. (1997); Czapliลski et al. (1996b)) were used to get the values presented in Table 11. As to the effective rate for transition of the $`d\mu {}_{}{}^{3}\mathrm{He}`$ complex from the state with the angular momentum $`J=1`$ to the state with $`J=0`$, it was calculated with allowance for the entire complicated branched chain of processes accompanying and competing with the rotational $`10`$ transition (see Table 11). The chain of these processes is considered in detail in Refs. Bogdanova et al. (1998); Bystritsky and Penโkov (1999); Bystritsky et al. (1999b, c). The effective rates of nuclear fusion in the $`d\mu {}_{}{}^{3}\mathrm{He}`$ complex found by us in runs I and II coincide within the measurement errors. This is also true for the $`d{}_{}{}^{3}\mathrm{He}`$ fusion rates $`\lambda _f^{J=0}`$ obtained by Eq. (37). A comparison of the measured $`\lambda _f^{J=0}`$ with the theoretical calculations show rather good agreement with Czapliลski et al. (1996b), a slight discrepancy with Refs. Penโkov (1997); Bogdanova et al. (1999) and considerable disagreement with Refs. Nagamine et al. (1989); Harley et al. (1989). The cause of this disagreement is not clear yet as also is not clear the discrepancy between $`\lambda _f^{J=0}`$ calculations in Refs. Nagamine et al. (1989); Penโkov (1997); Czapliลski et al. (1996b); Harley et al. (1989) (see Table 1). Note that the theoretical papers Refs. Nagamine et al. (1989); Penโkov (1997); Czapliลski et al. (1996b); Harley et al. (1989); Bogdanova et al. (1999) yield estimates with a different degree of approximation. A correct comparison of the experimental and theoretical $`\lambda _f^{J=0}`$ is possible only after carrying out some experiments with the $`\mathrm{D}_2+{}_{}{}^{3}\mathrm{He}`$ mixture ruling out model dependence on the effective rate of transition of the $`d\mu {}_{}{}^{3}\mathrm{He}`$ complex from the $`J=1`$ state to the $`J=0`$ state.
A comparison of the results of this paper with the experimental results Maev et al. (1999) reveals appreciable disagreement between them. The shortened form of presentation of the results Maev et al. (1999) does not allow us to find out sufficiently well the cause of this considerable disagreement. Note, however, some results of the intermediate calculations which, to our mind, disagree with the real estimates of the calculated quantities.
1. According to Ref. Maev et al. (1999), the fraction of the $`d\mu `$ atoms which were formed in the excited state under their experimental conditions and came to the ground state (per muon stop in the target) is $`C_{d\mu }=0.8`$. The quantity $`C_{d\mu }`$ is defined as
$$C_{d\mu }=\frac{1}{2}W_{p,d}(q_{1s}^{p\mu }+q_{1s}^{d\mu }),$$
(38)
where $`W_{p,d}`$ is the probability for direct muon capture by the HD molecule followed by formation of the muonic hydrogen atom or the excited $`d\mu `$ atom. $`q_{1s}^{p\mu }`$ and $`q_{1s}^{d\mu }`$ are the probabilities for the transition of the $`p\mu `$ and $`d\mu `$ atoms from the excited state to the $`1s`$ ground state. According to Refs. Tresch et al. (1998b); Gartner et al. (2000); Bystritsky et al. (2004); Tresch et al. (1999), under the Maev et. al experimental conditions the values of the quantities appearing in Eq. (38) were $`W_{p,d}=0.92`$ $`q_{1s}^{p\mu }=0.5`$, and $`q_{1s}^{d\mu }=0.8`$. Thus, as follows from our estimation, $`C_{d\mu }=0.6`$ and not 0.8 as stated.
2. The number of $`d\mu {}_{}{}^{3}\mathrm{He}`$ complexes formed in the course of data taking in their experiment was defined as
$$N_{d\mu {}_{}{}^{3}\mathrm{He}}=N_\mu C_{d\mu }\frac{\lambda _{d\mu {}_{}{}^{3}\mathrm{He}}}{\lambda _{d\mu }},$$
(39)
and correspond to $`N_{d\mu {}_{}{}^{3}\mathrm{He}}=(4.9\pm 0.4)\times 10^8`$.
According to our estimations, the quantities $`\lambda _{d\mu {}_{}{}^{3}\mathrm{He}}`$, $`\lambda _{d\mu }`$, $`\lambda _{pd\mu }`$ ($`pd\mu `$ molecule formation rate), and $`N_{d\mu {}_{}{}^{3}\mathrm{He}}`$ had the values $`\lambda _{d\mu {}_{}{}^{3}\mathrm{He}}=1.32\times 10^6\mathrm{s}^1`$ ($`\phi =0.0975`$ $`\mathrm{c}_{{}_{}{}^{3}\mathrm{He}}`$ = 0.056, $`\lambda _{d^3He}^0=2.42\times 10^8\mathrm{s}^1`$Bystritsky et al. (2004),
$`\lambda _{d\mu }`$ $``$ $`\lambda _0+\lambda _{d\mu {}_{}{}^{3}\mathrm{He}}\phi \mathrm{c}_{{}_{}{}^{3}\mathrm{He}}+\lambda _{pd\mu }\phi \mathrm{c}_p+\stackrel{~}{\lambda }_F\omega _d\phi \mathrm{c}_d`$ (40)
$``$ $`2.05\times 10^6\mathrm{s}^1`$
$`\lambda _{pd\mu }=5.6\times 10^6\mathrm{s}^1`$, which yields $`N_{d\mu {}_{}{}^{3}\mathrm{He}}3.7\times 10^8\text{s}^1`$ instead of $`(4.9\pm 0.4)\times 10^8\text{s}^1`$.
3. Their ionization chamber detection efficiency for protons from reaction (2aa) was defined as $`\epsilon =\epsilon _S\epsilon _\tau `$ and found to be $`\epsilon =0.082`$, where $`\epsilon _S=0.13`$ is the selection factor for events detected in compliance with certain amplitude and geometrical criteria, $`\epsilon _\tau =0.63`$ is the time factor to take of the fact that the detected events were analyzed in the time interval $`0.4t1.8\mu `$s. According to our estimation, $`\epsilon _\tau =e^{\lambda _{d\mu }t_1}e^{\lambda _{d\mu }t_2}=0.44`$, because under their experimental conditions the $`d\mu `$ disappearance rate is $`\lambda _{d\mu }2.05\times 10^6\text{s}^1`$, $`t_1=0.4\mu `$s, and $`t_2=1.8\mu `$s.
As can be seen, taking into account only the above items alone the upper limit of $`\stackrel{~}{\lambda }_f`$ is, to our mind, appreciably underestimated in the work of Maev et. al. Another cause of this underestimation might be the improper background subtraction procedure because they determined the background level using information from earlier experiments Balin et al. (1998) carried out under different conditions and at an experimental facility which was not completely analogous. In addition, it is slightly surprising that the background from muon capture by $`{}_{}{}^{3}\mathrm{He}`$ nuclei with the formation of protons in the energy region near 14.64 MeV is estimated at zero in Ref. Maev et al. (1999)(see <sup>2</sup><sup>2</sup>2According to Ref. Bystritsky et al. (2004), the fraction of protons from muon capture by the $`{}_{}{}^{3}\mathrm{He}`$ nucleus in the energy range $`14.314.64`$ MeV per $`\mu {}_{}{}^{3}\mathrm{He}`$ atom is $`W_{{}_{}{}^{3}He}^p=2\times 10^6`$.).
We believe that our $`\stackrel{~}{\lambda }_f`$ measurement results are reliable, which is confirmed by stable observation of nuclear fusion in both runs with the $`\mathrm{D}_2+{}_{}{}^{3}\mathrm{He}`$ mixture differing in density by a factor of about three. Nevertheless, as far as the experimental results obtained in this paper and in Ref. Maev et al. (1999) are concerned, the things are unfortunately uncertain and need clarifying.
There is a point important for comparison of the calculated $`\lambda _f^{J=0}`$ with the results of the previous experiments Maev et al. (1999) and this paper. Measurement of $`\lambda _f^{J=0}`$ is indirect because it is determined by Eq. (37) with the calculated effective rate for transition of the $`d\mu {}_{}{}^{3}\mathrm{He}`$ complex from the $`J=1`$ state to the $`J=0`$ state. Therefore, $`\lambda _f^{J=0}`$ is not uniquely defined and greatly depends on $`\stackrel{~}{\lambda }_{10}`$, which in turn is determined by the chain of processes accompanying and competing with the $`10`$ transition of the $`d\mu {}_{}{}^{3}\mathrm{He}`$ complex. To rule out this lack of uniqueness in determination of $`\lambda _f^{J=0}`$ and, in addition, to gain information on the effective $`10`$ transition rate $`\stackrel{~}{\lambda }_{10}`$ and the nuclear fusion rate $`\lambda _f^{J=1}`$ in the $`d\mu {}_{}{}^{3}\mathrm{He}`$ complex in the $`J=1`$ state, it is necessary, as proposed in Refs. Bystritsky and Penโkov (1999); Bystritsky et al. (1999b, c), to carry out an experiment with the $`\mathrm{D}_2+{}_{}{}^{3}\mathrm{He}`$ mixture at least at three densities in the range $`\phi =0.030.2`$, where not only protons from reaction (2aa) but also 6.85 keV $`\gamma `$ rays should be analyzed. Analysis of the results reported in this paper and in Ref. Maev et al. (1999) makes it possible to put forward some already obvious proposals as to getting unambiguous and precise information on important characteristics of $`\mu `$-molecular ($`\lambda _{d\mu {}_{}{}^{3}\mathrm{He}}`$, $`\stackrel{~}{\lambda }_{10}`$) and nuclear ($`\stackrel{~}{\lambda }_f`$, $`\lambda _f^{J=0}`$, $`\lambda _f^{J=1}`$) processes occurring in the $`\mathrm{D}_2+{}_{}{}^{3}\mathrm{He}`$ mixture. It is necessary to conduct experiments at no less than three densities of the ($`\mathrm{HD}+{}_{}{}^{3}\mathrm{He}`$) or ($`\mathrm{H}_2+\mathrm{D}_3(1\%)+{}_{}{}^{3}\mathrm{He}`$) mixture with detection of both protons from reaction (2aa) and 6.85 keV $`\gamma `$ rays, to increase at least three times the detection efficiency for protons $`\epsilon _p`$ and for muon decay electrons $`\epsilon _e`$ in comparison with the corresponding efficiencies in the present experiment.
###### Acknowledgements.
The authors would like to thanks R. Jacot-Guillarmod for his help during the conception of this experiment. We are thankful to V.F. Boreiko, A. Del Rosso, O. Huot, V.N. Pavlov, V.G. Sandukovsky, F.M. Penkov, C. Petitjean, L.A. Schaller and H. Schneuwly for they help during the construction of the experiments, the data taking period, and for very useful discussions. This work was supported by the Russian Foundation for Basic Research, Grant No. 01โ02โ16483, the Polish State Committee for Scientific Research, the Swiss National Science Foundation, and the Paul Scherrer Institute.
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# The Double Pulsar J0737โ3039: Testing the Neutron Star Equation of State
## 1 Introduction
The discovery of the first double pulsar, J0737โ3039 (Burgay et al. 2003; Lyne et al. 2004) consisting of two pulsars with spin periods of 22.7 ms (Pulsar A) and 2.77 s (Pulsar B) in a 2.4 hr orbit, has opened up a new window for testing fundamental physics under extreme conditions: not only will the system soon allow some of the best tests of General Relativity (Kramer et al. 2005), but also it is the first system where the physics of interacting magnetospheres from two pulsars can be studied (e.g. Arons et al. 2005). The orbit is decaying due to the loss of angular momentum by gravitational radiation, and the two neutron stars are expected to merge in only $`85`$Myr, a much shorter time than in any other double neutron star system.
One of the other interesting characteristics of J0737โ3039 is that the mass of the second neutron star to be formed (Pulsar B) is the lowest reliably measured mass for any neutron star to date: $`1.249\pm 0.001\text{M}_{}`$ (Kramer et al. 2005). Such a low mass may be an indication that the neutron star did not form in a standard iron-core-collapse supernova but in an electron-capture supernova (Nomoto 1984; Podsiadlowski et al. 2004)<sup>1</sup><sup>1</sup>1Note, however, that it is not entirely clear that electron-capture supernovae necessarily produce the lowest remnant masses, since iron cores with masses as low as $`1.15\text{M}_{}`$ may be unstable to collapse (see the discussion in Woosley, Heger & Weaver 2002).. These occur for ONeMg white dwarfs when the core density reaches a critical value at which electron captures (e-captures) onto Mg (and subsequently Ne) start, causing a loss of hydrostatic support in the core and triggering its collapse. One of the key aspects of an e-capture supernova is that the collapse takes places when the core density reaches a well-defined critical value ($`4.5\times 10^9`$g cm<sup>-3</sup>) which in turn occurs when the ONeMg core has grown to a well-defined critical mass ($`1.37\text{M}_{}`$; see ยง 3). Therefore, if Pulsar B indeed formed in an e-capture supernova, this would be the first instance for which the masses of the pre-collapse core and the post-collapse neutron star could both be determined, the former from a theoretical estimate of the critical mass for an e-capture supernova, the latter directly from the observed orbital parameters. Along with the pre-collapse gravitational mass, the corresponding baryon number can also be calculated. Since the loss of material during the formation of the neutron star is expected to be extremely small in this scenario (see ยง 3), this is also a good approximation to the baryon number of the neutron star. It is the purpose of this paper to demonstrate that both the observed orbital parameters of the system (in particular the low eccentricity) and the most likely evolutionary history of the system favour the formation of Pulsar B in an e-capture supernova and that comparison of its gravitational mass with the estimate obtained for the baryon number enables useful constraints to be placed on the neutron-star equation of state (EoS)<sup>2</sup><sup>2</sup>2A different test of the EoS of neutron-star matter can be derived from measuring the moment of inertia of Pulsar A from the effects of spin-orbit coupling, as proposed by Morrison et al. (2004)..
In ยง 2 we review the two most likely evolutionary channels that lead to systems like the double pulsar and show how this discussion supports the key assumption of Pulsar B having been formed in an e-capture supernova. In ยง 3 we provide a theoretical estimate for the pre-collapse core mass (with an estimate of the uncertainties) and in ยง 4 we demonstrate how the properties of the system can be used to constrain the neutron star EoS.
## 2 The Evolutionary History of J0737โ3039
There are two major evolutionary channels to form double neutron stars like the double pulsar which we shall refer to in this paper as the standard channel (e.g. Bhattacharya & van den Heuvel 1991) and the double-core channel (e.g. Brown 1995), respectively. In the standard channel (see the left panel in Fig. 1), J0737โ3039 originates from a massive binary in which the more massive star transfers its envelope to its companion star via stable Roche-lobe overflow (RLOF) before it collapses in a supernova (SN) explosion to form the first neutron star. If the SN explosion does not disrupt the system, the binary now consists of a neutron star and a massive main-sequence star, and the system will evolve into a high-mass X-ray binary (observationally it will initially look like a Be X-ray binary). As the secondary evolves, there will be a point at which it will fill its Roche lobe and start to transfer matter to the neutron star. Because of the large mass ratio of the system, this mass transfer is unstable and leads to a common-envelope (CE) and spiral-in phase, in which the neutron star spirals towards the centre of the massive companion inside the companionโs envelope. Provided that the CE phase does not lead to the complete merger of the two components and that the CE is ejected, the system evolves into a very close binary containing the helium core of the secondary and the neutron star. Depending on the mass of the helium star and the period of the system, another phase of mass transfer may occur, where the neutron star is spun up and becomes a fast โrecycledโ pulsar. Eventually the helium star collapses to form the second neutron star in the system.
In the double-core channel (see the right panel in Fig. 1), the binary components are very close in mass initially (within 5 โ 10 %), and the orbit is relatively wide, so that the primary only fills its Roche lobe after it has completed helium core burning (so-called Case C mass transfer) and has developed a CO core. At this stage, the secondary has already completed its hydrogen core burning phase and has evolved off the main sequence. Because of the high mass transfer rate, the accreting star expands to fill and ultimately overflow its Roche lobe, and the system is again expected to enter into a CE phase; but, in this case, the CE is formed from the combined envelopes of both stars. Inside the CE, there are the cores of the two stars, the more evolved one with a CO core and the less evolved He core of the secondary. The cores spiral-in inside the joint envelope until the envelope is ejected, leaving a very close binary consisting of the two evolved cores. The CO core soon collapses to form the first neutron star, leaving a binary consisting of a neutron star and a helium star. The further evolution is almost identical to the evolution in the standard scenario.
Because of the constraints on the initial mass ratio and the orbital separation, the double-core channel is expected to have a lower occurrence rate than the standard channel (by a factor of 2 to 10; Dewi, Podsiadlowski & Sena 2005; in preparation)<sup>3</sup><sup>3</sup>3However, if the spiral-in of a neutron star in a massive envelope always leads to hypercritical accretion onto the neutron star and its conversion into a black hole, as argued, e.g., by Brown (1995), the double-core channel would be the only one of these two channels that could produce double neutron star systems (also see Chevalier 1993)..
The formation of J0737โ3039 via the standard channel has been studied by Dewi & van den Heuvel (2004) and Willems & Kalogera (2004); both studies concluded that this system must have originated from a close helium starโneutron star (HeS-NS) binary where the system underwent mass transfer during the helium-star phase, spinning up the first-born neutron star in the process. Because the final stages of evolution in the standard and the double-core channel are essentially the same, i.e. the HeS-NS phase, the following discussion, which assumes the standard channel, also applies to the double-core channel.
For different reasons, Dewi & van den Heuvel (2004) and Willems & Kalogera (2004) found a similar pre-SN mass of the helium star progenitor of J0737โ3039. Willems & Kalogera (2004) took the threshold helium-star mass for the formation of a neutron star to be 2.1 $`\mathrm{M}_{}`$ (Habets 1986) as the lower limit. However, one should note that after the mass-transfer phase, the immediate pre-supernova mass can be as low as the Chandrasekhar mass ($`1.4\text{M}_{}`$). Dewi & van den Heuvel (2004) used 2.3 $`\mathrm{M}_{}`$ as the minimum possible helium-star mass at the time of the explosion, based on the assumption that lower-mass helium stars experience a further CE phase at the end of their evolution (Dewi et al. 2002; Dewi & Pols 2003). This lower limit on the pre-SN mass then required a minimum kick velocity of 70 $`\mathrm{km}\mathrm{s}^1`$, since the supernova mass loss on its own would produce a much larger post-supernova eccentricity than is consistent with the present orbital parameters of J0737โ3039.
However, the assumption of the occurrence of another CE phase in a lower-mass helium star is still an open question. A recent population synthesis study of the formation of double neutron stars (Dewi, Pols & van den Heuvel, 2005, in preparation) suggests that, particularly to explain the formation of J0737โ3039, it is more likely that this CE phase does not occur. In this case, the final pre-supernova mass can be much less than 2.3 $`\mathrm{M}_{}`$, indeed it can be as low as $`1.4\text{M}_{}`$, and no supernova kick is required to compensate for the mass loss; in particular it allows for the possibility that the supernova was symmetric (Dewi & van den Heuvel 2005), strongly favouring an e-capture SN for the second supernova.
### An electron-capture supernova to form Pulsar B?
The eccentricity of the double pulsar binary ($`e0.088`$ at present, most likely 0.11 โ 0.12 immediately after the second supernova; e.g. Burgay et al. 2003) is surprisingly low, much lower than one would expect if the system received a large supernova kick in the second supernova that formed Pulsar B. Indeed, such a low eccentricity can be most easily explained by a symmetric second supernova in which a moderate amount of mass is expelled from the system. In this case the eccentricity is given by $`e=\mathrm{\Delta }M/(M_\mathrm{A}+M_\mathrm{B})`$, where $`\mathrm{\Delta }M`$ is the mass lost in the supernova and $`M_\mathrm{A}`$ and $`M_\mathrm{B}`$ are the present masses of Pulsars A and B, respectively. Taking $`M_\mathrm{A}=1.338\text{M}_{}`$ and $`M_\mathrm{B}=1.249`$ (Kramer et al. 2005) and assuming a post-SN eccentricity $`e_0=0.12`$ then yields a pre-SN mass of the helium star of $`1.56\text{M}_{}`$. Such low-mass pre-SN helium stars typically form from HeS-NS binaries with initial helium stars of less than $`3\text{M}_{}`$ (see Dewi et al. 2002; Ivanova et al. 2003). This includes the mass range where helium stars are expected to end their evolution in an e-capture supernova (Nomoto 1984). This is also consistent with the speculation by Podsiadlowski et al. (2004) that e-capture supernovae may not produce large supernova kicks since the explosion may proceed on a timescale that is much shorter than the timescale on which the instabilities that produce large kicks can develop; this suggestion has received some theoretical support from recent core-collapse calculations (Scheck et al. 2004; H.-Th. Janka 2005 \[private communication\]). Since in any simple accretion model one would expect that the spin of Pulsar A would become aligned with the orbital momentum axis and since this orientation is not affected by a symmetric supernova, this may make the testable prediction that the post-SN misalignment angle between the spin of pulsar A and the orbital axis should be relatively small. This may indeed account for the surprising stability of the pulse shape of Pulsar A (Manchester et al. 2005).
### The second supernova kick and the space velocity of J0737โ3039
A low supernova kick velocity, as suggested by the low eccentricity, would typically also imply that the binary system should only have received a relatively small additional kick in the second supernova and that the system space velocity relative to the local standard of rest would not be much affected by it. Indeed, the space velocity of J0737โ3039 may provide an important constraint on its evolutionary history. Unfortunately, the situation regarding the value of this quantity is at present somewhat confused. Using interstellar scintillation to measure the transverse space velocity of the system, Ransom et al. (2004) determined a large system velocity of at least 140 km s<sup>-1</sup>. Subsequently, Coles et al. (2005) showed that the dispersion across the field was highly variable reducing the estimate of the scintillation velocity to a value as low as 66 km s<sup>-1</sup> and possibly even lower, since this value does not account for the motion of the Earth. More recently, Kramer et al. (2005) argued that the present limits on the proper motion of the system suggest a low transverse velocity of less than 30 km s<sup>-1</sup>, which would imply that it is statistically unlikely that the system received a large kick in the second supernova. On the basis of the original high estimate of the space velocity, Ransom et al. (2004) and Willems et al. (2004) concluded that the system should have received a large kick in the second supernova. However, they only considered the standard scenario in which the system is still fairly wide at the time of the first supernova; this implies that it cannot receive a large kick in the first supernova and remain bound (see Fig. 1). In contrast, in the double-core scenario where the system is already very tight at the time of the first supernova, the system is expected to receive a large systemic kick from the first supernova (of order 150 โ 400 km s<sup>-1</sup>: see Fig. 1 and Dewi et al. 2005, in preparation) and hence no additional kick from the second supernova would be required. As far as testing the EoS is concerned, it is not important whether the system velocity is low or high, since both can be understood within the framework of either of the two channels and a small second kick. In particular, we note how similar the final phases are in the two channels. A resolution of the issue of the systemโs space velocity could, however, provide a powerful discriminant between the standard channel and the double-core channel.
## 3 Electron-Capture Supernovae
An electron-capture supernova occurs when the central density of an ONeMg core reaches the threshold value $`\rho _{\mathrm{th}}`$ for electron captures on <sup>24</sup>Mg. This decreases the electron pressure and the electron fraction $`Y_e`$, lowering the Chandrasekhar mass and triggering the collapse of the core (Miyaji et al. 1980).
In order to estimate the uncertainties in the critical mass of the collapsing core, we performed a series of stellar structure calculations assuming non-rotating cores in hydrostatic equilibrium with a prescribed central density, homogeneous composition and a specified thermal profile. Since the heat released by the electron captures gives rise to a convective core, we adopted an isentropic thermal profile.
As our reference model we adopted the composition given by Gutiรฉrrez et al. (1996) with X(<sup>16</sup>O)=0.72, X(<sup>20</sup>Ne)=0.25, X(<sup>24</sup>Mg)=0.03, central density $`\rho _{\mathrm{th}}=4.5\times 10^9`$ g cm<sup>-3</sup> and a range of central temperatures from $`10^7`$ to $`10^9`$ K. We used the equation of state of Pols et al. (1995), assuming full ionization (i.e. we discard the ionization pressure term). Note that neglecting the Coulomb corrections in the equation of state would increase the total mass by $`3.79\times 10^2`$M for our reference composition.
We integrated the general relativistic (GR) equations of hydrostatic equilibrium out from the centre
$$\frac{\mathrm{d}P}{\mathrm{d}r}=\frac{Gm\rho }{r^2}\left(1+\frac{P}{\rho c^2}\right)\left(1+\frac{4\pi r^3P}{mc^2}\right)\left(1\frac{2Gm}{rc^2}\right)^1,$$
(1)
where
$$\frac{\mathrm{d}m}{\mathrm{d}r}=4\pi \rho r^2,$$
(2)
and
$$\frac{\mathrm{d}A}{\mathrm{d}r}=4\pi n_br^2\left(1\frac{2Gm}{rc^2}\right)^{\frac{1}{2}}.$$
(3)
Here $`m`$ and $`A`$ are the gravitational mass and baryon number enclosed within a sphere of radius $`r`$; $`\rho `$ is the density, including contributions from both the rest mass and the thermal energy; $`n_b`$ is the baryon number density. Rather than talking in terms of the baryon number $`A`$, which is a rather abstract quantity, it is convenient to convert this into a mass by multiplying by the atomic mass unit ($`931.50`$Mev/$`c^2`$), and we refer to this quantity as the baryonic mass. If we had ignored the general relativistic corrections, our estimate of the baryonic mass would have been increased by $`1.30\times 10^2`$M ($`1\%`$). This difference is larger than might have been expected; the reason, however, is analogous to that for the large difference between the Newtonian and GR values for the maximum mass of a neutron star (Oppenheimer & Volkoff 1939).
To check the validity of our procedure, we computed a model without the GR corrections to reproduce the case considered by Miyaji et al. (1987) and Nomoto (1987), where they used a composition $`X(^{16}`$O$`)=0.12`$, $`X(^{20}`$Ne$`)=0.76`$, $`X(^{24}`$Mg$`)=0.12`$, central density $`3.98\times 10^9`$ g cm<sup>-3</sup> and central temperature $`\mathrm{log}(T_c)=8.61`$. We find a total mass $`m=1.3754`$ M, which is very close to their published value of 1.375 M.
This composition is, however, no longer considered appropriate since the reaction rate <sup>12</sup>C($`\alpha ,\gamma `$)<sup>16</sup>O has been revised upwards (Dominguez, Tornambe & Isern 1993). This leads to a lower Ne and Mg abundance and a higher O abundance at the end of C burning than computed with the Fowler et al. (1975) rates (as was done in Miyaji et al. 1987). To estimate the uncertainties introduced by the composition, we investigated several different compositions found in the literature (Dominguez et al. 1993; Gutiรฉrrez et al. 1996; Ritossa, Garcรญa-Berro & Iben 1996; Gil-Pons & Garcรญa-Berro 2001) after the revision of the <sup>12</sup>C($`\alpha ,\gamma `$)<sup>16</sup>O rate. The composition is important in determining the Coulomb parameter and $`Y_e`$, and hence the magnitude of the Coulomb corrections. We use our reference model to give a lower bound on the critical mass (solid curve in Fig 2) and the composition $`X(^{16}`$O$`)=0.56`$, $`X(^{20}`$Ne$`)=0.29`$, $`X(^{24}`$Mg$`)=0.06`$, $`X(^{23}`$Na$`)=0.07`$ and $`X(^{12}`$C$`)=0.01`$, which mimics the composition of Gil-Pons & Garcรญa-Berro (2001) and has the lowest Coulomb correction, to give an upper bound to the critical mass (dashed curve in Fig 2).
Finally, we changed the threshold density $`\rho _{\mathrm{th}}`$ from $`4.5\times 10^9`$ to $`4\times 10^9`$ g cm<sup>-3</sup> in order to estimate the effect of the shift in the critical density due to the Coulomb corrections (the latter is the appropriate density without any Coulomb corrections; see Gutiรฉrrez et al. 1996 for a detailed discussion). This change decreases the critical mass by $`2\times 10^3`$M (dotted and dot-dashed lines in Fig 2).
In the previous studies, the central temperature after the onset of electron captures ranged from $`\mathrm{log}T_c=8.5`$ up to $`\mathrm{log}T_c=8.65`$. For this temperature range, Figure 2 then implies that the baryonic mass of the pre-collapse core should lie between 1.366 and 1.375$`\text{M}_{}`$ (the thin dotted curves) taking into account the uncertainties in temperature and composition. On the assumption that the loss of material during formation of the neutron star is negligible, this then gives the predicted range for the baryonic mass of the neutron star (which we refer to as $`M_0`$).
### Caveats
At this point, several caveats should be made about other effects which may systematically affect this estimate. First, we take the baryonic mass of the pre-collapse ONeMg core to be the same as that of the neutron star, neglecting any loss of material during the formation of the neutron star. In practice, some material may be ejected in the supernova in a neutrino-driven wind (Qian & Woosley 1996). However, because of the steep density gradient at the edge of the ONeMg core, we expect this mass loss to be small, probably less than a few times $`10^3\text{M}_{}`$, although it could potentially be as large as $`10^2\text{M}_{}`$ and the amount of mass loss itself depends on the EoS (Janka 2005, private communication). Second, an e-capture supernova may occur between carbon shell flashes when the central density increases significantly (although we note that this point is not yet fully resolved). Because of the discrete nature of the carbon flashes, this may introduce a natural variation in the critical mass from star to star because of variations of the thermal profile in the outer ONeMg core. In this context we note that the assumption of an isentropic profile is unlikely to be correct in the outer parts of the ONeMg core.
It is clear that these effects could significantly change our estimate for the baryonic mass of the neutron star (probably decreasing it) and therefore our present estimates should only be considered as preliminary. However, with the expected progress in simulating e-capture supernovae and calculating the evolution of their progenitors (with improved e-capture rates, a richer nuclear network, inclusion of accurate Coulomb corrections and using GR), we estimate that ultimately one should be able to pin-point the mass of the collapsing core to within $`2\times 10^3`$M.
## 4 Constraints on the Equation of State
On the hypothesis that the scenario presented in the previous sections is correct, we can use the information about the gravitational and baryonic masses of Pulsar B to place constraints on the EoS of neutron-star matter. These turn out to be quite interesting.
For any given EoS for neutron-star matter, one can calculate the relation between the gravitational mass and the baryonic mass (bearing in mind that the rotation speed is so low that taking the object to be spherical and non-rotating is an excellent approximation). The present gravitational mass (known from observations) and the baryonic mass (known from the stellar evolution calculations) then define an error box through which the relations calculated from the equations of state would need to pass in order to be consistent. In this section, we discuss how the constraints obtained in this way turn out. The observed gravitational mass ($`M_G=1.249\pm 0.001\text{M}_{}`$) and the calculated baryonic mass ($`M_0`$ in the range $`1.3661.375\text{M}_{}`$) specify the boundaries of the error box.
Results obtained by integrating the GR equations of hydrostatic equilibrium (Eqs 1 โ 3) for a range of EoSs are shown in the four panels of Figure 3, each corresponding to a particular class of equations of state (we use a modified form of the classification in the paper by Morrison et al. 2004). All of these EoSs give the maximum gravitational mass $`M_{\mathrm{max}}`$ for a non-rotating neutron star as being above $`1.5\text{M}_{}`$, in line with observations. Class I EoSs (top left-hand panel) come from non-relativistic many-body calculations with โrealisticโ potentials: APR98 (Akmal, Pandharipande & Ravenhall 1998), WFF88 (Wiringa, Fiks & Fabrocini 1988) and FPS (Lorenz, Ravenhall & Pethick 1993) include only nucleonic degrees of freedom, BJ74 (Bethe & Johnson 1974) and MOSZ74 (Mozskowski 1974) include also hyperonic components at the higher densities. Class II EoSs (top right-hand panel) use relativistic mean-field (or effective-field) approximations including hyperonic degrees of freedom: GLE210, GLE240 and GLE300 (Glendenning 2000), and HOF01 (Hofmann et al. 2001) are shown. We also include here one EoS representing a hybrid stellar model (nucleons + quarks): GLENHYB (Glendenning 2000). Class III EoSs (bottom left-hand panel) use non-relativistic phenomenological potentials of the Skyrme type (see Stone et al. 2003 and references therein). Class IV EoSs (lower right-hand panel) are for other phenomenological non-relativistic potentials: BPAL21 and BPAL31 (Prakash et al. 1997) have only nucleonic degrees of freedom, while BAL97 (Balberg & Gal 1997) includes hyperons at high density. All of these high-density EoSs were joined onto the Baym-Bethe-Pethick EoS (Baym, Bethe & Pethick 1971) at a density of $`1.4\times 10^{14}`$ g cm<sup>-3</sup> ($`0.080.09`$fm<sup>-3</sup>) and, in turn, this was joined onto the Baym-Pethick-Sutherland EoS (Baym, Pethick & Sutherland 1971) at $`4.2\times 10^{11}`$ g cm<sup>-3</sup> ($`2.5\times 10^4`$fm<sup>-3</sup>). By doing this, the inner and outer crust of the neutron star were treated in the same way for all of the EoSs, and so all of the differences seen result from differences in the treatment of the high-density matter. (Strange star models have not been included in Fig 3; they cover a very wide range and give curves passing both above and below the error box as well as curves passing through it.)
The most clear-cut result is that none of the class II models tested in this work give predictions in line with our constraint. For the other classes (I, III and IV), the situation is less clear and depends on the particular properties of each individual EoS. For the phenomenological Skyrme potentials, the EoSs give a wide range of predictions within the region delimited by SkI1 and BSk8. Those parametrisations giving $`M_G/M_0`$ curves passing through the error box give $`M_{\mathrm{max}}`$ between $`1.61.9\text{M}_{}`$. All of those for which $`M_{\mathrm{max}}>1.9\text{M}_{}`$ give curves passing above the error box while those for which $`M_{\mathrm{max}}<1.6\text{M}_{}`$ give curves passing below it. With reference to the discussion in Stone et al. (2003) and using the notation of that paper, we note that all of the Skyrme EoSs passing through the error box are type II parametrisations whereas all of those coming from type I parametrisations pass above it.
Apart from the comments made above, there is no simple general interpretation of the implications of our constraint for the physics behind the particular EoSs. It is important to recognise that our constraint (assuming that our overall scenario is correct) represents a necessary but not sufficient condition for choosing a suitable EoS for neutron star models. Additional observational information is needed, in particular concerning the neutron star maximum mass which, in combination with the present constraint, would give a more definitive criterion for the choice of physical model for the EoS. Also, the influence of different treatments for the inner and outer crust needs to be investigated before any final conclusion is drawn.
## 5 Conclusions
In this paper, we have demonstrated that the measured gravitational mass of Pulsar B in the double pulsar J0737โ3039 can be used to give a new test of the neutron-star EoS if one makes the all-important assumption that Pulsar B was formed in an electron-capture supernova. In this case, its baryonic mass can be estimated theoretically and comparison between this and its gravitational mass can then be used to constrain the EoS. We have re-constructed the possible evolutionary histories for J0737โ3039 in the main formation channels to support the hypothesis of Pulsar B having originated in an electron-capture supernova and have discussed possible tests of this hypothesis. Future refinements, both of the stellar evolution models (to better pin down the critical pre-collapse mass) and of electron-capture collapse models (to quantify the possible mass loss), should lead to a more stringent constraint which, when combined with other astrophysical EoS constraints, can provide new insight into the physics of neutron-star matter.
## Acknowledgements
We thank H.-Th. Janka and M. Kramer for very useful discussions and sharing the results of some of their unpublished work and C.M.Keil for supplying numerical data for EOS HOF01 in Figure 3. This work was in part supported by a European Research & Training Network on Type Ia Supernovae (HPRN-CT-20002-00303, PhP, PL), a Talent Fellowship (JDMD) from the Netherlands Organization for Scientific Research (NWO), EPSRC Grant 02300018 (WN), an Advanced Computing Grant from US DOE Scientific Discovery (JRS) and US DOE grant DE-FG02-94ER40834 (JRS).
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# Search for New Physics Using High Mass Tau Pairs from 1.96-TeV ๐โข๐ฬ Collisions
CDF Collaboration
(June 13, 2005)
## Abstract
We present the results of a search for anomalous resonant production of tau lepton pairs with large invariant mass, the first such search using the CDF II Detector in Run II of the Tevatron $`p\overline{p}`$ collider. Such anomalous production could arise from various new physics processes. In a data sample corresponding to 195 pb<sup>-1</sup> of integrated luminosity we predict 2.8$`\pm `$0.5 events from Standard Model background processes and observe 4. We use this result to set limits on the production of heavy scalar and vector particles decaying to tau lepton pairs.
At the Fermilab Tevatron $`p\overline{p}`$ collider, a number of nonโStandard-Model physics processes can lead to events with high-mass tau lepton pairs in the final state. Examples include the resonant production of Higgs scalars in two-Higgs-doublet models ref-twohiggs at large $`\mathrm{tan}\beta `$, the ratio of the vacuum expectation value of the two doublets. Two Higgs doublets are required, for example, in the Minimal Supersymmetric Standard Model ref-mssm , a favored candidate for extending the Standard Model. The heavy scalar and pseudoscalar Higgs bosons in this theory would decay to tau pairs about 9% of the time. Also, in supersymmetry, if R-parity is not conserved, heavy scalar neutrino production could have tau pair decay modes ref-rpvsnu . If there are heavy $`Z^{}`$ bosons, these could also produce high mass tau pairs in the final state, possibly even with enhanced tau couplings ref-zprime . With the large new data sample from Run II of the Tevatron it is thus of great interest to perform a generic search for high-mass tau pairs.
This Letter presents the results of a search for high-mass tau pairs performed using CDF II, the upgraded Collider Detector at Fermilab (CDF) ref-thesis . In 2002 and 2003 CDF recorded a data sample corresponding to 195 pb<sup>-1</sup> of integrated luminosity of $`p\overline{p}`$ collisions at a center of mass energy of 1.96 TeV. This is the first such search with the new high-statistics data sample ref-amy and new tau identification techniques.
Since the tau lepton decays to lighter leptons ($`e`$ or $`\mu `$) about 35% of the time, and to low-multiplicity hadronic states the rest of the time, this analysis selects events with one identified hadronic tau decay ($`\tau _h`$) and one other tau decay in the final state. Thus, there are three distinct final states, which we denote $`e\tau _h`$, $`\mu \tau _h`$, and $`\tau _h\tau _h`$. The main background to this search comes from Drell-Yan (DY) $`Z/\gamma ^{}\tau ^+\tau ^{}`$ production. Since we seek new particles with mass much larger than that of the $`Z`$, we use the observed rate for this background (at smaller tau pair masses) as a control sample, and define the signal region as that where the tau pairs have large visible invariant mass, with missing energy due to the neutrinos from the tau decays.
CDF II is a large general purpose detector with an overall cylindrical geometry surrounding the $`p\overline{p}`$ interaction region ref-cdf . The three-dimensional trajectories of charged particles produced in $`p\overline{p}`$ collisions are measured starting at radii of 1.5 cm with multiple layers of silicon microstrip detectors, and are measured at outer radii with an axial/stereo wire drift chamber (COT). The tracking system lies inside a uniform 1.4-T magnetic field produced by a superconducting solenoid, with the field oriented along the beam direction. Outside the solenoid lie the electromagnetic calorimeter and the hadronic calorimeters, which are segmented in pseudorapidity ($`\eta `$ref-eta and azimuth in a projective โtowerโ geometry. A set of strip/wire chambers (CES) located at a depth of six radiation lengths aids in reconstructing photons and electrons from the shower shape. Muons are identified by a system of drift chambers placed outside the calorimeter steel, which acts as an absorber for hadrons. The integrated luminosity of the $`p\overline{p}`$ collisions is measured to an accuracy of 6% using the Cerenkov Luminosity Counters ref-lum .
The $`e\tau _h`$ and $`\mu \tau _h`$ events of interest are recorded using triggers designed to select โlepton plus trackโ events: those with an $`e`$ or $`\mu `$ with transverse momentum ($`p_T`$) greater than 8 GeV/$`c`$ and another charged track with $`p_T>5`$ GeV/$`c`$ identified by the eXtremely Fast Tracker (XFT) portion of the trigger electronics ref-xft which reconstructs charged tracks in the COT. The efficiency of this trigger is measured using leptons from $`Z`$ boson decays, the $`\mathrm{{\rm Y}}`$ resonance, and photon conversions ref-NIM .
For selecting $`\tau _h\tau _h`$ events we use a trigger designed to select at least one hadronically decaying tau with $`E_T>20`$ GeV accompanied by at least 25 GeV missing energy in the plane transverse to the beam direction ($`/E_T`$). The tau is identified by matching an XFT track with $`p_T>5`$ GeV/$`c`$ to a calorimeter cluster. Data used from this trigger come from a sample corresponding to the first 72 pb<sup>-1</sup> of integrated luminosity recorded; this is less than that of the rest of the data used because of subsequent changes due to rate limitations.
Events selected by the triggers were recorded and processed later to reconstruct charged particle tracks, calorimeter clusters, and to identify electrons, muons, photons, jets, and $`/E_T`$. Electrons and muons are reconstructed using algorithms described in Ref. ref-cdf . Identification of hadronic decays of taus employs a novel โshrinking coneโ algorithm based on high-$`p_T`$ charged tracks in the silicon/COT system, and $`\pi ^0`$ candidates identified using the CES by matching strip clusters with wire clusters based on the energy in each.
The $`\tau _h`$ identification algorithm begins with a list of โseed tracksโ ranked in $`p_T`$, not yet used for another tau candidate, and having $`p_T>6`$ GeV/$`c`$. Then it finds the number of other tracks with $`p_T>1`$ GeV/$`c`$ and $`\pi ^0`$ candidates with at least 1 GeV whose momentum vector makes an angle of less than $`\alpha `$ with the seed track. The angle $`\alpha `$ is a function of $`E_{clu}`$, the energy in the calorimeter cluster associated with the seed track. The value of $`\alpha `$ is 10 or (5 GeV)/$`E_{clu}`$ radians, whichever is less. To allow for resolution effects, the value of $`\alpha `$ is not less than 100 mrad for $`\pi ^0`$ candidates, or 50 mrad for charged tracks. If any other tracks or $`\pi ^0`$ candidates have an angle greater than $`\alpha `$ but less than 30 to the seed track, or if the invariant mass calculated from the sum of all charged track and $`\pi ^0`$ candidate four-momenta exceeds 1.8 GeV/$`c^2`$, the $`\tau _h`$ candidate is rejected. Also, if there is at least 2 GeV of electromagnetic energy in calorimeter towers not part of the tau cluster, and whose centers have $`\mathrm{\Delta }R=\sqrt{(\mathrm{\Delta }\eta )^2+(\mathrm{\Delta }\varphi )^2}<0.4`$ from the tau seed track direction, the tau candidate is rejected. Candidates with momentum having an angle of less than 10 with that of of a previously identified $`e`$ or $`\mu `$ are rejected. The algorithm then considers further possible seed tracks, repeating the process until none remain.
The main challenge comes from the large production rate of hadronic jets, which can be misidentified as $`\tau _h`$. Using the selection described above, Figure 1 shows the efficiency for real hadronically decaying taus with $`|\eta |<1`$ to be reconstructed as $`\tau _h`$, using the simulation discussed below. The figure also shows the jet $`\tau _h`$ โfakeโ probability that hadronic jets are misidentified as hadronic tau decays. These jets, reconstructed in a cone size of $`\mathrm{\Delta }R=0.7`$, come from events recorded with triggers requiring various thresholds for calorimeter cluster energy.
To discriminate against background, for the $`e\tau _h`$ ($`\mu \tau _h`$) channel the electron (muon) must have a transverse energy of at least 10 GeV, the $`\tau _h`$ must have $`E_T>25`$ GeV, and the event must have $`/E_T>`$15 GeV. For the $`\tau _h\tau _h`$ channel, one $`\tau _h`$ must have $`E_T`$ greater than 25 GeV, and the other must have at least 10 GeV. The azimuthal angle between the $`/E_T`$ vector and the $`e`$ or $`\mu `$ (in $`e\tau _h`$ or $`\mu \tau _h`$ events) or the less-energetic of the two in $`\tau _h\tau _h`$ events must be less than 30.
For all events selected by the above cuts we calculate the โvisible massโ ($`m_{vis}`$) by adding the measured four-momenta of the two identified tau decay products in the event to the missing transverse energy four-momentum (for which the $`z`$ component is taken as zero), and then calculating the invariant mass of the sum. This quantity efficiently distinguishes between lower-mass production of tau pairs (mainly from $`Z`$ boson decays) and high-mass tau pairs from possible new massive resonant particle production.
The main source of events expected in the selected sample is DY production of $`Z/\gamma ^{}`$ decaying to lepton pair final states, and of these, tau pair production predominates. The production cross section times branching ratio to pairs of each charged lepton species for DY $`Z/\gamma ^{}`$ is assumed to be 250 pb ref-zxsec in the mass range 66-116 GeV/$`c^2`$. For the DY process and for the possible new physics processes discussed below, we simulate the production and decay using the PYTHIA 6.215 Monte Carlo program ref-pythia with CTEQ5L parton distribution functions (PDFโs) ref-CTEQ5L , with tau decays simulated by TAUOLA ref-tauola . Acceptance and resolution effects come from the full CDF II detector simulation.
The second largest source of events passing our selection criteria is hadronic jets which are misidentified as a $`\tau _h`$, for example from events with a $`W`$ boson decaying to a charged lepton and a neutrino plus a jet which passes the $`\tau _h`$ identification criteria. The estimated number of expected events comes from applying the jet $`\tau _h`$ โfakeโ rates to jets in events passing the trigger and other requirements, excluding the $`\tau _h`$ identification.
Various systematic uncertainties affect the predicted number of signal and background events. The largest is due to imperfect modeling of the tau identification efficiency. We perform a cross check of this efficiency using $`W\tau \nu `$ events recorded in the first 72 pb<sup>-1</sup>. Assuming a production cross section times branching ratio to $`\tau \nu `$ of 2688 pb ref-zxsec , this check yields a multiplicative factor of 0.97$`\pm `$0.10, which is incorporated into the acceptance calculation in the simulation. The 10% uncertainty in this factor, which affects each identified $`\tau _h`$ in the selected sample, includes the trigger efficiency uncertainty.
The uncertainties in the $`e`$ and $`\mu `$ identification and trigger efficiency, of 4% for $`e`$ and 5.5% for $`\mu `$, come from studies described elsewhere ref-thesis .
The jet $`\tau _h`$ fake background estimate has a 20% uncertainty reflecting the variation in the fake rate among the different trigger samples.
A 6% uncertainty due to imperfect modeling of the $`/E_T`$ comes from studies of transverse energy balancing in events with high energy jets recoiling against high energy photons.
Imperfect knowledge of the PDFโs leads to an 8% uncertainty in the DY and any new physics signal acceptances. The uncertainty is estimated from the variation of the acceptance using different PDF sets.
Table 1 summarizes the expected numbers of events by source for each channel, and shows the observed number of events in each search channel, for the control region dominated by $`Z`$ boson decay ($`m_{vis}<120`$ GeV$`/c^2`$). The observed number is in good agreement with that expected. This gives confidence that the estimated efficiencies and background rates are well understood, and we proceed to examine the signal region.
Table 2 shows, for each search channel, the numbers of events expected and the uncertainty for each background source in the signal region ($`m_{vis}>120`$ GeV/$`c^2`$). We observe four $`e\tau _h`$ events, and no $`\mu \tau _h`$ or $`\tau _h\tau _h`$ events. Given the uncertainties shown in the table, the observed number of events is in good agreement with that expected.
Figure 2 shows the distribution of visible mass in the signal and control regions, for the observed events and the predicted background. The distribution of the masses of the four events in the signal region is consistent with that expected from background.
Since we observe no significant excess rate of high-mass tau pair production, we determine upper bounds on the production cross section times branching ratio to tau pairs of hypothetical scalar and vector particles. As a general model for the acceptance for scalar particle production we use pseudoscalar Higgs boson ($`A`$) production, and for vector particle production we use a $`Z^{}`$ boson. The acceptance for both increases from near zero at masses of 100 GeV/$`c^2`$ to about 4% at high masses (500 GeV/$`c^2`$ or more).
To determine the upper bounds on the cross section times branching ratio we form a likelihood from the joint Poisson probability of all search channel results, and use a Bayesian method to incorporate the effects of systematic uncertainties, which are represented by truncated gaussian prior probability densities, including correlations. The likelihood is converted to a posterior probability density in the signal cross section using Bayes Theorem, assuming a prior in the signal rate which is uniform up to some high cutoff. The 95% CL upper limit comes from the integral of the posterior density.
Figure 3 shows the 95% CL upper bound on the cross section times branching ratio to tau pairs for scalar and vector particle production. Table 3 lists the upper limits on the production rate of scalar and vector particles as a function of mass. As an example of the sensitivity, these results would rule out a $`Z^{}`$ with Standard Model couplings having a mass less than 399 GeV/$`c^2`$, as indicated by the curve of cross section times branching ratio in the figure. The figure also shows the case of R-parity-violating scalar neutrino production and decay to tau pairs; this analysis, as an example, excludes a 377 GeV/$`c^2`$ scalar neutrino having coupling $`\lambda ^{}`$ to $`d\overline{d}`$ and branching ratio $`B`$ to tau pairs such that $`\lambda ^2B=`$0.01. In general the limits are readily interpreted within the context of new physics models in which new scalar or vector particles decay to tau lepton pairs.
We thank the Fermilab staff and the technical staffs of the participating institutions for their vital contributions. This work was supported by the U.S. Department of Energy and National Science Foundation; the Italian Istituto Nazionale di Fisica Nucleare; the Ministry of Education, Culture, Sports, Science and Technology of Japan; the Natural Sciences and Engineering Research Council of Canada; the National Science Council of the Republic of China; the Swiss National Science Foundation; the A.P. Sloan Foundation; the Bundesministerium fรผr Bildung und Forschung, Germany; the Korean Science and Engineering Foundation and the Korean Research Foundation; the Particle Physics and Astronomy Research Council and the Royal Society, UK; the Russian Foundation for Basic Research; the Comisiรณn Interministerial de Ciencia y Tecnologรญa, Spain; in part by the European Communityโs Human Potential Programme; and the Academy of Finland.
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# Lifetime of dynamic heterogeneity in strong and fragile kinetically constrained spin models
## 1 Introduction
The viscosity of supercooled liquids increases extremely rapidly when the temperature is reduced towards the glass temperature. It is firmly established that this dramatic slowing down is spatially heterogeneous. Local relaxation is widely distributed in time โ existence of broad stretched relaxations โ, but also in space โ existence of dynamic heterogeneity . The main physical aspect is that spatial fluctuations of local relaxations become increasingly spatially correlated when temperature decreases. Direct experimental investigations of the time and temperatures dependences of the relevant dynamic lengthscales at low temperature are however still missing.
To study dynamic heterogeneity, statistical correlators which probe more than two points in space and time have to be considered . For example, if one wants to study spatial correlations of the local dynamics one has to define a two-point, two-time correlator,
$$C_{2,2}(|ij|,t)=P_i(0,t)P_j(0,t)P_i(0,t)P_j(0,t),$$
(1)
where notations are adapted to lattice spin models. In equation (1), $`(i,j)`$ denote lattice sites, $`P_i(0,t)`$ quantifies the dynamics at site $`i`$ between times 0 and $`t`$ (autocorrelation or persistence functions), and brackets represent ensemble averages. The physical meaning of (1) is clear: given a spontaneous fluctuation of the two-time dynamics at site $`i`$, is there a similar fluctuation at site $`j`$? The quantity (1) has now been discussed both theoretically and numerically in some detail , generically revealing the existence of a growing spatial range of dynamic correlations in supercooled liquids accompanying an increasingly sluggish dynamics.
Logically, the next question is then: given spatial structures of the local relaxation between times 0 and $`t`$, what will this structure look like, say, between times $`t`$ and $`2t`$? In other words , how do dynamic heterogeneities evolve with time? This question is in fact simpler to address experimentally because no spatial resolution is needed and different experimental techniques can be devised: NMR , solvation dynamics , optical and dielectric hole-burning. In statistical terms, one wants to study a four-time correlation function of the general form
$$C_4(t_1,t_\mathrm{w},t_2)=P_i(0,t_1)P_i(t_1+t_\mathrm{w},t_1+t_\mathrm{w}+t_2),$$
(2)
which correlates dynamics between times 0 and $`t_1`$ and between times $`t_1+t_\mathrm{w}`$ and $`t_1+t_\mathrm{w}+t_2`$. Again the physical content of (2) is clear : given a dynamic fluctuation at site $`i`$ between in a certain time interval $`t_1`$, how long does it take for this fluctuation to be washed out? This leads to the general concept of a lifetime, $`\tau _{\mathrm{dh}}`$, for dynamic heterogeneity. While many investigations indicate that $`\tau _{\mathrm{dh}}`$ is in fact slaved to the alpha-relaxation time of the liquid, $`\tau _{\mathrm{dh}}\tau _\alpha `$, photobleaching experiments very close to the glass transition indicate that $`\tau _{\mathrm{dh}}`$ may become several orders of magnitude larger than $`\tau _\alpha `$, although in a surprisingly abrupt manner .
In this paper we study the lifetime of dynamic heterogeneity in kinetically constrained spin models of supercooled liquids . These models represent schematic coarse-grained models for the glass transition and provide a very efficient tool to study in detail many spatio-temporal aspects related to dynamic heterogeneity such as dynamic lengthscales , scaling , or decoupling phenomena . They are simple enough that analytical progress can be made and numerical simulations performed on a wide range of lengthscales and timescales, and yet rich enough that direct comparisons to both simulations and experiments can be made.
## 2 Models
Following previous works , we focus on two specific spin facilitated models in one spatial dimension, namely the one-spin facilitated Fredrickson-Andersen (FA) model and the East model that respectively behave as strong and fragile systems . These are probably the simplest models which incorporate the ideas that (i) mobility in supercooled liquids is both highly localized and sparse, as revealed by simulations ; (ii) a localized mobility very easily propagates to neighbouring regions, the dynamic facilitation concept. Detailed studies in spatial dimensions larger than one have shown that dimensionality does not play a relevant qualitative role , and justify therefore the present one-dimensional studies.
Both models are defined by the same non-interacting Hamiltonian, $`H=_in_i`$, expressed in terms of a mobility variable, $`n_i=1`$ when site $`i`$ is mobile, $`n_i=0`$ otherwise. Dynamic facilitation is incorporated at the level of the dynamic rules through kinetic constraints. In the FA model, the site $`i`$ can evolve with Boltzmann probability if at least one of its two neighbours is mobile, $`n_{i1}+n_{i+1}>0`$. In the East model the site $`i`$ can evolve only if its left neigbour is mobile, $`n_{i1}=1`$.
We have performed numerical simulations of both models using a continuous time Monte Carlo algorithm where all moves are accepted and the time is updated according to the corresponding statistical weight. Simulations have been performed โonlyโ over about 7 decades in time because extensive time averaging is required to accurately measure multi-time correlation functions such as equation (2).
## 3 Results
### 3.1 Dynamic filtering
There are several parameters involved in the four-time correlator (2) that need to be appropriately chosen. Since dynamic heterogeneity is more pronounced for times close to $`\tau _\alpha `$ it is sensible to first fix $`t_1=\tau _\alpha `$ and to study the remaining $`t_\mathrm{w}`$ and $`t_2`$ dependences. As a local dynamic correlator we first focus on the persistence function , $`P_i(0,t)=1`$ if spin $`i`$ has not flipped in the interval $`[0,t]`$, $`P_i(0,t)=0`$ otherwise. We also define the mean persistence, $`p(t)=P_i(t)`$, from which we measure $`\tau _\alpha `$ via $`p(\tau _\alpha )=e^1`$.
In figure 1 (left) we show the $`t_2`$ dependence of $`C_4(\tau _\alpha ,t_\mathrm{w},t_2)`$ for various $`t_\mathrm{w}`$ at $`T=0.3`$ in the FA model. We have normalized $`C_4`$ by $`p(\tau _\alpha )=e^1`$, its value at $`t_2=0`$. By definition, this function describes the persistence function in the interval $`[t_\mathrm{w}+\tau _\alpha ,t_\mathrm{w}+\tau _\alpha +t_2]`$ of those sites which had not flipped in the interval $`[0,\tau _\alpha ]`$, and were therefore slower than average. The first term in the correlator (2) plays the role of a dynamic filter , selecting a sub-population of sites which have an average dynamics different from the bulk. From earlier works studying the spatial correlator (1), it is known that those sites belong to compact clusters that represent the largest regions of space with no mobility defects at time 0 .
Immediately after filtering one expects therefore those slow regions to remain slow, as indeed observed in figure 1 for $`t_\mathrm{w}=0`$. When $`t_\mathrm{w}`$ increases, this selected population gradually forgets it was initially slow. When $`t_\mathrm{w}\mathrm{}`$, bulk dynamics is recovered,
$$\frac{C_4(t_1,t_\mathrm{w}\mathrm{},t_2)}{p(\tau _\alpha )}p(t_2),$$
(3)
as demonstrated by the full line in figure 1. In figure 1 (right) we also show the (logarithmic) distribution of relaxation times corresponding to the functions shown in the left panel, a representation sometimes preferred in experimental works . Both quantities are of course fully equivalent . It is clear from these figures that once a subset of sites has been dynamically selected the remaining relaxation is narrower than the bulk relaxation. In fact all persistence functions shown in figure 1 are well described by stretched exponentials. While $`\beta =1/2`$ is observed for the bulk dynamics, one finds $`\beta 0.83`$ at $`t_\mathrm{w}=0`$. Accordingly, distribution of relaxation times progressively broaden when $`t_\mathrm{w}`$ increases. These results are consistent with experimental observations.
In the FA model their interpretation is straightforward. Stretching in this model follows from an exponential distribution of distance between mobility defects . Dynamic filtering implies that this domain distribution is cut-off at small distance. Narrower lengthscale distributions directly imply narrower timescale distributions.
We have also investigated the effect of changing the โfilter efficiencyโ which in our case implies changing the duration of the filtering interval, $`[0,t_1]`$. While bulk distributions are found for $`t_1/\tau _\alpha 1`$ (weak filtering) distributions shift to larger times and become very narrow when $`t_1/\tau _\alpha `$ increases. In the following we work at constant filtering, $`t_1=\tau _\alpha `$.
### 3.2 โHomogeneousโ vs. โheterogeneousโ dynamics
The ability to select a sub-ensemble of sites that are slower than average is sometimes taken as a definition of dynamic heterogeneity , although lengthscales play no role in this view. That FA and East model display spatially heterogeneous dynamics is well-known, and the results of the previous section are therefore natural.
Another indicator of dynamic heterogeneity has been proposed based on the analysis of the four-time correlation (2). Consider the situation where $`t_\mathrm{w}=0`$, and $`t_1=t_2t/2`$. In that case, one studies a โthree-timeโ correlation
$$F_3(t)=P_i(0,t/2)P_i(t/2,t).$$
(4)
Two extreme behaviours can be expected for $`F_3(t)`$. (i) Dynamics in the intervals $`[0,t/2]`$ and $`[t/2,t]`$ are totally uncorrelated, and thus $`F_3(t)[p(t/2)]^2`$. (ii) Dynamics in the two intervals are strongly correlated, in the sense that those regions that survive filtering in $`[0,t/2]`$ are also those dominating the relaxation in the full interval $`[0,t]`$. In that case, $`F_3(t)p(t)`$. Scenarii (i) and (ii) have been termed โhomogeneousโ and โheterogeneousโ, respectively, although again lengthscales play no role in the distinction. Clearly, both estimates become equivalent when $`p(t)`$ decays exponentially.
Of course when studying the persistence function in the FA and East models, scenario (ii) strictly applies by definition, because $`P_i(0,t/2)P_i(t/2,t)=P_i(0,t)`$. In real materials, smoother dynamic functions are studied, directly defined from the particles positions instead of a mobility field. Our strategy is therefore to couple probe particles to our mobility field, see Refs. for technical details. From probe molecule displacements, $`\delta x(0,t)=x(t)x(0)`$, we define self-intermediate scattering functions, $`F_s(k,t)=\mathrm{cos}[k\delta x(0,t)]`$, and the analog of equation (4), $`F_3(k,t)=\mathrm{cos}[k\delta x(0,t/2)]\mathrm{cos}[k\delta x(t/2,t)]`$.
Our numerical results are presented in figure 2. Clearly the time dependence of $`F_3`$ closely follows the one of $`F_s(k,t)`$, in agreement with the โheterogeneousโ scenario described above. This is consistent with numerical results .
In the present approach, this result is a natural consequence of decoupling between structural relaxation and diffusion . At large wavevectors, $`k\pi `$, corresponding to distances of the order of the lattice spacing, $`F_s(k,t)`$ is dominated by the time distribution of the first jump of the probe molecule in the interval $`[0,t]`$, so that $`F_s(k,t)p(t)F_3(k,t)`$. At large distance, $`k<k^{}`$, Fickian diffusion holds , $`F_s(k,t)=\mathrm{exp}(k^2D_st)`$, and there is no distinction between homogeneous and heterogeneous relaxation. At intermediate wavectors, $`\pi >k>k^{}`$, the long-time decay of $`F_s(k,t)`$ is again dominated by the persistence time distribution, because the timescale it takes a molecule to make $`2\pi /k`$ steps is strongly dominated by the timescale to make the first step . This is just the condition for the heterogeneous scenario to hold, in agreement with figure 2. The characteristic wavevector separating the two regimes, $`k^{}(T)=1/\sqrt{\tau _\alpha D_s}`$, decreases when temperature decreases, opening a larger heterogeneous window; $`k^{}`$ also sets the upper limit of validity of Fickian diffusion in supercooled liquids .
### 3.3 Lifetime of dynamic heterogeneity
After dynamic filtering it takes some time for filtered distributions to reequilibrate towards the bulk relaxation, cf figure 1. To extract the typical lifetime of dynamic heterogeneity, $`\tau _{\mathrm{dh}}`$, we tried several procedures which all lead to similar results, based on how timescales (time decay of persistence functions or moments of the corresponding distributions) return to their equilibrium values. Following Refs. we also measured the integrated difference between filtered and bulk dynamics, $`\mathrm{\Delta }(t_\mathrm{w})_0^{\mathrm{}}๐t_2[C_4(\tau _\alpha ,t_\mathrm{w},t_2)/p(\tau _\alpha )p(t_2)]`$. From figure 1, we expect that $`\mathrm{\Delta }(t_\mathrm{w})`$ goes to 0 on a timescale $`\tau _{\mathrm{dh}}`$. In practice, we define $`\tau _{\mathrm{dh}}`$ as $`\mathrm{\Delta }(\tau _{\mathrm{dh}})/\mathrm{\Delta }(0)=e^1`$. In principle, $`\tau _{\mathrm{dh}}`$ depends on the filtering time, $`t_1`$, and on temperature, $`T`$.
We show in figure 3 results at various $`T`$ but constant filter efficiency, $`t_1=\tau _\alpha (T)`$, in East and FA models. While $`\tau _{\mathrm{dh}}`$ is set by $`\tau _\alpha `$ in the FA model (the tiny deviation observed in figure 3 is due to finite $`T`$ corrections which weaken when $`T`$ gets lower), this is not true in the East model where $`\tau _{\mathrm{dh}}`$ systematically grows faster than $`\tau _\alpha `$ at low $`T`$, as emphasized in the inset. Quantitatively a power law relationship, $`\tau _{\mathrm{dh}}\tau _\alpha ^{1+\zeta }`$, with $`\zeta 0.06`$, is a good description of the data, although alternative fitting formula could probably be used.
In the fragile case, $`\tau _{\mathrm{dh}}`$ can therefore be considered as an additional slow timescale characterizing the alpha-relaxation , on top of $`\tau _\alpha `$ and $`1/D_s`$ . The comparative study of FA and East models offers a possible physical interpretation. While both models display stretched relaxations, in the FA model stretching is constant, $`\beta =1/2`$, while $`\beta `$ increases linearly with $`T`$ in the East model . Therefore $`\tau _\alpha `$ represents the first moment of a distribution that becomes wider and wider when $`T`$ decreases. We attribute the small but systematic decoupling between $`\tau _{\mathrm{dh}}`$ and $`\tau _\alpha `$ to this broadening.
Unfortunately this decoupling does not quantitatively account for the results of photobleaching experiments which show that $`\tau _{\mathrm{dh}}/\tau _\alpha `$ increases strongly close to $`T_g`$ . In OTP, while $`\tau _\alpha `$ changes by about 1 decade when $`T`$ is changed from $`T_g+4`$ K to $`T_g+1`$ K, the ratio $`\tau _{\mathrm{dh}}/\tau _\alpha `$ changes by 2 orders of magnitude, so that $`\zeta 2`$. This value is much too large to be accounted for by the above results. Presumably, also, $`\beta `$ does not vary much on such a tiny temperature interval. Therefore, the present results cannot explain the experimental value $`\zeta 2`$ without invoking possible non-equilibrium effects due to the proximity of $`T_g`$. However we were able to predict instead a smaller, but definitely non-vanishing decoupling between the lifetime of dynamic heterogeneity and the alpha-relaxation time which could be detected in dynamic filtering experiments performed on a sufficiently large temperature window in fragile glass-formers.
## References
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# Introduction
## Introduction
Classical group theoretical analysis of differential equations whose foundations were laid by Sophus Lie over a hundred years ago finds increasing utilization in modern mathematical physics (see e.g. ). At the same time certain limits of the classical Lie approach become obvious that nevertheless do not allow full description of the symmetry of an equation under study . In particular, does not allow calculation of higher order symmetry operators that are widely used for calculation of reference frames admitting solution of equations in separated variables , in calculation of motion constants and in many other problems.
The present paper deals with investigation of non-Lie symmetry of the KleinโGordonโFock equation in $`(p+q)`$-dimensional Minkowsky space
$$L\phi \left(g^{\mu \nu }\frac{}{x_\mu }\frac{}{x_\nu }\varkappa ^2\right)\phi =0,$$
(1)
where $`\varkappa `$ is a real parameter,
$$g^{\mu \nu }=\{\begin{array}{cc}\hfill 0,& \mu \nu ,\hfill \\ \hfill 1,& \mu =\nu p,\hfill \\ \hfill 1,& p<\mu =\nu p+q,\hfill \end{array}$$
(5)
$`\phi =\phi (x_1,x_2,\mathrm{},x_{p+q})`$ is a function of $`p+q`$ variables.
A symmetry operator of equation (1) is an arbitrary operator $`Q`$ (linear, nonlinear, differential, integral) that transforms solutions of this equation into solutions, that is
$$L(Q\phi )=0,\text{if}L\phi =0$$
(6)
(see Section 1 below for more rigorous definition).
We will call a differential operator of a finite order $`n`$ being a symmetry operator of equation (1) a $`n`$-th order symmetry operator.
Description of the maximal (in the sense of Lie) symmetry of equation (1) may be reduced to finding of all linearly independent first order symmetry operators. Such operators are well-known, they form a basis of the Lie algebra of the generalized Poincarรฉ group $`P(p,q)`$ (for $`\varkappa 0`$) or for the conformal group in $`(p+q)`$-dimensional space (when $`\varkappa =0`$).
One of the main results of the present paper is calculation in explicit form of a complete set of $`n`$-th order symmetry operators of equation (1) for arbitrary $`n<\mathrm{}`$ and $`p+q4`$.
It is well-known that description of first order symmetry operators is based upon calculation of explicit form of the Killing vector that corresponds to the space of independent variables. We associate with higher order symmetry operators more complex fundamental objects that we call Killing tensors of rank $`j`$ and order $`s`$, with $`j,s=1,2,\mathrm{}`$ and conformal Killing tensors of rank $`j`$ and order $`s`$.
In this paper we give the definition of the mentioned tensors as a a complete set of linearly independent solutions of some overdetermined systems of PDE and find these tensors were found in explicit form for arbitrary fixed $`j`$ and $`s`$ in Minkowsky space of dimension $`p+q4`$. The results can be used in investigation of higher symmetries of a wide class of systems of partial differential equations of mathematical physics given in the same space, in particular, of relativistic and galilei-invariant wave equations.
Let us describe briefly arrangement of our presentation. Main definitions related to higher order symmetry operators, are adduced in Section 1, definition of Killing tensors of rank $`j`$ and order $`s`$ is given in Section 2, first order Killing tensors of rank $`j`$ and order $`s`$ in explicit form are found in Sections 3 and 4, conformal Killing tensors and Killing tensors of arbitrary rank and order are shown in Sections 6, 8 and 9. A complete set of $`n`$-th order symmetry operators for equation (1) with zero and non-zero โmassโ $`\varkappa `$ are adduced in Sections 5 and 7.
## 1 Symmetry operators of order $`๐`$
For the purpose of our study it is sufficient to consider only solutions of equation (1) defined on an open set $`D`$ of the four-dimensional manifold $`_{p+q}`$ consisting of points with co-ordinates $`(x_1,x_2,\mathrm{},x_{p+q})`$ and analytical with respect to real variables $`x_1,x_2,\mathrm{},x_{p+q}`$. The set of all such solutions forms a complex vector space that we designate by the symbol $`_0`$. Setting $`D`$ as fixed (e.g. assuming that $`D`$ coincides with $`_{p+q}`$), we will call $`_0`$ the set of solutions of equation (1).
Let us designate with $``$ a vector space of all complex-valued functions defined on $`D`$ and being real analytical, and with $`L`$ a linear differential operator (1) defined on $``$. Then $`L\psi `$ when $`\psi `$. At that $`_0`$ is such subspace of the vector space $``$ that coincides with zero-space (kernel) of the operator $`L`$.
Let $`๐_n`$ be a set (class) of differential operators of the order $`n`$ defined on $``$. Then a symmetry operator $`Q๐_n`$ of equation (1) is defined as follows.
Definition. A linear differential operator of order $`n`$
$$Q=\underset{i=0}{\overset{n}{}}Q_1,Q_i=H^{(a_1a_2\mathrm{}a_i)}\frac{^i}{x_{a_1}x_{a_2}\mathrm{}x_{a_i}},H^{(a_1a_2\mathrm{}a_i)}$$
(7)
is called a symmetry operator of equation (1) in the class $`๐_n`$ (or symmetry operator of order $`n`$) if
$$[Q,L]=\alpha _QL,\alpha _Q๐_{n1},$$
(8)
where $`[Q,L]=QLLQ`$ is the commutator of the operators $`Q`$ and $`L`$.
Relation (5) should be understood in the sense that operators in the right-hand and left-hand parts give the same acting on an arbitrary function $`\phi `$. Functions $`H^{a_1,a_2,\mathrm{},a_i}`$ to be determined are symmetric tensors of rank $`i`$. Hereinafter the parentheses enclose the set of symmetric indices.
It is easy to see that the relation (3) follows from (3) for each $`\phi _0`$. The reverse statement is also true: if an operator (4) satisfies the relation (3) for arbitrary $`\psi _0`$ then the condition (5) is satisfied for such operator with some operator $`\alpha _Q`$. In the case $`n=1`$ the symmetry operators defined above may be interpreted as generators of the symmetry group of the equation being considered . We can show that the set of symmetry operators $`Q๐_1`$ generates a Lie algebra, and corresponding finite transformations from the invariance group may be obtained by integration of the Lie equations .
Symmetry operators of order $`n>1`$ are not generators of a Lie algebra anymore and characterize generalized (non-Lie) symmetry of an equation under study. The problem of description of a complete set of $`n`$-th order for equation (1) be reduced to finding of the general solution of the operator equations (5).
## 2 Equations for coefficients of symmetry operators. <br>Killing tensors of rank $`๐`$ and order $`๐`$
For simplification of further calculations it is more convenient to present the operator $`Q`$ (4) as the sum of $`i`$-multiple anticommutators
$$Q=\underset{j=0}{\overset{n}{}}\widehat{Q}_j,$$
(9)
where
$$\widehat{Q}_j=[\mathrm{}[[F^{(a_1a_2\mathrm{}a_j)},\frac{}{x_{a_1}}]_+,\frac{}{x_{a_2}}]_+,\mathrm{}\frac{}{x_{a_j}}]_+,$$
(10)
$`[A,B]_+=AB+BA`$, $`F^{(a_1a_2\mathrm{}a_i)}`$ is a symmetric tensor of rank $`i`$. Expanding anticommutators and transferring differentiation operators to the righthand side, it is possible to reduce the expression (6) for the operator $`Q`$ to the form (4), and, vice versa, to write down any operator of the form (4) as (6).
We can use a similar representation for the operator $`\alpha _Q๐_{n1}`$
$$\alpha _Q=\underset{i=n2}{\overset{n1}{}}\widehat{\alpha }_i,\widehat{\alpha }_i=[\mathrm{}[[f^{a_1a_2\mathrm{}a_i},\frac{}{x_{a_1}}]_+,\frac{}{x_{a_2}}]_+,\mathrm{},\frac{}{x_{a_i}}]_+,$$
(11)
and write the product $`\alpha _QL`$ as
$$\alpha _QL\frac{1}{4}[[\alpha _Q,\frac{}{x_\mu }]_+,\frac{}{x^\mu }]_++\frac{1}{2}[\alpha _{Q_\mu },\frac{}{x_\mu }]_+,$$
(12)
where $`\alpha _{Q_\mu }=\frac{\alpha _Q}{x_\mu }`$.
Using the representations (6)โ(9) and taking into account that
$$[Q,L]=[Q_\mu ,\frac{}{x^\mu }]_+,Q_\mu =\frac{Q}{x_\mu },$$
(13)
it is possible to reduce the operator equation (5) to the system of equations for coefficients $`f^{a_1a_2\mathrm{}a_i}`$ and $`F^{a_1a_2\mathrm{}a_j}`$. In fact, substituting (6)โ(10) into (5) and putting equal coefficients at identical degrees of operators of differentiation, we obtain
$$^{(a_{j+1}}F^{a_1a_2\mathrm{}a_j)}=0,$$
(14)
$$f^{(a_1a_2\mathrm{}a_i)}0.$$
(15)
Here $`^{a_{j+1}}=\frac{}{x_{a_{j+1}}}`$, the round brackets contain symmetric indices (so symmetrization is implied (11)):
$$\frac{1}{j!}^{(a_{j+1}}F^{a_1a_2\mathrm{}a_j)}=^{a_{j+1}}F^{a_1a_2\mathrm{}a_j}+^{a_1}F^{a_{j+1}a_2\mathrm{}a_j}$$
$$+^{a_2}F^{a_1a_{j+1}\mathrm{}a_j}+\mathrm{}+^{a_j}F^{a_1a_2\mathrm{}a_{j+1}},$$
$`F^{a_1a_2\mathrm{}a_j}`$ is a symmetric tensor of the rank $`j`$.
If $`\varkappa =0`$, the equation for coefficients of the symmetry operator takes the following form:
$$^{(a_{j+1}}F^{a_1a_2\mathrm{}a_j)}=\delta ^{(a_ja_{j+1}}f^{a_1a_2\mathrm{}a_{j1})},$$
(16)
where $`F^{a_1a_2\mathrm{}a_j}`$ and $`f^{a_1a_2\mathrm{}a_{j1}}`$ are symmetric tensors with zero trace.
Convoluting equations (13) with respect to one pair of indices, we can eliminate the unknown functions $`f^{a_1a_2\mathrm{}a_{j1}}`$. As a result we get
$$^{(a_{j+1}}F^{a_1a_2\mathrm{}a_j)}\frac{j}{m+j1}^bF^{b(a_2a_3\mathrm{}a_j}g^{a_1a_{j+1})}=0,$$
(17)
$$f^{a_1a_2\mathrm{}a_{j1}}=\frac{j}{m+j1}^bF^{b(a_1a_2\mathrm{}a_{j1}}g^{a_ja_{j+1})},$$
(18)
where $`m=p+q`$ is dimension of the space of independent variables.
We see that the problem of description of symmetry operators of order $`n`$ for the equation (1) with $`\varkappa =0`$ appears to be equivalent to finding of the general solution of the system of partial differential equations given by the formula (11). This system is split with respect to the index $`j`$, as it splits into independent subsystems corresponding to $`j=0,1,\mathrm{},n`$. As it will be shown below, for complete description of the symmetry operators it is actually sufficient to solve only two such subsystems corresponding to $`j=n`$ and $`j=n1`$.
In the case $`j=1`$ the system (11) coincides with the Killing equations , and for $`j=2`$ it coincides with equations for the Killing tensor in the flat de Sitter space. The corresponding equations (14) determine conformal Killing vector and Killing tensor in the $`p+q`$-dimensional Minkowsky space.
We shall call functions $`F^{a_1a_2\mathrm{}a_j}`$ satisfying equations (11) (or (14)) Killing tensors (or conformal Killing tensors) or rank $`j`$ and order 1. The meaning of the term โorder 1โ (that we will omit sometimes) will be explained below.
The equations (11), (13) for a tensor of arbitrary rank were introduced (in the case $`j>2`$, without relation to any particular problem) in the paper . However, the general solution of these equations, as far as we are aware, was obtained in an explicit form only for $`j=1`$ and $`j=2`$ .
In the process of investigation of higher order symmetry operators admitted by systems of partial differential equations, we have to deal with more complicated equations for coefficients of such operators than those given by formulae (11) or (14). These equations include derivatives of the order $`s>1`$ and have the form
$$^{(a_{j+1}}^{a_{j+2}}\mathrm{}^{a_{j+s}}F^{a_1a_2\mathrm{}a_j)}=0,$$
(19)
where $`F^{a_1a_2\mathrm{}a_j}`$ is a symmetric tensor, and
$$\left[^{a_{j+1}}^{a_{j+2}}\mathrm{}^{a_{j+s}}\stackrel{~}{F}^{a_1a_2\mathrm{}a_j}\right]^{SL}=0,$$
(20)
where $`\stackrel{~}{F}^{a_1a_2\mathrm{}a_j}`$ is a symmetric tensor with zero trace, and the symbol $`[]^{SL}`$ designates the zero trace part of the tensor inside the square brackets (in our case it is a symmetric tensor of the rank $`R=j+s`$):
$$\left[G^{a_1a_2\mathrm{}a_R}\right]^{SL}=G^{a_1a_2\mathrm{}a_R}+\underset{d=1}{\overset{\{\frac{R}{2}\}}{}}(1)^dK_d\left(\underset{i=1}{\overset{d}{}}g^{a_{2i1}a_{2i}}\right)$$
$$\times F^{a_{2d+1}a_{2d+2\mathrm{}a_ib_1b_2b_3b_4\mathrm{}b_{2d1}b_{2d}}}{}_{g_{b_1b_2}g_{b_3b_4\mathrm{}}g_{b_{2d1}b_{2d}}}{}^{},$$
(21)
where $`\{\frac{R}{2}\}`$ is the integer part of the number $`\frac{R}{2}`$,
$$K_d=\frac{n!}{(n2d)!2^{d1}}\underset{i=1}{\overset{d}{}}\frac{1}{2(ni)+m2},$$
(22)
In the case $`s=1`$ the equations (16) and (17) can be reduced to equations (11) and (14) respectively.
We will call a symmetric tensor $`F^{a_1a_2\mathrm{}a_j}`$ satisfying equations (16) a Killing tensor of rank $`j`$ and order $`s`$.We will call a symmetric tensor $`\stackrel{~}{F}^{a_1a_2\mathrm{}a_j}`$ with zero trace satisfying equations (17) a conformal Killing tensor of rank $`j`$ and order $`s`$.
In Sections 3โ7 below we obtain the general solution of equations (11), (14) for arbitrary $`j`$ in the space of dimension $`p+q4`$. Equations (16), (17) are discussed in Sections 8, 9 where their general solution is found for $`p+q4`$ and arbitrary $`j`$ and $`s`$.
## 3 Reduction of equations for symmetry operators <br>to a system of linear algebraic equations
Let us start investigation of the system of equations (11) describing the Killing tensor of rank $`j`$ and order 1.
The system (11) may be written in the following symbolic form:
$$^{a_1a_2\mathrm{}a_{j+1}}=0,$$
(23)
where $`^{a_1a_2\mathrm{}a_{j+1}}`$ is a symmetric tensor of rank $`j+1`$ in $`m=p+q`$-dimensional space, and unknown functions are components of symmetric tensor of rank $`j`$ in $`m`$-dimensional space. Whence we can see that the system under investigation is overdetermined, including $`\left(\genfrac{}{}{0pt}{}{j+m}{j+1}\right)`$ equations for $`\left(\genfrac{}{}{0pt}{}{j+m1}{j}\right)`$ unknowns, $`\left(\genfrac{}{}{0pt}{}{b}{a}\right)=\frac{b!}{a!\left(ba\right)!}`$ designating binomial coefficients.
Following the general method for solving of overdetermined systems of partial differential equations , we consider the set of differential consequences of the system (11), obtained by differentiation of each term, $`k`$ times by $`x_{b_i}`$ $`(i=1,2,\mathrm{},k)`$. For each fixed $`k`$ such differential consequences are systems of linear homogeneous algebraic equations for derivatives
$$^{b_1}^{b_2}\mathrm{}^{b_k}^{a_{j+1}}F^{a_1a_2\mathrm{}a_j}F^{(a_1a_2\mathrm{}a_j,a_{j+1})b_1b_2\mathrm{}b_k}.$$
(24)
These systems have the form
$$F^{(a_1a_2\mathrm{}a_j,a_{j+1})b_1b_2\mathrm{}b_k}=0.$$
(25)
The system of equations (22) determines condition for vanishing of the tensor of rank $`j+k+1`$ symmetric with respect to $`j+1`$ indices $`a_1,a_2,\mathrm{},a_{j+1}`$ and with respect to $`k`$ indices $`b_1,b_2,\mathrm{},b_k`$, with unknown components of the tensor (21) of rank $`j+k+1`$ symmetric with respect to $`j`$ indices $`a_1,a_2,\mathrm{},a_j`$ and with respect to $`k+1`$ indices $`a_{j+1},b_1,\mathrm{},b_k`$. Whence we conclude that the corresponding numbers of equations $`(N_\text{e})`$ and of unknown variables $`(N_\text{u})`$ are given by the formulae
$$N_\text{e}=\left(\genfrac{}{}{0pt}{}{j+m}{j+1}\right)\left(\genfrac{}{}{0pt}{}{k+m1}{k}\right),N_\text{u}=\left(\genfrac{}{}{0pt}{}{j+m1}{j}\right)\left(\genfrac{}{}{0pt}{}{k+m}{k+1}\right),$$
(26)
where $`m=p+q`$ is dimension of the Minkowsky space where equations (11) are determined (that is the number of independent variables $`x_1,x_2,\mathrm{}`$ of the function $`F^{a_1a_2\mathrm{}a_j}`$).
According to (23)
$$N_\text{e}<N_\text{u},k<j,N_\text{e}=N_\text{u},k=j.$$
(27)
The formulae (23) allow calculation of the number of linearly independent solutions of equations (22), as the following statement is true:
Theorem 1. The system of linear algebraic equations (22) is not degenerate.
Proof of Theorem 1 is adduced below in Appendix.
We conclude from (24) in virtue of Theorem 1 that for $`k=j`$ the system of homogeneous linear algebraic equations (22) has only trivial solutions,
$$F^{a_1a_2\mathrm{}a_j,a_{j+1}b_1b_2\mathrm{}b_j}0.$$
Whence coefficients of the symmetry operator $`F^{a_1a_2\mathrm{}a_j}`$ are polynomials on $`x_a`$ $`(a=1,2,\mathrm{},m)`$ of order $`j`$. It follows from (23) that such polynomial contains $`N_j^m`$ arbitrary parameters, where
$$N_j^m=\underset{k=0}{\overset{j}{}}(N_\text{e}^kN_\text{u}^k)=\frac{1}{m}\left(\genfrac{}{}{0pt}{}{j+m1}{m1}\right)\left(\genfrac{}{}{0pt}{}{j+m}{m1}\right).$$
(28)
We see that equations (11) have $`N_j^m`$ linearly independent solutions that form a complete system. To find these solutions in explicit form it is necessary to find the general solution of the system of linear homogeneous equations (22) for arbitrary given $`j`$, $`m`$ and $`k<j`$, and then reconstruct polynomials $`F^{a_1a_2\mathrm{}a_j}`$ by found values of derivatives of the tensors $`F^{a_1\mathrm{}a_j,a_{j+1}b_1b_2\mathrm{}b_k}`$ (let us remind that indices after the comma designate derivatives with respect to the corresponding arguments). The general solution of equations (11) is adduced in Section 4 below.
## 4 Explicit form of Killing tensor of rank $`๐`$
According to the above proof, calculation of the explicit form of Killing tensor of rank $`j`$ is reduced to finding of the general solution of non-degenerate system of linear homogeneous algebraic equations given by the formula (22). Actual solution of this system with arbitrary given $`j`$ and $`m`$ is a rather difficult task that may be circumvented using the following observation.
Lemma 1. Let $`F^{a_1a_2\mathrm{}a_{j_0}}`$ be an arbitrary solution of the system (11) for $`j=j_0`$, and $`F^a`$ be a solution of the same system for $`j=1`$. Then the function
$$F^{a_1a_2\mathrm{}a_{j_0+1}}=F^{(a_1a_2\mathrm{}a_{j_0}}F^{a_{j_0+1})}$$
(29)
is a solution of the system (11) for $`j=j_0+1`$.
Proof is elementary and can be done by direct check.
Lemma 1 given an efficient algorithm for construction of solutions of equations (11). In fact, solutions of these equations for $`j=1`$ are well-known: they are Killing tensors , and a solution for arbitrary $`j`$ may be obtained from a solution for $`j=1`$ by successive application of the formula (26). If we manage to construct this way $`N_j^m`$ linearly independent solutions where $`N_j^m`$ is given by the formula (25), then such solutions form a complete system in virtue of Theorem 1.
Using the algorithm presented above we managed to obtain the general solution of equations (11) for $`m4`$ in the form
$$F^{a_1a_2\mathrm{}a_j}=g^{(a_{j1}a_j}F^{a_1a_2\mathrm{}a_{j2})}+f^{a_1a_2\mathrm{}a_j},$$
(30)
where $`F^{a_1a_2\mathrm{}a_{j2}}`$ is the general solution of equations (11) for $`j(j2)`$ depending on $`N_{j2}^m`$ arbitrary parameters, and $`f^{a_1a_2\mathrm{}a_j}`$ is a solution of equations (11) depending on $`N_j^mN_{j2}^m`$ arbitrary parameters.
The first addend in the right-hand part of the formula (27) corresponds to such symmetry operator (7) of order $`j`$ that on the set of solutions of equation (1) can be reduced to a symmetry operator of order $`j2`$. Explicit expressions for $`f^{a_1a_2\mathrm{}a_j}`$ corresponding to $`m4`$ are adduced below.
1. $`m=1`$. The corresponding tensor $`f^{a_1a_2\mathrm{}a_j}`$ can be reduced to a scalar not depending on the only variable.
2. $`m=2`$. Tensors $`f^{a_1a_2\mathrm{}a_j}`$ depend on two variables $`x_1`$ and $`x_2`$. The number of independent solutions, according to (25), is
$$N=N_j^2N_{j2}^2=2j+1.$$
(31)
Solutions are numbered by an integer number $`c`$ satisfying the condition
$$0cj,$$
(32)
and include for $`c=0`$ one, and for each $`c>0`$ two arbitrary parameters giving independent components of a symmetric zero trace tensor $`\lambda ^{a_1a_2\mathrm{}a_{jc}}`$ of rank $`jc`$. The explicit form of the corresponding solution $`f_c^{a_1a_2\mathrm{}a_j}`$ is given by the formula
$$f_c^{a_1a_2\mathrm{}a_j}=\epsilon \widehat{f}^{a_1a_2\mathrm{}a_j}+(1\epsilon )\widehat{f}^{(a_1a_2\mathrm{}a_{j1}}\epsilon ^{a^j)b}x_b,$$
(33)
where $`\epsilon ^{a^jb}`$ is the unit antisymmetric tensor, $`\epsilon =\frac{1}{2}[1+(1)^c]`$,
$$\widehat{f}^{a_1a_2\mathrm{}a_j}=\lambda ^{(a_1a_2\mathrm{}a_{jc}}\underset{\mu =0}{\overset{\{\frac{c}{2}\}}{}}\left(\underset{i=jc+1}{\overset{jc+2\mu }{}}x^{a_i}\right)^{}$$
$$\times \left(\underset{k=\{\frac{jc}{2}\}+\mu +1}{\overset{\mathrm{min}\{\frac{j}{2},\frac{j+1}{2}l\}}{}}g^{a_{2k+l}a_{2k})}\right)^{}(1)^\mu \left(\genfrac{}{}{0pt}{}{\{\frac{c}{2}\}}{\mu }\right)(x^2)^{\{\frac{c}{2}\}\mu },$$
$$\left(\underset{\lambda =A}{\overset{B}{}}f_\lambda \right)^{}=\{\begin{array}{cc}\underset{\lambda =A}{\overset{B}{}}f_\lambda ,\hfill & BA,\hfill \\ 1,\hfill & B<A,\hfill \end{array}x^2=x_1^2+x_2^2,l=(1)^{j+c+1},$$
(36)
and symmetrization over the indices $`a_1,a_2,\mathrm{},a_j`$ is implied.
3. $`m=3`$. The tensor $`f^{a_1a_2\mathrm{}a_j}`$ depends on three variables $`\stackrel{}{x}=(x_1,x_2,x_3)`$. The number of independent solutions is equal to
$$N=N_j^3N_{j2}^3=\frac{1}{3}(j+1)(2j^2+4j+3).$$
(37)
The solutions are numbered with pairs of integers $`c=(c_1,c_2)`$ satisfying the conditions
$$0c_12\left\{\frac{j}{2}\right\},0c_2ji\left\{\frac{c_1+1}{2}\right\},\epsilon _a=\frac{1}{2}[1+(1)^a],$$
(38)
and include for each $`c`$ the set $`2c_1+1`$ of arbitrary parameters giving independent components of a symmetric zero trace tensor $`\lambda ^{a_1a_2\mathrm{}a_{c_1}}`$ of rank $`c_1`$. Explicit forms of the corresponding solutions $`f_c^{a_1a_2\mathrm{}a_j}`$ are given by the formula
$$f_c^{a_1a_2\mathrm{}a_j}=\epsilon _{c_2}\widehat{f}_{c_1c_2}^{a_1a_2\mathrm{}a_j}+(1\epsilon _{c_2})\widehat{f}_{c_1c_2}^{b(a_1a_2\mathrm{}a_{j1}}\epsilon ^{a_j)bc}x_c,$$
(39)
where $`\epsilon ^{a^jbc}`$ is the unit antisymmetric tensor,
$$\widehat{f}_{c_1c_2}^{a_1a_2\mathrm{}a_j}=\underset{\mu }{}K_\mu \lambda _c^{\beta _\mu ,(A_\mu }\left(\underset{i=A_\mu +1}{\overset{A_\mu +L_\mu }{}}x^{a_i}\right)^{}\left(\underset{k=\{\frac{1}{2}(A_\mu +L_\mu )\}+1}{\overset{\mathrm{min}(\{\frac{j}{2}\},\{\frac{j+1}{2}\}l)}{}}g^{a_{2k}a_{2k+l}}\right)^{}(x^2)^{F_\mu }.$$
(40)
Here
$$\mu =(\mu _1,\mu _2,\mu _3,\mu _4,\mu _5),x^2=x_ax_bg^{ab},$$
$$K_\mu =(1)^{\mu _1+\mu _3+\mu _5}2\mu _3\frac{\{\frac{c_2}{2}\}!}{\mu _2!\mu _3!\mu _4!}\left(\genfrac{}{}{0pt}{}{\{\frac{c_1}{2}\}}{\mu _1}\right),$$
$$B_\mu =2\mu _2+\mu _3+\mu _5,A_\mu =jc_1B_\mu ,l=(1)^{c_2+j+1},$$
$$L_\mu =c_1+\mu _32\mu _1\mu _5,F_\mu =\mu _1+\mu _4,$$
(41)
and $`\lambda ^{B_\mu ,A_\mu }`$ is an arbitrary symmetric zero trace tensor of rank $`A_\mu +B_\mu `$ convoluted with $`B_\mu `$ vectors $`x_k`$:
$$\lambda ^{B_\mu ,A_\mu }=\lambda ^{b_1b_2\mathrm{}b_{B_\mu }a_1a_2\mathrm{}a_{A_\mu }}x_{b_1}x_{b_2}\mathrm{}x_{b_{B_\mu }}.$$
(42)
Summation in (35) is to be done over all possible nonnegative values of $`\mu `$ satisfying the conditions
$$0\mu _1\left\{\frac{c_1}{2}\right\},\mu _2+\mu _3+\mu _4=\left\{\frac{c_2}{2}\right\},0\mu _5\frac{1}{2}[1(1)^{c_1}].$$
(43)
As well as in the formula (31) symmetrization over the indices $`a_1,a_2,\mathrm{},a_j`$ in the right-hand part of (35) is implied.
4. $`m=4`$. The tensor $`f^{a_1a_2\mathrm{}a_j}`$ depends on four variables $`x=(x_1,x_2,x_3,x_4)`$. The number of independent solutions is equal to
$$N=N_j^4N_{j2}^4=\frac{1}{4!}(j+1)(j+2)(2j+3)(j^2+3j+4).$$
(44)
Solutions are numbered with triples of integers $`c=(c_1,c_2,c_3)`$ satisfying conditions (33) and (40):
$$0c_3j2\left\{\frac{c_1+1}{2}\right\}2c_2,$$
(45)
and include for each $`c`$ a set of $`N_c`$ arbitrary parameters, where
$$N_c=\{\begin{array}{cc}(c_2+2c_3+1)^2,\hfill & c_1=c_2+2c_3\hfill \\ 2(c_2+2c_3+1)(2c_1c_22c_3+1),\hfill & c_1c_2+2c_3,\hfill \end{array}$$
(48)
These parameters give independent components of an irreducible tensor
$$\lambda ^{a_1a_2\mathrm{}a_{R_1}[a_{R_1+1}b_1][a_{R_1+2}b_2]\mathrm{}[a_{R_1+R_2}b_{R_2}]},$$
where $`R_1=c_2+2c_3`$, $`R_2=c_1c_22c_3`$ (let us remind that an irreducible tensor of rank $`R_1+2R_2`$ has $`R_1`$ symmetric indices and $`R_2`$ symmetric pairs of antisymmetric indices, and convolution by any pair of indices and any triple of indices with completely antisymmetric tensor $`\epsilon _{\mu \nu \rho \sigma }`$ vanishes). Explicit expressions for the respective solutions are given by the formula (42):
$$f_c^{a_1a_2\mathrm{}a_j}=\underset{\mu }{}K_\mu \lambda ^{\beta _\mu ,(A_\mu ,D_c}\left(\underset{i=A_\mu +D_c+1}{\overset{A_\mu +D_c+L_\mu }{}}x^{a_i}\right)^{}$$
$$\times \left(\underset{k=\{\frac{1}{2}(A_\mu +D_c+L_\mu )\}+1}{\overset{\{\frac{j}{2}\}}{}}g^{a_{2k+l}a_{2k})}\right)^{}(x^2)^{F_\mu },$$
(49)
where $`\mu `$, $`x^2`$, $`K_\mu `$, $`B_\mu `$, $`F_\mu `$ are given by the formulae (36), (38) $`a,b=1,2,3,4`$, $`A_\mu =jc_1B_\mu D_c`$, $`D_c=c_3`$, $`l=(1)^{n+1}`$,
$$\lambda ^{B_\mu ,A_\mu ,D_c}=\lambda ^{b_1b_2\mathrm{}b_{B_\mu }a_1a_2\mathrm{}a_{A_\mu }[a_{A_\mu +1}d_1]\mathrm{}[a_{A_\mu +D_c}d_{D_c}]}$$
$$\times x_{b_1}x_{b_2}\mathrm{}x_{b_{B_\mu }}x_{d_1}x_{d_2}\mathrm{}x_{d_{D_c}},$$
(50)
symmetrization over the indices $`a_1,a_2,\mathrm{},a_j`$ is implied in the right-hand side of (42).
So, we have obtained the general solution of equations (11) in the space of dimension $`m4`$. One can verify by a direct check that the found solutions satisfy equations (11) and are linearly independent (it is not difficult to prove the latter considering $`i`$-fold convolutions of the found solutions with $`g^{kl},0i\{\frac{j}{2}\}`$). On the other side, these solutions form a complete system, as the number of arbitrary parameters they include is in compliance with the formula (25).
Let us also mention that we can present the general solution of equations (11) also in the form
$$F^{a_1a_2\mathrm{}a_j}=\underset{l=0}{\overset{j}{}}\lambda ^{a_1a_2\mathrm{}a_l[a_{l+1}b_1][a_{l+2}b_2]\mathrm{}[a_jb_{jl}]}x_{b_1}x_{b_2}\mathrm{}x_{b_{j_l}},$$
(51)
where $`\lambda ^{a_1a_2\mathrm{}a_l[a_{l+1}b_1]\mathrm{}[a_jb_{jl}]}`$ is a tensor, symmetric with respect to permutation of indices $`a_1,\mathrm{},a_j`$ and antisymmetric with respect to permutation of indices $`a_{l+j}`$ with $`b_i`$, $`1ijl`$, with vanishing convolution of this tensor over any three indices with $`\epsilon _{\mu \nu \rho \sigma }`$. The latter means that a cyclic permutation with respect to any triple of indices $`(a_k,a_{l+s},b_s)`$ gives zero, so the polynomial (44) admittedly satisfies equation (11). On the other side, the number of independent components of the tensor $`\lambda ^{a_1a_2\mathrm{}a_l[a_{l+1}b_1]\mathrm{}[a_jb_{jl}]}`$ for $`0lj`$ is exactly $`N_j^m`$ (25), so the formula (44) gives the general solution of equations (11). Decomposing tensors $`\lambda ^{a_1a_2\mathrm{}a_l[a_{l+1}b_1]\mathrm{}[a_jb_{jl}]}`$, $`0lj`$ into irreducible ones (that is having vanishing convolutions with respect to any pair of indices), we come to formulae (26)โ(43).
We formulate the above results as the following theorem.
Theorem 2. Equations (11) in a space of dimension $`m4`$ have $`N_j^m`$ linearly independent solutions. These solutions are polynomials of $`x_a`$ of degree $`j`$ and are given in explicit form by relations (26)โ(43).
The above theorem determines the explicit form of the Killing tensor of rank $`j`$ in a space of dimension $`m4`$.
## 5 Explicit form of symmetry operators $`๐ธ_๐`$ for $`๐\mathbf{}\mathrm{๐}`$
The above results allow presenting in explicit form of symmetry operators of order $`n`$ for equation (1) in $`m`$-dimensional space for arbitrary given $`n<\mathrm{}`$ and $`m4`$. For this purpose it is sufficient to look through all admissible values of $`c`$ given by the formulae (29), (33), (40) and construct in accordance to the formulae (26)โ(43) the corresponding expressions for Killing tensors of rank $`j`$ $`F_c^{a_1a_2\mathrm{}a_j}`$ (following (27), it is sufficient to restrict oneself with construction of $`f^{a_1a_2\mathrm{}a_j}`$), then substitute obtained expressions into (6), (7) and sum up over $`j`$ from $`0`$ to $`n`$.
In this Section we will realize this program for all $`n4`$ and $`m4`$, and write down in explicit form the corresponding symmetry operators.
Let us calculate the number of linearly independent symmetry operators of order $`n`$. It is equal, according to (27) to the number of linearly independent solutions of the system (11) for $`j=n,n1`$ or (see (25))
$$N(n,m)=N_m^n+N_m^{n1}=\frac{2n^2+2mn+m(m1)}{m(m1)}\left(\genfrac{}{}{0pt}{}{n+m2}{m2}\right)\left(\genfrac{}{}{0pt}{}{n+m1}{m2}\right),$$
(52)
In particular, for $`m=2,3,4`$
$$N(n,2)=(n+1)^2,$$
$$N(n,3)=\frac{1}{6}(n+1)(n+2)(n^2+3n+3),$$
$$N(n,4)=\frac{1}{72}(n+1)(n+2)^2(n+3)(n^2+4n+6).$$
(53)
Values of these numbers for $`n=1,2,3,4`$ are adduced in Table 1.
Table 1. Number of symmetry operators of order $`n`$
for equation (1) in $`m`$-dimensional space.
$`m\backslash n`$ $`1`$ $`2`$ $`3`$ $`4`$ $`2`$ $`4`$ $`9`$ $`16`$ $`25`$ $`3`$ $`7`$ $`26`$ $`70`$ $`155`$ $`4`$ $`11`$ $`60`$ $`225`$ $`665`$
Let us write down explicitly the corresponding solutions $`F_m^{(n)}=F_m^{a_1\mathrm{}a_n}`$ of equations (11)
$`m=1`$
$$n=1,2,3,4,F_1^{(n)}=\lambda _{1n};$$
$`m=2`$
$$n=0,F_2^{(0)}=\lambda _{20};$$
$$n=1,F_2^a=\lambda _{21}^a+\lambda _{21}\epsilon ^{ab}x_b;$$
$$n=2,F_2^{a_1a_2}=g^{a_1a_2}F_2^{(0)}+\lambda ^{a_1a_2}+\lambda ^{(a_1}\epsilon ^{a_2)^b}x_b+\lambda (g^{a_1a_2}x^2x^{a_1}x^{a_2});$$
(54)
$$n=3,F_2^{a_1a_2a_3}=g^{(a_1a_2}F_2^{a_3)}+\lambda _{(0,0)}^{a_1a_2a_3}+\lambda _{(0,1)}^{(a_1a_2}\epsilon ^{a_3)^b}x_b$$
$$+\lambda _{(0,2)}^{(a_1}(g^{a_2a_3)}x^2x^{a_2}x^{a_3)})+\lambda _{(0,3)}x^{(a_1}x^{a_2}\epsilon ^{a_3)^b}x_b;$$
(55)
$$n=4,F_2^{a_1a_2a_3a_4}=g^{(a_1a_2}F_2^{a_3a_4)}+\lambda _{(0,0)}^{a_1a_2a_3a_4}+\lambda _{(0,1)}^{(a_1a_2a_3}\epsilon ^{a_4)^b}x_b$$
$$+\lambda _{(0,2)}^{(a_1a_2}(x^{a_3}x^{a_4)}g^{a_3a_4)}x^2)+\lambda _{(0,4)}(x^{(a_1}x^{a_2}x^{a_3}x^{a_4)}$$
$$2x^{(a_1}x^{a_2}g^{a_3a_4)}+x^4g^{(a_1a_2}g^{a_3a_4)});$$
(56)
$`m=3`$
$$n=0,F_3^{(0)}=\lambda _3;$$
$$n=1,F_3^a=\lambda ^a+\lambda ^b\epsilon ^{abc}x_c;$$
$$n=2,F_3^{a_1a_2}=g^{a_1a_2}F_3^{(0)}+\lambda _{(0,0)}^{a_1a_2}+\lambda _{(0,1)}^{b(a_1}\epsilon ^{a_2)bc}x_c+\lambda _{(0,2)}^{a_1a_2}x^22\lambda _{(0,2)}^{b_1(a_1}x^{a_2)}x_{b_1}$$
$$+\lambda _{(0,2)}^{b_1b_2}x_{b_1}x_{b_2}g^{a_1a_2}+\lambda _{(1,0)}^{(a_1}x^{a_2)}\lambda _{(1,0)}^{b_1}x_{b_1}g^{a_1a_2}+\lambda _{(2,0)}x^{a_1}x^{a_2}$$
$$\lambda _{(2,0)}g^{a_1a_2}x^2;$$
(57)
$$n=3,F_3^{a_1a_2a_3}=g^{(a_1a_2}F_3^{a_3)}+\lambda _{(0,0)}^{a_1a_2a_3}+\lambda _{(0,1)}^{b(a_1a_2}\epsilon ^{a_3)bc}x_c+\lambda _{(0,2)}^{a_1a_2a_3}x^22\lambda _{(0,2)}^{b(a_1a_2}x^{a_3)}x_b$$
$$+\lambda _{(0,2)}^{b_1b_2(a_1}g^{a_2a_3)}x_{b_1}x_{b_2}+\lambda _{(0,3)}^{b(a_1a_2}\epsilon ^{a_3)bc}x_cx^22\lambda _{(0,3)}^{bd(a_1}x^{a_2}\epsilon ^{a_3)bc}x_cx_d$$
$$+\lambda _{(0,3)}^{db_1b_2}g^{(a_1a_2}\epsilon ^{a_3)dc}x_cx_{b_1}x_{b_2}+\lambda _{(1,0)}^{(a_1a_2}x^{a_3)}\lambda _{(1,0)}^{b(a_1}g^{a_2a_3)}x_b+\lambda _{(1,1)}^{b(a_1}x^{a_2}\epsilon ^{a_3)bc}x_c$$
$$+\lambda _{(2,0)}^{(a_1}x^{a_2}x^{a_3)}\lambda _{(1,1)}^{bd}g^{(a_1a_2}\epsilon ^{a_3)bc}x_cx_d+\lambda _{(2,1)}^bx^{(a_1}x^{a_2}\epsilon ^{a_3)bc}x_c$$
$$\lambda _{(2,0)}^{(a_1}g^{a_2a_3)}x^2\lambda _{(2,1)}^bg^{(a_1a_2}\epsilon ^{a_3)bc}x_cx^2;$$
$$n=4,F_3^{a_1a_2a_3a_4}=g^{(a_1a_2}F_3^{a_3a_4)}+\lambda _{(0,0)}^{a_1a_2a_3a_4}+\lambda _{(0,2)}^{a_1a_2a_3a_4}x^2+\lambda _{(0,1)}^{b(a_1a_2a_3}\epsilon ^{a_4)bc}x_c$$
$$2\lambda _{(0,2)}^{b(a_1a_2a_3}x^{a_4)}x_b+\lambda _{(0,2)}^{b_1b_2(a_1a_2}g^{a_3a_4)}x_{b_1}x_{b_2}+\lambda _{(0,3)}^{b(a_1a_2a_3}\epsilon ^{a_4)bc}x_cx^2$$
$$2\lambda _{(0,3)}^{bd(a_1a_2}x^{a_3}\epsilon ^{a_4)bc}x_cx_d+\lambda _{(0,4)}^{a_1a_2a_3a_4}x^4+\lambda _{(0,3)}^{db_1b_2(a_1}g^{a_3a_2}\epsilon ^{a_4)dc}x_cx_{b_1}x_{b_2}$$
$$4\lambda _{(0,4)}^{b(a_1a_2a_3}x^{a_4)}x_bx^2+4\lambda _{(0,4)}^{b_1b_2(a_1a_2}x^{a_3}x^{a_4)}x_{b_1}x_{b_2}+2\lambda _{(0,4)}^{b_1b_2(a_1a_2}g^{a_3a_4)}x_{b_1}x_{b_2}x^2$$
$$4\lambda _{(0,4)}^{b_1b_2b_3(a_1}x^{a_2}g^{a_3a_4)}x_{b_1}x_{b_2}x_{b_3}+\lambda _{(1,0)}^{(a_1a_2a_3}x^{a_4)}$$
$$+\lambda _{(0,4)}^{b_1b_2b_3b_4}x_{b_1}x_{b_2}x_{b_3}x_{b_4}g^{(a_1a_2}g^{a_3a_4)}\lambda _{(1,0)}^{b(a_1a_2}g^{a_3a_4)}x_b+\lambda _{(1,1)}^{b(a_1a_2}x^{a_3}\epsilon ^{a_4)bc}x_c$$
$$\lambda _{(1,1)}^{bd(a_1}g^{a_2a_3}\epsilon ^{a_4)bc}x_cx_d+\lambda _{(1,2)}^{(a_1a_2a_3}x^{a_4)}x^2\lambda _{(1,2)}^{b(a_1a_2}g^{a_3a_4)}x_bx^2$$
$$2\lambda _{(1,2)}^{b(a_1a_2}x^{a_3}x^{a_4)}x_b+3\lambda _{(1,2)}^{b_1b_2(a_1}x^{a_2}g^{a_3a_4)}x_{b_1}x_{b_2}+\lambda _{(2,0)}^{(a_1a_2}x^{a_3}x^{a_4)}$$
$$\lambda _{(1,2)}^{b_1b_2b_3}x_{b_1}x_{b_2}x_{b_3}g^{(a_1a_2}g^{a_3a_4)}\lambda _{(2,0)}^{(a_1a_2}g^{a_3a_4)}x^2+\lambda _{(2,1)}^{b(a_1}x^{a_2}x^{a_3}\epsilon ^{a_4)bc}x_c$$
$$\lambda _{(2,1)}^{b(a_1}g^{a_2a_3}\epsilon ^{a_4)bc}x_cx^2+\lambda _{(2,2)}^{(a_1a_2}x^{a_3}x^{a_4)}x^22\lambda _{(2,2)}^{b(a_1}x^{a_2}x^{a_3}x^{a_4)}x_b$$
$$+\lambda _{(2,2)}^{b_1b_2}x_{b_1}x_{b_2}x^{(a_1}x^{a_2}g^{a_3a_4)}\lambda _{(2,2)}^{(a_1a_2}g^{a_3a_4)}x^4+2\lambda _{(2,2)}^{b(a_1}x^{a_2}g^{a_3a_4)}x_bx^2$$
$$+\lambda _{(3,0)}^{(a_1}x^{a_2}x^{a_3}x^{a_4)}\lambda _{(2,2)}^{b_1b_2}x_{b_1}x_{b_2}x^2g^{(a_1a_2}g^{a_3a_4)}\lambda _{(3,0)}^{(a_1}x^{a_2}g^{a_3a_4)}x^2$$
$$\lambda _{(3,0)}^{b_1}x_{b_1}x^{(a_1}x^{a_2}g^{a_3a_4)}+\lambda _{(3,0)}^{b_1}x_{b_1}x^2g^{(a_1a_2}g^{a_3a_4)}+\lambda _{(4,0)}x^{(a_1}x^{a_2}x^{a_3}x^{a_4)}$$
$$2\lambda _{(4,0)}x^{(a_1}x^{a_2}g^{a_3a_4)}x^2+\lambda _{(4,0)}x^4g^{(a_1a_2}g^{a_3a_4)};$$
(58)
$`m=4`$
$$n=0,F_4^{(0)}=\lambda ;$$
$$n=1,F_4^a=\lambda _{(0,0,0)}^a+\lambda _{(0,0,1)}^{[ad_1]}x_{d_1};$$
$$n=2,F_4^{a_1a_2}=g^{a_1a_2}F_4^{(0)}+\lambda _{(0,0,0)}^{a_1a_2}+\lambda _{(0,0,1)}^{(a_1[a_2d_1])}x_{d_1}+\lambda _{(0,0,2)}^{[a_1d_1][a_2d_2]}x_{d_1}x_{d_2}$$
$$+\lambda _{(0,1,0)}^{a_1a_2}x^22\lambda _{(0,1,0)}^{b(a_1}x^{a_2)}x_b+\lambda _{(0,1,0)}^{b_1b_2}x_{b_1}x_{b_2}g^{a_1a_2}+\lambda _{(1,0,0)}^{(a_1}x^{a_2)}$$
$$\lambda _{(1,0,0)}^bx_bg^{a_1a_2}+\lambda _{(2,0,0)}x^{a_1}x^{a_2}\lambda _{(2,0,0)}g^{a_1a_2}x^2;$$
(59)
$$n=3,F_4^{a_1a_2a_3}=g^{(a_1a_2}F_4^{a_3)}+\lambda _{(0,0,0)}^{a_1a_2a_3}+\lambda _{(0,0,1)}^{(a_1a_2[a_3d_1])}x_{d_1}+\lambda _{(0,0,2)}^{(a_1[a_2d_1][a_3d_2])}x_{d_1}x_{d_2}$$
$$+\lambda _{(0,0,3)}^{[a_1d_1][a_2d_2][a_3d_3]}x_{d_1}x_{d_2}x_{d_3}2\lambda _{(0,1,0)}^{b_1(a_1a_2}x^{a_3)}x_{b_1}+\lambda _{(0,1,0)}^{a_1a_2a_3}x^2$$
$$+\lambda _{(0,1,0)}^{b_1b_2(a_1}g^{a_2a_3)}x_{b_1}x_{b_2}+\lambda _{(0,1,1)}^{(a_1a_2[a_3d_1])}x^2x_{d_1}2\lambda _{(0,1,1)}^{b_1(a_1[a_2d_1]}x^{a_3)}x_{b_1}x_{d_1}$$
$$+\lambda _{(0,1,1)}^{b_1b_2([a_1d_1]}g^{a_2a_3)}x_{b_1}x_{b_2}x_{d_1}+\lambda _{(1,0,0)}^{(a_1a_2}x^{a_3)}\lambda _{(1,0,0)}^{b_1(a_1}g^{a_2a_3)}x_{b_1}+\lambda _{(1,0,1)}^{(a_1[a_2d_1]}x^{a_3)}x_{d_1}$$
$$\lambda _{(1,0,1)}^{b_1([a_1d_1]}g^{a_2a_3)}x_{b_1}x_{d_1}+\lambda _{(2,0,0)}^{(a_1}x^{a_2}x^{a_3)}\lambda _{(2,0,0)}^{(a_1}g^{a_2a_3)}x^2+\lambda _{(2,0,1)}^{([a_1d_1]}x^{a_2}x^{a_3)}x_{d_1}$$
$$\lambda _{(2,0,1)}^{([a_1d_1]}g^{a_2a_3)}x^2x_{d_1};$$
(60)
$$n=4,F_4^{a_1a_2a_3a_4}=g^{(a_1a_2}F^{a_3a_4)}+\lambda _{(0,0,0)}^{a_1a_2a_3a_4}+\lambda _{(0,0,1)}^{(a_1a_2a_3[a_4d_1])}x_{d_1}$$
$$+\lambda _{(0,0,2)}^{(a_1a_2[a_3d_1][a_4d_2])}x_{d_1}x_{d_2}+\lambda _{(0,0,3)}^{(a_1[a_2d_1][a_3d_2][a_4d_3])}x_{d_1}x_{d_2}x_{d_3}$$
$$+\lambda ^{([a_1d_1][a_2d_2][a_3d_3][a_4d_4])}x_{d_1}x_{d_2}x_{d_3}x_{d_4}+\lambda _{(0,1,0)}^{a_1a_2a_3a_4}x^22\lambda _{(0,1,0)}^{b_1(a_1a_2a_3}x^{a_4)}x_{b_1}$$
$$+\lambda _{(0,1,0)}^{b_1b_2(a_1a_2}g^{a_3a_4)}x_{b_1}x_{b_2}+\lambda _{(0,1,1)}^{(a_1a_2a_3[a_4d_1])}x^2x_{d_1}2\lambda _{(0,1,1)}^{b_1(a_1a_2[a_3d_1]}x^{a_4)}x_{b_1}x_{d_1}$$
$$+\lambda _{(0,1,1)}^{b_1b_2(a_1[a_2d_1]}g^{a_3a_4)}x_{b_1}x_{b_2}x_{d_1}+\lambda _{(0,1,2)}^{(a_1a_2[a_3d_1][a_4d_2])}x^2x_{d_1}x_{d_2}$$
$$2\lambda ^{b_1(a_1[a_2d_1][a_3d_2]}x^{a_4)}x_{b_1}x_{d_1}x_{d_2}+\lambda ^{b_1b_2([a_1d_1][a_2d_2]}g^{a_3a_4)}x_{b_1}x_{b_2}x_{d_1}x_{d_2}$$
$$+\lambda _{(0,2,0)}^{a_1a_2a_3a_4}x^44\lambda _{(0,2,0)}^{b_1(a_1a_2a_3}x^{a_4)}x^2x_{b_1}+4\lambda _{(0,2,0)}^{b_1b_2(a_1a_2}x^{a_3}x^{a_4)}x_{b_1}x_{b_2}$$
$$+2\lambda _{(0,2,0)}^{b_1b_2(a_1a_2}g^{a_3a_4)}x^2x_{b_1}x_{b_2}4\lambda _{(0,2,0)}^{b_1b_2b_3(a_1}x^{a_2}g^{a_3a_4)}x_{b_1}x_{b_2}x_{b_3}$$
$$+g^{(a_1a_2}g^{a_3a_4)}\lambda _{(0,2,0)}^{b_1b_2b_3b_4}x_{b_1}x_{b_2}x_{b_3}x_{b_4}+\lambda _{(1,0,0)}^{a_1a_2a_3}x^{a_4)}\lambda _{(1,0,0)}^{b_1a_1a_2}g^{a_3a_4)}x_{b_1}$$
$$+\lambda _{(1,0,1)}^{(a_1a_2[a_3d_1]}x^{a_4)}x_{d_1}\lambda _{(1,0,1)}^{b_1(a_1[a_2d_2]}g^{a_3a_4)}x_{b_1}x_{d_1}+\lambda _{(1,0,2)}^{(a_1[a_2d_1][a_3d_2]}x^{a_4)}x_{d_1}x_{d_2}$$
$$\lambda _{(1,0,2)}^{b_1(a_1d_1][a_2d_2]}g^{a_3a_4)}x_{b_1}x_{d_1}x_{d_2}+\lambda _{(1,1,0)}^{(a_1a_2a_3}x^{a_4)}x^2\lambda _{(1,1,0)}^{b_1(a_1a_2}g^{a_3a_4)}x_{b_1}x^2$$
$$2\lambda _{(1,1,0)}^{b_1(a_1a_2}x^{a_4}x^{a_4)}x_{b_1}+3\lambda _{(1,1,0)}^{b_1b_2(a_1}x^{a_2}g^{a_3a_4)}x_{b_1}x_{b_2}\lambda _{(1,1,0)}^{b_1b_2b_3}g^{(a_1a_2}g^{a_3a_4)}x_{b_1}x_{b_2}x_{b_3}$$
$$+\lambda _{(2,0,0)}^{(a_1a_2}x^{a_3}x^{a_4}\lambda _{(2,0,0)}^{(a_1a_2}g^{a_3a_4)}x^2+\lambda _{(2,0,1)}^{(a_1[a_2d_1]}x^{a_3}x^{a_4)}x_{d_1}\lambda _{(2,0,1)}^{(a_1[a_2d_1]}g^{a_3a_4)}x^2x_{d_1}$$
$$+\lambda _{(2,0,2)}^{([a_1d_1][a_2d_2]}x^{a_3}x^{a_4)}x_{d_1}x_{d_2}\lambda _{(2,0,2)}^{([a_1d_1][a_2d_2]}g^{a_3a_4)}x^2x_{d_1}x_{d_2}+\lambda _{(2,1,0)}^{(a_1a_2}x^{a_3}x^{a_4)}x^2$$
$$2\lambda _{(2,1,0)}^{(b_1(a_1}x^{a_2}x^{a_3}x^{a_4)}x_{b_1}+\lambda _{(2,1,0)}^{b_1b_2}x^{(a_1}x^{a_2}g^{a_3a_4)}x_{b_1}x_{b_2}\lambda _{(2,1,0)}^{(a_1a_2}g^{a_3a_4)}x^4$$
$$+2\lambda _{(2,1,0)}^{b_1(a_1}x^{a_2}g^{a_3a_4)}x^2x_{b_1}\lambda _{(2,1,0)}^{b_1b_2}g^{(a_1a_2)}g^{a_3a_4)}x^2x_{b_1}x_{b_2}+\lambda _{(3,0,0)}^{(a_1}x^{a_2}x^{a_3}x^{a_4)}$$
$$\lambda _{(3,0,0)}^{b_1}x^{(a_1}x^{a_2}g^{a_3a_4)}x_{b_1}\lambda _{(3,0,0)}^{(a_1}x^{a_2}g^{a_3a_4}x^2+\lambda _{(3,0,0)}^{b_1}g^{(a_1a_2}g^{a_3a_4)}x^2x_{b_1}$$
$$+\lambda _{(4,0,0)}x^{a_1}x^{a_2}x^{a_3}x^{a_4}2\lambda _{(4,0,0)}x^{(a_1}x^{a_2}g^{a_3a_4)}x^2+\lambda _{(4,0,0)}g^{(a_1a_2}g^{a_3a_4)}x^4.$$
(61)
Substituting (47)โ(54) into (6), (7) and carrying differentiation operators to the right, we obtain explicit form of the corresponding symmetry operators. For $`n=1`$ we have a complete set of symmetry operators of the following form:
$$Q_1^a=P_a=i\frac{}{x^a},Q_1^{ab}=J_{ab}=x_aP_bx_bP_a.$$
(62)
We do not adduce explicit form of symmetry operators for $`n>1`$ because of the corresponding formulae being extremely cumbersome (in fact, these expressions are given by relations (6), (7), (47)โ(54)).
First order symmetry operators $`Q_1`$ adduced in (55), form a Lie algebra $`AP(p,q)`$, satisfying the following commutation relations:
$$[P_a,P_b]=0,[P_a,J_{bc}]=i(g_{ab}P_cg_{ac}P_b),$$
$$[J_{ab},J_{cd}]=i(g_{ac}J_{bd}+g_{bd}J_{ac}g_{ac}J_{bd}g_{bd}J_{ac}).$$
(63)
It is easy to notice using the representation (36) that symmetry operators of arbitrary order are polynomials of the operators (55). In other words, all symmetry operators of finite order of equation (1) belong to the enveloping algebra of the algebra $`AP(p,q)`$.
## 6 Explicit form of Killing tensor of arbitrary rank $`๐`$
Calculation of conformal Killing tensors of rank $`j`$ (that is construction of the general solution of equation (14)) may be done similarly to what was presented above in Sections 3โ5. Construction of such solution simplifies utilization of the result formulated in the following lemma.
Lemma 2. Let $`F^{a_1a_2\mathrm{}a_{j_0}}`$ be an arbitrary solution of the system (14) for $`j=j_0`$, and $`F^a`$ be a solution of this system for $`j=1`$. Then the function
$$F^{a_1a_2\mathrm{}a_{j_0+1}}=[F^{(a_1a_2\mathrm{}a_{j_0}}F^{a_{j_0}+1)}]^{SL},$$
(64)
where $`[]^{SL}`$ means the traceless part of the tensor in the square brackets (see (18)) is a solution of equations (14) for $`j=j_0+1`$.
Proof can be done by a direct check.
We adduce below without proof the general solution of equations (14)for $`m4`$ and arbitrary $`j`$.
By means of the reasoning similar to that in Section 3 we can show that in the two-dimensional space equations (14) are reduced to CauchyโRiemann equations, and corresponding symmetry operators are determined up to arbitrary analytical functions determining independent components of a symmetric traceless tensor $`F^{a_1a_2\mathrm{}a_j}`$ (there are two such components for $`j0`$ and one for $`j=0`$ (that is for the case when then tensor $`F^{a_1a_2\mathrm{}a_j}`$ is reduced to a scalar).
For $`m=3`$ the number of independent solutions of equations (14) is equal to
$$N_j^3=\frac{1}{3}(j+1)(2j+1)(2j+3).$$
(65)
Solutions are numbered by the pair of integers $`c=(c_1,c_2)`$ satisfying the conditions
$$0c_1j,0c_22c_1,$$
(66)
and for each $`c_1`$ contain $`(2c_1+1)`$ arbitrary parameters giving independent components of symmetric traceless tensor $`\lambda ^{a_1a_2\mathrm{}a_{c_1}}`$ of rank $`c_1`$. Explicit form of the corresponding solutions is given by the formula
$$F_{(c_1c_2)}^{a_1a_2\mathrm{}a_j}=\left[\epsilon _{c_2}f_{(c_1c_2)}^{a_1a_2\mathrm{}a_j}+(1\epsilon _{c_2})f_{(c_1c_2)}^{b(a_1a_2\mathrm{}a_{j1}}\epsilon ^{a_j)bc}x_c\right]^{SL},$$
(67)
where
$$f_{(c_1c_2)}^{a_1a_2\mathrm{}a_j}=\underset{m=0}{\overset{\{\frac{c_2}{2}\}}{}}(2)^m\left(\genfrac{}{}{0pt}{}{\{\frac{c_2}{2}\}}{m}\right)\lambda _{(c_1c_2)}^{b_1b_2\mathrm{}b_m(a_1a_2\mathrm{}a_{c_1m}}$$
$$\times x^{a_{c_1m+1}}x^{a_{c_1m+2}}\mathrm{}x^{a_j)}x_{b_1}x_{b_2}\mathrm{}x_{b_m}x^{2(\{\frac{c_2}{2}\}m)},$$
(68)
and the symbol $`[]^{SL}`$ means the traceless part of the corresponding tensor; see (18) for $`m=3`$.
For $`m=4`$ the number of independent solutions of equations (14) is equal to
$$N_j^4=\frac{1}{12}(j+1)^2(j+2)^2(2j+3).$$
(69)
The solutions are numbered by triples of integers $`c=(c_1,c_2,c_3)`$ satisfying the conditions
$$0c_1j,c_1c_2c_1,0c_3\left\{\frac{c_1|c_2|}{2}\right\},$$
(70)
and for each $`c`$ contain $`N_c`$ arbitrary parameters where
$$N_c=\{\begin{array}{cc}(c_1+1)^2,\hfill & c_1=|c_2|,\hfill \\ 2(|c_2|+2c_3+1)(2c_1|c_2|2c_3+1),\hfill & c_1|c_2|.\hfill \end{array}$$
(73)
These parameters determine independent components of an irreducible tensor of rank $`R=R_1+2R_2`$ where
$$R_1=|c_2|+2c_3,R_2=c_1|c_2|2c_3,$$
(74)
and explicit expressions for the corresponding solutions have the form
$$F_c^{(a_1a_2\mathrm{}a_j)}=[\underset{i=0}{\overset{m+c_3}{}}(1)^i\left(\genfrac{}{}{0pt}{}{m+c_3}{i}\right)(x^2)^i$$
$$\times \lambda ^{b_1b_2\mathrm{}b_{mi+c_3}(a_1a_2\mathrm{}a_{|c_2|m+i+c_3}[a_{|c_2|+im+1+c_3}d_1]\mathrm{}[a_{c_1m+ic_3}d_{c_1|c_2|2c_3}]}$$
$$\times x^{a_{c_1m+ic_3+1}}x^{a_{c_1m+ic_3+2}}\mathrm{}x^{a_j)}x_{b_1}x_{b_2}\mathrm{}x_{b_{mi+c_3}}$$
$$\times x_{d_1}x_{d_2}\mathrm{}x_{d_{c_1|c_2|2c_3}}]^{SL}.$$
(75)
Here $`\lambda ^{b_1\mathrm{}b_{mi+c_3a_1}\mathrm{}a_{|c_2|m+i+c_3}[a_{|c_2|+im+1+c_3}d_1]\mathrm{}[a_{c_1m+ic_3}d_{c_1|c_2|2c_3}]}`$ is an arbitrary irreducible tensor of rank $`R_1+2R_2`$ ($`R_1`$ and $`R_2`$ are given in (65)),
$$m=\{\begin{array}{cc}\hfill c_2,& c_2<0,\hfill \\ \hfill 0,& c_20,\hfill \end{array}$$
(78)
$`\left(\genfrac{}{}{0pt}{}{m+c_3}{i}\right)`$ is a binomial coefficient, and the symbol $`[]^{SL}`$ means the traceless part of the corresponding tensor; see (18), (19) for $`m=4`$. Symmetrisation is implied over the indices $`a_1,\mathrm{},a_j`$ in the righthand part (the sum over all possible permutations).
Thus, we have found the explicit form of the conformal tensor of rank $`j`$ for $`m4`$. The formula (66) determines the general form such tensor for arbitrary $`m>3`$, but at that the total number of independent solutions of equations (14) cannot be determined, in general, by the relation (62), but requires special calculation for each value of $`m`$.
Let us point out that the general solution of equations (14) for $`m>2`$ can be presented in the form
$$F^{a_1a_2\mathrm{}a_j}=[\underset{l,k=0}{\overset{j}{}}\underset{i=0}{\overset{jlk}{}}\lambda ^{b_1b_2\mathrm{}b_{jlki}(a_1a_2\mathrm{}a_{l+i}[a_{l+i+1}d_1]\mathrm{}[a_{l+i+k}d_k]}$$
$$\times (1)^i\left(\genfrac{}{}{0pt}{}{jlk}{i}\right)(x^2)^ix^{a_{l+k+i+1}}\mathrm{}x^{a_j)}x_{d_1}x_{d_2}\mathrm{}x_{d_k}x_{b_1}x_{b_2}\mathrm{}x_{b_{jlki}}],$$
(79)
where $`\lambda ^{b_1b_2\mathrm{}b_{jlki}a_1\mathrm{}a_{l+i}[a_{l+i+1}d_1]\mathrm{}[a_{l+i+k}d_k]}`$ is a tensor symmetric with respect to permutation of the indices $`b_1,\mathrm{},a_{l+j+k}`$ and antisymmetric with respect to permutation of the indices $`a_{l+i+f}`$ with $`d_f,f=1,2,\mathrm{},k`$, with convolution of this tensor by any three indices with an absolutely antisymmetric vanishing. Decomposing such tensor into irreducible tensors we come to formulae (58)โ(67) giving solutions of equations (14) for $`m=3,4`$.
## 7 Examples of solutions and symmetry operators <br>for $`๐\mathbf{}\mathrm{๐}`$
Let us adduce an explicit form of the solutions obtained and corresponding symmetry operators for $`m4`$ and $`n3`$. Quantities of such solutions in accordance to (58) and (62) are adduced in Table 2.
Table 2. Quantities of independent solutions of equations (14).
$`m\backslash j`$ $`1`$ $`2`$ $`3`$ $`3`$ $`10`$ $`35`$ $`84`$ $`4`$ $`15`$ $`84`$ $`300`$
Quantities of the corresponding symmetry operators of order $`n`$ can be obtained by summation of quantities of solutions from $`j=0`$ to $`j=n`$. We adduce the result in Table 3.
Table 3. Quantities of order $`n`$ symmetry operators of equation (1)
with $`\varkappa =0`$ in $`m`$-dimensional space.
$`m\backslash n`$ $`0`$ $`1`$ $`2`$ $`3`$ $`4`$ $`3`$ $`1`$ $`11`$ $`46`$ $`130`$ $`295`$ $`4`$ $`1`$ $`16`$ $`100`$ $`400`$ $`1225`$
For $`m=2`$ the number of solutions of equations (14) (and the number of the corresponding symmetry operators) is infinite as they are determined up to arbitrary functions.
Explicit expressions for all independent solutions of equations (14) for $`m4`$ and $`n3`$ are given by the following formulae $`(F^{(j)}=F^{(a_1a_2\mathrm{}a_j)})`$:
$`m=2`$
$$j=0,F^{(0)}=\phi ^0(x_1,x_2);$$
$$j>0,F^{11\mathrm{}1}=(\phi _j+\phi _j^{})+i(\xi _j+\xi _j^{});$$
$$F^{11\mathrm{}12}=i(\phi _j^{}\phi _j)+\xi _j\xi _j^{}.$$
(80)
Here $`\phi _j`$ and $`\xi _j`$ are arbitrary analytical functions of two variables $`x_1`$, $`x_2`$, and other components of the tensor $`F^{a_1a_2\mathrm{}a_j}`$ are expressed through (69) using properties of of zero trace and symmetry.
$`m=3`$
$$j=0,F^{(0)}=\lambda ;$$
$$j=1,F_{(0,0)}^a=\lambda _{(0,0)}x^a;F_{(1,0)}^a=\lambda _{(1,0)}^a;F_{(1,1)}^a=\epsilon _{abc}\lambda _{(1,1)}^bx^c;$$
$$F_{(1,2)}^a=2\lambda _{(1,0)}^bx_bx^a\lambda _{(1,2)}^ax^2;$$
(81)
$$j=2,F_{(0,0)}^{a_1a_2}=\lambda _{(0,0)}\left(x^{a_1}x^{a_2}\frac{1}{3}g^{a_1a_2}x^2\right);$$
$$F_{(1,0)}^{a_1a_2}=\lambda _{(1,0)}^{a_1}x^{a_2}+\lambda _{(1,0)}^{a_2}x^{a_1}\frac{2}{3}g^{a_1a_2}\lambda _{(1,0)}^bx^b;$$
$$F_{(1,1)}^{a_1a_2}=(x^{a_1}\epsilon ^{a_2}{}_{bc}{}^{}+x^{a_2}\epsilon ^{a_1}{}_{bc}{}^{})x^b\lambda _{(1,1)}^c;$$
$$F_{(1,2)}^{a_1a_2}=(x^{a_1}\lambda _{(1,2)}^{a_2}+x^{a_2}\lambda _{(1,2)}^{a_1})x^24x^{a_1}x^{a_2}\lambda _{(1,2)}^bx_b+\frac{2}{3}g^{a_1a_2}\lambda _{(1,2)}^bx_bx^2;$$
$$F_{(2,0)}^{a_1a_2}=\lambda _{(2,0)}^{a_1a_2};$$
$$F_{(2,1)}^{a_1a_2}=(\epsilon ^{a_1bc}\lambda _{(2,1)}^{ba_2}+\epsilon ^{a_2bc}\lambda _{(2,1)}^{ba_1})x_c;$$
$$F_{(2,2)}^{a_1a_2}=\lambda _{(2,2)}^{a_1a_2}x^2(x^{a_1}\lambda _{(2,2)}^{a_2b}+x^{a_2}\lambda _{(2,2)}^{a_1b})x_b+\frac{2}{3}g^{a_1a_2}\lambda _{(2,2)}^{bc}x_bx_c;$$
$$F_{(2,3)}^{a_1a_2}=2(x^{a_1}\epsilon ^{a_2}{}_{bk}{}^{}+x^{a_2}\epsilon ^{a_1}{}_{bk}{}^{})\lambda _{(2,3)}^{kd}x^bx_d(\epsilon ^{a_1}{}_{ck}{}^{}\lambda _{(2,3)}^{a_2k}+\epsilon ^{a_2}{}_{ck}{}^{}\lambda _{(2,3)}^{a_1k})x^cx^2;$$
$$F_{(2,4)}^{a_1a_2}=\lambda _{(2,4)}^{a_1a_2}x^42(x^{a_1}\lambda _{(2,4)}^{2c}+x^{a_2}\lambda _{(2,4)}^{a_1c}x_cx^2+4x^{a_1}x^{a_2}\lambda _{(2,4)}^{ca}x_cx_d;$$
(82)
$$j=3,F_{(0,0)}^{a_1a_2a_3}=\lambda _{(0,0)}\left(x^{a_1}x^{a_2}x^{a_3}\frac{1}{10}g^{(a_1a_2}x^{a_3)}x^2\right);$$
$$F_{(1,0)}^{a_1a_2a_3}=\lambda _{(1,0)}^{(a_1}x^{a_2}x^{a_3)}\frac{1}{5}(g^{(a_1a_2}\lambda _{(1,0)}^{a_3)}x^2+2g^{(a_1a_2}x^{a_3)}\lambda _{(1,0)}^bx_b);$$
$$F_{(1,1)}^{a_1a_2a_3}=x^{(a_1}x^{a_2}\epsilon ^{a_3)}{}_{bc}{}^{}x_{}^{b}\lambda _{(1,1)}^c\frac{1}{5}g^{(a_1a_2}\epsilon ^{a_3)}{}_{bc}{}^{}x_{}^{b}\lambda _{(1,1)}^c;$$
$$F_{(2,0)}^{a_1a_2a_3}=\lambda _{(2,0)}^{(a_1a_2}x^{a_3)}\frac{2}{5}g^{(a_1a_2}\lambda _{(2,0)}^{a_3)b}x_b;$$
$$F_{(2,1)}^{a_1a_2a_3}=x^{(a_1}\epsilon ^{a_2}{}_{bc}{}^{}\lambda _{(2,1)}^{a_3)c}x_b\frac{1}{5}g^{(a_1a_2}\epsilon ^{a_3)}{}_{bc}{}^{}\lambda _{(2,1)}^{cd}x^bx_d;$$
$$F_{(2,2)}^{a_1a_2a_3}=\lambda _{(2,2)}^{(a_1a_2}x^{a_3)}x^22x^{(a_1}x^{a_2}\lambda _{(2,2)}^{a_3)}x_b+\frac{4}{5}g^{(a_1a_2}x^{a_3)}\lambda _{(2,2)}^{bc}x_bx_c;$$
$$F_{(2,3)}^{a_1a_2a_3}=\epsilon ^{(a_1}{}_{bc}{}^{}x_{}^{a_2}(2x^{a_3)}\lambda _{(2,3)}^{bd}x_d\lambda _{(2,3)}^{a_3)b}x^2)x^c;$$
$$F_{(2,4)}^{a_1a_2a_3}=x^{(a_1}(\lambda _{(2,4)}^{a_2a_3)}x^44x^{a_2}\lambda _{(2,4)}^{a_3)b}x_bx^2+4x^{a_2}x^{a_3)}\lambda _{(2,4)}^{kl}x_kx_l)$$
$$\frac{2}{5}g^{(a_1a_2}(x^{a_3)}\lambda _{(2,4)}^{kl}x_kx_lx^2\lambda ^{a_3)b}x_bx^4);$$
$$F_{(3,0)}^{a_1a_2a_3}=\lambda _{(3,0)}^{a_1a_2a_3};$$
$$F_{(3,1)}^{a_1a_2a_3}=\epsilon ^{(a_3}{}_{bc}{}^{}\lambda _{}^{a_1a_2)b}x_c;$$
$$F_{(3,2)}^{a_1a_2a_3}=\lambda _{(3,2)}^{(a_1a_2a_3)}x^22x^{(a_3}\lambda _{(3,2)}^{a_1a_2)b}x_b+\frac{4}{5}g^{(a_1a_2}\lambda _{(3,2)}^{a_3)bc}x_bx_c;$$
$$F_{(3,3)}^{a_1a_2a_3}=\epsilon ^{(a_1}{}_{bc}{}^{}(\lambda _{(3,3)}^{a_2a_3)b}x^22x^{a_2}\lambda _{(3,3)}^{a_3)}x_bx_c+\frac{2}{5}g^{a_2a_3)}\lambda _{(3,3)}^{bcd}x_bx_cx_d);$$
$$F_{(3,4)}^{a_1a_2a_3}=\lambda _{(3,4)}^{(a_1a_2a_3)}x^44x^{(a_1}(\lambda _{(3,4)}^{a_2a_3)c}x_cx^2x^{a_2}\lambda _{(3,4)}^{a_3)bc}x_bx_c)$$
$$\frac{4}{5}g^{(a_1a_2}(2x^{a_3)}\lambda _{(3,4)}^{bcd}x_d\lambda _{(3,4)}^{a_3)bc}x^2)x_bx_c;$$
$$F_{(3,5)}^{a_1a_2a_3}=\epsilon ^{(a_1}(\lambda _{(3,5)}^{a_3a_2)b}x^44x^{a_2}\lambda _{(3,5)}^{a_3)bd}x_dx^2+4x^{a_2}x^{a_3)}\lambda _{(3,5)}^{bkl)b}x_kx_l)x^c;$$
$$F_{(3,6)}^{a_1a_2a_3}=\lambda _{(3,6)}^{(a_1a_2a_3)}x^66x^{(a_1}\lambda _{(3,6)}^{a_2a_3)c}x_cx^4$$
$$+12x^{(a_1}x^{a_2}\lambda _{(3,6)}^{a_3)bc}x_bx_cx^28x^{(a_1}x^{a_2}x^{a_3)}\lambda _{(3,6)}^{bcd}x_bx_cx_d;$$
(83)
$`m=4`$
$$j=0,F^{(0)}=\lambda ;$$
$$j=1,F_{(0,0,0)}^{(a)}=\lambda x^a;$$
$$F_{(1,1,0)}^{(a)}=\lambda ^bx^ax_b\lambda ^ax^2;$$
$$F_{(1,0,0)}^{(a)}=\lambda ^{[ad]}x_d;$$
$$F_{(1,1,0)}^{(a)}=\lambda ^a;$$
(84)
$$j=2,F_{(0,0,0)}^{a_1a_2}=\lambda _{(0,0,0)}\left(x^{a_1}x^{a_2}\frac{1}{4}g^{a_1a_2}x^2\right);$$
$$F_{(1,1,0)}^{a_1a_2}=2\lambda _{(1,1,0)}^bx_b\left(x^{a_1}x^{a_2}\frac{1}{4}g^{a_1a_2}x^2\right)\lambda _{(1,1,0)}^{a_1}x^{a_2}x^2$$
$$\lambda _{(1,1,0)}^{a_2}x^{a_1}x^2+\frac{1}{2}g^{a_1a_2}\lambda _{(1,1,0)}^cx_cx^2;$$
$$F_{(1,0,0)}^{a_1a_2}=\lambda _{(1,0,0)}^{[a_1d_1]}x^{a_2}x_{d_1}+\lambda _{(1,0,0)}^{[a_2d_1]}x^{a_1}x_{d_1};$$
$$F_{(1,1,0)}^{a_1a_2}=\frac{1}{2}g^{a_1a_2}\lambda _{(1,1,0)}^cx_c+\lambda _{(1,1,0)}^{a_1}x^{a_2}+\lambda _{(1,1,0)}^{a_2}x^{a_1};$$
$$F_{(2,2,0)}^{a_1a_2}=2\lambda _{(2,2,0)}^{b_1b_2}x^{a_1}x^{a_2}x_{b_1}x_{b_2}2\lambda _{(2,2,0)}^{b_1a_1}x^{a_2}x^2x_{b_1}$$
$$2\lambda _{(2,2,0)}^{b_1a_2}x^{a_1}x^2x_{b_1}+2\lambda _{(2,2,0)}^{a_1a_2}x^4+\frac{1}{2}g^{a_1a_2}\lambda _{(2,2,0)}^{b_1b_2}x^2x_{b_1}x_{b_2};$$
$$F_{(2,1,0)}^{a_1a_2}=\lambda _{(2,1,0)}^{b_1[a_1d_1]}x^{a_1}x_{b_1}x_{d_1}+\lambda _{(2,1,0)}^{b_1[a_2d_1]}x^{a_1}x_{b_1}x_{d_1}$$
$$\lambda _{(2,1,0)}^{a_1[a_2d_1]}x^2x_{d_1}\lambda _{(2,1,0)}^{a_2[a_1d_1]}x^2x_{d_1};$$
$$F_{(2,0,1)}^{a_1a_2}=\lambda _{(2,0,1)}^{ba_1}x^{a_2}x_b+\lambda _{(2,0,1)}^{ba_2}x^{a_1}x_b2\lambda _{(2,0,1)}^{a_1a_2}x^2\frac{1}{2}g^{a_1a_2}\lambda _{(2,0,1)}^{bc}x_bx_c;$$
$$F_{(2,0,0)}^{a_1a_2}=\lambda _{(2,0,0)}^{[a_1d_1][a_2d_2]}x_{d_1}x_{d_2};$$
$$F_{(2,1,0)}^{a_1a_2}=(\lambda _{(2,1,0)}^{a_1[a_2d]}+\lambda _{(2,1,0)}^{a_2[a_1d]})x_d;$$
$$F_{(2,2,0)}^{a_1a_2}=\lambda _{(2,2,0)}^{a_1a_2};$$
(85)
$$j=3,F_{(0,0,0)}^{a_1a_2a_3}=\lambda _{(0,0,0)}\left(x^{(a_1}x^{a_2}x^{a_3)}\frac{1}{2}g^{(a_1a_2}x^{a_3)}x^2\right);$$
$$F_{(1,1,0)}^{a_1a_2a_3}=\lambda _{(1,1,0)}^b\left(x^{(a_1}x^{a_2}x^{a_3)}\frac{1}{2}g^{(a_1a_2}x^{a_3)}x^2\right)x_b\lambda _{(1,1,0)}^{(a_1}x^{a_2}x^{a_3)}x^2$$
$$+\frac{1}{3}g^{(a_1a_2}x^{a_3)}\lambda _{(1,1,0)}^bx_bx^2+\frac{1}{6}g^{(a_1a_2}\lambda _{(1,1,0)}^{a_3)}x^4;$$
$$F_{(1,0,0)}^{a_1a_2a_3}=\lambda _{(1,0,0)}^{([a_1d]}x^{a_2}x^{a_3)}x_d\frac{1}{6}g^{(a_1a_2}\lambda _{(1,0,0)}^{[a_3d])}x^2x_d;$$
$$F_{(1,1,0)}^{a_1a_2a_3}=\lambda _{(1,1,0)}^{(a_1}x^{a_2}x^{a_3)}\frac{1}{3}g^{(a_1a_2}x^{a_3)}\lambda _{(1,1,0)}^bx_b\frac{1}{6}g^{(a_1a_2}\lambda _{(1,1,0)}^{a_3)}x^2;$$
$$F_{(2,2,0)}^{a_1a_2a_3}=\lambda _{(2,2,0)}^{b_1b_2}\left(x^{(a_1}x^{a_2}x^{a_3)}\frac{1}{2}g^{(a_1a_2}x^{a_3)}x^2\right)x_{b_1}x_{b_2}$$
$$2\lambda _{(2,2,0)}^{b(a_1}x^{a_2}x^{a_3)}x_bx^2+\frac{2}{3}g^{(a_1a_2}x^{a_3)}\lambda _{(2,2,0)}^{bc}x_bx_cx^2+\frac{1}{3}g^{(a_1a_2}\lambda _{(2,2,0)}^{a_3)b}x_bx^4$$
$$+\lambda _{(2,2,0)}^{(a_1a_2}x^{a_3)}x^4+\frac{1}{3}g^{(a_1a_2}\lambda _{(2,2,0)}^{a_3)b}x_bx^4;$$
$$F_{(2,1,0)}^{a_1a_2a_3}=\lambda _{(2,1,0)}^{b([a_1d]}x^{a_2}x^{a_3)}x_bx_d\lambda _{(2,1,0)}^{(a_1[a_2d]}x^{a_3)}x_dx^2$$
$$\frac{1}{6}g^{(a_1a_2}\lambda _{(2,1,0)}^{[a_3d])b}x^2x_bx_d+\frac{1}{6}g^{(a_1a_2}\lambda _{(2,1,0)}^{[a_3d])b}x^2x_bx_d;$$
$$F_{(2,0,0)}^{a_1a_2a_3}=\lambda _{(2,0,0)}^{([a_1d_1][a_2d_2]}x^{a_3)}x_{d_1}x_{d_2};$$
$$F_{(2,0,1)}^{a_1a_2a_3}=\lambda _{(2,0,1)}^{b(a_1}x^{a_2}x^{a_3)}x_b\frac{1}{3}g^{(a_1a_2}x^{a_3)}\lambda _{(2,0,1)}^{bc}x_bx_c$$
$$+\frac{1}{6}g^{(a_1a_2}\lambda _{(2,0,1)}^{a_3)b}x_bx^2\lambda _{(2,0,1)}^{(a_1a_2}x^{a_3)}x^2;$$
$$F_{(2,1,0)}^{a_1a_2a_3}=\lambda _{(2,1,0)}^{(a_1[a_2d]}x^{a_3)}x_d\frac{1}{6}g^{(a_1a_2}\lambda _{(2,1,0)}^{[a_3d])b}x_bx_d;$$
$$F_{(2,2,0)}^{a_1a_2a_3}=\lambda _{(2,2,0)}^{(a_1a_2}x^{a_3)}\frac{1}{3}g^{(a_1a_2}\lambda _{(2,2,0)}^{a_3)b}x_b;$$
$$F_{(3,3,0)}^{a_1a_2a_3}=\lambda _{(3,3,0)}^{b_1b_2b_3}\left(x^{(a_1}x^{a_2}x^{a_3)}x_{b_1}x_{b_2}x_{b_3}+\frac{1}{2}g^{(a_1a_2}x^{a_3)}x^2x_{b_1}x_{b_2}x_{b_3}\right)$$
$$\lambda _{(3,3,0)}^{(a_1a_2a_3)}x^63\lambda _{(3,3,0)}^{b_1b_2(a_1}x^{a_2}x^{a_3)}x^2x_{b_1}x_{b_2}+3\lambda _{(3,3,0)}^{b_1(a_1a_2}x^{a_3)}x^4x_{b_1}$$
$$\frac{1}{2}g^{(a_1a_2}\lambda _{(3,3,0)}^{a_3)b_1b_2}x_{b_1}x_{b_2}x^4;$$
$$F_{(3,2,0)}^{a_1a_2a_3}=\lambda _{(3,2,0)}^{b_1b_2([a_1d]}x^{a_2}x^{a_3)}x_{b_1}x_{b_2}x_d+\frac{1}{6}g^{(a_1a_2}\lambda _{(3,2,0)}^{[a_3d])b_1b_2}x^2x_{b_1}x_{b_2}x_d$$
$$2\lambda _{(3,2,0)}^{b(a_1[a_2d]}x^{a_3)}x_bx_dx^2+\lambda _{(3,2,0)}^{(a_1a_2[a_3d])}x_dx^4;$$
$$F_{(3,1,0)}^{a_1a_2a_3}=\lambda _{(3,1,0)}^{b_1([a_1d_1][a_2d_2]}x^{a_3)}x_bx_{d_1}x_{d_2}\lambda _{(3,1,0)}^{(a_1[a_2d_1][a_3d_2])}x_{d_1}x_{d_2}x^2;$$
$$F_{(3,0,0)}^{a_1a_2a_3}=\lambda _{(3,0,0)}^{[a_1d_1][a_2d_2][a_3d_3]}x_{d_1}x_{d_2}x_{d_3};$$
$$F_{(3,0,1)}^{a_1a_2a_3}=(\lambda _{(3,0,1)}^{b(a_1[a_2d]}x^{a_3)}\frac{1}{6}g^{(a_1a_2}\lambda _{(3,0,1)}^{[a_3d])bc}x_c)x_bx_d\lambda _{(3,0,1)}^{(a_1a_2[a_3d_1])}x_{d_1}x^2;$$
$$F_{(3,1,1)}^{a_1a_2a_3}=(\lambda _{(3,1,1)}^{b_1b_2(a_1}x^{a_2}x^{a_3)}\frac{1}{3}g^{(a_1a_2}x^{a_3)}\lambda _{(3,1,1)}^{b_1b_2c}x_c$$
$$+\frac{1}{2}g^{(a_1a_2}\lambda _{(3,1,1)}^{a_3)b_1b_2}x^2)x_{b_1}x_{b_2}2\lambda ^{b(a_1a_2}_{(3,1,1)}x^{a_3)}x_bx^2+\lambda ^{(a_1a_2a_3)}_{(3,1,1)}x^4;$$
$$F_{(3,1,0)}^{a_1a_2a_3}=\lambda _{(3,1,0)}^{(a_1[a_2d_1][a_3d_2]}x_{d_1}x_{d_2};$$
$$F_{(3,1,1)}^{a_1a_2a_3}=\lambda _{(3,1,1)}^{b(a_1a_2}x^{a_3)}x_b\frac{1}{3}g^{(a_1a_2}\lambda _{(3,1,1)}^{a_3)bc}x_bx_c\lambda _{(3,1,1)}^{(a_1a_2a_3)}x^2;$$
$$F_{(3,2,0)}^{a_1a_2a_3}=\lambda _{(3,2,0)}^{(a_1a_2[a_3d])}x_d;$$
$$F_{(3,3,0)}^{a_1a_2a_3}=\lambda _{(3,3,0)}^{a_1a_2a_3}$$
(86)
Here $`\lambda _{(\mathrm{})}`$, $`\lambda _{(\mathrm{})}^a`$, $`\lambda _{(\mathrm{})}^{a_1a_2}`$, $`\lambda _{(\mathrm{})}^{a_1a_2a_3}`$ are arbitrary symmetric tensors with zero trace.
Formulae (69)โ(75) give explicit form of all linearly independent conformal killing tensors of rank $`j3`$ in spaces of dimension $`p+q=2,3,4`$ (for $`p+q=2`$ $`j`$ is arbitrary). To obtain an explicit form of the corresponding operators it is sufficient to substitute (69)โ(75)into (6), (7) that is write down $`j`$-multiple anticommutators of $`F^{a_1a_2\mathrm{}a_j}`$ with $`\frac{}{x_{a_1}},\frac{}{x_{a_2}},\mathrm{},\frac{}{x_{a_j}}`$.
## 8 Killing tensors of rank $`๐`$ and order $`๐`$
Until now we considered solutions of equations (11), (14) that define Killing tensors (and conformal Killing tensors) of arbitrary rank $`j`$, but only of the first order. In this section we obtain explicit form of Killing tensors of rank $`j`$ and of arbitrary order $`s`$. Such tensors are determined as general solutions of equations (16).
The system of equations (16) is overdetermined including $`N_{js}^m`$ equations for $`\widehat{N}_{js}^m`$ unknown variables, where
$$N_{js}^m=\left(\genfrac{}{}{0pt}{}{j+s+m1}{m1}\right),\widehat{N}_{js}^m=\left(\genfrac{}{}{0pt}{}{j+m1}{m1}\right),m=p+q.$$
(87)
In the same way as it was done above in Section 3, we consider the set of differential consequences of the system under consideration that are obtained by $`k`$-multiple differentiation of every term of the equation by $`\frac{}{x_{a_1}},\frac{}{x_{a_2}},\mathrm{},\frac{}{x_{a_k}}`$. This set is a system of linear homogeneous algebraic equations of the following form:
$$F^{(a_1a_2\mathrm{}a_j,a_{j+1}a_{j+2}\mathrm{}a_{j+s})b_1b_2\mathrm{}b_k}=0,$$
(88)
where the following derivatives are unknown variables:
$$F^{a_1a_2\mathrm{}a_j,a_{j+1}a_{j+2}\mathrm{}a_{j+s}b_1b_2\mathrm{}b_k}^{a_{j+1}}^{a_{j+2}}\mathrm{}^{a_{j+s}}^{b_1}^{b_2}\mathrm{}^{b_k}F^{a_1a_2\mathrm{}a_j}.$$
(89)
The quantities of unknown variables $`N_\text{u}^k`$ and of equations $`N_\text{e}^k`$ are equal to
$$N_\text{u}^k=\left(\genfrac{}{}{0pt}{}{j+m1}{m1}\right)\left(\genfrac{}{}{0pt}{}{k+s+m1}{m1}\right),N_\text{e}^k=\left(\genfrac{}{}{0pt}{}{j+s+m1}{m1}\right)\left(\genfrac{}{}{0pt}{}{k+m1}{m1}\right),$$
(90)
so conditions (24) are also fulfilled.
It can be shown (see attachment) that the system (77) is non-degenerate, so it follows from (24) that
$$F^{a_1a_2\mathrm{}a_j,a_{j+1}a_{j+2}\mathrm{}a_{j+s}b_1b_2\mathrm{}b_j}0.$$
Whence we conclude that the Killing tensors of rank $`j`$ and order $`s`$ are polynomials of order $`j+s1`$. It follows from (77), (79) that such polynomial contains $`n_{js}^m`$ arbitrary parameters, where
$$n_{js}^m=\underset{i=0}{\overset{s1}{}}\left(\genfrac{}{}{0pt}{}{j+m1}{m1}\right)\left(\genfrac{}{}{0pt}{}{j+m1}{m1}\right)+\underset{k=0}{\overset{s1}{}}(N_\text{u}^kN_\text{e}^k)$$
$$=\frac{s}{m}\left(\genfrac{}{}{0pt}{}{j+m1}{m1}\right)\left(\genfrac{}{}{0pt}{}{j+s+m1}{m1}\right).$$
(91)
Here the first sum gives the number of independent solutions that have order by $`x`$ smaller than $`s`$. Such solutions can be written in the form
$$F_i^{a_1a_2\mathrm{}a_j}=\lambda _{b_1b_2\mathrm{}b_i}^{a_1a_2\mathrm{}a_j}x^{b_1}x^{b_2}\mathrm{}x^{b_i},i<s,$$
(92)
where $`\lambda _{b_1b_2\mathrm{}b_i}^{a_1a_2\mathrm{}a_j}`$ are numeric parameters with no limitations set by equations (77) (certainly the symmetry with respect to permutations of indices $`a_\lambda a_\lambda ^{}`$, $`b_\mu b_\mu ^{}`$, $`\lambda ,\lambda ^{}=1,2,\mathrm{},j`$, $`\mu ,\mu ^{}=1,2,\mathrm{},i`$.
In the case $`s=1`$ the formula (80) is reduced to (25). In particular, for $`m=2,3,4`$ we obtain from (80) that
$$n_{js}^2=\frac{1}{2}s(j+1)(j+s+1),$$
$$n_{js}^3=\frac{1}{12}s(j+1)(j+2)(j+s+1)(j+s+2),$$
$$n_{js}^4=\frac{1}{3!4!}s(j+1)(j+2)(j+3)(j+s+1)(j+s+2)(j+s+3),$$
(93)
Thus we have determined the number of linearly independent Killing tensors of rank $`j`$ and order $`s`$. To compute these tensors in explicit form we will use the following two lemmas.
Lemma 3. Let $`F_i^{a_1a_2\mathrm{}a_j}`$ be Killing tensors of rank $`j`$ and order $`s`$, and $`\phi `$ be a functions satisfying the equation
$$^\mu ^\nu \phi =0,\mu ,\nu =1,2\mathrm{}m.$$
(94)
Then the function
$$\stackrel{~}{F}^{a_1a_2\mathrm{}a_j}=\phi F^{a_1a_2\mathrm{}a_j}$$
(95)
is a Killing tensor of rank $`j`$ and order $`s`$.
Proof is reduced to direct check of the lemma statement that is $`(s+1)`$-multiple differentiation of (84) by $`^{a_{j+1}}`$, $`^{a_{j+1}}`$, โฆ, $`^{a_{j+s+1}}`$ and subsequent symmetrization of the obtained expression by $`a_1`$, $`a_2`$, $`a_{j+s+1}`$ using relations (16), (83).
Lemma 4. Let $`F^{a_1a_2\mathrm{}a_j}`$ be a Killing tensors of rank $`j`$ and order $`s`$. Then the convolution
$$\stackrel{~}{F}^{a_1a_2\mathrm{}a_{j1}}=F^{a_1a_2\mathrm{}a_j}x_{a_j}$$
(96)
is a Killing tensor of rank $`j+1`$ and order $`s+1`$.
Proof is similar.
The adduced Lemmas provide an effective algorithm for construction of Killing tensors of order $`s`$ from Killing tensors of order 1 found above in Section 4. The only difficulty in application of this algorithm is the need to sort out all linearly independent solutions of the system (16) (whose number is determined by the formulae (80), (82), as, generally speaking, there are more solutions of the form (84) than we need).
The general solution of equations (16) is determined in the following theorem.
Theorem 3. Equations (16) in the space of dimension $`m4`$ have $`n_{js}^m`$ linearly independent solutions where $`n_{js}^m`$ is given by the formula (82). These solutions have the form
$$F_{(s)}^{a_1a_2\mathrm{}a_j}=g^{(a_{j1}a_j}F_{(s)}^{a_1a_2\mathrm{}a_{j2})}+\epsilon _j\widehat{f}^{a_1a_2\mathrm{}a_j}$$
$$+\underset{d=1}{\overset{s}{}}x^{a_{j+1}}x^{a_{j+2}}\mathrm{}x^{a_{j+d1})},\epsilon _j=\frac{1}{2}[1+(1)^j],$$
(97)
where $`F^{a_1a_2\mathrm{}a_{j+d1}}`$ are Killing tensors of rank $`j+d1`$ and of order 1 whose explicit form is given by Theorem 2, $`F_{(s)}^{a_1a_2\mathrm{}a_{j2}}`$ are Killing tensors of rank $`j2`$ and of order $`s`$;
$$\widehat{f}^{a_1a_2\mathrm{}a_j}=\underset{\mu =0}{\overset{\frac{j}{2}1}{}}(1)^\mu \left(\genfrac{}{}{0pt}{}{\frac{j}{2}1}{\mu }\right)x^{(a_1}x^{a_2}\mathrm{}x^{a_{2\mu +1}}$$
$$\times g^{a_{2\mu +2}a_{2\mu +3}}\mathrm{}g^{a_{j2}a_{j1}}\lambda ^{[a_j),c]}x_c,$$
(98)
$`\lambda ^{[a_j,c]}`$ is an arbitrary antisymmetric tensor of rank $`2`$.
Proof. By virtue of Lemmas 3, 4 the function $`F_{(s)}^{a_1a_2\mathrm{}a_j}`$ given by formula (86) is a Killing tensor of rank $`j`$ and order $`s`$; the first term โ by definition, the second โ in accordance to Lemma 3 (being the product of of a Killing tensor of order 1 and $`\phi =\lambda ^\mu x_\mu `$), the third โ in accordance to Lemma 4 (each convolution with $`x_\mu `$ lowers the rank and increases the order of a Killing tensor, and we make the first convolution with the first-order tensors described above).
It is to some extent more difficult to make sure that formula (86) gives all linearly independent Killing tensors of order $`s`$. Proof of linear independence of all terms of the formula (86) is reduced to comparison of terms having the same order by $`x_{a_i}`$ using different convolutions by one, two etc. pairs of indices. Calculation of then umber of independent solutions given by formula (86) can be done easily by sorting through independent solutions for first-order tensors $`F^{a_1a_2\mathrm{}a_{j+d1}}`$ entering the last term (such tensors are described in Theorem 2, but it is necessary to restrict consideration with solutions corresponding to $`c_2>d1`$, as others give zero input into convolutions of (86)), and adding the number of solutions of the form (87). The result obtained is in compliance with the formula
$$N=n_{js}^mn_{j2s}^m,$$
where $`n_{js}^m`$ is the total number of solutions given by the formulae (82), $`n_{j2s}^m`$ is the number of solutions of the form $`g^{(a_{j1}a_j}F_s^{a_1a_2\mathrm{}a_{j2})}`$ that is also given in (82), $`N`$ is the total number of solutions under the summation sign and of solutions of the form (87). We omit the corresponding cumbersome calculations.
Formula (86) determines recurrent relations for calculation of explicit form of a Killing tensor of rank $`j`$ and order $`s`$ from a known tensor of order $`s`$ and rank $`j2`$. Such calculations may be easily checked starting from known killing tensors of order 1 and arbitrary rank, see Theorem 2.
Let us adduce as an example explicit expressions for Killing vectors of order $`s3`$ in three-dimensional space received from general relations (86):
$$s=1,F_{(1)}^a=\lambda ^a+\epsilon ^{abc}\eta _bx_c,$$
$$s=2,F_{(2)}^a=F_{(1)}^a+\lambda ^{ab}x_b+\lambda x^a+\epsilon ^{abc}\eta _{bd}x_cx^d+\xi ^ax^2x^a\xi ^bx_b,$$
$$s=3,F_{(3)}^a=F_{(2)}^a+\lambda ^{abc}x_bx_c+\stackrel{~}{\lambda }^bx_bx^a+\epsilon ^{abc}\eta _{bdl}x_cx^dx^l$$
$$+\xi ^{ab}x_bx^2x^a\xi ^{bc}x_bx_c+\epsilon ^{abc}x_b\xi _cx^2.$$
Here $`\epsilon ^{abc}`$ is a unit antisymmetric tensor, and other Greek letters designate arbitrary symmetric tensors with zero trace.
## 9 Conformal Killing tensors of rank $`๐`$ and order $`๐`$
Let us briefly discuss equations (17) describing conformal Killing tensors of rank $`j`$ and order $`s`$, and adduce without proof solutions of these equations for arbitrary $`j`$, $`s`$ and $`m4`$.
A constructive way for finding solutions of equations (17) is shown by the following statement that may be proven by direct check.
Lemma 5. Let $`F_s^{a_1a_2\mathrm{}a_j}`$ be a conformal Killing tensor of rank $`j`$ and order $`s`$, and $`\phi `$ be an arbitrary function satisfying the equation
$$^\mu ^\nu \phi =g^{\mu \nu }\lambda ,\lambda =\text{const}.$$
(99)
Then the function
$$\stackrel{~}{F}_{s+1}^{a_1a_2\mathrm{}a_j}=\phi F_s^{a_1a_2\mathrm{}a_j}$$
(100)
is a conformal Killing tensor of rank $`j`$ and order $`s+1`$.
We can show that quantities of linearly independent solutions of equations (17) for arbitrary $`j`$, $`s`$ and $`m=3,4`$ are given by formulae
$$m=3,\widehat{N}_{js}^3=\frac{s}{6}(2j+1)(2j+2s+1)(2j+s+1),$$
$$m=4,\widehat{N}_{js}^4=\frac{s}{12}(j+1)^2(j+s+1)^2(2j+2+s).$$
(101)
Using Lemma 5 we managed to construct $`N_{js}^m`$ linearly independent Killing tensors of rank $`j`$ and order $`s`$ (giving full system of solutions of equations (17)) in the following form:
$$F_s^{a_1a_2\mathrm{}a_j}=\underset{i=1}{\overset{s}{}}\left(F_i^{a_1a_2\mathrm{}a_j}(x^2)^{i1}+\underset{d=0}{\overset{si}{}}f_{i1d}^{a_1a_2\mathrm{}a_j}(x^2)^d\right).$$
(102)
Here $`F_i^{a_1a_2\mathrm{}a_j}`$ are conformal tensors of rank $`j`$ and order 1 given by formulae (60), (61) or (66) (the index โ$`i`$โ distinguishes independent solutions of (91) with various degrees of $`x^2`$), $`f_{i1d}^{a_1a_2\mathrm{}a_j}`$ are tensors of rank $`j`$ whose explicit form is adduced below.
In the case $`m=2`$ the number of conformal Killing tensors of order 1 appears to be infinite, see (69). The same formulae (69) give a general form of a conformal Killing tensor of arbitrary order $`s`$.
In the case $`m=3`$ functions $`f_{i1d}^{a_1a_2\mathrm{}a_j}`$ are characterized by an additional integer $`c`$,
$$0c2j,$$
(103)
and are determined up to arbitrary symmetric zero-trace tensor $`\stackrel{~}{\lambda }^{a_1a_2\mathrm{}a_R}`$ of rank $`R=2(c+i)1`$. Explicit form of these functions is given by by formula (93):
$$f_{id}^{a_1a_2\mathrm{}a_j}=\left[\epsilon _c\widehat{f}_{idc}^{a_1a_2\mathrm{}a_j}+(1\epsilon c)\widehat{f}^{b(a_1a_2\mathrm{}a_{j1}}\epsilon ^{aj)bc}x_c\right]^{SL},$$
(104)
where
$$\widehat{f}_{idc}^{a_1a_2\mathrm{}a_j}=\underset{n=0}{\overset{\{\frac{c}{2}\}}{}}(2)^n\left(\genfrac{}{}{0pt}{}{\{\frac{c}{2}\}}{n}\right)\stackrel{~}{\lambda }^{b_1b_2\mathrm{}b_{d+n}(a_1a_2\mathrm{}a_{jn}}$$
$$\times x^{a_{jn+1}}x^{a_{jn+2}}\mathrm{}x^{a_j)}x_{b_1}x_{b_2}\mathrm{}x_{b_{d+n}}x^{2(\{\frac{c}{2}\}n)},\epsilon _c=\frac{1}{2}[1+(1)^c]$$
(105)
with the symbol $`[]^{SL}`$ designating the zero trace part of the corresponding tensor, see (18), (19) for $`m=3`$, and the index $`d`$ is introduced for numeration of linearly independent solutions of (91) with different degrees of $`x^2`$.
In the case $`m=4`$ the functions $`f_{idc}^{a_1a_2\mathrm{}a_j}`$ are characterized by a pair of additional indices $`c=(c_1,c_2)`$
$$jc_1j,0c_2\left\{\frac{j|c_1|}{2}\right\}$$
(106)
and is determined up to arbitrary irreducible tensor $`\stackrel{~}{\lambda }^{a_1a_2\mathrm{}a_{R_1}[a_{R_1+1}b_1]\mathrm{}[a_{R_1+R_2}b_{R_2}]}`$ of rank $`R_1+2R_2`$ where
$$R_1=|c_1|+2c_2+i,R_2=j|c_1|2c_2.$$
(107)
Explicit form of these functions is given by the formula (97):
$$f_{id}^{a_1a_2\mathrm{}a_j}=[\underset{d=0}{\overset{n+c_2}{}}(1)^d\left(\genfrac{}{}{0pt}{}{n+c_2}{d}\right)(x^2)^d\lambda ^{b_1b_2\mathrm{}b_{n+c_2di}(a_1a_2\mathrm{}a_{|c_1|n+d+i+c_2+1}d_1]\mathrm{}}$$
$${}_{}{}^{\mathrm{}[a_{jn+i+dc_2}d_{j|c_1|2c_2}]}x_{}^{a_{jn+i+dc_2+1}}\mathrm{}x^{a_j}x_{b_1}\mathrm{}x_{b_{n+c_2di}}x_{d_1}\mathrm{}x_{d_{j|c_1|2c_2}}]^{SL},$$
(108)
where
$$n=\{\begin{array}{cc}\hfill c_1,& c_1<0,\hfill \\ \hfill 0,& c_10,\hfill \end{array}$$
(111)
and symmetrization is implied in the right-hand part by indices $`a_1,a_2,\mathrm{},a_j`$.
The formulae (91)โ(98) give in explicit form all linearly independent conformal Killing tensors of rank $`j`$ and order $`s`$ in space of dimension $`m=p+q4`$. In particular, conformal vectors of order $`s3`$ in three-dimensional space, in accordance to (91), (94) have the following form:
$$s=1,F_{(1)}^a=\lambda _{(1)}^a+\epsilon ^{abc}\eta _{(1)}^bx_c+\xi _{(1)}^ax^22x^a\xi _{(1)}^bx_b+\mu x^a;$$
$$s=2,F_{(2)}^a=F_{(1)}^a+\stackrel{~}{F}_{(1)}^ax^2+\lambda _{(2)}^{ab}x_b+\epsilon ^{abc}\eta _{(2)}^{bd}x_cx_d+\xi _{(2)}^{ab}x^2x_b2x^a\xi _{(2)}^{bc}x_bx_c;$$
$$s=3,F_{(3)}^a=F_{(2)}^a+x^4\stackrel{}{๐น}{}_{(1)}{}^{a}+x^2(\lambda _{(3)}^{ab}x_b+\epsilon ^{abc}\eta _{(3)}^{bd}x_cx_d+\xi _{(3)}^{ab}x_bx^22x^a\xi _{(3)}^{bc}x_bx_c)$$
$$+\lambda _{(3)}^{abc}x_bx_c+\epsilon ^{abc}\eta _{(3)}^{bdk}x_cx_dx_k+\epsilon _{(3)}^{abc}x_bx_cx^22x^a\xi _{(3)}^{bcd}x_bx_cx_d.$$
Here $`\epsilon ^{abc}`$ is the unit antisymmetric tensor, and other Greek letters designate arbitrary symmetric traceless tensors $`F_{(1)}^a`$, $`\stackrel{~}{F}_{(1)}^a`$ and $`\stackrel{}{๐น}_{(1)}^a`$ are first-order Killing vectors (generally speaking, different).
## 10 Conclusion
Let us sum up. We have defined the notion of Killing tensor of rank $`j`$ and order $`s`$ and of conformal Killing tensor of rank $`j`$ and order $`s`$. These tensors are defined as general solutions of equations (16) or (17) that in the case $`s=1`$ coincide with generally accepted equations for Killing tensors and conformal Killing tensors, see e.g. .
We limit ourselves with investigation of equations (16) and (17) in flat de Sitter space, and generalization of these equations for for the case of spaces with non-zero curvature requires replacement of $`^{a_i}`$ for covariant derivatives.
Equations (16) and (17) are natural generalizations of the Killing equations and arise in description of higher-order symmetry operators. In the present paper we show relation of these equations (for first-order tensors) with higher-order symmetry operators of KleinโGordonโFock equation, see Sections 2, 5. Equations for Killing tensors and conformal Killing tensors of order $`s<1`$ arise in problems of description of symmetry operators of order $`s`$ for systems of partial differential equations โ in particular, for the Maxwell equations . We have found in explicit form all non-equivalent Killing tensors of rank $`j`$ and order $`s`$ in the space of dimension $`p+q`$ for arbitrary $`j`$ and $`s`$ and $`p+q4`$. Limitation by dimension of space is based, on the one hand, on practical reasons (the absolute majority of equations of mathematical physics being the field of research interests of the authors, have dimension $`m4`$ with respect to independent variables), and, on the other side, on difficulties that had not been overcome to the moment in proof of non-degeneracy of systems of algebraic equations for coefficients of Killing tensors in spaces of arbitrary dimension, see attachment. At that formulae (42), (66), (86) giving solutions of equations (16), (17) for $`p+q=4`$, probably give the general solution of these equations for arbitrary $`p+q4`$.
The found general solutions for Killing tensors and conformal Killing tensors of arbitrary order and rank may find quite large use in description of symmetry operators of systems of partial differential equations. In this paper using these solutions we found full set of symmetry operators of arbitrary finite order for KleinโGordonโFock equations with zero and non-zero mass.
## Appendix A Non-degeneracy of systems <br>of equations for coefficients of Killing tensors
We will adduce proof of Theorem 1 stating non-degeneracy of system of linear algebraic equations (22). As it is cumbersome we adduce it in abridged form.
The main difficulty of the analysis of system (22) is the need to do it for arbitrary value of $`j`$, that is for a system of arbitrary fixed dimension given by formula (20).
Let us consider equations (22) for $`m=4`$, at that equations for $`m<4`$ will be included into the analysis as particular cases. Indices $`a_1,a_2,\mathrm{},a_{j+1}`$ and $`b_1,b_2,\mathrm{},b_k`$ with $`kj`$ independently take values from 1 to 4. At that, as it is easy to notice, the system (22) splits at non-linked subsystems $`M(s_1,s_2,s_3,s_4)`$, where $`s_l`$ $`(l=1,2,3,4)`$ gives the number of indices having the value $`l`$. It is obvious that
$$s_1+s_2+s_3+s_4=j+1+k,$$
(A.1)
so $`0s_lj+k+1`$.
System (22) is non-degeneracy iff all its subsystems $`M(s_1,s_2,s_3,s_4)`$ are non-degeneracy. Without loss of generality, for arbitrary subsystem $`M(s_1,s_2,s_3,s_4)`$ we can put
$$s_1s_2s_3s_4,$$
(A.2)
other cases can be reduced to (A.2) by renumeration of variables.
Let us prove non-degeneracy of an arbitrary subsystem $`M(s_1,s_2,s_3,s_4)`$.
We designate by the symbol $`n_l`$ $`(l=1,2,3,4)`$ the number of indices of unknown variable $`F^{a_1a_2\mathrm{}a_j,a_{j+1}b_1b_2\mathrm{}b_k}`$ present on the left of the comma and equal to $`l`$, and by the symbol $`m_l`$ โ the number of indices after the comma that are equal to $`l`$. Obviously, the following should be satisfied,
$$m_c+n_c=s_c,n_1+n_2+n_3+n_4=j,m_1+m_2+m_3+m_4=k+1,$$
(A.3)
so out of eight numbers $`n_l`$ and $`m_l`$ only three will be linearly independent (see (A.1)). Let us choose the following numbers as independent: $`n_1`$, $`n_2`$ and $`n_3`$, then the triple $`(n_1,n_2,n_3)`$ will completely determine a vector $`F^{a_1a_2\mathrm{}a_j,a_{j+1}b_1b_2\mathrm{}b_k}`$ from the subsystem $`M(s_1,s_2,s_3,s_4)`$. Using for such vector the designation $`F(n_1,n_2,n_3)`$ and considering relations (A.1)โ(A.3), we can write any equation (22) from the subsystem $`M(s_1,s_2,s_3,s_4)`$ in one of the following forms:
$$(n_3+1)F(0,0,n_3)+(jn_3)F(0,0,n_3+1)=0,$$
$$\mathrm{max}\{s_3k1,1\}n_3s_31;$$
(A.4)
$$(n_2+1)F(0,n_2,n_3)+n_3F(0,n_2+1,n_31)+(jn_2n_3)F(0,n_2+1,n_3)=0,$$
$$\mathrm{max}\{0,s_1+s_2k1\}n_2s_21,$$
$$\mathrm{max}\{0,s_1+s_2+s_3k1n_2\}n_3\mathrm{min}\{s_3,jn_2\};$$
(A.5)
$$(n_1+1)F(n_1,n_2,n_3)+n_2F(n_1+1,n_21,n_3)+n_3F(n_1+1,n_2,n_31)$$
$$+(jn_1n_2n_3)F(n_1+1,n_2,n_3)=0,$$
$$\mathrm{max}\{0,s_1k1\}n_1s_11,\mathrm{max}\{0,s_1+s_2k1n_1\}n_2s_2,$$
$$\mathrm{max}\{0,s_1+s_2+s_3k1n_1n_2\}n_3\mathrm{min}\{s_3,jn_1n_2\}.$$
(A.6)
When $`s_1>0`$ (the case $`s_1=0`$ is considered below) there are three possibilities:
$$1.s_1+s_2<k+1,$$
$$2.s_1+s_2k+1,s_1<k+1,$$
$$3.s_1k+1.$$
(A.7)
Let us consider these possibilities one by one.
In the case 1 the system under consideration is given by the formulae (A.4)โ(A.6).
Let us present the vector $`F(n_1,n_2,n_3)`$ in the form of a column whose components are numbered by the index
$$F(n_1,n_2,n_3)=(F(n_1,n_2,\stackrel{~}{n}_3),F(n_1,n_2,\stackrel{~}{n}_3+1)\mathrm{}F(n_1,n_2,\widehat{n}_3))^T,$$
(A.8)
where $`\stackrel{~}{n_3}`$ and $`\widehat{n_3}`$ are minimal and maximal values of $`n_3`$, and each vector $`F(n_1,n_2,\stackrel{~}{n}_3+k)`$, in its turn, will be regarded as a column whose components are numbered by the index $`n_2`$, $`0n_2s_2`$. Then it is possible to write equations (A.4)โ(A.6) in the matrix form:
$$AF=0,$$
(A.9)
where
$$A=\left(\begin{array}{ccccc}B_0& & & & \\ E_1& B_1& & & \\ & 2E_2& B_2& & \\ & & \mathrm{}& \mathrm{}& \\ & & & s_1E& B_{s_1}\end{array}\right),$$
(A.15)
$$B_l=\left(\begin{array}{ccccc}D_l& & & & \\ E_{l1}& D_{l+1}& & & \\ & 2E_{l2}& D_{l+2}& & \\ & & \mathrm{}& \mathrm{}& \\ & & & s_2E_{ls_2}& D_{l+s_2}\end{array}\right).$$
(A.21)
Here $`E_k`$ and $`E_{lf}`$ ($`l,k=1,2,\mathrm{},s_1`$, $`f=1,2,\mathrm{},s_2)`$ are unit matrices whose number of rows coincides with the number of rows of the adjacent matrices on the right (the number of columns of $`B_k`$ ($`D_l`$) coincides with the number of rows of $`B_{k+1}`$ ($`D_{l+1}`$)), $`D_l`$ being the matrices whose explicit form is determined below. Namely, with $`s_1+s_2+s_3<k+1`$:
$$D_l=\left(\begin{array}{ccccc}jl+1& & & & \\ 1& jl& & & \\ & 2& jl1& & \\ & & \mathrm{}& \mathrm{}& \\ & & & s_3& jl+1s_3\end{array}\right)$$
(A.27)
with $`s_1+s_2+s_3k+1`$
$$D_l=\left(\begin{array}{ccccc}a_1& ja_1l+1& & & \\ & a_1+1& ja_1l& & \\ & & \mathrm{}& \mathrm{}& \\ & & & a_2& ja_2l+1\end{array}\right),$$
(A.32)
where $`a_1=\mathrm{max}\{0,s_1+s_2+s_3kl\}`$, $`a_2=\mathrm{min}\{s_3,jl+1\}`$.
At that in the cases $`a_1=0`$ or (and) $`s_3jl+1`$ in the matrix (A.12) the first or (and) last column should be crossed out.
Note A.1. It can be shown that there are always will be less of matrices $`D_l`$ of the form (A.11) (or (A.12) $`a_10`$ $`ja_1l+10`$) than of the matrices (A.12) with $`a_1=ja_2l+10`$.
Our task is to prove that all rows of the matrix $`A`$ (A.10) are linearly independent.
Writing this matrix in the equivalent form
$$A^{}=\left(\begin{array}{ccccc}E_1& B_1& & & \\ & 2F_2& B_2& & \\ & & \mathrm{}& \mathrm{}& \\ & & & s_1E_{s_1}& B_{s_1}\\ B_0& & & & \end{array}\right)$$
and subjecting $`A^{}`$ to the transformation $`A^{}A^{\prime \prime }=VA^{}W`$ that does not change the rank, with $`V`$ and $`W`$ being reversible matrices of the form
$$V=\left(\begin{array}{ccccc}E_1& & & & \\ & E_2& & & \\ & & \mathrm{}& & \\ & & & E_{s_1}& \\ B_0& B_0B_1& \mathrm{}& (1)^{s_1+1}\frac{1}{s_1!}B_0B_1B_2\mathrm{}B_{s_1}& E_0\end{array}\right),$$
$$W=\left(\begin{array}{ccccccc}E_1& B_1& \frac{1}{2!}B_1B_2& \frac{1}{3!}B_1B_2B_3& \mathrm{}& (1)^{s_1}\frac{1}{(s_11)!}B_1B_2\mathrm{}B_{s_11}& \\ & E_2& B_2& \frac{1}{2!}B_2B_3& \mathrm{}& (1)^{s_11}\frac{1}{(s_12)!}B_2B_3\mathrm{}B_{s_11}& \\ & & \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \\ & & & & E_{s_1}& B_{s_11}& \\ & & & & & E_{s_1}& \end{array}\right),$$
we get
$$A^{\prime \prime }=\left(\begin{array}{ccccc}E_1& & & & \\ & E_2& & & \\ & & \mathrm{}& & \\ & & & E_{s_1}& \\ & & & & \frac{1}{s_1!}B_0B_1B_2\mathrm{}B_{s_1}\end{array}\right)$$
(A.38)
and proof of linear independence of the rows of the matrix (A.10) is reduced to proof of linear independence of the rows of the matrix
$$\widehat{B}^{s_1}=B_0B_1\mathrm{}B_{s_1}.$$
(A.39)
Lemma A.1. The matrix $`(A.14)`$ can be split into blocks $`\widehat{B}_{lf}^{s_1}`$ including the same number of rows as the matrix $`D_{s_1+l1}`$, $`1ls_{2+1}`$, $`1fs_{2+1}`$. These blocks have the following form:
$$B_{lf}^{s_1}=\{\begin{array}{c}\widehat{0},s_2lf<0\text{or}s_1+2lfs_2,\hfill \\ \left(\genfrac{}{}{0pt}{}{s_1+1}{lf}\right)P_{lf}^{l1}D_{l1}D_l\mathrm{}D_{s_1+f1},0lfs_1,\hfill \\ \left(\genfrac{}{}{0pt}{}{s_1+1}{s_1+1}\right)P_{s_1+1}^{l1}E,lf=s_1+1.\hfill \end{array}$$
(A.43)
Here $`D_{l1},D_l,\mathrm{}`$ are matrices from $`(A.10)`$, the general form of which is given by the formulae (A.11), (A.12), $`\widehat{0}`$ and $`E_k`$ are zero and unity matrices of corresponding dimensions, $`\left(\genfrac{}{}{0pt}{}{m}{k}\right)=\frac{m!}{k!(mk)!}`$, $`P_m^k=\frac{k!}{(km)!}`$.
Proof can be done by induction, by successive investigation of the products $`B_0B_1`$, $`B_0B_1B_2,\mathrm{}`$.
Lemma A.2. The matrix (A.14) by finite number of elementary transformations can be reduced to the form
$$\widehat{B}^{s_1}=\left(\begin{array}{cccccccc}& & & & & & & D_0D_1\mathrm{}D_{s_1+s_2}\\ & & & & & & \hfill \mathrm{}\text{}& \\ & & & & & \hfill D_{s_11}D_{s_1}\mathrm{}D_{s_2+1}\text{}& & \\ & & & & \hfill D_{s_1}D_{s_1+1}\mathrm{}D_{s_2}\text{}& & & \\ & & & \hfill E_{s_2s_1}\text{}& & & & \\ & & \mathrm{}& & & & & \\ & E_2& & & & & & \\ E_1& & & & & & & \end{array}\right)$$
(A.52)
Proof. We can describe elementary transformations mentioned in the lemma in the following way:
1. Let us go from the matrix (A.14) to the matrix $`\mathrm{\Pi }_d`$ (with $`d`$ successively taking values $`1,2,\mathrm{},s_2s_1`$), using the following algorithm:
1) Represent the matrix (A.14) in the block form (A.15) and operate with โrowsโ with the number $`l`$, including blocks $`\widehat{B}_{l1}^{s_1},\widehat{B}_{l2}^{s_1},\mathrm{},\widehat{B}_{ls_2+1}^{s_1}`$.
2) Multiply $`(t+1)`$-th โrowโ of the matrix $`\widehat{B}^{s_1}`$ by the matrix $`D_{t1}`$, and the โrowโ of the matrix $`\widehat{B}^{s_1}`$ with the number $`t`$ โ by the number $`s_1t+2`$ and subtract the latter from the former. Write the obtained result instead of the โrowโ with the number $`t`$. Perform this operation successively with all โrowsโ for $`t=1,2,\mathrm{},s_1`$.
3) Multiply $`(s_1+2)`$-th โrowโ of the matrix $`\widehat{B}^{s_1}`$ by $`D_{s_1}\mathrm{}D_{s_2}\mathrm{}D_{s_1+d1}`$ and subtract from the obtained result the โrowโ with the number $`s_1+1`$ multiplied by $`d`$. Write the obtained result instead of the โrowโ with the number $`s_1`$, leaving the remaining โrowsโ unchanged.
4) As a result of the described transformations $`\widehat{B}^{s_1}\mathrm{\Pi }_d^{}`$, where $`\mathrm{\Pi }_d^{}`$ is the matrix having in the first column the only non-vanishing element $`(\mathrm{\Pi }_d^{})_{s_1+\mathrm{2\hspace{0.17em}1}}`$. It is possible to get by means of elementary transformations that all elements of the $`s_1+2`$-th โrowsโ would also vanish (except $`(\mathrm{\Pi }_d^{})_{s_1+\mathrm{2\hspace{0.17em}1}}`$).
5) Let us put the $`s_1+2`$-th โrowโ to the lowest position. As a result we get the matrix $`\mathrm{\Pi }_{s_1s_2}`$ of the following form:
$$\mathrm{\Pi }_{s_2s_1}=\left(\begin{array}{ccccc}& & & & M\\ & & & E_{s_2s_1}& \\ & & \mathrm{}& & \\ & E_2& & & \\ E_1& & & & \end{array}\right),$$
(A.58)
where $`M`$ is a matrix that can be split into blocks of the following form:
$$M_{lf}=\{\begin{array}{c}0,s_1lf<0,\hfill \\ \left(\genfrac{}{}{0pt}{}{s_2+1}{lf}\right)P_{l1}^{lf}D_{l1}D_l\mathrm{}D_{f+s_2+1},0lfs_1,\hfill \end{array}(f,l)=1,2,\mathrm{},s_{1+1}.$$
2. Let us simplify the matrix $`M`$ by using successively the adduced algorithm for $`q=0,1,\mathrm{},s_11`$, and for each value of $`q`$ โ for $`k=1,2,\mathrm{},s_1q`$.
1) Let us multiply the $`(k+1)`$-th โrowโ of the matrix $`M`$ by the matrix $`D_{k1}`$, and the $`k`$-th column โ by the number $`s_1k+2`$, and subtract the former from the latter, leaving other โrowsโ unchanged.
2) Perform this operation successively for all $`k=1,2,\mathrm{},s_1q`$, simplifying the matrices obtained at each step by means of elementary transformations vanishing all elements of โrowsโ except one that was the only non-vanishing in its column.
3) Perform operations 1), 2) successively for all $`q=0,1,\mathrm{},s_11`$.
As a result we come to the matrix (A.16). The lemma is proved.
Lemma A.3. Let the matrix $`D=d_{ab}`$, $`a=1,2,\mathrm{},s`$, $`b=1,2\mathrm{},r`$, $`s<r`$ have the rang $`s`$, with all minors $`D`$ being positive, and the matrix $`B`$ have the form
$$B=\left(\begin{array}{ccccc}b_{11}& & & & \\ b_{21}& b_{22}& & & \\ & b_{32}& b_{33}& & \\ & & \mathrm{}& \mathrm{}& \\ & & & b_{r1r2}& b_{r1r1}\\ & & & & b_{rr1}\end{array}\right),$$
(A.65)
where $`b_{kk}>0`$ and $`b_{k+1k}>0`$, $`k=1,2,\mathrm{},r1`$. Then the matrix $`C=DB`$ also has the rank $`s`$, and all minors $`C`$ are positive.
Proof is reduced to direct utilization of the BinetโCauchy formula representing spinors of the matrix product $`DB`$ via sum of the products of minors of the matrix $`D`$ by minors of the matrix $`B`$. As a result each minor of the matrix $`C`$ can be represented as the sum of positive values. The Lemma is proved.
By virtue of the Lemma A.2 proof of linear independence of the matrix $`A`$ (A.10) is reduced to proof of linear independence of rows of the matrix $`๐^d`$,
$$๐^dD_{s_1d}D_{s_1d+1}\mathrm{}D_{s_2+d},d=0,1,2,\mathrm{},s_1,$$
(A.66)
where $`D_{s_{1d}},D_{s_1d+1},\mathrm{}`$ are matrices of the form (A.11) or (A.12). By virtue of Note A.1 the number of rows of each matrix (A.19) does not exceed the number of its columns. Considering successively the products $`D_{s_1d}D_{s_1d+1},D_{s_1d}D_{s_1d+1}D_{s_1d+2},\mathrm{}`$ and using each time either the Silvester inequality or Lemma A.3, it is not difficult to show that the rank of the matrix (A.19) coincides with the number of its rows, and, whence, all rows of the matrix $`A`$ (A.10) are linearly independent.
We have proved non-degeneracy of of the system (A.4)โ(A.6) for the case 1 from (A.7). In the case 2 when $`s_1+s_2k+1`$, the system (A.4)โ(A.6) is reduced to equations (A.5), (A.6) that can also be written in matrix form (A.9) where $`A`$ is given by (A.10), but the blocks $`B_l`$ have a new form. Namely, with $`s_2<k+1`$
$$B_l=\left(\begin{array}{ccccc}(s_1+s_2kl)E_0& D_0& & & \\ & (s_1+s_2kl+1)E_1& D_1& & \\ & & \mathrm{}& \mathrm{}& \\ & & & s_2E_{ks_1+l}& D_{ks_1+l}\end{array}\right),$$
(A.71)
$`l=0,1,\mathrm{},s_1+s_2k1`$, and
$$B_l=\left(\begin{array}{cccc}D_{l+ks_1s_2}& & & \\ E_1& D_{l+ks_1s_2+1}& & \\ & \mathrm{}& \mathrm{}& \\ & & s_2E_{s_2}& D_{l+ks_1}\end{array}\right),$$
(A.76)
if $`s_1+s_2kls_1`$. If $`s_2k`$, then all matrices $`B_l`$, $`0ls_1`$, are given by the formula (A.20).
Explicit form for the matrices $`D_l`$ is given by relations (A.22)โ(A.23):
for $`s_2<k`$
$$D_l=\left(\begin{array}{ccccc}a_1& jla_1& & & \\ & a_{1+1}& jla_11& & \\ & & \mathrm{}& \mathrm{}& \\ & & & a_2& jla_2\end{array}\right),$$
(A.81)
where $`a_1=\mathrm{max}\{0,s_1+s_2+s_3kl1\}`$, $`a_2=\mathrm{min}\{s_3,jl\}`$, $`l=0,1,\mathrm{}kd_1`$,
$$D_l=\left(\begin{array}{cccc}a_1& \hfill jl+k+1s_1s_2a_1\text{}& & \\ & \hfill a_1+1\text{}& \hfill jl+ks_1s_2a_1\text{}& \\ & & \hfill \mathrm{}\text{}& \mathrm{}\\ & & \hfill a_2\text{}& jl+ks_1s_2a_2\end{array}\right),$$
(A.86)
where $`a_1=\mathrm{max}\{0,s_3l\}`$, $`a_2=\mathrm{min}\{s_3,jl+k+1s_1s_2\}`$, $`l=ks_1+1,ks_1+2,\mathrm{},k`$.
If $`a_1=0`$ then first columns in (A.22) and (A.23) should be crossed out, if $`s_3jl`$, then the last column in (A.22) should be crossed out, and with $`s_3jl+k+1s_1s_2`$ it is necessary to cross out the last column in (A.23).
In the case $`s_2k`$ the matrices $`D_l`$ are given by the formula (A.22) for all $`l`$.
Our task is to prove linear independence of rows of the matrix $`A`$ determined by relations (A.10), (A.20)โ(A.23). Transforming this matrix to the form (A.13) we reduce this problem again to investigation of the matrix (A.14) that in our case can be split into blocks of the form
$$\widehat{B}_{lf}^{S_1}=\{\begin{array}{c}\widehat{0},s_2lf<s_1+s_2k\text{or}ks_2+2lfs_2,\hfill \\ \left(\genfrac{}{}{0pt}{}{s_1+1}{l+s_1+s_2kf}\right)P_{l+s_1+s_2kf}^{l+s_1+s_2k1}D_{l1}D_l\mathrm{}D_{ks_21f},\hfill \\ s_1s_2+klfks_2,\hfill \\ \left(\genfrac{}{}{0pt}{}{s_1+1}{s_1+1}\right)P_{s_1+1}^{l+s_1+s_2k1}E,lf=k+1s_2,\hfill \end{array}$$
(A.91)
where $`1lk+1s_2`$, $`1fs_2+1`$.
Further proof is done in full analogy with the proof for the case 1.
Let us consider now the third case from (A.7). The corresponding system of equations (A.4)โ(A.6) is reduced to equations (A.6). Writing these equations in the matrix form (A.9) we come to the corresponding matrix $`A`$ of the following form:
$$A=\left(\begin{array}{ccccc}(s_1k)E_0& B_0& & & \\ & (s_1k+1)E_1& B_1& & \\ & & \mathrm{}& \mathrm{}& \\ & & & s_1E_k& B_k\end{array}\right),$$
(A.96)
where
$$B_r=\left(\begin{array}{ccccc}(s_2R)\stackrel{~}{E}_0& D_0& & & \\ & (s_2R+1)\stackrel{~}{E}_1& D_1& & \\ & & \mathrm{}& \mathrm{}& \\ & & & s_2\stackrel{~}{E}_R& D_R\end{array}\right),R=0,1,\mathrm{},k,$$
$$D_l=\left(\begin{array}{ccccc}(s_3l)& \hfill j+k+1s_1s_2s_3\text{}& & & \\ & \hfill s_3l+1\text{}& j+ks_1s_2s_3& & \\ & & \mathrm{}& \mathrm{}& \\ & & & s_3& j+k+1s_1s_2s_3l\end{array}\right),$$
$`l=0,1,\mathrm{},k`$.
All rows of the matrix $`A`$ (A.25) are evidently linearly independent.
Thus we had proved that the system of equations (A.4)โ(A.6) is non-degenerate in all cases listed in the formulae (A.7). Thus the system (72) in the case $`m=4`$ is non-degenerate.
Considering only such systems of equations (A.4)โ(A.6) that correspond to $`S_1=0`$ we get a full set of non-linked subsystems of the system (72) for $`m=3`$, and in the case $`s_1=s_2=0`$ we come to full set of non-linked subsystems of the system (72) for $`m=2`$. Consequently non-degeneracy of the system (72) for $`m=2`$ and $`m=3`$ follows from the adduced proof as a particular case.
Similarly (but with involvement of somewhat more cumbersome calculations) it is possible to prove non-degeneracy of the systems of linear algebraic equations (74) for coefficients of Killing tensors of rank $`j`$ and order $`s`$.
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# 1 Introduction
## 1 Introduction
This paper deals with the factorization properties of deep inelastic scattering (DIS) in the region of phase space where $`1x\mathrm{\Lambda }_{\mathrm{QCD}}/Q`$, with $`Q`$ the large energy carried by the virtual photon. In this kinematical region the final-state jet carries an energy of order $`Q`$, but has a small invariant mass $`p_x^2=Q^2(1x)/xQ\mathrm{\Lambda }_{\mathrm{QCD}}`$. We assume that perturbation theory is valid at both the hard scale $`Q^2`$ and the jet scale $`Q\mathrm{\Lambda }_{\mathrm{QCD}}`$. The invariant mass of the target proton defines a third, non-perturbative scale $`M_P^2\mathrm{\Lambda }_{\mathrm{QCD}}^2`$. In similar cases in inclusive $`B`$ decay it is possible to derive factorization formulas which separate the physics from the three scales $`Q^2Q\mathrm{\Lambda }_{\mathrm{QCD}}\mathrm{\Lambda }_{\mathrm{QCD}}^2`$ into a convolution of the generic form
$$HJS.$$
(1)
The functions $`H`$ and $`J`$ are perturbatively calculable hard and jet functions depending on fluctuations at the scales $`Q^2`$ and $`Q\mathrm{\Lambda }_{\mathrm{QCD}}`$ respectively, and $`S`$ is a non-perturbative function containing physics at the low-energy scale $`\mathrm{\Lambda }_{\mathrm{QCD}}^2`$. The symbol $``$ stands for a convolution. Our goal is to use effective field theory methods to establish whether such a factorization formula can be derived for deep inelastic scattering in the large-$`x`$ limit.
Recent studies of perturbative factorization in $`B`$ decay have relied heavily on soft-collinear effective theory (SCET) . These include many applications to inclusive decay, both at leading order and including power corrections . For inclusive $`B`$ decay these proofs are rather straightforward. Inclusive decay deals with interactions between hard-collinear particles fluctuating at the jet scale $`m_b\mathrm{\Lambda }_{\mathrm{QCD}}`$ with soft particles fluctuating at the non-perturbative scale $`\mathrm{\Lambda }_{\mathrm{QCD}}^2`$. The leading-order Lagrangian interactions between soft and hard-collinear particles can be decoupled by field redefinitions involving Wilson lines . After integrating out hard fluctuations in a first step of matching, the factorization of the SCET matrix elements into a convolution of jet and soft functions is more or less a natural consequence of this decoupling at the level of the Lagrangian.
Applications of SCET to exclusive decay are considerably more complicated . Exclusive processes typically involve both soft and collinear particles fluctuating at the scale $`\mathrm{\Lambda }_{\mathrm{QCD}}^2`$, in addition to hard-collinear fluctuations which are integrated out in the first step of a two-step matching procedure. It has been argued that a low-energy theory of soft and collinear particles contains a third mode, referred to as soft-collinear . This follows from an analysis of loop diagrams with soft and collinear external lines by the method of regions . This soft-collinear โmessenger modeโ has the special property that it can interact with both soft and collinear particles without taking them far off shell. These modes introduce an additional, highly non-perturbative soft-collinear scale $`\mathrm{\Lambda }_{\mathrm{QCD}}^3/Q`$. To prove factorization formulas of the type in (1), one must show that the effects of this fourth scale are irrelevant to the low-energy matrix elements defining the soft functions $`S`$. This has been emphasized in . In $`B`$ decay the soft-collinear scale is relevant at the endpoints of convolution integrals linking non-perturbative soft and collinear functions, so the soft-collinear field has often been associated with endpoint divergences in these integrals .
In this paper we show that the soft-collinear mode is relevant to an analysis of DIS at large $`x`$. Near the endpoint, DIS involves the three widely separated scales $`Q^2(1x)Q^2\mathrm{\Lambda }_{\mathrm{QCD}}^2`$. Our main finding is that we cannot correlate the two small scales by the definition $`\lambda ^2(1x)\mathrm{\Lambda }_{\mathrm{QCD}}/Q`$ without introducing a fourth scale, $`\mathrm{\Lambda }_{\mathrm{QCD}}^3/QQ^2\lambda ^6`$. The appearance of this fourth scale is associated with the soft-collinear mode. For values of $`x`$ satisfying $`1x\mathrm{\Lambda }_{\mathrm{QCD}}/Q`$, the low-energy matrix element defining the parton distribution function involves fluctuations at both the collinear and soft-collinear scales. An attempt to use effective field theory methods to prove a factorization formula such as (1) leads instead to an expression
$$H\left(\frac{Q^2}{\mu ^2}\right)J(\frac{Q^2(1x)}{\mu ^2},\frac{Q\omega _{sc}}{\mu ^2})f(\frac{\mathrm{\Lambda }_{\mathrm{QCD}}^2}{\mu ^2},\frac{\mathrm{\Lambda }_{\mathrm{QCD}}^2\omega _{sc}}{Q\mu ^2}),$$
(2)
where $`\omega _{sc}\mathrm{\Lambda }_{\mathrm{QCD}}`$ is a convolution variable. Since the parton distribution function $`f`$ contains a non-perturbative dependence on the large energy $`Q`$, factorization is spoiled.
The organization of this paper is as follows. In Section 2 we define our power counting and identify the relevant SCET fields by applying the method of regions to a representative loop diagram. Section 3 deals with matching the QCD Lagrangian and electromagnetic current onto a version of SCET which accounts for these momentum regions. In Section 4 we show with a tree-level example that the parton distribution function is sensitive to soft-collinear effects, and discuss this further in Section 5 with a one-loop calculation. In Section 6 we summarize the implications of the soft-collinear mode on factorization. We compare our results with previous work in Section 7 and conclude in Section 8.
## 2 Power counting and momentum regions
Deep inelastic scattering involves the scattering of an energetic virtual photon with a large invariant mass $`q^2=Q^2`$ off a proton with momentum $`P`$ to form a hadronic jet carrying momentum $`p_x`$ and an invariant mass $`p_x^2=Q^2(1x)/x`$, where
$$x=\frac{q^2}{2Pq}=\frac{Q^2}{2Pq}.$$
(3)
We are interested in the region of phase space where the hadronic jet carries a large energy of order $`Q`$, but has a small invariant mass on the order of the jet scale $`Q\mathrm{\Lambda }_{\mathrm{QCD}}`$. More precisely, we work in the kinematic region where $`p_x^2Q\mathrm{\Lambda }_{\mathrm{QCD}}Q^2(1x)`$. This correlates the two small scales $`1x\mathrm{\Lambda }_{\mathrm{QCD}}/Q`$. We make this explicit in the effective theory by introducing an expansion parameter $`\lambda ^2(1x)\mathrm{\Lambda }_{\mathrm{QCD}}/Q`$. We then calculate the cross section as a double series in the perturbative coupling constant and $`\lambda `$. In terms of $`\lambda `$ the invariants $`Pp_xQ^2,p_x^2Q^2\lambda ^2,`$ and $`P^2Q^2\lambda ^4`$ define three widely separated scales $`Q^2Q^2\lambda ^2Q^2\lambda ^4`$. In this paper we investigate whether we can derive a factorization formula which separates the physics from these scales.
Our analysis relies on soft-collinear effective theory. Unlike in applications of SCET to $`B`$ decay, there is no natural Lorentz frame in which to describe the scattering process. We find the Breit frame most convenient for what follows. In terms of two light-like vectors $`n_\pm `$ satisfying $`n_+n_{}=2`$, the components of the photon momentum $`q^\mu `$ in the Breit frame are given by $`(n_+q,q_{},n_{}q)=(Q,0,Q)`$. If the proton momentum is $`P=(Q/x,0,M_p^2x/Q)`$, then at leading order in $`\lambda `$ the jet momentum $`p_x=q+P`$ is given by $`p_x=(Q(1x)/x,0,Q)`$ and satisfies $`p_x^2=Q^2(1x)/x`$. We will refer to momenta with the scaling $`p_cQ(1,\lambda ^2,\lambda ^4)`$ as collinear, and momenta with the scaling $`p_{\overline{hc}}Q(\lambda ^2,\lambda ,1)`$ as anti-hard-collinear. With this terminology, the proton momentum is collinear and the final-state jet momentum is anti-hard-collinear. Above and in the rest of the paper we work in the reference frame where the transverse components of the external momenta vanish.
To construct the effective theory we must first identify the momentum regions that produce on-shell singularities in loop diagrams. SCET fields are then introduced to reproduce the effects of these momentum regions. The relevant momentum regions depend on the choice of $`x`$. In the kinematical regime where $`1x\mathrm{\Lambda }_{\mathrm{QCD}}/Q`$, we find that we must consider hard, anti-hard-collinear, collinear, and soft-collinear regions.
The appearance of soft-collinear instead of soft modes will have important consequences in our analysis. Before we begin, it is useful to explain their origin in simple terms. For the final-state jet to be anti-hard-collinear requires that $`n_+P+n_+q=n_+p_xQ\lambda ^2`$. This is possible only if $`n_+P=n_+q+\omega `$, where $`\omega `$ is a residual momentum scaling as $`\omega Q\lambda ^2`$ and $`n_+q=Q`$ is a large kinematic piece. This is similar to HQET, where the $`b`$-quark momentum is $`m_bv+k_s`$, with $`k_s`$ a soft residual momentum $`k_sm_b(\lambda ^2,\lambda ^2,\lambda ^2)`$ and $`m_bv`$ a large kinematic piece. For a collinear particle, however, the residual momentum cannot be soft because $`n_{}k_sQ\lambda ^2`$, while for a collinear momentum $`n_{}p_cQ\lambda ^4`$. The simplest possibility is that the residual momentum scales as $`Q(\lambda ^2,\lambda ^3,\lambda ^4)`$, a scaling to which we will refer as soft-collinear. We will show below that this is indeed the scaling which is relevant in loop diagrams. This leads us to interpret the soft-collinear mode as the residual momentum of a collinear field.
We will now make these observations more rigorous by analyzing a loop diagram using the method of regions , similarly to . As a simplification, we begin with the scalar version of the triangle diagram shown in Figure 1. This allows us to identify the relevant momentum regions without complications related to Dirac algebra. The external lines carry a collinear momentum $`p_p`$ and an anti-hard-collinear momentum $`p_x=p_p+q`$. We set all masses to zero, and regularize IR divergences by keeping the external lines off shell by an amount $`p_x^2Q^2\lambda ^2`$ and $`p_p^2P^2Q^2\lambda ^4`$. The integral in the full theory is given by
$`I`$ $`=`$ $`{\displaystyle [dL]\frac{1}{(L+p_x)^2}\frac{1}{(L+p_p)^2}\frac{1}{L^2}}`$ (4)
$`=`$ $`{\displaystyle \frac{1}{Q^2}}\left[\mathrm{ln}{\displaystyle \frac{p_x^2}{Q^2}}\mathrm{ln}{\displaystyle \frac{p_p^2}{Q^2}}+{\displaystyle \frac{\pi ^2}{3}}\right],`$
where we have defined the measure as
$$[dL]=i16\pi ^2\left(\frac{\mu ^2e^{\gamma _E}}{4\pi }\right)^ฯต\frac{d^dL}{(2\pi )^d},$$
(5)
and expanded the result to leading order in $`\lambda `$. At leading order $`Q^2=n_+p_pn_{}p_x`$.
We seek to reproduce this result by the method of regions. This strategy splits the loop integration into contributions from momentum regions according to the scaling of their light-cone components with $`\lambda `$. The integrand is expanded as appropriate for the particular momentum region before evaluating the integral. Once all relevant regions are identified, their sum reproduces the full theory result. We start with the hard region, where the loop momentum scales as $`LQ(1,1,1).`$ The expanded integral is
$`I_h`$ $`=`$ $`{\displaystyle [dL]\frac{1}{(L^2+n_{}p_xn_+L)}\frac{1}{(L^2+n_+p_pn_{}L)}\frac{1}{L^2}}`$ (6)
$`=`$ $`{\displaystyle \frac{1}{Q^2}}\left[{\displaystyle \frac{1}{ฯต^2}}{\displaystyle \frac{1}{ฯต}}\mathrm{ln}{\displaystyle \frac{Q^2}{\mu ^2}}+{\displaystyle \frac{1}{2}}\mathrm{ln}^2{\displaystyle \frac{Q^2}{\mu ^2}}{\displaystyle \frac{\pi ^2}{12}}\right].`$
We have regularized additional divergences with dimensional regularization in $`d=42ฯต`$ dimensions.
The integral $`I_h`$ contains logarithms depending on the hard scale $`Q^2`$. This is a generic feature: the result for a given region always involves logarithms at that momentum scale. For this reason we need to consider the anti-hard-collinear and collinear regions, since these integrals can depend on $`p_x^2`$ and $`p_p^2`$. For the anti-hard-collinear region, where the loop momentum scales as $`LQ(\lambda ^2,\lambda ,1)`$, we find
$`I_{\overline{hc}}`$ $`=`$ $`{\displaystyle [dL]\frac{1}{(L+p_x)^2}\frac{1}{(n_{}Ln_+p_p)}\frac{1}{L^2}}`$ (7)
$`=`$ $`{\displaystyle \frac{1}{Q^2}}\left[{\displaystyle \frac{1}{ฯต^2}}+{\displaystyle \frac{1}{ฯต}}\mathrm{ln}{\displaystyle \frac{p_x^2}{\mu ^2}}{\displaystyle \frac{1}{2}}\mathrm{ln}^2{\displaystyle \frac{p_x^2}{\mu ^2}}+{\displaystyle \frac{\pi ^2}{12}}\right].`$
For the collinear region, where $`LQ(1,\lambda ^2,\lambda ^4)`$, we have
$`I_c`$ $`=`$ $`{\displaystyle [dL]\frac{1}{(n_+Ln_{}p_x)}\frac{1}{(L+p_p)^2}\frac{1}{L^2}}`$ (8)
$`=`$ $`{\displaystyle \frac{1}{Q^2}}\left[{\displaystyle \frac{1}{ฯต^2}}+{\displaystyle \frac{1}{ฯต}}\mathrm{ln}{\displaystyle \frac{p_p^2}{\mu ^2}}{\displaystyle \frac{1}{2}}\mathrm{ln}^2{\displaystyle \frac{p_p^2}{\mu ^2}}+{\displaystyle \frac{\pi ^2}{12}}\right].`$
Taking the sum of the regions considered so far does not reproduce the result for the full integral (4). It is easy to check that $`II_hI_{\overline{hc}}I_c`$ contains logarithms depending on $`p_x^2p_p^2/Q^2Q^2\lambda ^6`$. This is taken into account by including the soft-collinear region, where $`LQ(\lambda ^2,\lambda ^3,\lambda ^4)`$. This region gives
$`I_{sc}`$ $`=`$ $`{\displaystyle [dL]\frac{1}{(n_+Ln_{}p_x+p_x^2)}\frac{1}{(n_{}Ln_+p_p+p_p^2)}\frac{1}{L^2}}`$ (9)
$`=`$ $`{\displaystyle \frac{1}{Q^2}}\left[{\displaystyle \frac{1}{ฯต^2}}+{\displaystyle \frac{1}{ฯต}}\mathrm{ln}{\displaystyle \frac{Q^2\mu ^2}{p_x^2p_p^2}}+{\displaystyle \frac{1}{2}}\mathrm{ln}^2\left({\displaystyle \frac{p_x^2p_p^2}{Q^2\mu ^2}}\right)+{\displaystyle \frac{\pi ^2}{4}}\right].`$
Adding $`I_h+I_{\overline{hc}}+I_c+I_{sc}`$, we see that the poles cancel, and that we recover the result for the full integral given in (4). We will construct a version of SCET which accounts for these momentum regions in the next section.
Note that the soft and hard-collinear regions are needed in applications of SCET to $`B`$ decay, but are not needed here. The soft region, where $`LQ(\lambda ^2,\lambda ^2,\lambda ^2)`$, is irrelevant because
$`I_s`$ $`=`$ $`{\displaystyle [dL]\frac{1}{(n_+Ln_{}p_x+p_x^2)}\frac{1}{(n_{}Ln_+p_p)}\frac{1}{L^2}}`$ (10)
$`=`$ $`{\displaystyle \frac{1}{Q^2}}{\displaystyle [dL]\frac{1}{(n_+L+n_+p_x)}\frac{1}{(n_{}L)}\frac{1}{L^2}}=0.`$
To derive the second line we used $`p_x^2=n_+p_xn_{}p_x`$ (recall $`p_x_{}=0`$), and then that scaleless integrals vanish in dimensional regularization. The hard-collinear integrand, where $`LQ(1,\lambda ,\lambda ^2)`$, is also scaleless and vanishes.
While it may be possible to eliminate the soft-collinear scale $`\mathrm{\Lambda }_{\mathrm{QCD}}^3/Q`$ by introducing an IR regulator to cut off momentum regions with virtuality smaller than the QCD scale $`\mathrm{\Lambda }_{\mathrm{QCD}}^2`$, we find it more convenient to keep the collinear quarks off shell by an amount $`p_p^2P^2`$ and use dimensional regularization. In our end analysis we will adopt the philosophy of , and interpret any sensitivity of low-energy matrix elements to the soft-collinear mode as a breakdown of factorization.
We should emphasize that all results are frame independent. It is also possible to carry out the analysis in the target rest frame, where the proton momentum is soft. We can identify the scaling of the light-cone components of the momentum regions in the rest frame by performing a Lorentz boost to this frame, which amounts to rescaling $`n_\pm `$. The components of a generic momentum change according to $`(n_+p,p_{},n_{}p)(n_+p\lambda ^2,p_{},n_{}p\lambda ^2)`$. The correspondence between the two frames is given by
| | | Breit Frame | | Rest Frame |
| --- | --- | --- | --- | --- |
| | hard | $`Q(1,1,1)`$ | $``$ | $`Q(\lambda ^2,1,\frac{1}{\lambda ^2})`$ |
| | anti-hard-collinear | $`Q(\lambda ^2,\lambda ,1)`$ | $``$ | $`Q(\lambda ^4,\lambda ,\frac{1}{\lambda ^2})`$ |
| | collinear | $`Q(1,\lambda ^2,\lambda ^4)`$ | $``$ | $`Q(\lambda ^2,\lambda ^2,\lambda ^2)`$ |
| | soft-collinear | $`Q(\lambda ^2,\lambda ^3,\lambda ^4)`$ | $``$ | $`Q(\lambda ^4,\lambda ^3,\lambda ^2)`$ |
Although the individual light-cone components of the momentum regions scale differently in the two frames, the number of regions is the same. Moreover, the result for each region depends on invariants at that scale and is therefore frame independent. This can be seen from the explicit results, or by noticing that each integrand is invariant under the simultaneous rescalings of $`n_\pm `$ shown above. In the effective theory this is referred to as reparameterization invariance (RPI) .
## 3 Matching onto SCET
This section deals with matching the QCD Lagrangian and electromagnetic current onto SCET. Our eventual goal is to examine the factorization properties of the hadronic tensor using effective field theory methods. In inclusive processes all QCD effects are contained in the hadronic tensor, which is given by the spin-averaged matrix element between proton states
$$W^{\mu \nu }=\frac{1}{\pi }\mathrm{Im}P|T^{\mu \nu }|P,$$
(11)
where the current correlator $`T^{\mu \nu }`$ is defined through the time-ordered product
$$T^{\mu \nu }=id^4ze^{iqz}\mathrm{T}\left\{J^\mu (z)J^\nu (0)\right\}.$$
(12)
Here $`J^\mu `$ is the electromagnetic current, and $`q`$ is the momentum of the incoming photon. We will evaluate the correlator in effective field theory by separating the contributions from the momentum regions identified in the previous section, namely
| hard | $`Q(1,1,1)`$ |
| --- | --- |
| anti-hard-collinear | $`Q(\lambda ^2,\lambda ,1)`$ |
| collinear | $`Q(1,\lambda ^2,\lambda ^4)`$ |
| soft-collinear | $`Q(\lambda ^2,\lambda ^3,\lambda ^4)`$ |
We calculate the hadronic tensor using a two-step matching procedure familiar from applications of SCET to inclusive $`B`$ decay in the shape-function region . In the first step, we match the QCD Lagrangian and electromagnetic current onto SCET by integrating out fluctuations at the hard scale $`Q^2`$ and introducing effective theory fields for the regions $`p_{\overline{hc}},p_c,p_{sc}`$. The Lagrangian can be derived exactly, and will be discussed in the next sub-section. The current, on the other hand, receives corrections from fluctuations at the hard scale. These corrections can be absorbed into a hard Wilson coefficient, which we will calculate at one loop in Section 3.2. In a second step of matching we evaluate the hadronic tensor (11) using the SCET Lagrangian and current. In this step of matching we integrate out fluctuations at the hard-collinear scale $`Q\mathrm{\Lambda }_{\mathrm{QCD}}`$ and match onto the parton distribution function. We discuss this at tree level in Section 4 and at one loop in Section 5.
### 3.1 SCET Lagrangian
The QCD Lagrangian for light quarks contains no hard scale and the SCET Lagrangian can be derived exactly . For the case at hand, we have
$$_{\mathrm{QCD}}_{c+sc}+_{\overline{hc}+sc}+_{sc}+_{YM},$$
(13)
where $`_{c+sc}`$ contains the collinear Lagrangian as well as interactions with the soft-collinear gluon field, and analogously for $`_{\overline{hc}+sc}`$. There is no interaction term $`_{c+\overline{hc}}`$ for processes where the initial and final states contain only one type of collinear field . The soft-collinear Lagrangian $`_{sc}`$ can be found in , and the Yang-Mills Lagrangian for each sector is the same as in QCD.
The Lagrangian $`_{c+sc}`$ can be derived using the methods of , as was done in . The result for the leading-order Lagrangian $`_{c+sc}`$ is
$`_{c+sc}=\overline{\xi }_c(in_{}D_{c+sc}+(iD/_c_{}m_q){\displaystyle \frac{1}{in_+D_c}}(iD/_c_{}+m_q)){\displaystyle \frac{n/_+}{2}}\xi _c,`$ (14)
where $`iD_{c+sc}^\mu =i^\mu +gA_c^\mu +gA_{sc}^\mu `$. In interactions between collinear and soft-collinear fields the soft-collinear fields are multipole expanded and depend on $`z_{}^\mu =(n_+z)n_{}^\mu /2`$. We have omitted a pure glue interaction term, which will not be needed here. We can derive the Lagrangian $`_{\overline{hc}+sc}`$ by making the replacements $`n_{}n_+`$ and $`\varphi _c\varphi _{\overline{hc}}`$ in the expressions above. The result is
$`_{\overline{hc}+sc}=\overline{\xi }_{\overline{hc}}(in_+D_{\overline{hc}+sc}+iD/_{\overline{hc}_{}}{\displaystyle \frac{1}{in_{}D_{\overline{hc}}}}iD/_{\overline{hc}_{}}){\displaystyle \frac{n/_{}}{2}}\xi _{\overline{hc}}.`$ (15)
We have again omitted a pure glue interaction term. In interactions between anti-hard-collinear and soft-collinear fields the soft-collinear fields must be multipole expanded and depend only on $`z_+^\mu =(n_{}z)n_+^\mu /2`$.
We have included a collinear quark mass $`m_q\mathrm{\Lambda }_{\mathrm{QCD}}Q\lambda ^2`$ in the leading-order Lagrangian $`_{c+sc}`$ above. We are free to include such a mass without changing the regions analysis. In fact, keeping the collinear momentum off shell by an amount $`p_p^2Q^2\lambda ^4`$ effectively gave such a scale to the collinear line, adding an actual mass just changes $`p_p^2p_p^2m_q^2`$ in the collinear propagator. This does not eliminate soft-collinear effects. We checked this claim by modifying the scalar triangle integral to include a mass $`m_qQ\lambda ^2`$ for the collinear line and confirmed that, at least to one loop, the regions analysis is unchanged. We have no proof that the regions analysis is unchanged beyond one loop, and in the following calculations we will always set $`m_q=0`$ for simplicity.
A property of the Lagrangians crucial for factorization proofs is that the soft-collinear fields can be decoupled from the anti-hard-collinear and collinear fields through field redefinitions involving Wilson lines . We introduce the Wilson lines
$`S_{sc}(z)`$ $`=`$ $`\mathrm{P}\mathrm{exp}\left(ig{\displaystyle _{\mathrm{}}^0}๐sn_{}A_{sc}(z+sn_{})\right)`$ (16)
$`S_{\overline{sc}}(z)`$ $`=`$ $`\mathrm{P}\mathrm{exp}\left(ig{\displaystyle _{\mathrm{}}^0}๐sn_+A_{sc}(z+sn_+)\right)`$ (17)
along with similar objects $`W_c`$ and $`W_{\overline{hc}}`$, where the soft-collinear fields are replaced by collinear or anti-hard-collinear fields, and $`n_+n_{}`$. After making the field redefinitions
$`\xi _c`$ $`=`$ $`S_{sc}\xi _c^{(0)},A_c=S_{sc}A_c^{(0)}S_{sc}^{},W_c=S_{sc}W_c^{(0)}S_{sc}^{},`$ (18)
$`\xi _{\overline{hc}}`$ $`=`$ $`S_{\overline{sc}}\xi _{\overline{hc}}^{(0)},A_{\overline{hc}}=S_{\overline{sc}}A_{\overline{hc}}^{(0)}S_{\overline{sc}}^{},W_{\overline{hc}}=S_{\overline{sc}}W_{\overline{hc}}^{(0)}S_{\overline{sc}}^{},`$
the fields with the superscript $`0`$ no longer interact with the soft-collinear fields. This factorization of soft-collinear fields at the level of the Lagrangians does not guarantee the factorization of the current correlator (12), however, because the effects may reappear in time-ordered products with the external currents .
### 3.2 SCET current at one loop
Having obtained the relevant SCET Lagrangian, we now consider the one-loop matching of the electromagnetic current onto its effective field theory expression. This was done previously in and we agree with the results obtained there. We will repeat the calculation to show how logarithms related to the soft-collinear mode are essential to the analysis.
At leading order the matching of the electromagnetic current onto SCET takes the form
$$\overline{\psi }_c(z)\gamma ^\mu \psi _{\overline{hc}}(z)๐s๐t\stackrel{~}{C}(s,t,\mu )(\overline{\xi }_cW_c)(z+sn_+)\gamma ^\mu (W_{\overline{hc}}^{}\xi _{\overline{hc}})(z+tn_{}).$$
(19)
As in , we consider a single quark flavor with unit charge. The convolution arises because $`n_+p_c`$ and $`n_{}p_{\overline{hc}}`$ are on the order of the hard scale, so the operator can be non-local by an amount $`1/Q`$ in these directions. Setting $`z`$ to zero and using translational invariance, the current can be written as
$$C(n_+P_cn_{}P_{\overline{hc}},\mu )(\overline{\xi }_cW_c)(0)\gamma ^\mu (W_{\overline{hc}}^{}\xi _{\overline{hc}})(0),$$
(20)
where the Fourier-transformed coefficient function is
$$C(n_+P_cn_{}P_{\overline{hc}},\mu )=๐s๐t\stackrel{~}{C}(s,t)e^{i(sn_+P_ctn_{}P_{\overline{hc}})},$$
(21)
and $`P_{c,\overline{hc}}`$ are momentum operators. In our case these are both $`Q`$ so we have $`C(Q^2,\mu )`$.
To calculate the one-loop matching conditions we take the difference of the QCD result from that evaluated in SCET. The QCD graph is the same as in Figure 1 but evaluated with the Feynman rules of QCD. We find it useful to break up the QCD result into contributions from each momentum region, as we did with the scalar triangle. The matching conditions are related only to the hard region. For the QCD result we find
$`I_{\mathrm{QCD}}`$ $`=`$ $`{\displaystyle \frac{C_F\alpha _s}{4\pi }}\gamma ^\mu \left[{\displaystyle \frac{1}{ฯต_{\mathrm{UV}}}}\mathrm{ln}{\displaystyle \frac{Q^2}{\mu ^2}}2\mathrm{ln}{\displaystyle \frac{p_p^2}{Q^2}}\mathrm{ln}{\displaystyle \frac{p_x^2}{Q^2}}2\mathrm{ln}{\displaystyle \frac{p_p^2}{Q^2}}2\mathrm{ln}{\displaystyle \frac{p_x^2}{Q^2}}{\displaystyle \frac{2\pi ^2}{3}}\right]`$ (22)
$`=`$ $`I_h+I_{\overline{hc}}+I_c+I_{sc},`$
where
$`I_h`$ $`=`$ $`{\displaystyle \frac{C_F\alpha _s}{4\pi }}\gamma ^\mu \left[{\displaystyle \frac{2}{ฯต^2}}+{\displaystyle \frac{2}{ฯต}}\left(\mathrm{ln}{\displaystyle \frac{Q^2}{\mu ^2}}2\right)+{\displaystyle \frac{1}{ฯต_{\mathrm{UV}}}}\mathrm{ln}^2{\displaystyle \frac{Q^2}{\mu ^2}}+3\mathrm{ln}{\displaystyle \frac{Q^2}{\mu ^2}}+{\displaystyle \frac{\pi ^2}{6}}8\right],`$ (23)
$`I_{\overline{hc}}`$ $`=`$ $`{\displaystyle \frac{C_F\alpha _s}{4\pi }}\gamma ^\mu \left[{\displaystyle \frac{2}{ฯต^2}}{\displaystyle \frac{2}{ฯต}}\left(\mathrm{ln}{\displaystyle \frac{p_x^2}{\mu ^2}}1\right)+\mathrm{ln}^2{\displaystyle \frac{p_x^2}{\mu ^2}}2\mathrm{ln}{\displaystyle \frac{p_x^2}{\mu ^2}}{\displaystyle \frac{\pi ^2}{6}}+4\right],`$ (24)
$`I_c`$ $`=`$ $`{\displaystyle \frac{C_F\alpha _s}{4\pi }}\gamma ^\mu \left[{\displaystyle \frac{2}{ฯต^2}}{\displaystyle \frac{2}{ฯต}}\left(\mathrm{ln}{\displaystyle \frac{p_p^2}{\mu ^2}}1\right)+\mathrm{ln}^2{\displaystyle \frac{p_p^2}{\mu ^2}}2\mathrm{ln}{\displaystyle \frac{p_p^2}{\mu ^2}}{\displaystyle \frac{\pi ^2}{6}}+4\right],`$ (25)
$`I_{sc}`$ $`=`$ $`{\displaystyle \frac{C_F\alpha _s}{4\pi }}\gamma ^\mu \left[{\displaystyle \frac{2}{ฯต^2}}+{\displaystyle \frac{2}{ฯต}}\mathrm{ln}{\displaystyle \frac{p_x^2p_p^2}{Q^2\mu ^2}}\mathrm{ln}^2{\displaystyle \frac{p_x^2p_p^2}{Q^2\mu ^2}}{\displaystyle \frac{\pi ^2}{2}}\right].`$ (26)
We have expanded all results to leading order in $`\lambda `$, and used that $`Q^2=n_+p_pn_{}p_x`$ at this order. We must supplement these graphs with the wave-function renormalization for off-shell quarks, which gives a contribution
$$I_w=\frac{1}{2}\frac{C_F\alpha _s}{4\pi }\gamma ^\mu \left[\frac{1}{ฯต_{\mathrm{UV}}}1+\mathrm{ln}\frac{p_i^2}{\mu ^2}\right]$$
(27)
for each external quark line. The UV poles cancel in the sum $`I_h+I_w`$, as required by current conservation.
The next step is to evaluate the SCET diagrams in Figure 2. Evaluating the graphs in the figure using the Feynman rules of SCET reproduces the result for the QCD regions calculation. By this we mean that Figure 2(a) evaluates to $`I_{sc}`$, Figure 2(b) to $`I_{\overline{hc}}`$, and Figure 2(c) to $`I_c`$. This just confirms that we have constructed the effective theory correctly. The wave-function graphs in the effective theory are the same as (27) .
The difference between the two theories is that the hard integral $`I_h`$ is absent in SCET. Its finite part is taken into account by a hard matching coefficient, and its infinite part is reproduced by a renormalization factor $`Z_J`$ applied to the bare current. Including the tree-level contribution, the matching coefficient is therefore
$$C(Q^2,\mu )=1+\frac{C_F\alpha _s}{4\pi }\left[\mathrm{ln}^2\frac{Q^2}{\mu ^2}+3\mathrm{ln}\frac{Q^2}{\mu ^2}+\frac{\pi ^2}{6}8\right],$$
(28)
and the renormalization factor is
$$Z_J=1+\frac{C_F\alpha _s}{4\pi }\left[\frac{2}{ฯต^2}\frac{3}{ฯต}+\frac{2}{ฯต}\mathrm{ln}\frac{Q^2}{\mu ^2}\right].$$
(29)
The hard coefficient $`C(Q^2,\mu )`$ and the renormalization factor $`Z_J`$ depend on the hard scale $`Q^2`$. For the infinite counter terms, this is possible only after a cancellation between logarithms that occurs when adding the anti-hard-collinear, collinear, and soft-collinear graphs. That logarithms of UV origin related to the soft-collinear field are needed to ensure that the renormalization factor $`Z_J`$ depends only on the hard scale $`Q^2`$ was first noted in , in a slightly different context. On the other hand, no such cancellation occurs for the finite terms, where the sum of the anti-hard-collinear, collinear, and soft collinear graphs still contains logarithms at each scale. We will see in Sections 5 and 6 that the matrix element of the current correlator shares this property, and that the logarithms related to the soft-collinear scale cause problems for factorization.
## 4 Matching onto parton distributions at tree level
Matching onto the intermediate theory has absorbed the effects of hard fluctuations into a short-distance Wilson coefficient. This leaves the three widely separated scales $`Q\mathrm{\Lambda }_{\mathrm{QCD}}\mathrm{\Lambda }_{\mathrm{QCD}}^2\mathrm{\Lambda }_{\mathrm{QCD}}^3/Q`$. We now examine the factorization properties of the hadronic tensor (11). To achieve a perturbative factorization of the form (1), we would need to show that the soft-collinear scale $`\mathrm{\Lambda }_{\mathrm{QCD}}^3/Q`$ is irrelevant. We could then perform a second and final step of matching at the scale $`Q\mathrm{\Lambda }_{\mathrm{QCD}}`$, and identify the associated matching coefficient with the jet function $`J`$. The low-energy matrix element would define a parton distribution function $`f`$ characterized by fluctuations at the collinear scale $`\mathrm{\Lambda }_{\mathrm{QCD}}^2`$ only, which would be linked to $`J`$ in a convolution integral. The purpose of this section is to demonstrate with a tree-level example that this is impossible, by showing that soft-collinear effects do not decouple from the low-energy matrix element. In fact, the jet function is linked to the parton distribution function by a convolution variable related to the soft-collinear scale. We will argue that a full separation of scales would require integrating out the collinear scale $`\mathrm{\Lambda }_{\mathrm{QCD}}^2`$ in a third step of matching, which however cannot be done perturbatively.
At leading order, the current correlator (12) is given by the time-ordered product
$$T^{\mu \nu }=id^4ze^{iqz}\mathrm{T}\left\{\overline{\chi }_c(z)\gamma ^\mu \chi _{\overline{hc}}(z)\overline{\chi }_{\overline{hc}}(0)\gamma ^\nu \chi _c(0)\right\},$$
(30)
where we have defined the fields
$$\chi _cW_c^{}\xi _c,\chi _{\overline{hc}}W_{\overline{hc}}^{}\xi _{\overline{hc}},$$
(31)
which are manifestly gauge invariant under anti-hard-collinear and collinear gauge transformations. At tree level and to lowest order in $`g`$ the hard Wilson coefficient $`C(Q^2,\mu )`$ and the Wilson lines $`W`$ are unity. We perform a second step of matching by integrating out the anti-hard-collinear fields. This is done at the scale $`Q\mathrm{\Lambda }_{\mathrm{QCD}}`$, which we treat as perturbative. To do this at tree level, we first perform the decoupling redefinition (18) on the anti-hard-collinear fields (and immediately drop the superscript (0)), and then contract the anti-hard-collinear fields into a propagator. This is represented by the Feynman diagram in Figure 3(a). The anti-hard-collinear propagator is given in momentum space by
$$0|\xi _{\overline{hc}}(z)_{a\alpha }\overline{\xi }_{\overline{hc}}(0)_{b\beta }|0=\frac{d^4L}{(2\pi )^4}e^{iLz}\frac{in_{}L}{L^2+i0}\left(\frac{n/_+}{2}\right)_{\alpha \beta }\delta _{ab}.$$
(32)
This forces $`z`$ to scale as an anti-hard-collinear quantity, and we need to perform the multipole expansion accordingly. This is in general different from the multipole expansion in SCET Lagrangian interactions, because the photon injects a large external momentum $`q`$ into the diagram. In particular, since $`z`$ scales as anti-hard-collinear, the collinear and soft-collinear fields can depend only on $`z_+^\mu =(n_{}z/2)n_+^\mu `$. The result for the current correlator is then
$`T^{\mu \nu }`$ $`=`$ $`{\displaystyle d^4ze^{iqz}\overline{\xi }_c(z_+)\gamma ^\mu S_{\overline{sc}}(z_+)\frac{n/_+}{2}S_{\overline{sc}}^{}(0)\gamma ^\nu \xi _c(0)}`$ (33)
$`{\displaystyle \frac{d^4L}{(2\pi )^4}e^{iLz}\frac{n_{}L}{n_{}Ln_+L+L_{}^2+i0}}`$
The soft-collinear and collinear fields do not depend on $`n_+z`$ or $`z_{}`$, so we can perform these integrations. We then have
$`T^{\mu \nu }`$ $`=`$ $`{\displaystyle \frac{d(n_{}z)}{2}\frac{d(n_+L)}{2\pi }e^{in_+Ln_{}z/2}\frac{n_{}q}{n_{}qn_+L+i0}}`$ (34)
$`e^{in_+qn_{}z/2}\overline{\xi }_c(z_+)S_{\overline{sc}}(z_+)S_{\overline{sc}}^{}(0)\gamma ^\mu {\displaystyle \frac{n/_+}{2}}\gamma ^\nu \xi _c(0).`$
Even though $`n_+p_cQ`$ and $`n_+p_{sc}Q\lambda ^2`$ we cannot set the argument of the soft-collinear Wilson line $`S_{\overline{sc}}(z_+)`$ to zero. This is because $`n_+p_c+n_+qQ\lambda ^2`$, so we need to keep $`n_+p_{sc}Q\lambda ^2`$ in the $`n_+LQ\lambda ^2`$ component of the anti-hard-collinear propagator. We will discuss this further below. For now, we simply note that soft-collinear effects do not decouple even at leading order in the $`1/Q`$ expansion.
In order to calculate the hadronic tensor (11) we now take the matrix element of the current correlator between proton states. We define a parton distribution function through the spin-averaged matrix element
$`P|\overline{\chi }_c(tn_+)S_{\overline{sc}}(tn_+)S_{\overline{sc}}^{}(0)\gamma ^\mu {\displaystyle \frac{n/_+}{2}}\gamma ^\nu \chi _c(0)|P=\stackrel{~}{f}(t)\mathrm{tr}\left[{\displaystyle \frac{n/_{}}{2}}\gamma ^\mu {\displaystyle \frac{n/_+}{2}}\gamma _\nu \right](n_+q).`$ (35)
The factor of $`n_+q=Q+๐ช(Q\lambda ^2)`$ preserves manifest boost invariance. Although not necessary for tree-level matching, we have reinserted the Wilson lines $`W_c`$ in order to define a gauge invariant hadronic matrix element. The Fourier transformed function is
$$\stackrel{~}{f}(t)=๐\omega e^{i\omega t}f(\omega ).$$
(36)
Inserting this into (33), the hadronic tensor becomes
$$W^{\mu \nu }=\frac{1}{\pi }\mathrm{Im}๐\omega f(\omega )\frac{Q}{n_+q\omega +i0}\mathrm{tr}\left[\frac{n/_{}}{2}\gamma ^\mu \frac{n/_+}{2}\gamma ^\nu \right].$$
(37)
As written, (37) obscures the power counting in the effective theory. We have used a delta function to eliminate $`n_+LQ\lambda ^2`$, so we must have $`n_+q\omega Q\lambda ^2`$. This requires that $`\omega =n_+q+\omega _{sc}`$, where $`\omega _{sc}Q\lambda ^2`$. The large component of the collinear momentum simply balances that of the incoming photon momentum, it is the residual soft-collinear component $`\omega _{sc}`$ which controls the dynamics. We can make this transparent by integrating out the collinear scale $`\mathrm{\Lambda }_{\mathrm{QCD}}^2`$ and matching onto a low-energy theory defined at the soft-collinear scale. This is of course not possible because QCD is already strongly coupled at the collinear scale, and we will revisit this point below. At tree level, however, such a matching is trivial: we simply introduce a new field by writing $`e^{iQn_{}z/2}\xi _c(z)=\xi _{Q,sc}(z)`$, so that $`\xi _{Q,sc}`$ carries the residual momentum $`\omega _{sc}`$.<sup>1</sup><sup>1</sup>1This field is similar to the SCET field $`\xi _{Q,n}`$ defined in the label formalism , the difference being that the residual momentum is soft-collinear instead of ultra-soft. Figure 3(b) shows the tree-level diagram for the current correlator evaluated using the $`\xi _{Q,sc}`$. We calculate the correlator using the same steps as before, and define a distribution function $`f_{sc}`$ analogously to (35), but in terms of $`\xi _{Q,sc}`$ instead of $`\chi _c`$. We furthermore choose $`n_+q=Q+n_+p_x`$, with $`n_+p_x=Q(1x)`$. The result is
$$W^{\mu \nu }=๐\omega _{sc}J(n_+p_x\omega _{sc})f_{sc}(\omega _{sc})\mathrm{tr}\left[\frac{n/_{}}{2}\gamma ^\mu \frac{n/_+}{2}\gamma ^\nu \right],$$
(38)
where
$$J(n_+p_x\omega _{sc})=\frac{1}{\pi }\mathrm{Im}\frac{Q}{n_+p_x\omega _{sc}+i0}=Q\delta (n_+p_x\omega _{sc}).$$
(39)
This completes the tree-level matching calculation. We could have obtained this same result in the free-quark decay picture by calculating the diagram in Figure 3 and taking the imaginary part. In the free-quark picture at tree level we can interpret $`f_{sc}(\omega _{sc})=\delta (\omega _{sc})`$, so that the convolution $`Jf_{sc}`$ reproduces the result for the diagram. We went through the extra step of defining the parton distribution function in terms of a hadronic matrix element in order to draw some parallels between (38) and the factorization formula derived for inclusive $`B`$ decay in the shape-function region, where one finds a convolution of the form
$$๐\omega _sJ(n_+p_x\omega _s)S(\omega _s),$$
(40)
with $`n_+p_x=m_b(1x)\mathrm{\Lambda }_{\mathrm{QCD}}`$. The function $`J(\omega _s)`$ is a perturbatively calculable jet function containing physics at the scale $`m_b\mathrm{\Lambda }_{\mathrm{QCD}}`$, and $`S(\omega _s)`$ is a shape function containing non-perturbative effects at the scale $`\mathrm{\Lambda }_{\mathrm{QCD}}^2`$. It is defined by the HQET matrix element
$$S(\omega _s)=๐te^{it\omega _s}\overline{B}_v|\overline{h}_v(tn_{})h_v(0)|\overline{B}_v.$$
(41)
The crucial difference between $`B`$ decay and DIS is that the heavy-quark field $`h_v`$ carries a residual momentum $`\omega _s`$ which is soft. Matching can be stopped at the perturbative scale $`m_b\mathrm{\Lambda }_{\mathrm{QCD}}`$. This should be compared with (38), where the convolution involves the soft-collinear residual momentum $`\omega _{sc}`$. To isolate the physics at this low scale requires an extra step of matching at the scale $`p_p^2\mathrm{\Lambda }_{\mathrm{QCD}}^2`$. We did this above in order to define $`f_{sc}`$, but it is important to understand that this was only a formal manipulation. Since we cannot do this matching perturbatively, we must always lump the collinear and soft-collinear effects together into one non-perturbative function. We will emphasize in Section 6 that any sensitivity of the parton distribution function to the soft-collinear scale signals a breakdown of factorization. However, to explain how effective field theory could in principle be used to separate all the scales, we end this section by considering a fictitious QCD where perturbation theory is valid at the collinear scale. In this fictitious theory we can remove collinear fluctuations when matching the electromagnetic current onto SCET. This matching takes the form
$$\overline{\psi }_c\gamma ^\mu \psi _{\overline{hc}}C(Q^2,\mu )D_c(p_p^2,\mu )\overline{\xi }_{Q,sc}\gamma ^\mu \chi _{\overline{hc}}.$$
(42)
The matching coefficient $`D_c(p_p^2,\mu )`$ reproduces the effects of collinear loop diagrams, and could be obtained at one loop from the finite part of our expressions in Section 3.2, see (23-26). To consider such a (fictitious) matching of the SCET current will be useful in some of the discussion in the next two sections.
## 5 Matching onto parton distributions at one loop
In this section we examine the one-loop corrections to the current correlator (12), and interpret the results in terms of the effective theory. The relevant one-loop diagrams are shown in Figure 4. Note that the graph in Figure 4(e) containing collinear exchange, as well as graphs 4(h) and 4(i), are not actually SCET graphs. In these graphs the short-dashed propagator is hard, not anti-hard-collinear. It was our intention to remove all hard fluctuations in the first step of matching, but we have clearly not done so. Although these graphs are power suppressed by a factor of $`p_x^2/Q^2\lambda ^2`$, we find it awkward to generate power-suppressed graphs from the leading-order Feynman rules of the effective theory. A formal solution to this problem is to remove the collinear scale when matching the current, as in (42). We cannot do this matching perturbatively, but we are interested only in the sum of collinear and soft-collinear graphs, which can equally well be written in this way. This simplifies the book-keeping, because after taking this step the graphs 4(e), 4(h) and 4(i) no longer exist, so the short-dashed propagator is always anti-hard-collinear. Note that we have not drawn box diagrams related to gluon distributions. These are power suppressed, either because the intermediate propagator is hard, analogously to 4(h), or because they involve insertions of soft-collinear quark fields (not to be confused with $`\xi _{Q,sc}`$), which are absent from the leading order SCET Lagrangian .
We now give results for the remaining diagrams, which we calculate using the free-quark picture. As before, we keep the external collinear quarks off shell by an amount $`p_p^2`$ when performing the matching. We work in Feynman gauge. With this choice of gauge graphs 4(d) and 4(e) vanish, as do the parts of 4(b), 4(f) involving soft-collinear exchange, since $`n_\pm ^2=0`$. We suppress the Dirac structure, which is always the same as in the tree-level expression (35) after summing over spins, and also the Wilson coefficients $`C^2(Q^2,\mu )`$, which appear as a multiplicative factor.
The non-vanishing anti-hard-collinear graphs add up to
$`T_{\overline{hc}}={\displaystyle \frac{C_F\alpha _s}{4\pi }}{\displaystyle \frac{Q^2}{p_x^2}}\{[{\displaystyle \frac{1}{ฯต}}1+\mathrm{ln}{\displaystyle \frac{p_x^2}{\mu ^2}}]`$
$`+[{\displaystyle \frac{4}{ฯต^2}}+{\displaystyle \frac{4}{ฯต}}(1\mathrm{ln}{\displaystyle \frac{p_x^2}{\mu ^2}})+2\mathrm{ln}^2{\displaystyle \frac{p_x^2}{\mu ^2}}4\mathrm{ln}{\displaystyle \frac{p_x^2}{\mu ^2}}{\displaystyle \frac{\pi ^2}{3}}+8]\},`$ (43)
the collinear graphs (including wave-function graphs) evaluate to
$`T_c={\displaystyle \frac{C_F\alpha _s}{4\pi }}{\displaystyle \frac{Q^2}{p_x^2}}\{[{\displaystyle \frac{1}{ฯต}}1+\mathrm{ln}{\displaystyle \frac{p_p^2}{\mu ^2}}]`$
$`+[{\displaystyle \frac{4}{ฯต^2}}+{\displaystyle \frac{4}{ฯต}}(1\mathrm{ln}{\displaystyle \frac{p_p^2}{\mu ^2}})+2\mathrm{ln}^2{\displaystyle \frac{p_p^2}{\mu ^2}}4\mathrm{ln}{\displaystyle \frac{p_p^2}{\mu ^2}}{\displaystyle \frac{\pi ^2}{3}}+8]\},`$ (44)
and the soft-collinear graphs give
$$T_{sc}=\frac{C_F\alpha _s}{4\pi }\frac{Q^2}{p_x^2}\left[\frac{4}{ฯต^2}+\frac{4}{ฯต}\mathrm{ln}\frac{p_x^2p_p^2}{Q^2\mu ^2}2\mathrm{ln}^2\frac{p_x^2p_p^2}{Q^2\mu ^2}\pi ^2\right].$$
(45)
The $`1/ฯต`$ terms in the sum of all diagrams is subtracted by the current renormalization factor $`Z_J^2`$ given in (29). This is possible only after the same cancellation between logarithms in the divergent pieces of the anti-hard-collinear, collinear, and soft-collinear graphs that we observed when matching the SCET current. No such cancellation occurs in the finite pieces, where logarithms at each scale remain. We can interpret the finite parts as the one-loop corrections to a convolution of functions characterizing the physics at the various scales. This takes the form
$$\frac{1}{\pi }\mathrm{Im}\left(T_{\overline{hc}}+T_c+T_{sc}\right)=J^{(1)}\left[S^{(0)}f_{sc}^{(0)}\right]+J^{(0)}\left[S^{(1)}f_{sc}^{(0)}+S^{(0)}f_{sc}^{(1)}\right].$$
(46)
The superscript refers to the $`n`$-loop correction to each function. We have defined a function $`S=D_c^2`$ (see (42)), which takes into account collinear effects, and grouped the sum of the collinear and soft-collinear corrections inside the square brackets. In this step of matching we want to obtain the one-loop correction to the jet function $`J`$. To do this rigorously, we would first need to calculate the renormalized expression for the object $`\left[Sf_{sc}\right]`$, using the free-quark picture. This calculation would require a more precise formulation of an effective theory defined at the soft-collinear scale. We will not go through this exercise here, but rather assume that we can construct a low-energy theory that properly accounts for the IR physics related to the collinear and soft-collinear fields. The difference between this low-energy theory and SCET is that the anti-hard-collinear fields are absent, so the the matching function $`J`$ is given by imaginary part of the finite piece of $`T_{\overline{hc}}`$ in (5). This imaginary part is singular at $`p_x^2=0`$ and must be interpreted in terms of distributions to be integrated against a smooth function $`F(p_x^2)`$. To stay in the region where the SCET treatment is valid requires a cut on $`p_x^2`$, so we find it convenient to express the results in terms of star distributions, which are defined as
$`{\displaystyle _0^z}๐xF(x)\left({\displaystyle \frac{1}{x}}\right)_{}^{[u]}`$ $`=`$ $`{\displaystyle _0^z}๐x{\displaystyle \frac{F(x)F(0)}{x}}+F(0)\mathrm{ln}{\displaystyle \frac{z}{u}},`$
$`{\displaystyle _0^z}๐xF(x)\left({\displaystyle \frac{\mathrm{ln}(x/u)}{x}}\right)_{}^{[u]}`$ $`=`$ $`{\displaystyle _0^z}๐x{\displaystyle \frac{F(x)F(0)}{x}}\mathrm{ln}{\displaystyle \frac{z}{u}}+{\displaystyle \frac{F(0)}{2}}\mathrm{ln}^2{\displaystyle \frac{z}{u}}.`$ (47)
In terms of these distributions, we find
$`J^{(1)}`$ $``$ $`\left[S^{(0)}f_{sc}^{(0)}\right]={\displaystyle \frac{1}{\pi }}\mathrm{Im}T_{\overline{hc}}`$ (48)
$`=`$ $`Q^2{\displaystyle \frac{C_F\alpha _s}{4\pi }}\left[\left(7\pi ^2\right)\delta (p_x^2)3\left({\displaystyle \frac{1}{p_x^2}}\right)_{}^{[\mu ^2]}+4\left({\displaystyle \frac{\mathrm{ln}(p_x^2/\mu ^2)}{p_x^2}}\right)_{}^{[\mu ^2]}\right].`$
In the limit $`x1`$ the star distribution is related to the plus distribution, and (48) agrees with a corresponding expression in . This matching function also appears in inclusive $`B`$ decay in the shape-function region .
Having obtained an expression for the one-loop jet function, we end this section by taking a closer look at the low-energy physics relevant to the parton distribution function. If the hadronic tensor obeyed the factorization formula (1), then the low-energy physics would depend on the collinear scale $`\mathrm{\Lambda }_{\mathrm{QCD}}^2`$ only. However, we have shown that the low-energy theory contains a product of collinear and soft-collinear functions, what we called $`\left[Sf_{sc}\right]`$ above, and that this object contains logarithms at both the collinear and soft-collinear scales. The soft-collinear scale $`\mathrm{\Lambda }_{\mathrm{QCD}}^3/Q`$ depends on the large energy $`Q`$, so not all of the $`Q`$ dependence has been factorized into the hard coefficient $`C(Q^2,\mu )`$. We will explain the consequences of this in the next section.
## 6 Soft-collinear effects and factorization
In this section we consolidate our results concerning the factorization of the hadronic tensor. To summarize, we found that for values of $`x`$ satisfying $`1x\mathrm{\Lambda }_{\mathrm{QCD}}/Q`$, DIS involves four scales
$$Q^2Q\mathrm{\Lambda }_{\mathrm{QCD}}\mathrm{\Lambda }_{\mathrm{QCD}}^2\mathrm{\Lambda }_{\mathrm{QCD}}^3/Q.$$
(49)
The relevance of the soft-collinear scale $`\mathrm{\Lambda }_{\mathrm{QCD}}^3/Q`$ makes it impossible to derive a factorization formula of the type (1). To clarify this, we find it useful to first consider a fictitious version of QCD, where the collinear scale $`\mathrm{\Lambda }_{\mathrm{QCD}}^2`$ is perturbative and the soft-collinear scale is non-perturbative. In this fictitious QCD, we can derive a factorization formula by matching onto a low-energy theory defined at the soft-collinear scale. To do so, we split up the initial-state parton momentum as $`p_p=Qn_{}/2+p_{sc}`$, where the soft-collinear residual momentum satisfies $`n_+p_{sc}Q(1x)\mathrm{\Lambda }_{\mathrm{QCD}}`$ and $`n_{}p_{sc}M_P^2/Q\mathrm{\Lambda }_{\mathrm{QCD}}^2/Q`$. This treats the parton as a massless on-shell collinear quark carrying momentum $`Qn_{}/2`$, which receives a residual momentum $`p_{sc}`$ through interactions with soft-collinear partons. Beyond tree level the factorization formula contains a convolution between $`f_{sc}`$ and $`D_c`$, in addition to that between $`f_{sc}`$ and $`J`$. This is because $`n_{}p_{sc}n_{}p_c`$, so the $`n_{}p_{sc}`$ momentum can be distributed between the collinear and soft-collinear fields, just as the $`n_+p_{sc}n_+p_x`$ momentum can be distributed between the anti-hard-collinear and soft-collinear fields. Writing the mass scales (49) in terms of $`Q`$, $`p_{sc}`$, and $`p_x`$ we find a factorization formula of the form
$$WH\left(\frac{Q^2}{\mu ^2}\right)J\left(\frac{Qn_+p_xQn_+p_{sc}}{\mu ^2}\right)f_{sc}\left(\frac{n_{}p_{sc}n_+p_{sc}}{\mu ^2}\right)S\left(\frac{Qn_{}p_{sc}}{\mu ^2}\right),$$
(50)
where $`n_+p_{sc}`$ and $`n_{}p_{sc}`$ are convolution variables. The hard function $`H`$ and the soft function $`S`$ are related to the Wilson coefficients arising when matching the SCET current as in (42), $`H=C^2`$ and $`S=D_c^2`$. The jet function $`J`$ is calculated as explained in Section 5, and the soft-collinear function $`f_{sc}`$ is defined by the spin-averaged matrix element in (35), but with $`\chi _c\xi _{Q,sc}`$. The jet function $`J`$ and the collinear function $`S`$ are not linked directly through a convolution. Instead, they are linked to each other only through a mutual convolution with the function $`f_{sc}`$. Although this formula is qualitatively different from (1), we could in principle derive the renormalization group equations for this effective theory, and use them to resum all large logarithms involving the ratio $`\mathrm{\Lambda }_{\mathrm{QCD}}/Q`$. This scenario has been mentioned in , in analogy with techniques used for the off-shell Sudakov form factor .
Our derivation of (50) was based on an effective field theory approach that integrated out the larger scales until reaching the smallest scale, which is soft-collinear. In real QCD it is not possible to use perturbation theory at the scale $`\mathrm{\Lambda }_{\mathrm{QCD}}^2`$. This obligates us to stop the matching procedure at the jet scale $`Q\mathrm{\Lambda }_{\mathrm{QCD}}`$ and lump the collinear and soft-collinear effects into one non-perturbative function. We have seen that some cancellations occur between the sum of the infinite parts of the collinear and soft-collinear graphs, but this does not occur in the finite pieces defining the matrix elements. The hadronic tensor therefore takes the form
$$WH\left(\frac{Q^2}{\mu ^2}\right)J\left(\frac{Qn_+p_xQn_+p_{sc}}{\mu ^2}\right)f(\frac{\mathrm{\Lambda }_{\mathrm{QCD}}^2}{\mu ^2},\frac{\mathrm{\Lambda }_{\mathrm{QCD}}^2n_+p_{sc}}{Q\mu ^2}),$$
(51)
where we have inserted the physical scaling $`n_{}p_{sc}\mathrm{\Lambda }_{\mathrm{QCD}}^2/Q`$. The notation makes clear that the parton distribution function $`f`$ contains physics at both the collinear and soft-collinear scales. We cannot match perturbatively at the scale $`\mathrm{\Lambda }_{\mathrm{QCD}}^2`$, so we have no way of deriving a low-energy theory that would allow us to resum logarithms at the soft-collinear scale, and large logarithms depending on $`\mathrm{\Lambda }_{\mathrm{QCD}}/Q`$ remain. In other words, the parton distribution function contains a non-perturbative dependence on the large energy $`Q`$. This is different from both (1) and (50). We conclude that a perturbative factorization of scales is not possible in this region of phase space.
## 7 Comparison with previous work
### 7.1 Diagrammatic Approach
Factorization formulas for deep inelastic scattering near the endpoint have been derived using diagrammatic methods in . It seems that the effective field theory calculation leads us to different conclusions concerning the perturbative factorization of scales. The differences can be traced directly to the soft-collinear mode. In turn, we found that the soft-collinear mode is relevant in a very specific region of phase space, where $`1x`$ is correlated with $`\mathrm{\Lambda }_{\mathrm{QCD}}/Q`$ through the relation $`1x\mathrm{\Lambda }_{\mathrm{QCD}}/Q\lambda ^2`$. To the best of our knowledge, such a power counting has not been implemented within the diagrammatic approach, where one takes the limit $`1x0`$ without making the above-mentioned correlation. To understand the significance of this, recall that the effective field theory approach led us to split the parton distribution function into two parts according to
$$fS\left(\frac{Qn_{}p_{sc}}{\mu ^2}\right)f\left(\frac{n_{}p_{sc}n_+p_{sc}}{\mu ^2}\right)S\left(\frac{\mathrm{\Lambda }_{\mathrm{QCD}}^2}{\mu ^2}\right)f\left(\frac{\mathrm{\Lambda }_{\mathrm{QCD}}^2(1x)}{\mu ^2}\right).$$
(52)
A very similar observation has been made in the diagrammatic approach, where $`S`$ and $`f`$ are related to $`\varphi `$ and $`V`$ . The function $`f`$ is linked to the jet function $`J`$ by the convolution variable $`n_+p_{sc}`$. Boost invariance and dimensional analysis require that this enter the parton distribution function in the combination $`n_{}p_{sc}n_+p_{sc}/\mu ^2`$. From this alone it is apparent that the parton distribution function involves fluctuations at two scales, as shown above. The second scale depends on $`1x`$, and need not be soft-collinear. One sees this clearly from (9). For generic values of $`p_x^2Q^2(1x)`$, the soft-collinear region is replaced by an $`x`$-dependent soft region scaling as $`(Q(1x),\mathrm{\Lambda }_{\mathrm{QCD}}\sqrt{(1x)},\mathrm{\Lambda }_{\mathrm{QCD}}^2/Q`$). The function $`f`$ is associated with the vacuum matrix element of a Wilson loop built out of gauge fields with this scaling. As long as $`n_+p_{sc}Q(1x)Q\lambda _D`$, with $`\lambda _D`$ numerically small but still $`๐ช(1)`$, then both the collinear modes and these additional soft modes are parametrically of the order $`\mathrm{\Lambda }_{\mathrm{QCD}}^2`$. Formulas derived with the diagrammatic approach are valid within this particular large-$`x`$ limit. We emphasize that this is a different large-$`x`$ limit than that considered in our work. The non-factorizable soft-collinear effects studied here emerge for values of $`x`$ satisfying $`1x\mathrm{\Lambda }_{\mathrm{QCD}}/Q`$. To apply effective field theory methods in the most straightforward way requires that we make this correlation, because only then can we calculate the results as an expansion in a single small parameter $`\lambda .`$ This power counting for $`1x`$ also ensures that we avoid the resonance region, where $`1x\mathrm{\Lambda }_{\mathrm{QCD}}^2/Q^2`$. The failure of factorization for $`1x\mathrm{\Lambda }_{\mathrm{QCD}}/Q`$ suggests that the most useful application of SCET to DIS in the endpoint region might instead use a multi-scale approach to study the limit $`1x\mathrm{\Lambda }_{\mathrm{QCD}}/Q`$ more carefully. This would involve replacing the soft-collinear modes by the $`x`$-dependent soft modes identified above, carefully re-deriving the factorization formula for the large-$`x`$ limit obtained within the diagrammatic approach , and studying power corrections in terms of SCET operators. This could make use of techniques similar to those developed for the multi-scale operator expansion in inclusive $`B`$ decay .
### 7.2 SCET based approach
The first application of SCET to DIS can be found in , which is however limited to the standard OPE region and has little overlap with our work. In Manohar carried out a SCET analysis of DIS at large $`x`$, also using a two-step matching procedure. In the first step, the author matched QCD in the Breit frame onto a version of SCET involving hard-collinear fields interacting with anti-hard-collinear fields via soft gluon exchange. This differs from the version of SCET used here, which involves collinear fields interacting with anti-hard-collinear fields via soft-collinear gluon exchange. While our two approaches differ conceptually, our results for the anomalous dimension and hard matching coefficient of the SCET current agree. The results are the same because the leading-order Lagrangians $`_{c+sc},_{\overline{hc}+sc}`$ are of the same form as $`_{hc+s},_{\overline{hc}+s}`$. The author used these results to derive some interesting consequences for the anomalous dimension of the SCET current. We disagree on some points concerning the calculation of the hadronic tensor in the second step of matching. The major difference is that found that the effects of soft gluon exchange are irrelevant to the low-energy matrix element defining the parton distribution function (in the Breit frame). Using the translation between the leading-order Lagrangians given above, this would imply the irrelevance of soft-collinear effects, which we did not observe here. This also contradicts the results for the large-$`x`$ limit derived in the diagrammatic approach, where the low-energy matrix element splits into a product of collinear and soft functions, often called $`\varphi `$ and $`V`$ .
## 8 Conclusions
We used soft-collinear effective theory to examine the factorization properties of deep inelastic scattering in the region of phase space where $`1x\mathrm{\Lambda }_{\mathrm{QCD}}/Q`$. An analysis of loop diagrams in the Breit frame showed that the appropriate effective theory includes anti-hard-collinear, collinear, and soft-collinear fields. We found that soft-collinear effects ruin perturbative factorization. An attempt to use SCET to prove a perturbative factorization formula yields instead an expression where the low-energy matrix element defining the parton distribution function contains a non-perturbative dependence on the large energy $`Q`$. It is therefore impossible to separate the three scales $`Q^2Q\mathrm{\Lambda }_{\mathrm{QCD}}\mathrm{\Lambda }_{\mathrm{QCD}}^2`$ in terms of a factorization formula. These complications related to the soft-collinear mode are similar to those found in a SCET analysis of the heavy-to-light form factors relevant to exclusive $`B`$ meson decay . They do not appear in an analysis of factorization for inclusive $`B`$ decay in the shape-function region, where the presence of a heavy quark ensures that soft instead of soft-collinear fields are relevant to the effective theory construction.
Our conclusions are true as long as $`1x`$ is correlated with $`\mathrm{\Lambda }_{\mathrm{QCD}}/Q`$ through the relation $`1x\mathrm{\Lambda }_{\mathrm{QCD}}/Q`$. If $`1x`$ is numerically small but still larger than $`\mathrm{\Lambda }_{\mathrm{QCD}}/Q`$, the standard large-$`x`$ factorization formula derived within the diagrammatic approach is valid. As $`1x`$ approaches the endpoint, however, non-factorizable soft-collinear effects emerge. It would be interesting to use a multi-scale effective field theory approach to carefully re-derive the large-$`x`$ factorization formula using SCET, quantify power corrections in terms of SCET operators, and more carefully study the limit $`1x\mathrm{\Lambda }_{\mathrm{QCD}}/Q`$.
## Acknowledgements
I am grateful to Thomas Becher, Thorsten Feldmann, Thomas Mannel, and Matthias Neubert for useful discussions and comments on the manuscript. This work was supported by the DFG Sonderforschungsbereich SFB/TR09 โComputational Theoretical Particle Physicsโ.
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# Quantum phase transitions of spin chiral nanotubes
## 1 Motivation
These days the experimental techniques have made a lot of progress and good experimentalists have synthesized the nanoscale magnets which realize the models that had been discussed only theoretically before. Among them are the BIP-TENO as the spin-$`1`$ two-leg ladder , the compound \[(CuCl<sub>2</sub>tachH)<sub>3</sub>Cl\]Cl<sub>2</sub> (tach=cis,trans-1,3,5-triamino-cyclohexane) as the triangular spin nanotube , and the oxygen molecules adsorbed on the inner surface of the porous material as a many-leg spin tube . The spin ladders and tubes have attracted attention in the course of the extensive studies on the low-dimensional magnets, part of which started from the discussions on the Haldane gap of the one-dimensional spin chains. The appearance of the spin excitation gap is one of the interesting macroscopic quantum phenomena. In spin-$`1/2`$ systems, it occurs when the frustration in the spin-spin interaction or the lattice dimerization is introduced. A good example of the former case is the triangular spin nanotube whose spontanous dimerized and gapped ground state had been theoretically investigated and now got a strong possibility of an experimental realization in the near future . During the course of studies on the low-dimensional systems, the huge amount of studies on the carbon nanotubes were done in the 90โs, and the striking properties on the electric conductivity of the chiral nanotubes were found out. Namely, a carbon nanotube can be either a metal or a semiconductor, either of which is determined by the chiral vector .
With the great progress in experiments and theories, it is now interesting in both respects to ask, if the carbon atoms are replaced by some magnetic ions on the chiral nanotubes and the antiferromagnetic superexchange interactions between them are introduced, are there any relations between the chiral vector and the magnetic properties analogous to the carbon nanotubes ? The real experiments should be done in the near future. As for the numerical experiments, we take the spin-$`1/2`$ Heisenberg antiferromagnets on the chiral nanotubes made from the honeycomb strips as the simplest systems that the quantum effects bring about interesting properties. Remarkably, these models always enable us to perform numerical experiments by the quantum Monte Carlo method without suffering from the notorious negative sign problem . This is not generally the case for the spin nanotubes because they sometimes contain frustration in the spin-spin interactions . The details are explained in the next section. We investigate the ground state properties of the spin chiral nanotubes with several kinds of chiral vectors.
## 2 Model
Our model is the spin-$`1/2`$ antiferromagnetic Heisenberg tube. The tube is made from a strip of the honeycomb lattice. An example of the strip is shown in Fig. 1. The honeycomb strip can be seen as the coupled dimerized chains in the strongest dimerization limit with the configuration of the dimers in the anti-phase way along the interchain direction. Using the terminology of the low-dimensional strongly correlated systems, we call the bond along the dimerized chains as the โlegโ and that along the interchain direction as the โrungโ. The Hamiltonian is written as follows.
$`H`$ $`=J_x{\displaystyle \underset{k=1}{\overset{n}{}}}{\displaystyle \underset{j=1}{\overset{L}{}}}\left[{\displaystyle \frac{1+(1)^{j+k}}{2}}\right]๐_{j,k}๐_{j+1,k}`$ (1)
$`+J_y{\displaystyle \underset{j=1}{\overset{L}{}}}\left({\displaystyle \underset{k=1}{\overset{n1}{}}}๐_{j,k}๐_{j,k+1}+๐_{j,n}๐_{j+m,1}\right)`$
Here $`๐_{j,k}`$ is the spin-$`1/2`$ operator with $`j`$ ($`k`$) specifying the position on the leg (rung). The number of legs is $`n`$, and the number of sites along the leg is $`L`$. We impose the periodic boundary condition in both of the rung and the leg directions such that $`๐_{L+1,k}๐_{1,k}`$ and $`๐_{j,n+1}๐_{j+m,1}`$. If $`m`$ is finite, chirality is introduced in the tube structure. The coupling constants $`J_x`$ and $`J_y`$ are positive. We set the $`x`$-axis along the dimerized chains. The ground state property in the thermodynamic limit, $`L\mathrm{}`$, is investigated.
In our spin nanotube, the same structure as that of the carbon nanotube is kept when the imposed periodic boundary condition along the rung direction satisfies the condition that $`(n+m)`$ is an even number. In this case, there is no frustration between the spin-spin interactions and we can perform numerical experiments by the quantum Monte Carlo method without suffering from the negative sign problem . The chiral vector of the tube is defined by $`p๐_1+q๐_2`$, where the lattice vectors $`๐_1`$ and $`๐_2`$ are shown in Fig. 1. In terms of our parameters, it is written as $`(p,q)=((m+n)/2,(mn)/2)`$.
## 3 Method and Results
We determine the quantum critical points of the model and discuss their distribution in our parameter space. The positions of them are specified by the parameter $`RJ_y/(J_x+J_y)`$ which connects continuously the strong leg coupling limit at $`R=0`$ to the strong rung coupling limit at $`R=1`$. The phase diagram is drawn on the $`R`$ axis. We note that the isotropic nanotube is realized on the point $`R=0.5`$. The phase boundaries are determined by the the expectation value of the Lieb-Schultz-Mattis slow-twist operator , which is used to detect the bond order specific to the gapped states. For our model, it is defined as follows.
$$z_L=\mathrm{exp}\left[i\frac{2\pi }{L}\underset{j=1}{\overset{L}{}}j\left(\underset{k=1}{\overset{n}{}}S_{j,k}^z\right)\right]_0$$
(2)
Here $`<>_0`$ means the expectation value in the ground state. The quantum critical points are given by the thermodynamic limit of the zeros of $`z_L`$. The observables including $`z_L`$ are evaluated by the quantum Monte Carlo method with the continuous-time loop algorithm , which is actually a finite-temperature algorithm, for several finite sizes. Assuming the Lorentz invariance , the temperature $`T`$ is fixed to the value of $`T=1/L`$. The critical point in both of the limits $`L,T\mathrm{}`$ i.e. the thermodynamic limit at zero temperature, is determined by the crossing point of $`z_L`$โs calculated for several $`L`$โs. An example is shown in Fig. 2 for the chiral nanotube with $`(n,m)=(3,1)`$. We can identify the quantum critical point at $`R_\mathrm{c}=0.38\pm 0.01`$. Thus determined quantum critical points for several chiral nanotubes are shown in Fig. 3.
## 4 Discussions
The characteristic of the resultant distribution of the quantum critical points in Fig. 3 is the difference with respect to the parity of the number of the fully dimerized chains, $`n`$. For the nanotube with even (odd) $`n`$, the number of the quantum critical points is seen to be $`n/2`$ ($`(n1)/2`$). This reflects the well-known even-odd effect in the ground state of uniform (i.e. non-dimerized) spin ladders that is expected to hold also to the spin tubes without frustration. In this effect, the even-legged uniform systems have an excitation gap while the odd-legged ones do not.
We sum up with a few remarks on the comparison of the spin nanotubes with the carbon ones. We note that $`n`$ is written in terms of the chiral vector, $`(p,q)`$, as $`n=pq`$. The topology of the ground-state phase diagram of the spin nanotubes is determined by the parity of $`n`$, in contrast to the carbon nanotubes where the electric property on the point $`R=0.5`$ is known by asking whether $`(pq)`$ is a multiple of $`3`$ or not. So far we have only determined the quantum critical points of the spin chiral nanotubes and discussed the topology of the ground-state phase diagram with respect to the chiral vector. Further investigations are now in progress.
## Acknowledgements
The authors (MM and TS) thank Prof. Masaki Mito and Prof. Hiroyuki Nojiri for valuable discussions. One of the authors (MM) would like to express gratitude to Mr. Naoki Kobayashi, Mr. Takashi Mesaki, and Prof. Riichiro Saito for introducing him the physics of carbon nanotubes and to Mr. Shinsei Ryu for useful comments. The loop algorithm codes for the present calculations are based on the library โLOOPER version 2โ developed by one of the authors (ST) and Dr. K. Kato and the codes for parallel simulations are based on the library โPARAPACK version 2โ also developed by one of the authors (ST). The numerical calculations for the present work were done on the SGI 2800 at the Supercomputer center in the Institute for Solid State Physics, University of Tokyo and on the HPC Server at the Condensed Matter and Statistical Physics Theory Group, Tohoku University.
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# Phase diagram of two-species Bose-Einstein condensates in an optical lattice
## I Introduction
Since the realization of Bose-Einstein condensates (BECs) in dilute atomic gases, a number of interesting experiments have been conducted to investigate multispecies Bose gases, in which two or more states of condensates exist together in a magnetic or optical trap . Recently, vortex states have been obtained in a two-species Bose gas . Progress in the experiments exploring dilute mixtures of quantum gases has stimulated intensive research on the properties of mixed Bose gases at zero temperature and finite temperature as well .
The BEC trapped in an optical lattice exhibites a novel feature, namely the quantum phase transition between a Mott-insulator and a superfluid . Such quantum phase transition has attracted considerable attention in recent years. As a matter of fact, atomic gas of bosons in BEC subjected to a lattice potential which is turned on smoothly can be kept in the superfluid phase as long as the atom-atom interactions are small comparing with the tunnel coupling. In this regime, the kinetic energy is dominant in the total energy of the boson system. With an increase of the potential depth of the optical lattice, it is getting more and more difficult for bosons to tunnel from one site to the other, and finally the system attends an insulator phase above a critical value of the potential depth. In this case, the phase coherence is absent and the number of boson atoms in each lattice site becomes the same. The system possesses a Mott-insulator behavior. Various approaches have been proposed to understand theoretically the quantum phase transition and to determine the phase diagram as a function of BEC parameters .
Motivated by both the experimental and theoretical progress, we in the present paper study the phase diagram for superfluid and insulator phases of two-species BECs in a one-dimensional (1D) optical lattice and the property of persistent current as well. The paper is organized as follows. In Sec. II, the exact macroscopic wave functions of the condensates which are not in the tight-binding regime are constructed by solving the coupled nonlinear Schrรถdinger equations. In Sec. III, the phase diagram is determined analytically according to the macroscopic wave functions of the condensates, i.e., the order parameters. Finally, we summarize our results in Sec. IV.
## II The exact macroscopic wave functions
In this 1D geometry, the confinement along the radial direction is so tight that the trap frequency $`\omega _0`$ along the radial direction is much greater than the mean-field interaction energy. At low temperatures, the dynamics of the atoms in the radial direction is essentially โfrozen,โ with all the atoms occupying the ground state of the harmonic trap with the wave function that
$$\varphi _0(y,z)=\sqrt{\frac{1}{\pi l_0^2}}\mathrm{exp}\left[\left(y^2+z^2\right)/2l_0^2\right]\text{.}$$
(1)
Here the extension of the wave function in the radial direction is given by the length scale $`l_0\sqrt{\mathrm{}/m\omega _0}`$ of a harmonic oscillator, where $`m`$ is the mass of the atoms.
In the mean-field regime, $`l_0`$ is much greater than the radius of the interatomic potential $`R_e`$. The scattering of atoms in this effective 1D system is thus still a process of three-dimension. According to ref. , the effective coupling constant in this 1D system is
$$g_{1D}=\frac{2\mathrm{}^2}{m}\frac{a}{l_0\left(l_0Ca\right)}\text{,}$$
(2)
where $`a`$ is the s-wave scattering length and $`C`$ is a numerical constant of the order unit. The term $`Ca`$ in Eq. (2) is negligible for $`l_0>>R_e`$. In this limit, the expression for $`g_{1D}`$ is the same as that obtained by averaging over the radial wave function (1),
$$g_{1D}=g_{3D}_{\mathrm{}}^{\mathrm{}}_{\mathrm{}}^{\mathrm{}}๐y๐z\varphi _0^4(y,z)=\frac{2\mathrm{}^2a}{ml_0^2}\text{.}$$
(3)
We use this expression in the rest of this paper.
We consider the two-species BECs in a 1D periodic potential. The energy functional is seen to be
$`E[\psi _1,\psi _2]`$ $`=`$ $`{\displaystyle }dx\{{\displaystyle \underset{i=1,2}{}}\left[{\displaystyle \frac{\mathrm{}^2}{2m_i}}\right|{\displaystyle \frac{\psi _i}{x}}|^2+V_i(x)|\psi _i|^2`$ (5)
$`+{\displaystyle \frac{\mathrm{}^2a_i}{m_il_i^2}}|\psi _i|^4]+{\displaystyle \frac{2\mathrm{}^2a_{12}}{\sqrt{m_1m_2}l_1l_2}}|\psi _1|^2|\psi _2|^2\}\text{,}`$
where $`\psi _i`$, $`m_i`$, and $`l_i=\sqrt{\mathrm{}/m_i\omega _0}`$ are the macroscopic wave functions of the condensates, the mass, and the harmonic-oscillator lengths in the radial direction of the $`i`$th species $`\left(i=1,2\right)`$, respectively. $`a_1`$, $`a_2`$, and $`a_{12}`$ denote the s-wave scattering lengths between same-species and interspecies collisions. $`V_i(x)`$ are the periodic potentials,
$$V_i(x)=V_{0,i}\text{sn}^2(k_Lx,k)\text{,}$$
(6)
with $`V_{0,i}`$ denoting the magnitude of potentials, where $`k_L=2\pi /\lambda `$ is the wavevector of the laser light and $`\lambda `$ is the wavelength, corresponding to a lattice period $`d=\lambda /2`$. sn$`(k_Lx,k)`$ is the Jacobian elliptic sine function with modulus $`k`$ $`\left(0k1\right)`$. In the limit $`k=0`$, the Jacobian elliptic sine reduces to the sinusoid function and thus $`V(x)`$ possesses a standard form of the standing light wave. For values of $`k<0.9`$, the potential is virtually indistinguishable from a standing light wave. Finally, for $`k1`$, $`V(x)`$ becomes an array of well-separated hyperbolic secant potential barriers or wells.
The governing equations of the trapped BECs are obtained in terms of the variational procedure ,
$$i\mathrm{}\frac{\psi _i}{t}=\frac{\delta E}{\delta \psi _i^{}}\text{,}$$
(7)
which leads to the coupled nonlinear Schrรถdinger equations
$`i\mathrm{}{\displaystyle \frac{\psi _1}{t}}`$ $`=`$ $`{\displaystyle \frac{\mathrm{}^2}{2m_1}}{\displaystyle \frac{^2\psi _1}{x^2}}+{\displaystyle \frac{2\mathrm{}^2a_1}{m_1l_1^2}}|\psi _1|^2\psi _1`$ (9)
$`+{\displaystyle \frac{2\mathrm{}^2a_{12}}{\sqrt{m_1m_2}l_1l_2}}|\psi _2|^2\psi _1+V_1(x)\psi _1\text{,}`$
$`i\mathrm{}{\displaystyle \frac{\psi _2}{t}}`$ $`=`$ $`{\displaystyle \frac{\mathrm{}^2}{2m_2}}{\displaystyle \frac{^2\psi _2}{x^2}}+{\displaystyle \frac{2\mathrm{}^2a_{12}}{\sqrt{m_1m_2}l_1l_2}}|\psi _1|^2\psi _2`$ (11)
$`+{\displaystyle \frac{2\mathrm{}^2a_2}{m_2l_2^2}}|\psi _2|^2\psi _2+V_2(x)\psi _2\text{.}`$
For the case of weakly coupled condensates in an optical lattice , the wave function $`\psi `$ can be decomposed as a sum of wave functions localized in each well of the periodic potential (tight-binding approximation) with the assumption relying on the fact that the height of the interwell barrier is much higher than the chemical potential. We, however, do not restrict ourselves to the low-energy case and look for the global condensate wave functions of excitations: $`\psi _i(x,t)=\varphi _i(x)\mathrm{exp}\left(i\mu _it/\mathrm{}\right)`$, where $`\mu _i`$ $`(i=1,2)`$ are the chemical potentials. Thus the spatial wave functions satisfy the stationary coupled nonlinear Schrรถdinger equations that
$`\mu _1\varphi _1`$ $`=`$ $`{\displaystyle \frac{\mathrm{}^2}{2m_1}}{\displaystyle \frac{^2\varphi _1}{x^2}}+{\displaystyle \frac{2\mathrm{}^2a_1}{m_1l_1^2}}|\varphi _1|^2\varphi _1`$ (13)
$`+{\displaystyle \frac{2\mathrm{}^2a_{12}}{\sqrt{m_1m_2}l_1l_2}}|\varphi _2|^2\varphi _1+V_1(x)\varphi _1\text{,}`$
$`\mu _2\varphi _2`$ $`=`$ $`{\displaystyle \frac{\mathrm{}^2}{2m_2}}{\displaystyle \frac{^2\varphi _2}{x^2}}+{\displaystyle \frac{2\mathrm{}^2a_{12}}{\sqrt{m_1m_2}l_1l_2}}|\varphi _1|^2\varphi _2`$ (15)
$`+{\displaystyle \frac{2\mathrm{}^2a_2}{m_2l_2^2}}|\varphi _2|^2\varphi _2+V_2(x)\varphi _2\text{.}`$
With the general form of spatial wave functions $`\varphi _i(x)`$ written as $`\varphi _i(x)=r_i(x)\mathrm{exp}\left[i\phi _i\left(x\right)\right]`$, Eq. (13) can be separated as real and imaginary parts. We then integrate once for the imaginary part and obtain the first-order differential equations for the phases $`\phi _i(x)`$,
$$\phi _i^{^{}}(x)=\frac{\alpha _i}{r_i^2(x)}\text{,}$$
(16)
where parameters $`\alpha _i`$ $`(i=1,2)`$ are constants of integration to be determined. Substituting Eq. (16) into the real part obtained from Eq. (13) and integrating again, we find
$`\left(r_1r_1^{^{}}\right)^2`$ $`=`$ $`{\displaystyle \frac{2a_1}{l_1^2}}r_1^6{\displaystyle \frac{2m_1\mu _1}{\mathrm{}^2}}r_1^4+\beta _1r_1^2\alpha _1^2`$ (19)
$`+{\displaystyle \frac{4a_{12}\sqrt{m_1}}{\sqrt{m_2}l_1l_2}}r_1^2{\displaystyle r_2^2d\left(r_1^2\right)}`$
$`+{\displaystyle \frac{2m_1}{\mathrm{}^2}}r_1^2{\displaystyle V_1\left(x\right)d\left(r_1^2\right)\text{,}}`$
$`\left(r_2r_2^{^{}}\right)^2`$ $`=`$ $`{\displaystyle \frac{2a_2}{l_2^2}}r_2^6{\displaystyle \frac{2m_2\mu _2}{\mathrm{}^2}}r_2^4+\beta _2r_2^2\alpha _2^2`$ (22)
$`+{\displaystyle \frac{4a_{12}\sqrt{m_2}}{\sqrt{m_1}l_1l_2}}r_2^2{\displaystyle r_1^2d\left(r_2^2\right)}`$
$`+{\displaystyle \frac{2m_2}{\mathrm{}^2}}r_2^2{\displaystyle V_2\left(x\right)d\left(r_2^2\right)\text{,}}`$
where $`\beta _i`$ $`(i=1,2)`$ denote additional constants of integration.
We then construct the solutions as
$$r_i^2(x)=A_i\text{sn}^2(k_Lx,k)+B_i\text{,}$$
(23)
where the constants $`B_i`$ $`(i=1,2)`$ determine the mean amplitudes and act as the dc offsets for the numbers of the condensed atoms , and the parameters $`A_i`$ $`(i=1,2)`$ are to be determined.
Substituting Eq. (23) into Eq. (19) and using identities of Jacobian elliptic functions, we obtain eight equations for the parameters $`\alpha _i`$, $`\beta _i`$, $`\mu _i`$, and $`A_i`$. Eliminating $`\beta _i`$, we find
$`A_1`$ $`=`$ $`{\displaystyle \frac{\frac{\sqrt{m_1}l_1l_2a_{12}}{\sqrt{m_2}}\left(m_2V_{0,2}\mathrm{}^2k_L^2k^2\right)a_2l_1^2\left(m_1V_{0,1}\mathrm{}^2k_L^2k^2\right)}{2\mathrm{}^2\left(a_1a_2a_{12}^2\right)}}\text{,}`$ (24)
$`A_2`$ $`=`$ $`{\displaystyle \frac{\frac{\sqrt{m_2}l_1l_2a_{12}}{\sqrt{m_1}}\left(m_1V_{0,1}\mathrm{}^2k_L^2k^2\right)a_1l_2^2\left(m_2V_{0,2}\mathrm{}^2k_L^2k^2\right)}{2\mathrm{}^2\left(a_1a_2a_{12}^2\right)}}\text{,}`$ (26)
and
$`\alpha _1^2`$ $`=`$ $`B_1k_L^2\left[{\displaystyle \frac{k^2}{A_1}}B_1^2+\left(1+k^2\right)B_1+A_1\right]\text{,}`$ (27)
$`\alpha _2^2`$ $`=`$ $`B_2k_L^2\left[{\displaystyle \frac{k^2}{A_2}}B_2^2+\left(1+k^2\right)B_2+A_2\right]\text{,}`$ (28)
and
$`\mu _1`$ $`=`$ $`{\displaystyle \frac{\mathrm{}^2k_L^2}{2m_1}}(1+k^2+{\displaystyle \frac{6a_1}{l_1^2k_L^2}}B_1+{\displaystyle \frac{4a_{12}\sqrt{m_1}}{l_1l_2k_L^2\sqrt{m_2}}}B_2`$ (30)
$`+{\displaystyle \frac{2a_{12}\sqrt{m_1}}{l_1l_2k_L^2\sqrt{m_2}}}{\displaystyle \frac{A_2}{A_1}}B_1+{\displaystyle \frac{m_1V_{0,1}}{\mathrm{}^2k_L^2}}{\displaystyle \frac{B_1}{A_1}})\text{,}`$
$`\mu _2`$ $`=`$ $`{\displaystyle \frac{\mathrm{}^2k_L^2}{2m_2}}(1+k^2+{\displaystyle \frac{6a_2}{l_2^2k_L^2}}B_2+{\displaystyle \frac{4a_{12}\sqrt{m_2}}{l_1l_2k_L^2\sqrt{m_1}}}B_1`$ (32)
$`+{\displaystyle \frac{2a_{12}\sqrt{m_2}}{l_1l_2k_L^2\sqrt{m_1}}}{\displaystyle \frac{A_1}{A_2}}B_2+{\displaystyle \frac{m_2V_{0,2}}{\mathrm{}^2k_L^2}}{\displaystyle \frac{B_2}{A_2}})\text{.}`$
For $`k=0`$, $`\text{sn}(k_Lx,0)=\mathrm{sin}(k_Lx)`$, the solutions reduce to
$$\psi _i(x,t)=\sqrt{A_i^0\mathrm{sin}^2\left(k_Lx\right)+B_i}\mathrm{exp}\left\{i\left[\phi _i^0\left(x\right)\mu _i^0t/\mathrm{}\right]\right\}\text{,}$$
(33)
where
$`A_1^0`$ $`=`$ $`{\displaystyle \frac{\sqrt{m_1m_2}a_{12}l_1l_2V_{0,2}m_1a_2V_{0,1}l_1^2}{2\mathrm{}^2\left(a_1a_2a_{12}^2\right)}},`$ (34)
$`A_2^0`$ $`=`$ $`{\displaystyle \frac{\sqrt{m_2m_1}a_{12}l_1l_2V_{0,1}m_2a_1V_{0,2}l_2^2}{2\mathrm{}^2\left(a_1a_2a_{12}^2\right)}}.`$ (35)
The phases $`\phi _i^0\left(x\right)`$ $`(i=1,2)`$ are determined by nonlinear equations
$$\mathrm{tan}\left[\phi _i^0\left(x\right)\right]=\pm \sqrt{1+\frac{A_i^0}{B_i}}\mathrm{tan}\left(k_Lx\right)$$
(36)
and
$`\mu _1^0`$ $`=`$ $`{\displaystyle \frac{\mathrm{}^2k_L^2}{2m_1}}(1+{\displaystyle \frac{6a_1}{l_1^2k_L^2}}B_1+{\displaystyle \frac{4a_{12}\sqrt{m_1}}{l_1l_2k_L^2\sqrt{m_2}}}B_2`$ (38)
$`+{\displaystyle \frac{2a_{12}\sqrt{m_1}}{l_1l_2k_L^2\sqrt{m_2}}}{\displaystyle \frac{A_2^0}{A_1^0}}B_1+{\displaystyle \frac{m_1V_{0,1}}{\mathrm{}^2k_L^2}}{\displaystyle \frac{B_1}{A_1^0}})\text{,}`$
$`\mu _2^0`$ $`=`$ $`{\displaystyle \frac{\mathrm{}^2k_L^2}{2m_2}}(1+{\displaystyle \frac{6a_2}{l_2^2k_L^2}}B_2+{\displaystyle \frac{4a_{12}\sqrt{m_2}}{l_1l_2k_L^2\sqrt{m_1}}}B_1`$ (40)
$`+{\displaystyle \frac{2a_{12}\sqrt{m_2}}{l_1l_2k_L^2\sqrt{m_1}}}{\displaystyle \frac{A_1^0}{A_2^0}}B_2+{\displaystyle \frac{m_2V_{0,2}}{\mathrm{}^2k_L^2}}{\displaystyle \frac{B_2}{A_2^0}})\text{.}`$
The constants $`A`$ and $`B`$ are related by restrictions such that $`B_1A_1^0`$ for $`A_1^0<0`$ and $`B_10`$ for $`A_1^0>0`$; $`B_2A_2^0`$ for $`A_2^0<0`$ and $`B_20`$ for $`A_2^0>0`$.
## III The phase diagram
The average particle number densities $`n_i`$ for the two species are obtained as
$`n_i`$ $`=`$ $`{\displaystyle \frac{1}{L}}{\displaystyle _0^L}\left|\psi _i(x,t)\right|^2๐x`$ (41)
$`=`$ $`{\displaystyle \frac{1}{h\pi }}{\displaystyle _0^{h\pi }}\left[A_i^0\mathrm{sin}^2\left(x^{^{}}\right)+B_i\right]๐x^{^{}}\text{,}`$ (42)
where $`x^{^{}}=k_Lx`$ and $`L=hd`$ denotes the length of the optical lattice with $`h=1,2,3,\mathrm{}`$ This leads to
$$B_i=n_i\frac{A_i^0}{2}\text{.}$$
(43)
Then the macroscopic wave functions of the condensates Eq. (33) can exist only when
$`n_1`$ $``$ $`{\displaystyle \frac{|A_1^0|}{2}}={\displaystyle \frac{\left|a_{12}V_{0,2}a_2V_{0,1}\right|}{4\mathrm{}\omega _0\left|a_1a_2a_{12}^2\right|}}\text{,}`$ (44)
$`n_2`$ $``$ $`{\displaystyle \frac{|A_2^0|}{2}}={\displaystyle \frac{\left|a_{12}V_{0,1}a_1V_{0,2}\right|}{4\mathrm{}\omega _0\left|a_1a_2a_{12}^2\right|}}\text{.}`$ (45)
The condensate atom currents can be evaluated from the usual definition, $`j=(\mathrm{}/m)`$Im$`\left[\psi ^{}\left(\psi /x\right)\right]`$ , with the exact wave functions Eq. (33) which are seen to be travelling matter waves. The result is
$`j_i`$ $`=`$ $`\pm {\displaystyle \frac{\mathrm{}k_L}{m_i}}\sqrt{B_i(B_i+A_i^0)}`$ (46)
$`=`$ $`\pm {\displaystyle \frac{\mathrm{}k_L}{m_i}}\sqrt{n_i^2{\displaystyle \frac{(A_i^0)^2}{4}}}\text{,}`$ (47)
which are independent of space-time variables and therefore persistent currents. We may demand that the wave functions Eq. (33) satisfy the periodic boundary condition $`\psi _i(x,t)=\psi _i(x+L,t)`$ which is naturally fulfilled as the total length $`L`$ is an integer times the lattice constant $`d`$ ($`L=hd,h=1,2,3,\mathrm{}`$). These periodic solutions in 1D space with spatial period $`L`$ are equivalent to the solutions in a ring of circumstance $`L`$. The persistent currents then can be viewed as in the optical lattice ring.
It is found from Eq. (46) that the persistent currents are valid only for conditions in which the number density of atoms is greater than critical values, $`n_i|A_i^0|/2`$, i.e., when the macroscopic wave functions of the condensates exist. These persistent currents are similar to the 1D Frรถhlich superconductivity induced by the traveling lattice wave and can be controlled by adjusting the barriers height of the periodic potentials and parameters of the bosonic atoms. The currents increase with the decrease of $`|A_i^0|`$ and approach the asymptotic maximum values $`j_{i,\mathrm{max}}=\pm \mathrm{}n_ik_L/m_i`$ when $`|A_i^0|`$ become vanishingly small. With the recent progress made on confinement of atoms in the light-induced as well as the magnetic-field-induced atom waveguides , the persistent currents may be observed experimentally in the future.
The energy spectrum for the two species is obtained as
$`\mu _1^0`$ $`=`$ $`E_{R,1}+{\displaystyle \frac{V_{0,1}}{2}}+2\mathrm{}\omega _0(a_1n_1+a_{12}n_2)\text{,}`$ (48)
$`\mu _2^0`$ $`=`$ $`E_{R,2}+{\displaystyle \frac{V_{0,2}}{2}}+2\mathrm{}\omega _0(a_2n_2+a_{12}n_1)\text{,}`$ (49)
where $`E_{R,i}=\mathrm{}^2k_L^2/2m_i`$ are the recoil energy of an atom absorbing one of the lattice phonons .
In this paper, we consider the two species both with repulsive interactions, namely, $`a_1>0`$, $`a_2>0`$, $`a_{12}>0`$. Then the macroscopic wave functions of the condensates can exist only when
$`\mu _1^0`$ $``$ $`E_{R,1}+{\displaystyle \frac{V_{0,1}}{2}}+{\displaystyle \frac{a_1\left|a_{12}V_{0,2}a_2V_{0,1}\right|}{2\left|a_1a_2a_{12}^2\right|}}`$ (51)
$`+{\displaystyle \frac{a_{12}\left|a_{12}V_{0,1}a_1V_{0,2}\right|}{2\left|a_1a_2a_{12}^2\right|}}\text{,}`$
$`\mu _2^0`$ $``$ $`E_{R,2}+{\displaystyle \frac{V_{0,2}}{2}}+{\displaystyle \frac{a_2\left|a_{12}V_{0,1}a_1V_{0,2}\right|}{2\left|a_1a_2a_{12}^2\right|}}`$ (53)
$`+{\displaystyle \frac{a_{12}\left|a_{12}V_{0,2}a_2V_{0,1}\right|}{2\left|a_1a_2a_{12}^2\right|}}\text{.}`$
The macroscopic wave functions of the condensates $`\psi _i(x,t)`$ are complex functions defined as the expectation value of the boson field operators: $`\psi _i(x,t)\widehat{\mathrm{\Psi }}_i(x,t)`$, which have the meaning of order parameters and characterize the off-diagonal long-range behavior of the one-particle density matrix $`\rho _i(x^{^{}},x,t)=\widehat{\mathrm{\Psi }}_i^+(x^{^{}},t)\widehat{\mathrm{\Psi }}_i(x,t)`$ . So the condensates can be kept in the superflud phase only when the macroscopic wave functions exist. Otherwise, the phase coherence and the currents vanish and therefore the condensates are in the insulator phase. Then we obtain four cases of phases for two-species BECs in a 1D optical lattice as follows.
Case 1. The two species are both in the superfluid phase, namely
$`n_1`$ $``$ $`{\displaystyle \frac{\left|a_{12}V_{0,2}a_2V_{0,1}\right|}{4\mathrm{}\omega _0\left|a_1a_2a_{12}^2\right|}}\text{,}`$ (54)
$`n_2`$ $``$ $`{\displaystyle \frac{\left|a_{12}V_{0,1}a_1V_{0,2}\right|}{4\mathrm{}\omega _0\left|a_1a_2a_{12}^2\right|}}\text{,}`$ (55)
which is labeled as SS in the phase diagram, Fig. 1.
Case 2. Species 1 is in the superfluid phase while species 2 is in the insulator phase,
$`n_1`$ $``$ $`{\displaystyle \frac{\left|a_{12}V_{0,2}a_2V_{0,1}\right|}{4\mathrm{}\omega _0\left|a_1a_2a_{12}^2\right|}}\text{,}`$ (56)
$`n_2`$ $``$ $`{\displaystyle \frac{\left|a_{12}V_{0,1}a_1V_{0,2}\right|}{4\mathrm{}\omega _0\left|a_1a_2a_{12}^2\right|}}\text{,}`$ (57)
labeled as SI.
Case 3. Species 1 is in the insulator phase and species 2 is in the superfluid phase,
$`n_1`$ $``$ $`{\displaystyle \frac{\left|a_{12}V_{0,2}a_2V_{0,1}\right|}{4\mathrm{}\omega _0\left|a_1a_2a_{12}^2\right|}}\text{,}`$ (58)
$`n_2`$ $``$ $`{\displaystyle \frac{\left|a_{12}V_{0,1}a_1V_{0,2}\right|}{4\mathrm{}\omega _0\left|a_1a_2a_{12}^2\right|}}\text{,}`$ (59)
labeled as IS.
Case 4. The two species are both in the insulator phase,
$`n_1`$ $``$ $`{\displaystyle \frac{\left|a_{12}V_{0,2}a_2V_{0,1}\right|}{4\mathrm{}\omega _0\left|a_1a_2a_{12}^2\right|}}\text{,}`$ (60)
$`n_2`$ $``$ $`{\displaystyle \frac{\left|a_{12}V_{0,1}a_1V_{0,2}\right|}{4\mathrm{}\omega _0\left|a_1a_2a_{12}^2\right|}}\text{,}`$ (61)
labeled as II.
The quantum phases of the condensates can be determined by all parameters of BECs and optical lattice as shown above. In Fig. 1, we show the phase diagram with the various same-species s-wave scattering lengths and equal particle number density $`n_1=n_2=n`$ and the magnitudes of potentials ($`V_{0,1}=V_{0,2}=V_0`$ ) for simplicity. Thus the conditions for the two species in the superfluid phase are given by
$$\frac{V_0}{4n\mathrm{}\omega _0}\frac{\left|a_1a_2a_{12}^2\right|}{\left|a_{12}a_2\right|}$$
(62)
and
$$\frac{V_0}{4n\mathrm{}\omega _0}\frac{\left|a_1a_2a_{12}^2\right|}{\left|a_{12}a_1\right|}\text{,}$$
(63)
respectively.
Component separation in two-species BECs has been predicted by means of mean-field theory and observed in experiments when the relation of the scattering lengths that $`a_{12}>\sqrt{a_1a_2}`$ is fulfilled. From the above conditions Eqs. (62), (63) and the phase diagram Fig. 1, we find that the larger values of $`a_{12}`$ favor the superfluid phase in the two-species mixture and the component separation according to the experimental observation. Particularly when the interspecies scattering length approaches the value of the same-species such that $`a_{12}=a_2`$ or $`a_1`$, the conditions Eqs. (62), 63) result in the superfluid phase independent of the potential magnitude $`V_0`$. One should not be surprised by this result since we consider the case in which the chemical potential is always higher than the potential magnitude $`V_0`$ seen from Eq. (51).
## IV Conclusion
In conclusion, the exact macroscopic wave functions of two-species BECs in an optical lattice beyond the tight-binding approximation are studied. The phase diagram is determined analytically according to the order parameters, and persistent currents in an optical lattice ring are obtained explicitly in terms of the exact wave functions, which are seen to be traveling matter waves.
## V Acknowledgments
This work was supported by the NSF of China under Grants No. 10475053, No. 60490280, No. 90406017, and No. 90403034.
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# Tensor products of ๐ญโข๐ฐโข๐ฉโข(2|2) representations
## 1 Introduction
This short note is devoted to the representation theory of the Lie superalgebra $`๐ญ๐ฐ๐ฉ(2|2)`$. The latter describes symmetries of several important physical systems, ranging from strings moving in an $`AdS_3`$ background to the quantum Hall effect (see e.g. ). The questions we address here, however, are purely mathematical and we shall only comment very briefly on our motivation from physics, leaving more concrete applications to a future publication.
The so-called A-series of simple Lie superalgebras consists of $`๐ฐ๐ฉ(n|m)`$ with $`nm`$ and $`๐ญ๐ฐ๐ฉ(n|n)`$.<sup>1</sup><sup>1</sup>1The Lie superalgebras $`๐ญ๐ฐ๐ฉ(n|n)`$ are obtained from $`๐ฐ๐ฉ(n|n)`$ when we remove the 1-dimensional center. Finite dimensional representations of these Lie superalgebras with diagonal Cartan elements have been constructed and investigated extensively (see e.g. for some early papers). Most of this work focused on irreducible representations, disregarding that even the very simplest Lie superalgebras from the A-series possess plenty of non-trivial indecomposables (see e.g. ). In fact, the Lie superalgebras $`๐ฐ๐ฉ(n|m)`$ with $`n,m>1`$ admit so many of them that they cannot even be classified . Irreducible representations, on the other hand, are rather easy to list (see e.g. and references therein). These fall into two different classes known as typical (long) and atypical (short) representations . More general indecomposables may be regarded as composites of the latter.
Investigations of tensor products for the A-series cannot avoid dealing with indecomposable representations. In fact, it is well known that the product of two irreducibles is often not fully reducible and what is even worse: indecomposables of Lie superalgebras do not form an ideal in the fusion ring. Consequently, it would e.g. not be possible to determine how many times a given irreducible representation appears in a higher tensor product if we only knew the number of irreducibles in the two-fold products of irreducible representations, simply because the fusion paths leading to irreducible representations may pass through indecomposables. In other words, the study of tensor products for representations of Lie superalgebras cannot be consistently truncated to irreducible representations. General results show, however, that there is a preferable class of so-called projective representations which gives rise to an ideal in the fusion ring. This contains the typical representations along with certain maximal indecomposable composites of atypicals. The latter are known as projective covers of atypical representations.
The indecomposable representations that emerge in products of irreducibles add a lot of interesting novel structure to the study of tensor products, but they certainly also bring a lot of extra difficulties. Even though tensor products of Lie superalgebras from the A-series have been investigated in the literature (see e.g. for some early work), most of the existing work excludes degenerate cases in which indecomposables appear as one of the factors or in the decomposition of the product. More extensive results on degenerate products seem to be restricted to a few simple examples, including $`๐ฐ๐ฉ(2|1)`$ and $`๐ค๐ฉ(1|1)`$ (see e.g. ). And even in these cases, a complete treatment was only found recently .
In this note we shall study all finite dimensional tensor products for $`๐ญ๐ฐ๐ฉ(2|2)`$ between typical and atypical representations as well as their projective covers. We shall see explicitly how the products involving at least one projective representation can be decomposed into a sum of typicals and projective covers (propositions 1, 2 and 3 below). On the other hand, a new family of indecomposables arises in tensor products of atypical representations (see proposition 4). The next section contains an introduction to the three most important types of $`๐ญ๐ฐ๐ฉ(2|2)`$ representations. Section 3 is then devoted to our new results on tensor products of these representations. A complete list of our findings is provided in subsection 3.1 along with a few additional comments. Our claims are supported by two central ideas which are explained in subsection 3.2 before we demonstrate how they work together in our computation of tensor products. Since the full calculations are quite cumbersome and not very illuminating, we shall illustrate the key steps in two representative examples rather than attempting to present a general proof. In the concluding section, we will discuss at least briefly our motivations from physics and some potential applications.
## 2 Representations of $`๐ญ๐ฐ๐ฉ(2|2)`$
In this first section we shall mainly review known results about some finite dimensional representations of $`๐ค=๐ญ๐ฐ๐ฉ(2|2)`$. We shall provide a complete list of irreducible representations, both typical and atypical, explain how they are constructed from the so-called Kac modules and we shall describe the projective covers of all atypical representations.
### 2.1 The Lie superalgebra $`๐ญ๐ฐ๐ฉ(2|2)`$
The Lie superalgebra $`๐ค=๐ญ๐ฐ๐ฉ(2|2)`$ possesses six bosonic generators $`K^{ab}=K^{ba}`$ with $`a,b=1,\mathrm{},4`$. In addition, there are eight fermionic generators that we denote by $`S_\alpha ^a`$ with the new index $`\alpha =1,2`$ and where $`a`$ assumes the same values as for the bosonic generators. The relations of $`๐ค`$ are given by
$$\begin{array}{cc}\hfill [K^{ab},K^{cd}]& =i\left[\delta ^{ac}K^{bd}\delta ^{bc}K^{ad}\delta ^{ad}K^{bc}+\delta ^{bd}K^{ac}\right]\hfill \\ \hfill [K^{ab},S_\gamma ^c]& =i\left[\delta ^{ac}S_\gamma ^b\delta ^{bc}S_\gamma ^a\right]\hfill \\ \hfill [S_\alpha ^a,S_\beta ^b]& =\frac{1}{2}ฯต_{\alpha \beta }ฯต^{abcd}K^{cd}.\hfill \end{array}$$
(2.1)
Here, $`ฯต_{\alpha \beta }`$ and $`ฯต^{abcd}`$ denote the usual $`ฯต`$-symbols with two and four indices, respectively, and a summation over repeated indices is implied.
Note that the even subalgebra $`๐ค^{(0)}=๐ญ๐ฐ๐ฉ(2|2)^{(0)}`$ is isomorphic to $`๐ฐ๐ฉ(2)๐ฐ๐ฉ(2)`$. The odd part $`๐ค^{(1)}`$ of $`๐ญ๐ฐ๐ฉ(2|2)`$ is spanned by the eight fermionic generators. We split the latter into two sets of four generators
$$๐ค_+^{(1)}=\text{span}\{S_2^a\},๐ค_{}^{(1)}=\text{span}\{S_1^a\}.$$
As indicated by the subscript $`\pm `$, we shall think of the fermionic generators $`S_1^a`$ as annihilation operators and of $`S_2^a`$ as creation operators.
Let us furthermore recall that the group SL(2,$``$) acts on $`๐ค`$ through outer automorphisms. For an element $`u=(u_{\alpha }^{}{}_{}{}^{\beta })`$ SL(2,$``$) the latter read
$$\gamma _u(K^{ab})=K^{ab},\gamma _u(S_\alpha ^a)=u_{\alpha }^{}{}_{}{}^{\beta }S_\beta ^a,$$
(2.2)
Consistency with the defining relations of our Lie superalgebra is straightforward to check. It only uses the fact that $`\text{det}(u)=1`$.
### 2.2 Finite dimensional irreducible representations
The irreducible finite dimensional representations of $`๐ค`$ are labeled by pairs $`j_1,j_2`$ with $`j_i=0,1/2,1,\mathrm{}`$. All these representations are highest weight representations and they are uniquely characterized by the highest weights $`(j_1,j_2)`$ of the corresponding even subalgebra $`๐ค^{(0)}๐ฐ๐ฉ(2)๐ฐ๐ฉ(2)`$. We shall see in a moment how such representations may be constructed explicitly.
For reasons that we shall understand below, the irreducible representations of $`๐ค`$ fall into two classes. So-called typical representations appear for $`j_1j_2`$. We shall denote them as $`[j_1,j_2]`$. Their dimension is given by
$$\mathrm{dim}[j_1,j_2]=d_{[j_1,j_2]}=16(2j_1+1)(2j_2+1).$$
(2.3)
Representations with labels $`j=j_1=j_2`$ are atypical, since their dimension is smaller than a naive application of formula (2.3) would suggest. For these irreducible representations we shall employ the symbol $`[j]`$ and one finds that
$$\mathrm{dim}[j]=d_{[j]}=16j(j+1)+2.$$
The formula holds for $`j0`$. The representation $`[0]`$ is the trivial one-dimensional representation. Let us also point out that the 14-dimensional atypical representation $`[1/2]`$ is the adjoint representation of $`๐ญ๐ฐ๐ฉ(2|2)`$ .
The irreducible representations $`[i,j]`$ and $`[j]`$ of the Lie superalgebra $`๐ค`$ can be restricted to the even subalgebra $`๐ค^{(0)}`$. With respect to this restricted action they decompose according to
$`[j]|_{๐ค^{(0)}}`$ $``$ $`(j+{\scriptscriptstyle \frac{1}{2}},j{\scriptscriptstyle \frac{1}{2}})2(j,j)(j{\scriptscriptstyle \frac{1}{2}},j+{\scriptscriptstyle \frac{1}{2}})(\text{for }j>0),`$ (2.4)
$`[i,j]|_{๐ค^{(0)}}`$ $``$ $`(i,j)\left[2(0,0)2({\scriptscriptstyle \frac{1}{2}},{\scriptscriptstyle \frac{1}{2}})(0,1)(1,0)\right].`$ (2.5)
Here and in the following, the pairs $`(i,j)`$ denote irreducible representations of the even subalgebra. Note that these decomposition formulas are consistent with our expressions for the dimension of irreducible representations.
The irreducible representations of $`๐ค`$ possess one property that will become very important later on: They admit an implementation of the outer automorphisms $`\gamma _u`$, see (2.2). To make this more precise, let us introduce the symbol $`\pi `$ for the representation that sends elements $`X๐ญ๐ฐ๐ฉ(2|2)`$ to linear maps on the representation space $`_\pi `$. If we compose a representation $`\pi `$ with an automorphism $`\gamma `$ of the Lie superalgebra, we obtain a new representation $`\pi \gamma `$ on the same graded vector space. In general, this new representation differs from $`\pi `$. If $`\pi `$ is one of our finite dimensional irreducible representations $`[i,j]`$ or $`[j]`$ and $`\gamma `$ one of the automorphisms in (2.2), however, then the new representation is equivalent to the original one, i.e. for every $`u`$ SL(2,$``$) there exists an invertible linear map $`U_\pi :_\pi _\pi `$ such that
$$\pi \gamma _u(X)=U_\pi \pi (X)U_\pi ^1\text{ for all }X๐ญ๐ฐ๐ฉ(2|2).$$
(2.6)
Let us stress that the map $`uU_\pi `$ defines a representation of the subgroup SL(2,$``$) of outer automorphisms on the representation space $`_\pi `$.
### 2.3 Kac modules and irreducible representations
It is useful for us to discuss briefly how the irreducible representations we have listed above can be constructed. The idea is rather standard: We begin with an irreducible highest weight representation $`(j_1,j_2)`$ of the even subalgebra $`๐ค^{(0)}`$. We declare that the corresponding representation space $`V_{(j_1,j_2)}`$ is annihilated by $`S_1^a`$ and then generate a so-called Kac module $`[j_1,j_2]`$ by application of the raising operators $`S_2^a`$,
$$[j_1,j_2]:=\mathrm{Ind}_{๐ค^{(0)}๐ค_{}^{(1)}}^๐คV_{(j_1,j_2)}=๐ฐ(๐ค)_{๐ค^{(0)}๐ค_{}^{(1)}}V_{(j_1,j_2)}.$$
Here, we have extended the $`๐ค^{(0)}`$ module $`V_{(j_1,j_2)}`$ to a representation of $`๐ค^{(0)}๐ค_{}`$ by setting $`S_1^aV_{(j_1,j_2)}=0`$. Note that we can apply at most four fermionic generators to the states in $`V_{(j_1,j_2)}`$. Therefore, the dimension of this Kac module is given by
$$\mathrm{dim}[j_1,j_2]=16(2j_1+1)(2j_2+1).$$
The Kac module $`[j_1,j_2]`$ is irreducible whenever $`j_1j_2`$. In these generic cases it agrees with the typical representation. For $`j_1=j_2=j`$, however, the associated Kac module turns out to be reducible but not fully reducible, i.e. it cannot be written as a direct sum of irreducible representations. If $`j1`$ the structure of the Kac module can be encoded in the following chain
$$[j,j]:[j][j+{\scriptscriptstyle \frac{1}{2}}][j{\scriptscriptstyle \frac{1}{2}}][j],$$
(2.7)
or, equivalently, in a planar diagram in which one direction refers to the spin $`j`$ of the atypical constituents,
$$[j,j]:\text{}$$
(2.8)
Since pictures of this type will appear frequently throughout this text, let us pause here for a moment and explain carefully how to decode their information. We read the diagram (2.8) from right to left. The rightmost entry in our chain contains the so-called socle of the indecomposable representation, i.e. the largest fully reducible invariant submodule we can find. In the case of our Kac module, the socle happens to be irreducible and it is given by the atypical representation $`[j]`$. If we divide the Kac module by the submodule $`[j]`$, we obtain a new indecomposable representation of our Lie superalgebra. Its diagram is obtained from the one above by removing the last entry and arrow. The socle of this quotient is a direct sum of the two atypical representations $`[j\pm 1/2]`$. It is rather obvious how to iterate this procedure until the entire indecomposable representation is split up into floors with only direct sums of irreducible representations appearing on each floor.
There are two special cases for which the decomposition of the Kac module does not follow the generic pattern. These are the cases $`j=0`$ and $`j=1/2`$,
$`[0,0]:`$ $`[0][{\scriptscriptstyle \frac{1}{2}}][0],`$ (2.9)
$`[{\scriptscriptstyle \frac{1}{2}},{\scriptscriptstyle \frac{1}{2}}]:`$ $`[{\scriptscriptstyle \frac{1}{2}}][1][0][0][{\scriptscriptstyle \frac{1}{2}}].`$ (2.10)
Let us note that our formula for the dimension of atypical representations follows directly from the decomposition of the corresponding Kac modules.
### 2.4 Projective covers of atypical representations
As we have seen in the last subsection, atypical representations can be extended into larger indecomposables. Kac modules are only one example of such composites and we shall indeed see several others as we proceed. Among them, however, there is one special class, the so-called projective covers $`๐ซ_๐ค(j)`$. By definition, these are the largest indecomposables whose socle consists of a single atypical representation $`[j]`$. General results imply that such a maximal indecomposable extension of $`[j]`$ exists and is unique. In case of $`j3/2`$, the structure of $`๐ซ_๐ค(j)`$ is encoded in the following diagram
$`๐ซ_๐ค(j):`$ $`[j]2[j+{\scriptscriptstyle \frac{1}{2}}]2[j{\scriptscriptstyle \frac{1}{2}}][j+1]4[j][j1]`$
$`2[j+{\scriptscriptstyle \frac{1}{2}}]2[j{\scriptscriptstyle \frac{1}{2}}][j].`$
Note that $`๐ซ_๐ค(j)`$ contains an entire Kac module as proper submodule. In this sense, the Kac modules are extendable. We also observe one rather generic feature of projective covers: they are built up from different Kac modules in a way that resembles the pattern in which Kac modules are constructed out of irreducibles (see eq. (2.8)). One may see this even more clearly if $`๐ซ_๐ค(j)`$ is displayed as a 2-dimensional diagram in which the additional direction keeps track of the spin $`j`$ of the atypical constituents $`[j]`$,
$$๐ซ_๐ค(j):\text{}$$
(2.12)
We will continue to switch between such planar pictures and diagrams of the form (2.4). The remaining cases $`j=0,1/2,1`$ have to be listed separately. When $`j=1`$ the picture is very similar only that we have to insert $`2[0]`$ in place of $`[j1]`$,
$$๐ซ_๐ค(1):[1]2[{\scriptscriptstyle \frac{3}{2}}]2[{\scriptscriptstyle \frac{1}{2}}][2]4[1]2[0]2[{\scriptscriptstyle \frac{3}{2}}]2[{\scriptscriptstyle \frac{1}{2}}][1].$$
(2.13)
The projective cover of the atypical representation $`[1/2]`$ is obtained from the generic case by the formal substitution $`2[j1/2]3[0]`$,
$$๐ซ_๐ค({\scriptscriptstyle \frac{1}{2}}):[{\scriptscriptstyle \frac{1}{2}}]2[1]3[0][{\scriptscriptstyle \frac{3}{2}}]4[{\scriptscriptstyle \frac{1}{2}}]2[1]3[0][{\scriptscriptstyle \frac{1}{2}}].$$
(2.14)
Finally, the projective cover $`๐ซ_๐ค(0)`$ of the trivial representation is given by,
$$๐ซ_๐ค(0):[0]3[{\scriptscriptstyle \frac{1}{2}}]2[1]6[0][{\scriptscriptstyle \frac{1}{2}}][0].$$
(2.15)
The reader is invited to convert the last three formulas into planar pictures. This concludes our list of the projective covers. The representations $`๐ซ_๐ค(j)`$ will arise as important building blocks in the decomposition of tensor products in the next section. Together, typical representations and the projective covers of atypicals form the subset of so-called projective representations. What makes this class particularly interesting is its behavior under tensor products. In fact, it is rather well-known that projective representations of a Lie superalgebra form an ideal in the fusion ring. We will see this very explicitly in the concrete decomposition formulas for tensor products below.
## 3 Tensor products of $`๐ญ๐ฐ๐ฉ(2|2)`$ representations
This section contains the central results of this work, i.e. explicit formulas for the decomposition of tensor products between all the finite dimensional irreducible representations and projectives we have introduced in the previous section. We begin by stating our main results. Then we shall explain the two key ideas that are needed in the proof. Finally, we analyze two very representative examples.
### 3.1 Tensor products of irreducible representations
Before we start listing our results, we would like to introduce one object that will enable us to determine the contributions from typical representations in most of the tensor products below. We shall denote by $`\pi _๐ค`$ a map which associates a direct sum of typical representations to any finite dimensional representation of the bosonic subalgebra. On the irreducible representations of $`๐ค^{(0)}`$ it gives
$$\pi _๐ค(i,j):=\{\begin{array}{cc}[i,j]\hfill & \text{ for }ij\hfill \\ 0\hfill & \text{ for }i=j.\hfill \end{array}$$
(3.1)
This prescription is extended to any direct sum of irreducibles by linearity. $`\pi _๐ค`$ enters e.g. in the following formula for the tensor product of typical and atypical representations.
###### Proposition 1 (Tensor product of typical and atypical representations).
The tensor product of an atypical representation $`[n]`$ with a typical representation $`[i,j]`$ is given by
$$[n][i,j]=\pi _๐ค\left([n]|_{๐ค^{(0)}}(i,j)\right)\text{ for }i+j.$$
When $`i+j`$, on the other hand, the tensor product can also contain projective covers,
$$[n][i,j]=\pi _๐ค\left([n]|_{๐ค^{(0)}}(i,j)\right)\underset{l=p}{\overset{q}{}}๐ซ_๐ค(l)$$
(3.2)
where $`p=\mathrm{max}(|ni|,|nj|)`$ and $`q=\mathrm{min}(n+i,n+j)`$. The decomposition of the irreducible representation $`[n]`$ into representations of the even subalgebra $`๐ค^{(0)}`$ was spelled out in eq. (2.4).
It is interesting to observe that the tensor product of the representation $`[n]`$ with the typical representation $`[0,2n]`$ contains a single projective cover $`๐ซ_๐ค(n)`$. Proposition 2 may be employed to determine tensor products between typical and any indecomposable representation. The prescription requires introducing a new map $`๐ฎ_๐ค`$ which sends an indecomposable representation $``$ to a sum of atypicals,
$$๐ฎ_๐ค()=\underset{j}{}[:[j]][j].$$
(3.3)
Here, the symbol $`[:[j]]`$ denotes the total number of atypical representations $`[j]`$ in the decomposition series of $``$. In the case where $``$ is one of the projective covers, for example, these numbers can be read off from eqs. (2.4)-(2.15). With this notation, the tensor product between a typical representation $`[i,j]`$ and the indecomposable $``$ reads,
$$[i,j]๐ฎ_๐ค()[i,j].$$
(3.4)
As we have anticipated, we may now employ proposition 1 to decompose the tensor product on the right hand side. Thereby, we can determine e.g. the fusion of typical representations and projective covers.
###### Proposition 2 (Tensor product of typical representations).
The tensor product of two atypical representations $`[i_1,j_1]`$ and $`[i_2,j_2]`$ is given by
$$[i_1,j_1][i_2,j_2]=\pi _๐ค\left([i_1,j_1]|_{๐ค^{(0)}}(i_2,j_2)\right)\text{ for }i_1+i_2+j_1+j_2.$$
If $`i_1+i_2+j_1+j_2`$, on the other hand, projective covers may appear in the decomposition,
$`[i_1,j_1][i_2,j_2]`$ $`=`$ $`\pi _๐ค\left([i_1,j_1]_{๐ค^{(0)}}(i_2,j_2)\right)2{\displaystyle \underset{m=2p}{\overset{2q}{}}}๐ซ_๐ค({\scriptscriptstyle \frac{m}{2}})`$
$`\{\begin{array}{cc}\delta _{p,q+1}๐ซ_๐ค(q+{\scriptscriptstyle \frac{1}{2}})\hfill & ,p>q\hfill \\ \left(1\delta _{j_1+j_2}^{i_1+i_2}\right)๐ซ_๐ค(q+{\scriptscriptstyle \frac{1}{2}})\left(1\delta _{|j_1j_2|}^{|i_1i_2|}\right)๐ซ_๐ค(p{\scriptscriptstyle \frac{1}{2}})\delta _{p,0}๐ซ_๐ค(0)\hfill & ,pq.\hfill \end{array}`$
The decomposition of typical representations $`[i,j]`$ into irreducibles of $`๐ค^{(0)}`$ appears in eq. (2.5). We have also introduced the parameters $`p`$ and $`q`$ by $`p=\mathrm{max}(|i_1i_2|,|j_1j_2|)`$ and $`q=\mathrm{min}(i_1+i_2,j_1+j_2)`$. Note that the last term in the first line subtracts one copy of the projective cover $`๐ซ_๐ค(0)`$ whenever $`p`$ vanishes.
At this point we are able to decompose all tensor products which involve at least one typical factor. Our next task is to analyze the products in which at least one factor is a projective cover $`๐ซ_๐ค(j)`$.
###### Proposition 3 (Tensor product of atypical representations and projective covers).
The tensor product of an atypical representation $`[n]`$ with a projective cover $`๐ซ_๐ค(j)`$ with $`j0`$ and $`n0`$ is given by
$`[n]๐ซ_๐ค(j)`$ $`=`$ $`\pi _๐ค\left([n]|_{๐ค^{(0)}}H_{(j)}\right)2{\displaystyle \underset{m=2|nj|}{\overset{2(n+j)}{}}}P_๐ค({\scriptscriptstyle \frac{m}{2}})\delta _{n,j}๐ซ_๐ค(0)`$ (3.6)
where $`H_{(j)}=2(j,j)(j+{\scriptscriptstyle \frac{1}{2}},j+{\scriptscriptstyle \frac{1}{2}})(j{\scriptscriptstyle \frac{1}{2}},j{\scriptscriptstyle \frac{1}{2}}).`$ (3.7)
The decomposition of atypical representations into modules of $`๐ค^{(0)}`$ can be found in eq. (2.4). In the special case of $`j=0`$ one obtains
$$[n]๐ซ_๐ค(0)=\pi _๐ค\left([n]|_{๐ค^{(0)}}H_{(0)}\right)4๐ซ_๐ค(n)$$
Here we have to insert the representation $`H_{(0)}=2(0,0)2({\scriptscriptstyle \frac{1}{2}},{\scriptscriptstyle \frac{1}{2}})`$ of the bosonic subalgebra in the argument of $`\pi _๐ค`$.
From proposition 3 we may compute the tensor product of a projective cover with any other indecomposable through the following extension of formula (3.4) to projective covers of atypical representations,
$$๐ซ_๐ค(j)๐ฎ_๐ค()๐ซ_๐ค(j)$$
(3.8)
for all indecomposables $``$. The symbol $`๐ฎ`$ has been introduced in equation (3.3) above such that it associates to $``$ the direct sum of irreducibles that appear in its decomposition series.
Note that so far we were able to express all tensor products through a direct sum of typical representations and projective covers. This is consistent with the before-mentioned general result that projective representations form an ideal in the fusion ring of representations. There remains, however, one more family of tensor products to be determined. These are the tensor products between two atypical representations. Not surprisingly, we shall encounter a new set of indecomposables in such tensor products. These possess the following form
$`\pi _{ij}^{\mathrm{indec}}`$ $`=`$ $`{\displaystyle \underset{k=|ij|}{\overset{i+j1}{}}}[k+{\scriptscriptstyle \frac{1}{2}}]{\displaystyle \underset{k=|ij|}{\overset{i+j}{}}}2[k]{\displaystyle \underset{k=|ij|}{\overset{i+j1}{}}}[k+{\scriptscriptstyle \frac{1}{2}}](ij)`$ (3.9)
$`\pi _{jj}^{\mathrm{indec}}`$ $`=`$ $`{\displaystyle \underset{k=0}{\overset{2j1}{}}}[k+{\scriptscriptstyle \frac{1}{2}}]3[0]{\displaystyle \underset{k=1}{\overset{2j}{}}}2[k]{\displaystyle \underset{k=0}{\overset{2j1}{}}}[k+{\scriptscriptstyle \frac{1}{2}}].`$ (3.10)
Note that the second line is essentially a special case of the first except that we replaced the term $`2[jj]`$ by $`3[0]`$. When rewritten in terms of our 2-dimensional diagrams, these representations read,
$$\pi _{ij}^{\mathrm{indec}}:\text{}\pi _{jj}^{\mathrm{indec}}:\text{}$$
Proposition 4: (Tensor product of atypical representations) The tensor product of two atypical representations $`[i]`$ and $`[j]`$ is given by
$$[i][j]=\pi _๐ค\left((i,i)(j,j)\right)\delta _{i,j}[0]\pi _{ij}^{\mathrm{indec}}.$$
(3.11)
The formula holds for $`i,j=1/2,1,3/2,`$. Tensor products of the trivial representation $`[0]`$ with any other atypical representation are obvious.
Let us note that for $`i=j`$ one copy of the atypical representation $`[0]`$ appears as a summand, i.e. it is not part of the single indecomposable representation that contains all the other atypical building blocks. We have also calculated a few tensor products between the new indecomposable representations $`\pi _{ij}`$ and atypicals. Such products turn out to generate further indecomposables with a structure that is similar to the one of $`\pi _{ij}`$ but involves different multiplicities. It seems within reach to fully analyze the fusion ring that is generated by irreducibles, but since the applications we have in mind do not require such an exhaustive study, we have not investigated these issues much further.
### 3.2 Proof of the decomposition formulas - general ideas
There are two main ideas that enter into the proof of the above formulas. To begin with, we shall exploit systematically that the outer automorphisms are implementable not only in the irreducibles of $`๐ค`$ but also in all their tensor products. This will ultimately organize all the typical subrepresentations and, more importantly, the elements of the composition series of indecomposables into multiplets of the group SL(2,$``$) of outer automorphisms.
In a second step we then restrict the action of $`๐ค`$ to the action of an embedded $`๐ฅ=๐ฐ๐ฉ(2|1)`$. Knowledge about the representation theory of the latter (see also ) along with a few new results from will then allow us to uniquely determine the structure of indecomposable $`๐ค`$ representations that appear in the tensor products of irreducibles and projectives. The consistency of our results has also been verified by checking associativity of the tensor products involving three representations of irreducible or projective type with labels equal or below $`4`$ on a computer.
#### Implementation of outer automorphisms.
We have stated above that the action of the SL(2,$``$) outer automorphisms of $`๐ค`$ can be implemented in all of its irreducible representations. Now we will argue that implementability of outer automorphisms is respected by the operation of tensor products, by restriction to the socle and by quotients. This implies that outer automorphisms may be implemented separately in all representations that appear in the decomposition of tensor products of irreducible representations. In this sense, such representations are rather special.
###### Lemma 1.
Suppose that the action of outer automorphisms can be implemented in two representations $`\pi `$ and $`\pi ^{}`$ of the Lie superalgebra $`๐ค`$ and that these implementations respect the $`_2`$ gradings of the representation spaces. Then it can also be implemented in the tensor product $`\pi \pi ^{}`$.
###### Proof:.
We denote the implementations of the outer automorphism $`\gamma `$ in the representations $`\pi `$ and $`\pi ^{}`$ by $`U_\pi (\gamma )`$ and $`U_\pi ^{}(\gamma )`$, respectively. Then the implementation in the tensor product is trivially given by $`U_{\pi \pi ^{}}(\gamma )=U_\pi (\gamma )U_\pi ^{}(\gamma )`$. โ
###### Lemma 2.
Let $`\pi `$ be a representation of $`๐ค`$ in which the outer automorphism $`\gamma `$ is implemented by $`U_\pi (\gamma )`$. Let furthermore $`\pi ^{}`$ be a subrepresentation that is invariant under the action of $`U_\pi (\gamma )`$, i.e. $`U_\pi (\gamma )_\pi ^{}_\pi ^{}`$. Then the action of $`\gamma `$ is implementable in the factor representation $`\pi /\pi ^{}`$.
###### Proof:.
The statement is obvious. โ
###### Lemma 3.
Suppose that the action of an outer automorphism $`\gamma `$ can be implemented in a representation $`\pi `$ of the Lie superalgebra $`๐ค`$ by the implementation map $`U_\pi (\gamma ):_\pi _\pi `$ and let $`\pi ^{}`$ be the socle of $`\pi `$. Then $`U_\pi (\gamma )|__\pi ^{}`$ is an implementation of $`\gamma `$ in $`\pi ^{}`$, i.e. $`U_\pi (\gamma )_\pi ^{}_\pi ^{}`$.
###### Proof:.
By definition, the socle $`_\pi ^{}`$ is the maximal semisimple submodule of $`_\pi `$. Its image $`U_\pi (\gamma )_\pi ^{}_\pi `$ carries an action of the Lie superalgebra. The corresponding subrepresentation is given by $`\pi ^{}\gamma `$ and hence it decomposes into a sum of irreducibles just like $`\pi ^{}`$ itself. In other words, the subspace $`U_\pi (\gamma )_\pi ^{}`$ is a semisimple submodule of $`_\pi `$. From the maximality of the socle we therefore conclude $`_\pi ^{}=U_\pi (\gamma )_\pi ^{}`$. โ
While the previous lemmas hold in full generality for all Lie superalgebras, the next one relies on special properties of the Lie superalgebra $`๐ญ๐ฐ๐ฉ(2|2)`$ and its class of SL(2,$``$) automorphisms $`\gamma _u`$ as defined in (2.2).
###### Lemma 4.
Let $`\pi `$ be a representation of $`๐ค=๐ญ๐ฐ๐ฉ(2|2)`$ in which an outer automorphism $`\gamma _u`$ of the form (2.2) is implemented by $`U_\pi `$ and let $`\pi ^{}`$ be an irreducible subrepresentation on the space $`_\pi ^{}`$. Then $`\pi `$ defines a subrepresentation on the space $`U_\pi _\pi ^{}`$ and the latter is isomorphic to $`\pi ^{}`$.
###### Proof.
For any vector $`v_\pi ^{}`$ we find $`\pi (x)U_\pi v=U_\pi \pi \gamma ^1(x)vU_\pi _\pi ^{}`$, hence $`U_\pi _\pi ^{}`$ is invariant. To see that the resulting representation is isomorphic to the original one we just need to realize a) that it is also irreducible, b) that the dimension agrees and c) that the weights and weight multiplicities coincide. While b) is obvious, c) holds because $`\gamma _u`$ acts trivially on the bosonic generators and especially on the Cartan elements. Assume now that there exists a proper subrepresentation on the space $`_{\pi ^{\prime \prime }}U_\pi _\pi ^{}`$. Then, using the arguments of the first line above, we find a proper subrepresentation of the original representation $`\pi ^{}`$ on the space $`U_\pi ^1_{\pi ^{\prime \prime }}_\pi ^{}`$, contradicting its irreducibility. โ
As simple as these statements are, they will be rather useful for our analysis of tensor products. They imply in particular, that all the indecomposable representations that arise in tensor products of irreducibles, and hence all the representations we discuss in this note, allow for an implementation of the outer automorphisms. This insight is particularly useful when we analyze the internal structure of indecomposable composites. In fact, if the action of outer automorphisms is implementable in an indecomposable, then it is implementable in its socle and the associated factor representation. Hence, the transformation properties under the group SL(2,$``$) respect the structure of such indecomposable representations, i.e. they organize each floor of their decomposition diagram into SL(2,$``$) multiplets of irreducible representations. For the indecomposables that appear in this text, the multiplicities are displayed in appendix A.
#### Decomposition with respect to $`๐ฐ๐ฉ(2|1)`$.
In this paragraph we shall study a particular embedding of $`๐ฅ=๐ฐ๐ฉ(2|1)`$ into $`๐ค`$ and explain how the irreducible representations of $`๐ค`$ decompose into representations of $`๐ฅ`$. We do not intend to present a complete introduction into $`๐ฐ๐ฉ(2|1)`$ here, but restrict to a short list of relevant notations. More details can be found in the standard literature (see e.g. ) and in .
The even subalgebra $`๐ฅ^{(0)}`$ of $`๐ฐ๐ฉ(2|1)`$ is given by the sum $`๐ค๐ฉ(1)๐ฐ๐ฉ(2)`$. Consequently, the Kac modules are labeled by pairs $`(b,j)`$ where $`b`$ and $`j=1/2,1,\mathrm{}`$. We shall denote these representations by
$$\{b,j\}:=\mathrm{Ind}_{๐ฅ^{(0)}๐ฅ_{}^{(1)}}^๐ฅV_{(b1/2,j1/2)}=๐ฐ(๐ฅ)_{๐ฅ^{(0)}๐ฅ_{}^{(1)}}V_{(b1/2,j1/2)}.$$
The shift of the labels in our notations will turn out to be rather convenient in the following. We denote the Kac modules of $`๐ฅ`$ by brackets $`\{,\}`$ to distinguish them from those of $`๐ค`$. For $`b\pm j`$, the Kac modules $`\{b,j\}`$ give rise to irreducible representations of dimension $`d_{\{b,j\}}=8j`$. Kac modules $`\{b,j\}`$ with $`b=\pm j`$ are indecomposable. Their structure can be encoded in the following short diagram
$$\{\pm j,j\}:\{j\}_\pm \{j{\scriptscriptstyle \frac{1}{2}}\}_\pm .$$
Here, $`\{j\}_\pm `$ denote the $`(4j+1)`$dimensional irreducible atypical representations of $`๐ฅ`$.
The last type of representations that we shall need below are the projective covers $`๐ซ_๐ฅ^\pm (j)`$ of the atypical representations $`\{j\}_\pm `$. For $`j0`$, their structure is encoded in the following picture
$$๐ซ_๐ฅ^\pm (j):\{j\}_\pm \{j+{\scriptscriptstyle \frac{1}{2}}\}_\pm \{j{\scriptscriptstyle \frac{1}{2}}\}_\pm \{j\}_\pm .$$
(3.12)
These spaces are $`16j+4`$ dimensional as one can easily check by adding up the dimensions of the atypical composition series. The projective cover of the trivial representation $`\{0\}=\{0\}_\pm `$ is an $`8`$-dimensional representation that is given by
$$๐ซ_๐ฅ(0):\{0\}\{{\scriptscriptstyle \frac{1}{2}}\}_+\{{\scriptscriptstyle \frac{1}{2}}\}_{}\{0\}.$$
(3.13)
Results for tensor products of these representations are spelled out in appendix B. Tensor products of irreducible representations were originally computed in and results for tensor products of irreducible with the projective representations $`๐ซ_๐ฅ^\pm (n)`$ can be found in .
What we shall need in the following are explicit formulas for the decomposition of irreducible $`๐ค=๐ญ๐ฐ๐ฉ(2|2)`$ representations with respect to the subalgebra $`๐ฅ=๐ฐ๐ฉ(2|1)`$. Obviously, these depend on the explicit choice of the embedding. Here we shall consider the case in which the central element $`Z๐ฅ^{(0)}`$ of the even part in $`๐ฐ๐ฉ(2|1)`$ is identified with the Cartan element of the first $`๐ฐ๐ฉ(2)`$ subalgebra of $`๐ค^{(0)}`$. With this choice we find
$`[n]|_๐ฅ`$ $``$ $`\{n\}_+\{n\}_{}\left([n]\right)`$ (3.14)
where $`\left([n]\right)={\displaystyle \underset{b=n+1/2}{\overset{n1/2}{}}}\{b,n+{\scriptscriptstyle \frac{1}{2}}\}.`$
Below we shall think of $``$ as a map that sends any sum of atypical $`๐ค`$-representations to a sum of typical representations for $`๐ฅ`$. For typical representations $`[i,j]`$ of $`๐ค`$ we shall first assume that $`j>i`$. The other case will be treated below.
$$[i,j]|_๐ฅ\underset{b=i}{\overset{i}{}}\left(\{b,j+1\}(1\delta _{j,0})\{b,j\}\{b{\scriptscriptstyle \frac{1}{2}},j+{\scriptscriptstyle \frac{1}{2}}\}\{b+{\scriptscriptstyle \frac{1}{2}},j+{\scriptscriptstyle \frac{1}{2}}\}\right).$$
(3.15)
Note that with our assumption $`i<j`$ only typical representations occur on the right hand side. The formulas also holds true for $`i>j`$ and $`i+j`$ a half-integer. When $`i>j`$ and $`i+j`$ integer, on the other hand, the formula must be modified and the decomposition turns out to contain some of the projective covers $`๐ซ_๐ฅ^\pm (j)`$ of atypical representations,
$`[i,j]|_๐ฅ`$ $``$ $`{\displaystyle \underset{b=i}{\overset{i}{}}}{}_{}{}^{}(\{b,j+1\}(1\delta _{j,0})\{b,j\}\{b{\scriptscriptstyle \frac{1}{2}},j+{\scriptscriptstyle \frac{1}{2}}\}\{b+{\scriptscriptstyle \frac{1}{2}},j+{\scriptscriptstyle \frac{1}{2}}\})`$ (3.16)
$`{\displaystyle \underset{\nu =\pm }{}}\left(๐ซ_๐ฅ^\nu (j)๐ซ_๐ฅ^\nu (j+1/2)\right)`$
where the symbol $`^{}`$ instructs us to omit all terms that would formally be associated with atypical representations. In addition, we shall agree throughout this paper to replace the sum $`๐ซ_๐ฅ^+(0)๐ซ_๐ฅ^{}(0)`$ by $`๐ซ_๐ฅ(0)`$ whenever it appears. This is purely formal and has no meaning in terms of representations. Note that the left hand side only contains projective representations. Let us finally spell out the decomposition formulas for the projective covers and the indecomposables $`\pi ^{\text{indec}}`$ which we defined in eq. (3.9, 3.10),
$`๐ซ_๐ค(j)|_๐ฅ`$ $``$ $`{\displaystyle \underset{\nu =\pm }{}}\left(๐ซ_๐ฅ^\nu (j+{\scriptscriptstyle \frac{1}{2}})\mathrm{\hspace{0.17em}2}๐ซ_๐ฅ^\nu (j)๐ซ_๐ฅ^\nu (|j{\scriptscriptstyle \frac{1}{2}}|)\right)๐ฎ\left(๐ซ_๐ค(j)\right)`$ (3.17)
$`\pi _{ij}^{\mathrm{indec}}|_๐ฅ`$ $``$ $`{\displaystyle \underset{\nu =\pm }{}}\left(\{|ij|\}_\nu \{i+j\}_\nu {\displaystyle \underset{p=|ij|+\frac{1}{2}}{\overset{i+j+\frac{1}{2}}{}}}๐ซ_๐ฅ(p)\right)๐ฎ(\pi _{ij}^{\mathrm{indec}})`$ (3.18)
These formulas also hold for the special cases of $`j=0`$ and $`i=j`$ if we agree to replace $`\{0\}_+\{0\}_{}`$ by $`\{0\}`$. Once more such a replacement is purely formal and has no meaning in terms of representations. We note that it is rather easy to infer the projective covers $`๐ซ_๐ฅ`$ in the two decomposition formulas from our planar pictures for $`๐ซ_๐ฅ`$ and the indecomposables $`\pi ^{\text{indec}}`$. This concludes the presentation of the background material that is needed in the proof of our decomposition formulas for tensor products of $`๐ค`$.
### 3.3 Proof of the decomposition formulas - examples
Rather than trying to go through the general proof of our formulas we would like to illustrate the main ideas in two rather representative examples. These are the tensor product of the atypical representation $`[1]`$ with itself and with $`[0,2]`$.
#### The tensor product $`[1][1]`$.
Let us begin by collecting a few results on the 34-dimensional atypical representation $`[1]`$ of $`๐ค`$. With respect to the embedded $`๐ฅ`$, this representations decomposes as follows
$$[1]|_๐ฅ\{1\}_+\{1\}_{}\{{\scriptscriptstyle \frac{1}{2}},{\scriptscriptstyle \frac{3}{2}}\}\{{\scriptscriptstyle \frac{1}{2}},{\scriptscriptstyle \frac{3}{2}}\}.$$
Using standard results (see appendix B) about tensor products of irreducible representations of $`๐ฅ`$, we obtain the following decomposition formula for the tensor product $`[1][1]`$ in terms of representations of $`๐ฅ`$,
$$\left([1][1]\right)|_๐ฅ2\{0\}2๐ซ_๐ฅ(0)\underset{\nu =\pm }{}\left(\{2\}_\nu 3๐ซ_๐ฅ^\nu ({\scriptscriptstyle \frac{1}{2}})๐ซ_๐ฅ^\nu (1)2๐ซ_๐ฅ^\nu ({\scriptscriptstyle \frac{3}{2}})\right)+\mathrm{}$$
(3.19)
where the dots stand for a sum of typical $`๐ฅ`$ representations. The latter will not play any role for the following analysis.
It is also easy to find the typical $`๐ค`$ representations in the tensor product of $`[1]`$ with itself,
$`\left([1][1]\right)^{\mathrm{typ}}`$ $`=`$ $`[0,1][0,2][1,0][1,2][2,0][2,1]`$
$`\left([1][1]\right)^{\mathrm{typ}}|_๐ฅ`$ $``$ $`2๐ซ_๐ฅ(0){\displaystyle \underset{\nu =\pm }{}}\left(2๐ซ_๐ฅ^\nu ({\scriptscriptstyle \frac{1}{2}})๐ซ_๐ฅ^\nu (1)๐ซ_๐ฅ^\nu ({\scriptscriptstyle \frac{3}{2}})\right)+\mathrm{}`$ (3.20)
When we passed to the second line, we have inserted the results from our decomposition formulas for typical $`๐ค`$ representations with respect to the embedded $`๐ฅ`$. Once more, we have omitted all typical $`๐ฅ`$ representations.
A comparison of eqs. (3.19) and (3.20) gives the following intermediate result,
$$\left([1][1]\left([1][1]\right)^{\mathrm{typ}}\right)|_๐ฅ2\{0\}\underset{\nu =\pm }{}\left(\{2\}_\nu ๐ซ_๐ฅ^\nu ({\scriptscriptstyle \frac{1}{2}})๐ซ_๐ฅ^\nu ({\scriptscriptstyle \frac{3}{2}})\right)+\mathrm{}$$
The atypical $`๐ฅ`$ representations that remain must come from the decomposition of the $`๐ค`$ indecomposables that appear in the tensor product of $`[1]`$. A short glance on our decomposition formulas (3.14) and (3.18) confirms that this is indeed the case.
We can actually use this example to derive the structure of the indecomposable representation $`\pi _{11}`$ from the information on its restriction to $`๐ฅ`$. In fact, it is rather easy to see that the composition series of $`\pi _{11}`$ contains the following list of atypical representations, each displayed with a multiplicity that refers to its transformation properties under the action of the SL(2,$``$) outer automorphisms<sup>2</sup><sup>2</sup>2The subscripts may be determined from the SL(2,$``$) transformation properties of the bosonic multiplets in the involved tensor factors. The latter are listed in appendix A.
$$[0]_1,[0]_3,2[{\scriptscriptstyle \frac{1}{2}}]_1,[1]_2,2[{\scriptscriptstyle \frac{3}{2}}]_1,[2]_2.$$
A moment of thought reveals that all but the $`[0]_1`$ representation must be part of one single indecomposable in order to be able to recover our knowledge about the decomposition with respect to $`๐ฅ`$. It is at this point where our assignment of multiplicities becomes crucial. In fact, our knowledge from the $`๐ฅ`$ embedding would e.g. have been consistent with including two of the four trivial $`๐ญ๐ฐ๐ฉ(2|2)`$ representations into the indecomposable. But since there is no doublet, we were forced to include the triplet. Hence, there is only a single irreducible representation $`[0]`$ left. In other words, the tensor product $`[1][1]`$ contains only one true invariant. The presence of additional states which transform trivially but sit in the indecomposables has to be contrasted with the case of ordinary simple Lie algebras where the tensor product of a representation with its conjugate contains precisely one such state.
#### The tensor product $`[1][0,2]`$.
With respect to the embedded $`๐ฅ๐ค`$, the typical $`๐ค`$ representation $`[0,2]`$ decomposes as follows
$$[0,2]|_๐ฅ\{0,3\}\{0,2\}\{{\scriptscriptstyle \frac{1}{2}},{\scriptscriptstyle \frac{5}{2}}\}\{{\scriptscriptstyle \frac{1}{2}},{\scriptscriptstyle \frac{5}{2}}\}.$$
Using again standard results about tensor products of irreducible representations of $`๐ฅ`$ we obtain the following decomposition formula for the tensor product $`[1][0,2]`$ in terms of representations of $`๐ฅ`$,
$$\left([1][0,2]\right)|_๐ฅ\underset{\nu =\pm }{}\left(๐ซ_๐ฅ^\nu ({\scriptscriptstyle \frac{1}{2}})2๐ซ_๐ฅ^\nu (1)๐ซ_๐ฅ^\nu ({\scriptscriptstyle \frac{3}{2}})\right)+\mathrm{}$$
(3.21)
where the dots stand for a sum of typical $`๐ฅ`$ representations as in the last example. One can easily find the typical $`๐ค`$ representations in the tensor product of $`[1]`$ with $`[0,2]`$,
$`\left([1][0,2]\right)^{\mathrm{typ}}`$ $`=`$ $`[{\scriptscriptstyle \frac{1}{2}},{\scriptscriptstyle \frac{3}{2}}][{\scriptscriptstyle \frac{1}{2}},{\scriptscriptstyle \frac{5}{2}}][{\scriptscriptstyle \frac{1}{2}},{\scriptscriptstyle \frac{7}{2}}]2[1,2]2[1,3][{\scriptscriptstyle \frac{3}{2}},{\scriptscriptstyle \frac{5}{2}}].`$ (3.22)
We observe that none of these representations contributes any indecomposable upon restriction to the embedded $`๐ฅ`$. Hence, we obtain
$$\left([1][0,2]\left([1][0,2]\right)^{\mathrm{typ}}\right)|_๐ฅ\underset{\nu =\pm }{}\left(๐ซ_๐ฅ^\nu ({\scriptscriptstyle \frac{1}{2}})2๐ซ_๐ฅ^\nu (1)๐ซ_๐ฅ^\nu (3/2)\right)+\mathrm{}$$
(3.23)
According to eq. (3.17), this particular sum of projective $`๐ฅ`$ representations comes from the decomposition of the projective cover $`๐ซ_๐ค(1)`$. Hence, we conclude that the latter appears in the tensor product of $`[1]`$ with $`[0,2]`$, in agreement with our proposition 1.
Before we conclude, we would like to point out that the we can read off the internal structure of $`๐ซ_๐ค(1)`$ from the information on its restriction to $`๐ฅ`$. In fact one can see rather easily that the composition series of the indecomposables in the tensor product of $`[1]`$ and $`[0,2]`$ contains the following list of atypical representations, each displayed with a multiplicity that refers to its transformation properties on the action of the outer automorphisms,<sup>3</sup><sup>3</sup>3The relevant data that allow to determine the mutiplicities can be found in appendix A.
$$[0]_2,2[{\scriptscriptstyle \frac{1}{2}}]_2,[1]_3,3[1]_3,2[{\scriptscriptstyle \frac{3}{2}}]_2,[2]_1.$$
Our planar picture for $`๐ซ_๐ค(1)`$ provides the unique pattern in which we can form a composite from these constituents that is consistent with the information (3.23) on the restriction to the subalgebra $`๐ฅ`$.
## 4 Conclusions and Outlook
In this note we have succeeded to decompose all tensor products between finite dimensional irreducible and projective representations of $`๐ญ๐ฐ๐ฉ(2|2)`$. Whereas tensor products involving at least one projective representation were shown explicitly to stay within the class of projectives, we have constructed a new family of finite dimensional indecomposable representations that appear in the tensor product of atypicals. Preliminary investigations show that tensor products of these new indecomposables $`\pi _{ij}`$ with atypicals generate yet another family of representations whose structure resembles the one of $`\pi _{ij}`$, though with different multiplicities of the involved atypical building blocks. Since the applications we have in mind only require tensor products in which at least one factor is projective, we have not pushed our investigations further into this direction. We believe, however, that results can be obtained using the techniques we have developed above. Even though indecomposables of $`๐ญ๐ฐ๐ฉ(2|2)`$ cannot be classified, it may well be possible to classify all those representations that arise in multiple tensor products of irreducibles. In fact, according to section 3.2, the latter admit an implementation of the SL(2,$``$) outer automorphisms and therefore they form a rather distinguished sub-class of representations.
The techniques we have used here may also be applied to other Lie superalgebras of the A-series, in particular to $`๐ญ๐ฐ๐ฉ(n|n)`$. A promising approach would be to address $`๐ฐ๐ฉ(n|1)`$ first and then to advance to non-trivial second label. Obviously, the structure of the representation theory, becomes much richer for larger Lie superalgebras, in particular because multiply atypical representations can occur (see, e.g., and references therein). Some partial results in this direction will be published elsewhere.
We finally want to sketch at least one concrete physics problem to which we hope to apply the rather mathematical results of this note. During recent years, non-linear $`\sigma `$-models on supergroups and supercosets have surfaced in a variety of distinct problems and in particular through studies of string theory in certain RR backgrounds . Many specific and important properties of these models, such as e.g. the possible existence of conformal invariance even in the absence of a Wess-Zumino term, originate from peculiar features of the underlying Lie superalgebra .
Since the isometries of $`AdS_3\times S^3`$ are generated by two copies of the even subalgebra $`๐ฐ๐ฉ(2,)๐ฐ๐ฒ(2)`$ of the non-compact real form $`๐ญ๐ฐ๐ฒ(1,1|2)`$ of $`๐ญ๐ฐ๐ฉ(2|2)`$, it is not hard to believe that the corresponding $`\sigma `$-models enter the description of strings in an $`AdS_3`$ background . In fact, for strings moving in the presence of a pure NSNS background field the physics is described by a WZW model for $`๐ญ๐ฐ๐ฉ(2|2)`$. This theory possesses a holomorphic and antiholomorphic $`๐ญ๐ฐ๐ฉ(2|2)`$ current symmetry and it may be solved exactly after decoupling the bosons and the fermions, using results established in and references therein. After the marginal deformation which arises from turning on a RR background field, however, the local (worldsheet) symmetries of the system are reduced drastically and so far no solution has been found, in spite of the significant interest in such models (see also e.g. ).
In a first step one may hope to determine the exact spectra of theories with $`AdS_3\times S^3`$ target as a function of the strength of the RR flux. Results in imply that states which transform according to the same representation of the remaining global $`๐ญ๐ฐ๐ฉ(2|2)๐ญ๐ฐ๐ฉ(2|2)`$ symmetry experience the same energy shift as we switch on the RR background field. This motivates to classify all string states according to their behavior under the action of $`๐ญ๐ฐ๐ฉ(2|2)`$. Such states arise by application of creation operators on certain โground statesโ in the theory. At the WZW-point, these creation operators are the negative modes of the $`๐ญ๐ฐ๐ฉ(2|2)`$ currents. The latter transform in the adjoint representation of $`๐ญ๐ฐ๐ฉ(2|2)`$. Hence, in order to determine the transformation properties of excited states, we must control tensor powers (and the symmetric parts therein) of the adjoint representation. This is exactly where the results of the present note feed into studies of strings in $`AdS_3`$. Such an analysis is beyond the scope of this work but we plan to come back to these issues in a forthcoming publication.
Acknowledgment: It is a pleasure to thank Gleb Arutyunov, Jerome Germoni, Hubert Saleur, Paul Sorba and Anne Taormina for many useful discussions. This work was partially supported by the EU Research Training Network grants โEuclidโ, contract number HPRN-CT-2002-00325, โSuperstring Theoryโ, contract number MRTN-CT-2004-512194, and โForcesUniverseโ, contract number MRTN-CT-2004-005104. TQ is supported by a PPARC postdoctoral fellowship under reference PPA/P/S/2002/00370 and partially by the PPARC rolling grant PPA/G/O/2002/00475. We are grateful for the kind hospitality at the ESI during the workshop โString theory in curved backgroundsโ which stimulated the present work.
## Appendix A Appendix: SL(2,$``$) multiplicities
As it is mentioned in the main text, the implementation of the outer automorphisms helps a lot in understanding the decomposition of tensor products. In this appendix we would like to explain how the SL(2,$``$) action organizes representations into multiplets and list explicit results for all the indecomposables we are interested in.
Let us recall that the SL(2,$``$) automorphisms act trivially on the bosonic subalgebra. Hence, each representation in which these automorphisms are implemented may be decomposed into a sum of
$$(i,j)V^J(i,j)_n\text{ where }n=2J+1,$$
$`(i,j)`$ are representations of the even subalgebra and $`V^J`$ carries an action of the automorphism group SL(2,$``$). Since all our representations are assumed to be finite dimensional, the same must be true for $`V^J`$. This means that only SL(2,$``$) representations with half-integer spin $`J`$ and dimension $`n=2J+1`$ can arise.
With the previous remarks in mind we can now move ahead and analyse how the various representations that appeared in the main text decompose into the building blocks $`(i,j)_n`$. We shall restrict our explicit lists here to the irreducible representations. Let us begin with the generic typical irreducible representations for $`i,j>1/2`$:
$`[i,j]=\begin{array}{ccc}& & (i+1,j)_1\hfill \\ & (i+{\scriptscriptstyle \frac{1}{2}},j+{\scriptscriptstyle \frac{1}{2}})_2\hfill & \\ & & (i,j+1)_1\hfill \\ & (i+{\scriptscriptstyle \frac{1}{2}},j{\scriptscriptstyle \frac{1}{2}})_2\hfill & \\ (i,j)_3\hfill & & (i,j)_1\hfill \\ & (i{\scriptscriptstyle \frac{1}{2}},j+{\scriptscriptstyle \frac{1}{2}})_2\hfill & \\ & & (i1,j)_1\hfill \\ & (i{\scriptscriptstyle \frac{1}{2}},j{\scriptscriptstyle \frac{1}{2}})_2\hfill & \\ & & (i,j1)_1.\hfill \end{array}`$ (A.10)
In case one of the labels $`i,j`$ is equal to $`1/2`$, the generic decomposition gets reduced. When $`j=1/2`$ and $`i>1/2`$ we find
$`[i,{\scriptscriptstyle \frac{1}{2}}]=\begin{array}{ccc}& (i+{\scriptscriptstyle \frac{1}{2}},1)_2\hfill & (i+1,{\scriptscriptstyle \frac{1}{2}})_1\hfill \\ & & \\ & (i+{\scriptscriptstyle \frac{1}{2}},0)_2\hfill & (i,{\scriptscriptstyle \frac{3}{2}})_1\hfill \\ (i,{\scriptscriptstyle \frac{1}{2}})_3\hfill & & \\ & (i{\scriptscriptstyle \frac{1}{2}},1)_2\hfill & (i,{\scriptscriptstyle \frac{1}{2}})_1\hfill \\ & & \\ & (i{\scriptscriptstyle \frac{1}{2}},0)_2\hfill & (i1,{\scriptscriptstyle \frac{1}{2}})_1.\hfill \end{array}`$ (A.18)
Obviously, the case of $`i=1/2`$ and $`j>1/2`$ is analogous. The series of representations with $`j=0,i>1/2,`$ possesses an even shorter picture
$`[i,0]=\begin{array}{ccc}& & (i+1,0)_1\hfill \\ & (i+{\scriptscriptstyle \frac{1}{2}},{\scriptscriptstyle \frac{1}{2}})_2\hfill & \\ (i,0)_3\hfill & & (i,1)_1\hfill \\ & (i{\scriptscriptstyle \frac{1}{2}},{\scriptscriptstyle \frac{1}{2}})_2\hfill & \\ & & (i1,0)_1.\hfill \end{array}`$ (A.24)
The last typical representation that remains to be treated is the case of $`i=1/2`$ and $`j=0`$ for which one finds
$`[{\scriptscriptstyle \frac{1}{2}},0]=\begin{array}{ccc}& (1,{\scriptscriptstyle \frac{1}{2}})_2\hfill & ({\scriptscriptstyle \frac{3}{2}},0)_1\hfill \\ ({\scriptscriptstyle \frac{1}{2}},0)_3\hfill & & \\ & (0,{\scriptscriptstyle \frac{1}{2}})_2\hfill & ({\scriptscriptstyle \frac{1}{2}},1)_1.\hfill \end{array}`$ (A.28)
Now we can turn to the irreducible atypical representations with $`j1/2`$,
$`[j]=\begin{array}{cc}& (j+{\scriptscriptstyle \frac{1}{2}},j{\scriptscriptstyle \frac{1}{2}})_1\hfill \\ (j,j)_2\hfill & \\ & (j{\scriptscriptstyle \frac{1}{2}},j+{\scriptscriptstyle \frac{1}{2}})_1.\hfill \end{array}`$ (A.32)
The representation $`[0]`$ is trivially given by $`(0,0)_1`$. This concludes our list of irreducible representations. Similarly, we could now analyse all those indecomposables which allow for an implementation of the SL(2,$``$) automorphisms. Let us stress, that the requirement of implementability is not fulfilled for the Kac modules $`[j,j]`$. Hence, the above formulas should only be used for $`ij`$.
In subsection 3.2 we argued that the atypical constituents of the indecomposables $`๐ซ_๐ค(j)`$ and $`\pi _{ij}`$ are organized in multiplets of SL(2,$``$). For such multiplets we shall employ the symbol $`[j]_m`$. Note that the decomposition of $`[j]_m`$ in terms of $`๐ค^{(0)}sl(2)`$ representations is obtained from the corresponding decomposition of $`[j]`$ (see above) by tensoring with the SL(2,$``$) representation $`V^I`$ where $`m=2I+1`$. In case of the projective covers $`๐ซ(j),j0,`$ one finds the following structure,
$$๐ซ_๐ค(j):\text{}$$
(A.33)
By $`[j]_{3,1}`$ we mean that there is one triplet $`[j]_3`$ and one singlet $`[j]_1`$. This result can be verified by the explicit decomposition of the tensorproduct between $`[1/2]`$ and $`[j+1/2,j1/2]`$ in which, up to typicals, exactly one projective cover $`๐ซ(j)`$ appears.
$`๐ซ(0)`$ has to be treated seperately. Its structure is encoded in a picture of the form
$$๐ซ_๐ค(0):\text{}$$
(A.34)
Finally, we also want to list the multiplicities for the indecomposables $`\pi _{ij}`$ that arise in tensor products of atypicals. For these representations we find
$$\pi _{ij}^{\mathrm{indec}}:\text{}\pi _{jj}^{\mathrm{indec}}:\text{}$$
## Appendix B Appendix: Tensor products for $`๐ฐ๐ฉ(2|1)`$
In order to list results on the tensor products of $`๐ฐ๐ฉ(2|1)`$ representations, we would like to introduce a map $`\pi _๐ฅ`$ which sends representations of the bosonic subalgebra $`๐ฅ^{(0)}`$ to typical representations of $`๐ฅ`$. Its action on irreducibles is given by
$$\pi _๐ฅ(b{\scriptscriptstyle \frac{1}{2}},j{\scriptscriptstyle \frac{1}{2}})=\{\begin{array}{cc}\{b,j\}\hfill & \text{ for }b\pm j,\hfill \\ 0\hfill & \text{ for }b=\pm j.\hfill \end{array}$$
(B.35)
The map $`\pi _๐ฅ`$ may be extended to a linear map on the space of all finite dimensional representations of $`๐ฅ^{(0)}`$.
The first tensor product we would like to display is the one between two typical representations . In our new notations, the decomposition is given by
$`\{b_1,j_1\}\{b_2,j_2\}`$ $`=`$ $`\pi _๐ฅ\left((b_1{\scriptscriptstyle \frac{1}{2}},j_1{\scriptscriptstyle \frac{1}{2}})\{b_2,j_2\}|_{๐ฅ^{(0)}}\right)`$ (B.40)
$`\{\begin{array}{cc}๐ซ_๐ฅ(\pm |b_1+b_2|{\scriptscriptstyle \frac{1}{2}})& \text{for}b_1+b_2=\pm (j_1+j_2)\hfill \\ ๐ซ_๐ฅ^\pm (|b_1+b_2|)๐ซ_๐ฅ^\pm (|b_1+b_2|{\scriptscriptstyle \frac{1}{2}})& \text{for}b_1+b_2\pm \{|j_1j_2|+1,\mathrm{},j_1+j_21\}\hfill \\ ๐ซ_๐ฅ(\pm |b_1+b_2|)& \text{for}b_1+b_2=\pm |j_1j_2|.\hfill \end{array}`$
Note that neither $`j_1`$ nor $`j_2`$ can vanish so that the three cases listed above are mutually exclusive. The first term computes all the typical representations that appear in the tensor product. All it requires is the decomposition of typical $`๐ฅ`$ representations into irreducibles of the bosonic subalgebra,
$$\{b,j\}|_{๐ฅ^{(0)}}=(b,j)(b+{\scriptscriptstyle \frac{1}{2}},j{\scriptscriptstyle \frac{1}{2}})(b{\scriptscriptstyle \frac{1}{2}},j{\scriptscriptstyle \frac{1}{2}})(b,j1).$$
and a computation of tensor products for representations of $`๐ฅ^{(0)}=๐ค๐ฉ(1)๐ฐ๐ฉ(2)`$ which presents no difficulty. The outcome is then converted into a direct sum of typical representations through our map $`\pi _๐ฅ`$.
Tensor products of typical with atypical representations can also be found in Marcuโs paper. The results are
$`\{b_1,j_1\}\{j_2\}_\pm `$ $`=`$ $`\pi _๐ฅ\left((b_1{\scriptscriptstyle \frac{1}{2}},j_1{\scriptscriptstyle \frac{1}{2}})\{j_2\}_\pm |_{๐ฅ^{(0)}}\right)`$
$`\{\begin{array}{cc}๐ซ_๐ฅ^{}(|b_1\pm j_2|{\scriptscriptstyle \frac{1}{2}})\hfill & \text{for}b_1\pm j_2\{|j_1j_2|+1,\mathrm{},j_1+j_2\}\hfill \\ ๐ซ_๐ฅ^\pm (|b_1\pm j_2|)\hfill & \text{for}b_1\pm j_2\pm \{|j_1j_2|,\mathrm{},j_1+j_21\}.\hfill \end{array}`$
This formula may be used to determine the tensor product of typical representations with projective covers and other indecomposables. These tensor products are simply given by
$$\{b,j\}\{b,j\}๐ฎ_๐ฅ().$$
(B.42)
Here, the maps $`๐ฎ_๐ฅ`$ is defined such that it sends the indecomposable representation $``$ to a direct sum of its atypical building blocks, i.e.
$$๐ฎ_๐ฅ()=\underset{\pm j}{}[:\{j\}_\pm ]\{j\}_\pm .$$
(B.43)
Here, $`[:\{j\}_\pm ]`$ counts how many times the atypical representation $`\{j\}_\pm `$ appears in the composition series of $``$. In the special case of $`=๐ซ_๐ฅ^\pm (j)`$, these numbers may be read off from the diagrams (3.12) and (3.13).
Having gone through the entire list of products which involve at least one typical factor, we would like to turn to the fusion of projective covers with any other type of representation. The tensor product between a projective cover $`๐ซ_๐ค^\pm (j)`$ and an atypical representation $`\{n\}_\pm `$ with $`n>0`$ is given by
$`๐ซ_๐ฅ^\pm (j)\{n\}`$ $`=`$ $`\pi \left(H_j^\pm \{n\}|_{๐ฅ^{(0)}}\right)๐ซ_๐ฅ(\pm j+n)`$
where $`H_j^\pm =(\pm j{\scriptscriptstyle \frac{1}{2}},j{\scriptscriptstyle \frac{1}{2}})(\pm (j+{\scriptscriptstyle \frac{1}{2}}){\scriptscriptstyle \frac{1}{2}},j)`$ (B.44)
and $`H_0^\pm =H_0=(0,0)(1,0)`$. We also set $`\{|n|\}_\pm =\{\pm |n|\}`$. As before, we can exploit this formula further to determine all tensor products between a projective cover and an indecomposable composite of atypical representations. The corresponding formula employs our symbol $`๐ฎ_๐ฅ`$ (see eq. (B.43)),
$$๐ซ_๐ฅ^\pm (j)๐ซ_๐ฅ^\pm (j)๐ฎ_๐ฅ().$$
(B.45)
As an application, we can spell out the tensor product between two projective covers $`๐ซ_๐ฅ^\pm (j_1),j_10,`$ and $`๐ซ_๐ฅ(j_2)=๐ซ_๐ฅ^{\text{sign}j_2}(|j_2|)`$,
$`๐ซ_๐ฅ^\pm (j_1)๐ซ_๐ฅ(j_2)`$ $`=`$ $`\pi _๐ฅ\left(H_{j_1}^\pm ๐ซ_๐ฅ(j_2)|_{๐ฅ^{(0)}}\right)`$
$`๐ซ_๐ฅ(\pm j_1+j_2+{\scriptscriptstyle \frac{1}{2}})2๐ซ_๐ฅ(\pm j_1+j_2)๐ซ_๐ฅ(\pm j_1+j_2{\scriptscriptstyle \frac{1}{2}}).`$
The $`๐ค^{(0)}`$ modules $`H_j^\pm `$ were defined in eq. (B.44). In the argument of $`\pi _๐ฅ`$ the product $``$ refers to the fusion between representations of the bosonic subalgebra $`๐ฅ^{(0)}=๐ค๐ฉ(1)๐ฐ๐ฉ(2)`$.
We are finally lacking a formula for the tensor product of two atypical representations. According to , such products are given by
$`\{j_1\}_\pm \{j_2\}_\pm `$ $`=`$ $`\{j_1+j_2\}_\pm {\displaystyle \underset{j=|j_1j_2|}{\overset{j_1+j_21}{}}}\{\pm (j_1+j_2+{\scriptscriptstyle \frac{1}{2}}),j+{\scriptscriptstyle \frac{1}{2}}\},`$ (B.47)
$`\{j_1\}_+\{j_2\}_{}`$ $`=`$ $`\left\{|j_1j_2|\right\}_{\text{sign}(j_1j_2)}{\displaystyle \underset{j=|j_1j_2|+1}{\overset{j_1+j_2}{}}}\{j_1j_2,j\}.`$ (B.48)
Proofs for all these formulas can be found in and in our recent paper .
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# A pseudo-potential analog for zero-range photoassociation and Feshbach resonance
## Abstract
A zero-range approach to atom-molecule coupling is developed in analogy to the Fermi-Huang pseudo-potential treatment of atom-atom interactions. It is shown by explicit comparison to an exactly-solvable finite-range model that replacing the molecular bound-state wavefunction with a regularized delta-function can reproduce the exact scattering amplitude in the long-wavelength limit. Using this approach we find an analytical solution to the two-channel Feshbach resonance problem for two atoms in a spherical harmonic trap.
Coupling between atoms and molecules in quantum-degenerate gases is an ever-present aspect of ultracold atomic physics. Feshbach resonances (FR) InoAndSte98 are now routinely used for control over atomic interactions CorClaRob00 and the formation of molecular Bose-Einstein condensates JocBarAlt03 ; GreRegJin03 ; XuMukAbo03 ; DurVolMar04 . Laser-induced photoassociation (PA) is also widely employed TheThaWin04 ; WynFreHan00 , having the advantage of control over the coupling strength ProPicJun03 . While a zero-range approach to atom-atom collisions has long been a cornerstone of BEC theory, an analog to the Fermi-Huang pseudo-potential approach FerHua has yet to be formulated to treat multichannel free-bound coupling in ultra-cold atomic gases.
In the long-wavelength limit, the energy-dependence of the scattering phase-shift for atomic collisions takes a universal form, with all information about the details of the interaction potential contained in a single parameter, the scattering length. As a result, the full interaction potential can be replaced by a regularized delta-function pseudo-potential, which yields the correct scattering amplitude up to a third-order correction in the ratio of the effective range to the incident wavelength. In this Letter we formulate an analogous approach to atom-molecule coupling, by replacing the bound-state wavefunction with the zero-range object which correctly reproduces the long-wavelength scattering amplitude. The resulting model contains no divergences and does not require a momentum cut-off. It is likely that this model will play an important role in understanding the role played by atom-atom correlations in FR and PA physics, particularly in the strong-coupling regime, where such effects play a dominant role.
The first zero-range model for BEC atom-molecule coupling, proposed by Heinzen and coworkers, replaced the bound-state wavefunction with a delta-function HenWynDru00 . This approach was shown by Holland and coworkers to contain a UV divergence when pair correlations were taken into account HolParWal00 , thus limiting its applicability. Holland and coworkers demonstrated that this divergence could be removed via a momentum cut-off and re-normalized detuning. As we will see, this approach fails in the presence of a background scattering length.
We begin our analysis by considering a pair of atoms described by a relative wavefunction $`\varphi _j(๐ซ,t)`$, where $`j=1,2`$ corresponds to an internal spin state. The eigenstates of this system obey the Schrรถdinger equation,
$$E\varphi _j(๐ซ)=\frac{\mathrm{}^2}{2\mu }^2\varphi _j(๐ซ)+\underset{k}{}V_{jk}(๐ซ,t)\varphi _k(๐ซ),$$
(1)
where $`E`$ is the energy eigenvalue, $`\mu `$ is the reduced mass and $`V_{jk}(๐ซ)`$ is the inter-atom potential. For our model system we assume that the first channel sees a flat potential, $`V_{11}(๐ซ)=0`$. The second channel sees a spherical-well potential of depth $`V_0`$ and radius $`w`$, $`V_{22}(๐ซ)=U_0V_0U(wr)`$, where $`U_0`$ is the continuum threshold energy and $`U(x)`$ is the unit-step function. In the absence of coupling terms, i.e. for $`V_{12}(๐ซ)=0`$, the spectrum of the second channel consists of a continuum of states above the threshold energy, $`U_0`$, and a discrete set of bound states with energies between $`U_0`$ and $`U_0V_0`$. The bound-states are all of the form
$$\psi _b(๐ซ)=\{\begin{array}{cc}๐ฉ_b\frac{e^{r/a_b}}{r}& :r>w\\ ๐ฉ_b\frac{e^{w/a_b}}{\mathrm{sin}(k_bw)}\frac{\mathrm{sin}(k_br)}{r}& :r<w\end{array}$$
(2)
where $`a_b`$ and $`k_b`$ satisfy the equations $`\frac{\mathrm{}^2}{2\mu }\left[k_b^2+1/a_b^2\right]=V_0`$ and $`\mathrm{cot}(k_bw)=1/(k_ba_b)`$, and $`๐ฉ_b`$ is determined by normalization. The bound state energies are $`E_b=U_0\mathrm{}^2/(2\mu a_b^2)`$.
We proceed by first expanding the second channel wavefunction, $`\varphi _2(๐ซ)`$, onto its bare eigenstates under the simplifying assumptions that only a single bound state is near-resonantly coupled to the first channel so that all other states may be neglected. We assume the interaction potential has the form $`V_{12}(๐ซ,t)\frac{\mathrm{}^2G}{\mu }e^{i\omega t}`$. Taking $`E=\mathrm{}^2k^2/(2\mu )`$ then leads to an eigenvalue problem for a continuum coupled to a single bound state,
$`{\displaystyle \frac{1}{2}}\left[k^2+^2\right]\varphi _1(๐ซ)`$ $`=`$ $`G\psi _b(๐ซ)c`$ (3)
$`{\displaystyle \frac{1}{2}}\left[k^22\mathrm{\Delta }\right]c`$ $`=`$ $`G^{}{\displaystyle d^3r\psi _b^{}(๐ซ)\varphi _1(๐ซ)},`$ (4)
where $`c`$ is the probability amplitude for the atom pair to be in the bound state, and $`\mathrm{\Delta }=\frac{\mu }{\mathrm{}^2}(U_0\omega )\frac{1}{2a_b^2}`$ is the detuning away from the atom-molecule resonance at $`k=0`$. The coupling constant $`G`$ will depend on the details of the atom-molecule coupling scheme.
Our goal is now to solve this eigenvalue problem, under the boundary conditions
$`\underset{r\mathrm{}}{lim}\varphi _1(๐ซ)`$ $`=`$ $`{\displaystyle \frac{e^{ikr}}{r}}+f{\displaystyle \frac{e^{ikr}}{r}}`$ (5)
$`\underset{r0}{lim}r\varphi _1(๐ซ)`$ $`=`$ $`0,`$ (6)
in order to determine the scattering amplitude $`f=f(k)`$. The solution can be obtained via the ansatz
$$\varphi _1(๐ซ)=\{\begin{array}{cc}\frac{e^{ikr}}{r}+f\frac{e^{ikr}}{r}+\frac{2Ga_b^2}{1+(a_bk)^2}c\psi _b(๐ซ)& :r>w\\ \beta \frac{\mathrm{sin}(kr)}{r}+\frac{2G}{k^2K_b^2}c\psi _b(๐ซ)& :r<w\end{array}.$$
(7)
This ansatz explicitly satisfies (3), as well as the boundary conditions (5-6). Equation (4), together with the continuity equations $`\varphi _1(w^+)=\varphi _1(w^{})`$ and $`\varphi _1(w^+)=\varphi _1(w^{})`$ can then be used to determine the three unknowns $`f`$, $`c`$, and $`\beta `$. These equations are linear in the three unknowns, and can be thus solved in a straightforward manner.
The long-wavelength limit requires that $`1/k`$ be large compared to the size of the bound-state. As the size of the bound-state is $`w+a_b`$, this is equivalent to the limits $`kw1`$ and $`ka_b1`$. For our model potential the condition $`K_b>1/w`$ is always satisfied, so that $`k/K_b`$ is a small parameter as well. Expanding the scattering amplitude $`f(k)`$ in terms of these small parameters then yields
$$f(k)=\frac{k^22\delta ik\frac{|\chi |^2}{\pi }}{k^22\delta +ik\frac{|\chi |^2}{\pi }}+O[\epsilon ^3],$$
(8)
where $`\epsilon \{kw,ka_b,k/K_b\}`$, and we have introduced the light-shifted detuning
$`\delta `$ $`=`$ $`\mathrm{\Delta }8\pi |G|^2๐ฉ_b^2e^{2w/a}[{\displaystyle \frac{e^{2w/a}}{๐ฉ_b^2K_b^2}}[1+{\displaystyle \frac{1}{K_b^2a_b^2}}]{\displaystyle \frac{a_b^3}{2}}`$ (9)
$`+[1+{\displaystyle \frac{1}{K_b^2a_b^2}}]^2a_b^2(a_b+w)],`$
and the effective coupling constant
$$\chi =4\pi G๐ฉ_be^{w/a}a_b(a_b+w)\left[1+\frac{1}{K_b^2a_b^2}\right].$$
(10)
The important point here is that all of the details of the potential can be absorbed into effective detuning and coupling constants.
We now consider a zero-range model in which the bound-state wavefunction $`\psi _b(๐ซ)`$ in (3) is replaced by a regularized delta-function, $`G\psi _b(๐ซ)\chi \delta ^3(๐ซ)\frac{}{r}r`$. In addition, the detuning $`\mathrm{\Delta }`$ is replaced by the light-shifted detuning $`\delta `$ and the coupling constant $`G`$ is replaced by the effective coupling constant $`\chi `$. The Schrรถdinger equation for this model is given by
$`{\displaystyle \frac{1}{2}}\left[k^2+^2\right]\varphi _1(๐ซ)`$ $`=`$ $`\chi \delta ^3(๐ซ)c`$ (11)
$`{\displaystyle \frac{1}{2}}\left[k^22\delta \right]c`$ $`=`$ $`\chi ^{}{\displaystyle d^3r\delta ^2(๐ซ)\frac{}{r}r\varphi _1(๐ซ)}.`$ (12)
This problem can be solved by making use of the ansazt $`\varphi _1(๐ซ)=\frac{e^{ikr}}{r}+f\frac{e^{ikr}}{r}`$ and the identity $`^2\frac{1}{r}=4\pi \delta ^3(๐ซ)`$. The scattering amplitude is readily found to be
$$f(k)=\frac{k^22\delta ik\frac{|\chi |^2}{\pi }}{k^22\delta +ik\frac{|\chi |^2}{\pi }},$$
(13)
which agrees with the result (8) up to a correction of third-order in the small parameters $`kw`$, $`ka_b`$, and $`k/K_b`$. Thus the zero-range model (11-12) will reproduce correctly the long-wavelength atom-molecule quantum-dynamics of our model potential.
Second quantization of this model yields the Hamiltonian
$`\widehat{}={\displaystyle \frac{\mathrm{}^2}{4m}}{\displaystyle d^3r\left[\widehat{\psi }^{}(๐ซ)^2\widehat{\psi }(๐ซ)+\frac{1}{2}\widehat{\mathrm{\Psi }}^{}(๐ซ)^2\widehat{\mathrm{\Psi }}(๐ซ)\right]}`$
$`+{\displaystyle \frac{\mathrm{}^2\chi }{\sqrt{2}m}}{\displaystyle d^3Rd^3r\widehat{\mathrm{\Psi }}^{}(๐)\delta ^3(๐ซ)\frac{}{r}r\widehat{\psi }(๐+\frac{๐ซ}{2})\widehat{\psi }(๐\frac{๐ซ}{2})}`$
$`+H.c.`$ (14)
where $`\widehat{\psi }(๐ซ)`$ is the annihilation operator for an atom of mass $`m=2\mu `$, and $`\widehat{\mathrm{\Psi }}(๐ซ)`$ is the annihilation operator for a molecule of mass $`2m`$. The system of equations (11-12) can be derived from this Hamiltonian via the 2-atom quantum state
$`|\mathrm{\Psi }`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}{\displaystyle d^3Rd^3r\mathrm{\Phi }(๐)\varphi (๐ซ)\widehat{\psi }^{}(๐+\frac{๐ซ}{2})\widehat{\psi }^{}(๐\frac{๐ซ}{2})|0}`$ (15)
$`+`$ $`c{\displaystyle d^3R\mathrm{\Phi }(๐)\widehat{\mathrm{\Psi }}^{}(๐)|0},`$
where $`\mathrm{\Phi }(๐)`$ is an arbitrary center-of-mass wavefunction and $`|0`$ is the vacuum state. The Hamiltonian (A pseudo-potential analog for zero-range photoassociation and Feshbach resonance) should form the basis of any field-theoretical description of zero-range atom-molecule coupling.
As an example, we now solve the problem of two bosonic atoms in a spherical harmonic oscillator (with frequency $`\omega _{trap}`$) with both s-wave collisions and coupling to a bound state in a second channel. With $`E=\mathrm{}\omega _{trap}(\nu _n+3/2)`$, $`\delta \mathrm{}\omega _{trap}\delta `$, and using harmonic oscillator units, the time-independent Schrรถdinger equation can be written as
$`\left[\nu _n+{\displaystyle \frac{1}{2}}^2{\displaystyle \frac{1}{2}}r^2+{\displaystyle \frac{3}{2}}\right]\varphi _n(๐ซ)=2\pi \left({\displaystyle \frac{a}{\lambda }}\right)\delta ^3(๐ซ){\displaystyle \frac{}{r}}r\varphi _n(๐ซ)`$ (16)
$`+\pi ^{3/4}\mathrm{\Omega }\delta ^3(๐ซ)c_n`$
$`\left[\nu _n\delta \right]c_n=\pi ^{3/4}\mathrm{\Omega }{\displaystyle d^3r\delta ^3(๐ซ)\frac{}{r}r\varphi _n(๐ซ)},`$ (17)
where $`n`$ is an integer label for each quantum level (the lowest energy level corresponding to $`n=0`$), $`a`$ is the background scattering length, $`\lambda `$ is the harmonic oscillator length of the trap, and $`\mathrm{\Omega }=\lambda ^2\pi ^{3/4}\chi `$. The normalized eigenfunctions are found to be BusEngRza98
$$\varphi _n(๐ซ)=\frac{\mathrm{\Omega }}{2\pi ^{3/4}}\frac{\mathrm{\Gamma }[\frac{\nu _n}{2}]}{\beta (\nu _n,\frac{a}{\lambda })}c_ne^{r^2/2}U(\frac{\nu _n}{2},\frac{3}{2},r^2),$$
(18)
$$c_n=\left[1+\frac{\mathrm{\Omega }^2\sqrt{\pi }}{2}\frac{\mathrm{\Gamma }[\frac{\nu _n}{2}]}{\mathrm{\Gamma }[\frac{\nu _n+1}{2}]}\frac{[\psi (\frac{\nu _n}{2})\psi (\frac{\nu _n+1}{2})]}{\beta ^2(\nu _n,\frac{a}{\lambda })}\right]^{1/2},$$
(19)
where $`\beta (\nu ,x)=12x\mathrm{\Gamma }[\frac{\nu }{2}]/\mathrm{\Gamma }[\frac{\nu +1}{2}]`$, $`U(a,b,z)`$ is the confluent hypergeometric function and $`\psi (z)`$ is the polygamma function AbrSte65 . The eigenvalues $`\{\nu _n\}`$ are determined by the characteristic equation
$$\delta =\nu _n\frac{\mathrm{\Omega }^2\sqrt{\pi }\mathrm{\Gamma }[\frac{\nu _n}{2}]}{\mathrm{\Gamma }[\frac{\nu _n+1}{2}]\beta (\nu _n,\frac{a}{\lambda })},$$
(20)
where there is an apparently non-trivial relation $`|c_n|^2=d\nu _n/d\delta `$. It is straightforward to show that the spectrum of eigenvalues will agree exactly with those of a single-channel system with the energy-dependent effective scattering length
$$a_{eff}(\nu )=a+\frac{\lambda }{2}\frac{\sqrt{\pi }\mathrm{\Omega }^2}{(\nu \delta )},$$
(21)
which is the familiar Feshbach Resonance result. The only difference between the true atom-molecule eigenstates and the equivalent single-channel states with scattering length $`a_{eff}`$, is the presence of the bare-molecule population, $`|c_n|^2`$. From a series expansion of (18) the $`1/r`$ part of $`\varphi _n(๐ซ)`$ is found to be $`\frac{\mathrm{\Omega }}{2\pi ^{1/4}\beta (\nu _n,\frac{a}{\lambda })}\frac{c_n}{r}`$. Only for $`a=0`$ is this term independent of $`\nu _n`$, so that it can be removed via a renormalized detuning HolParWal00 .
On resonance we have $`\nu _n=\delta `$ and $`|a_{eff}|\mathrm{}`$. A careful analysis shows that this requires $`\nu _n=2n1`$ and $`c_n0`$. Thus the eigenvalues are driven to odd-integer or โfermionizedโ values, for which the regular part of $`\varphi _n(๐ซ)`$ vanishes at $`๐ซ=0`$. Inserting this result into Eq. (19) gives an analytic expression for the on-resonance molecular fraction,
$$|c_n|^2=\frac{1}{1+\alpha _n\mathrm{\Omega }^2},$$
(22)
where $`\alpha _n=\frac{(2n)!!}{(2n1)!!}\pi /2`$. For the low lying levels we have $`\alpha _0=\pi /2`$, $`\alpha _1=\pi `$ and $`\alpha _2=4\pi /3`$.
The energy-dependence in the effective scattering length is critical to understanding the cross-over between the weak-coupling and strong coupling regimes. The requirement for a significant deviation from the bare-trap spectrum is $`a_{eff}/\lambda 1`$. Obtaining this condition via Feshbach resonance requires $`\delta =\stackrel{~}{\nu }_n\pm \sqrt{\pi }\mathrm{\Omega }^2/2`$. If this width is smaller than the level spacing, only a single level can be near-resonant for a given detuning. In this weak-coupling regime, $`\mathrm{\Omega }^21`$, the spectrum consists of a series of avoided crossings between the bare molecular level and the uncoupled eigenstates of the โopenโ channel. At each avoided crossing there will be strong mixing between a single trap level and the molecular state. Sweeping the detuning can select which trap level is resonantly coupled to the molecular state. This is illustrated in Figure 1, where we have plotted the eigenvalue spectrum as a function of the detuning for the case $`a=.3\lambda `$ and $`\mathrm{\Omega }=.2`$. The dotted lines show the uncoupled ($`\mathrm{\Omega }=0`$) eigenvalues. The shifts in the asymptotic values of the energy levels from the bare trap spectrum ($`\nu _n=2n`$) are due to the presence of s-wave collisions. The asymptotic state at $`.7`$ is the bound state of the โopenโ channel, which is an eigenstate of the trap plus pseudo-potential system.
In the strong coupling regime, defined as $`\mathrm{\Omega }^21`$, the width of the resonance is much larger than the trap level-spacing, hence many levels can be resonant simultaneously. Thus the low-lying levels all lie very close to their on-resonance values of $`\nu _n=2n1`$. This is illustrated in Figure 2, which shows the eigenvalue spectrum as a function of detuning for the case $`\mathrm{\Omega }=10`$ and $`a=0`$. In this regime Eq. (22) is a good estimate for the molecular fraction, showing that the molecular amplitude decreases dramatically with increasing coupling strength. To understand this effect, we simply make the reasonable assumption that in the strong-coupling limit all quasi-resonant levels are mixed with equal amplitudes. For $`\mathrm{\Omega }^21`$, the number of near resonant levels is $`N_{levels}\mathrm{\Omega }^2`$. If we equate the probability for any given bare-state to the total probability divided by the approximate number of levels we arrive at $`|c|^21/\mathrm{\Omega }^2`$, which agrees well with Eq. (22).
In Figure 3 we plot $`a_{eff}`$ and $`|c_n|^2`$ versus detuning for several cases of interest. In Figs 3a and 3b we show the weak-coupling case $`\mathrm{\Omega }=.2`$ and $`a=.3`$ for levels $`n=1`$ and $`n=2`$ respectively. The $`n=1`$ case shows a sweep (right to left) from the lowest โunboundโ state into the bound state in the โopenโ channel. The $`n=2`$ case shows a transfer from one โunboundโ state to another. As the level is swept through resonance we see a broad feature in the molecular fraction $`|c_n|^2`$, whose maximum value is slightly larger than the on-resonance value (22) and occurs to the right of the resonance. Figures 3c and 3d show the intermediate case $`\mathrm{\Omega }=1`$ and $`a=0`$ for levels $`n=0`$ and $`n=1`$. We see in the $`n=1`$ case that the molecular fraction is significantly reduced compared to the weak-coupling regime. Lastly, in Figures 3e and 3f we see the strong-coupling case $`\mathrm{\Omega }=10`$ and $`a=0`$, for levels $`n=0`$ and $`n=1`$. We see that in the strong coupling regime, the scattering length can be tuned from $`\mathrm{}`$ to $`+\mathrm{}`$, with a negligible bare-molecular component.
In conclusion, we see that the effects of pair-correlations play a major role in atom-molecule coupling, resulting in the appearance of a $`1/r`$ singularity in the relative wavefunction together with a corresponding decrease in the bare-molecule population. This suggests that for molecule formation it is best to have a weak coupling, while for manipulation of atomic interactions, e.g. for BCS pairing of fermions HolKokChi01 ; RegGreJin04 , a strong coupling will remove the corresponding bare-molecule population. In FR the free-space coupling strength is predetermined by atomic properties, hence $`\mathrm{\Omega }`$ can only be increased by decreasing the trap size. In PA, however, the coupling strength is readily increased by increasing the laser intensity. This suggests that laser-induced photoassociation may have a significant advantage over Feshbach Resonance for tuning atom-atom interactions.
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# Non-vanishing of the twisted cohomology on the complement of hypersurfaces
## 1. Introduction
Let $`V_1,\mathrm{},V_m`$ be hypersurfaces in the complex projective space $`^n`$ of dimension $`n`$ and let $``$ be a complex local system of rank one over the complement $`M=^n_{i=1}^mV_i`$. In \[A, KN, Ch, Di\], we know vanishing theorems of the twisted cohomology that is the cohomology with the coefficients in the local system $``$ over $`M`$ as follows. Under the generic situation for $`V_1,\mathrm{},V_m`$, if $``$ is non-trivial and generic, then the twisted cohomology on $`M`$ vanishes except in the highest dimension:
$$H^k(M,)=0\text{ for }kn.$$
In particular, vanishing theorems for the case of hyperplanes was found in \[Ko, Yu, CDO, Ka2\] (cf. \[ESV, STV\]). Recently, arrangements of hyperplanes with non-vanishing twisted cohomology were studied well and many examples were found (cf. \[CS, Fa, LY, Ka3\]). \[Yu2\] implied that most of them consist of special elements of pencils in $`^2`$. In this paper we generarize it to hypersurefaces in the general dimension. We shall show that hypersurfaces that is the support of some divisors in some linear system have non-vanishing twisted cohomologies besides the top dimension.
Let $`\mathrm{\Omega }_M`$ denote the sheaf of germs of holomorphic forms on $`M`$ and let $`๐ช_M=\mathrm{\Omega }_M^0`$. Let $`D_1,\mathrm{},D_m`$ be effective divisors on $`^n`$ such that the support of $`D_i`$ is $`V_i`$ for $`i`$ and $`D_i`$ and $`D_j`$ are linearly equivalent for $`ij`$. In this case, $`D_i`$โs have the same degree (see \[H\]). Let $`\lambda =(\lambda _1,\mathrm{},\lambda _m)`$ be a complex weight with $`_{i=1}^m\lambda _i=0`$. For $`ij`$, we have $`D_iD_j`$ is the divisor of some rational function $`f_{ij}`$. Fix $`j`$ and define the global one-form $`\omega _\lambda =_{ij}\lambda _id\mathrm{log}f_{ij}`$. This is independent on the choice of $`j`$ and then we can denote
$$\omega _\lambda =\underset{i=1}{\overset{m}{}}\lambda _id\mathrm{log}D_i\mathrm{\Gamma }(M,\mathrm{\Omega }_M^1),\underset{i=1}{\overset{m}{}}\lambda _i=0.$$
We can define the flat connection $`_\lambda =d+\omega _\lambda :๐ช_M\mathrm{\Omega }_M^1`$ and define the local system $`_\lambda `$ by its kernel. Since $`M`$ is a Stein manifold, we have
$$H^k(M,_\lambda )H^k(\mathrm{\Gamma }(M,\mathrm{\Omega }_M),_\lambda )$$
and $`H^k(M,_\lambda )=0`$ for $`k>n`$ (cf. \[De\]).
###### Theorem 1.
Let $`1<n<sm`$. Let $`D_1,\mathrm{},D_m`$ be effective divisors on $`^n`$ with same degree and let $`M`$ be the complement of their supports. Suppose
* $`D_1,\mathrm{},D_s`$ are elements of some linear system $`\mathrm{\Lambda }`$ on $`^n`$ with dimension $`n1`$.
* There exist a base point $`P`$ of $`\mathrm{\Lambda }`$ such that $`D_1\mathrm{}D_n`$ is normal crossing locally at $`P`$ and $`D_i`$ does not pass through $`P`$ for $`i=s+1,\mathrm{},m`$.
* $`D_1,\mathrm{},D_s`$ are in general position as points in $`\mathrm{\Lambda }=^{n1}`$.
If $`\lambda `$ is a non-trivial weight, that is, $`\lambda ^m`$, such that $`_{i=1}^s\lambda _i=0`$ and $`\lambda _i=0`$ for $`i=s+1,\mathrm{},m`$, then we have
$$dimH^{n1}(M,_\lambda )\left(\genfrac{}{}{0pt}{}{s2}{n1}\right).$$
Moreover, if $`s<m`$ then we have
$$dimH^n(M,_\lambda )\left(\genfrac{}{}{0pt}{}{s2}{n1}\right).$$
###### Remark.
If the support of $`_{i=1}^mD_i`$ contains a hyperplane $`H`$, we can consider $`M`$ as the complement of affine hypersurfaces in $`^n=^nH`$. The setting of the affine hypersurfaces case was found in \[KN, Ki\] (cf. Section 5).
###### Remark.
Let $`D_1,\mathrm{},D_m`$ be divisors of degree $`d`$ defined only by hyperplanes. Let $`๐`$ be the set of irreducible components of the support of $`_{i=1}^mD_i`$. Then $`๐`$ is an arrangement of hyperplanes in $`^n`$ and $`M=^n_{H๐}H`$. If $`D_1,\mathrm{},D_m`$ are divisors of degree one, then they are defined by hyperplanes and this result is known (cf. \[Fa\]).
## 2. Preliminary: Rational forms
Let $`V`$ be a vector space of dimension $`\mathrm{}`$ over a field $`K`$ of characteristic zero. Let $`S=K[V]`$ be the symmetric algebra of $`V^{}`$ and let $`F=K(V)`$ be the quotient field of $`S`$. We consider $`S`$ as the polynomial algebra and $`F`$ as the field of rational functions on $`V`$. Let $`\mathrm{\Omega }(V)=_{p=0}^{\mathrm{}}\mathrm{\Omega }^p(V)`$ be the exterior algebra of $`F`$-vector space $`FV^{}`$ and let $`d`$ be the usual differential. When we choose a basis $`x_1,\mathrm{},x_{\mathrm{}}`$ of $`V^{}`$, we have $`S=K[x_1,\mathrm{},x_{\mathrm{}}]`$, $`F=K(x_1,\mathrm{},x_{\mathrm{}})`$,
$$df=\underset{i=1}{\overset{\mathrm{}}{}}\frac{f}{x_i}x_i=\underset{i=1}{\overset{\mathrm{}}{}}\frac{f}{x_i}dx_i\text{ for }fF,$$
$`\mathrm{\Omega }^0(V)=F`$, $`\mathrm{\Omega }^1(V)=FV^{}=Fdx_1\mathrm{}Fdx_{\mathrm{}}`$, and, $`\mathrm{\Omega }^p(V)=_{i_1<\mathrm{}<i_p}Fdx_{i_1}\mathrm{}dx_{i_p}`$. For $`p2`$ and $`\omega _1,\mathrm{},\omega _p\mathrm{\Omega }^1(V)`$, define
$$\mathrm{\Delta }[\omega _1:\mathrm{}:\omega _p]:=\underset{k=1}{\overset{p}{}}(1)^{k1}\omega _1\mathrm{}\widehat{\omega _k}\mathrm{}\omega _p.$$
###### Lemma 2.
Let $`p2`$ and $`\omega _1,\mathrm{},\omega _p\mathrm{\Omega }^1(V)`$.
1. For a permutation $`\sigma `$ of $`\{1,\mathrm{},p\}`$, we have
$$\mathrm{\Delta }[\omega _{\sigma (1)}:\mathrm{}:\omega _{\sigma (p)}]=\mathrm{sign}(\sigma )\mathrm{\Delta }[\omega _1:\mathrm{}:\omega _p].$$
2. If $`2jp2`$, then
$`\mathrm{\Delta }[\omega _1:\mathrm{}:\omega _p]=`$ $`\mathrm{\Delta }[\omega _1:\mathrm{}:\omega _j]\omega _{j+1}\mathrm{}\omega _p`$
$`+(1)^j\omega _1\mathrm{}\omega _j\mathrm{\Delta }[\omega _{j+1}:\mathrm{}:\omega _p].`$
3. $`\mathrm{\Delta }[\omega _1:\mathrm{}:\omega _p]=(\omega _1\omega _2)\mathrm{\Delta }[\omega _2:\mathrm{}:\omega _p]`$.
4. $`\mathrm{\Delta }[\omega _1:\mathrm{}:\omega _p]=(1)^{p1}(\omega _1\omega _2)(\omega _2\omega _3)\mathrm{}(\omega _{p1}\omega _p)`$.
5. $`\mathrm{\Delta }[\omega _1:\mathrm{}:\omega _p]=(\omega _2\omega _1)(\omega _3\omega _1)\mathrm{}(\omega _p\omega _1)`$.
6. $`\omega _1\mathrm{\Delta }[\omega _1:\mathrm{}:\omega _p]=\omega _1\omega _2\mathrm{}\omega _p`$.
Since (6), if $`\mathrm{\Delta }[\omega _1:\mathrm{}:\omega _p]=0`$ then $`\omega _1,\mathrm{},\omega _p`$ is $`F`$-linearly dependent. However, the inverse is not true in general. The rational function $`fF`$ is said to be homogeneous of degree $`d`$, if $`f=g_1/g_2`$ with homogeneous polynomials $`g_1`$, $`g_2`$ and $`d=\mathrm{deg}g_1\mathrm{deg}g_2`$.
###### Lemma 3.
Assume $`p2`$ and $`f_1,\mathrm{},f_pFK`$ are homogeneous of same degree. Then $`\mathrm{\Delta }[df_1/f_1:\mathrm{}:df_p/f_p]=0`$, if and only if, $`df_1/f_1\mathrm{}df_p/f_p=0`$.
###### Proof.
Let $`\omega _1,\mathrm{},\omega _p\mathrm{\Omega }^1(V)`$. By the direct computation, we can obtain
* $`\mathrm{\Delta }[\omega _1:\mathrm{}:\omega _p]=0`$, if and only if, there exists $`(g_1,\mathrm{},g_p)F^p\{(0,\mathrm{},0)\}`$ such that $`g_1+\mathrm{}+g_p=0`$ and $`g_1\omega _1+\mathrm{}+g_p\omega _p=0`$.
Take $`\omega _i=df_i/f_i`$ and $`g_i=h_if_i`$. Then we get
* For $`p2`$ and $`f_1,\mathrm{},f_pFK`$, we have $`\mathrm{\Delta }[df_1/f_1:\mathrm{}:df_p/f_p]=0`$, if and only if, there exists $`(h_1,\mathrm{},h_p)F^p\{(0,\mathrm{},0)\}`$ such that $`h_1f_1+\mathrm{}+h_pf_p=0`$ and $`h_1df_1+\mathrm{}+h_pdf_p=0`$.
When $`f_i`$โs are homogeneous of same degree, by using the Euler derivation, $`h_1df_1+\mathrm{}+h_pdf_p=0`$ induces $`h_1f_1+\mathrm{}+h_pf_p=0`$. โ
In particular, if $`c_1f_1+\mathrm{}+c_pf_p=0`$ for some $`(c_1,\mathrm{},c_p)K^p\{(0,\mathrm{},0)\}`$, then $`\mathrm{\Delta }[df_1/f_1:\mathrm{}:df_p/f_p]=0`$. If $`f_1^{n_1}\mathrm{}f_p^{n_p}=1`$ for some non-zero integers $`n_1,\mathrm{},n_p`$ then $`\mathrm{\Delta }[df_1/f_1:\mathrm{}:df_p/f_p]=0`$ also.
###### Remark.
$`\mathrm{\Delta }`$ is a natural generalization of the linear derivation on Orlik-Solomon Algebras (\[OT\]), which is in the degree one case.
## 3. Proofs
###### Proof of Theorem 1.
Let $`[x_0:x_1:\mathrm{}:x_n]`$ be homogeneous coordinates of $`^n`$. We can assume $`D_i`$ is given by a homogeneous polynomial $`F_i(x)`$ of degree $`d`$. So we can write $`\omega _\lambda =_{i=1}^m\lambda _idF_i/F_i`$. It is easy to check that $`dF_i/F_idF_j/F_j`$ and $`\omega _\lambda `$ are global forms. For $`1i_1,\mathrm{},i_pm`$, define a holomorphic form on $`M`$ by
$$\eta [i_1,\mathrm{},i_p]:=\mathrm{\Delta }[\frac{dF_{i_1}}{F_{i_1}}:\mathrm{}:\frac{dF_{i_p}}{F_{i_p}}].$$
Since Lemma 2 (4), if it is not zero then $`\eta [i_1,\mathrm{},i_p]`$ is a global ($`p1`$)-form. By (A1) and (A3), $`F_1,\mathrm{},F_n`$ becomes a basis of the vector subspace of $`\mathrm{\Gamma }(^n,๐ช(d))`$ defining the linear system $`\mathrm{\Lambda }`$. So we can write $`F_j=a_{1j}F_1+\mathrm{}+a_{nj}F_n`$ for some constant $`a_{ij}`$ and define the $`n\times s`$-matrix $`A=(a_{ij})`$. By (A3), any $`n\times n`$-minor of $`A`$ is not zero. Due to Lemma 3, we have $`\eta [i_1,\mathrm{},i_{n+1}]=0`$ for $`1i_1<\mathrm{}<i_{n+1}s`$. Because we can write $`\omega _\lambda =_{ii_1}\lambda _i(dF_i/F_idF_{i_1}/F_{i_1})`$, using Lemma 3, we have $`\omega _\lambda \eta [i_1,\mathrm{},i_n]=0`$ and then $`_\lambda (\eta [i_1,\mathrm{},i_n])=0`$ for $`1i_1<\mathrm{}<i_ns`$. We note that, for $`1i_1<\mathrm{}<i_ns`$, by (A2) and (A3), we have $`dF_{i_1}/F_{i_1}\mathrm{}dF_{i_n}/F_{i_n}0`$ and, by Lemma 3, $`\eta [i_1,\mathrm{},i_n]0`$. Thus, $`\eta [i_1,\mathrm{},i_n]`$ is a $`_\lambda `$-closed ($`n1`$)-form for $`1i_1<\mathrm{}<i_ns`$.
By (A2), we take a local neighborhood $`U`$ and coordinates $`x=(x_1,\mathrm{},x_n)`$ at $`P`$ such that $`D_i`$ is defined by $`x_i=0`$ for $`i=1,\mathrm{},n`$. Let $`\alpha _j(x)=a_{1j}x_1+\mathrm{}+a_{nj}x_n`$ and $`H_j=\{\alpha _j(x)=0\}`$. So we get the central arrangement $`๐=\{H_j\}_{1js}`$ of hyperplanes in $`^nU`$, ($`_{i=1}^sH_i`$ is the origin). Let $`M(๐)`$ denote the complement of an arrangement $`๐`$. Then we have $`H^k(UM,_\lambda |_{UM})=H^k(M(๐),\stackrel{~}{}_\lambda )`$ where $`\stackrel{~}{}_\lambda `$ is the rank one local system on $`M(๐)`$ whose monodromy around the hyperplane $`H_j`$ is $`\mathrm{exp}(2\pi \sqrt{1}\lambda _j)`$. Since $`\lambda `$ is non-trivial and $`_{i=1}^s\lambda _i=0`$, without loss of generality, we may assume that $`\lambda _1`$ and $`\lambda _s`$ are not integers. Now, choosing $`H_1๐`$, we get the deconing $`๐๐`$ (see \[OT\]), which is an arrangement of $`s1`$ affine hyperplanes in $`^{n1}H_1`$. Note that $`M(๐)M(๐๐)\times ^{}`$ by the restriction of the Hopf bundle. Since any $`n\times n`$-minor of $`A`$ is not zero, $`๐`$ is generic and $`๐๐`$ is in general position (\[OT\]). On the other hand, for the integer weight $`k^n`$ with $`_{i=1}^sk_i=0`$, we know that the local system $`\stackrel{~}{}_\lambda `$ is equivalent to the local system $`\stackrel{~}{}_{\lambda +k}`$ associated to the integer shift weight $`\lambda +k`$ (see \[OT2\]). By shifting a weight if necessary, we can assume that $`\lambda (\{0\})^n`$. Since $`H^k(M(๐),\stackrel{~}{}_\lambda )H^k(M(๐๐),\stackrel{~}{}_\lambda )H^{k1}(M(๐๐),\stackrel{~}{}_\lambda )`$ (cf. \[Fa\]), in this case, the following is known.
###### Lemma 4 (cf. \[Ha, KN, Ki, Ka\]).
Let $`๐=\{H_j=\{\alpha _j=0\}:1js\}`$ be a generic arrangement of hyperplanes in $`^n`$ and let $`M(๐)`$ be its complement. For a complex weight $`\lambda =(\lambda _1,\mathrm{},\lambda _s)`$ such that $`\lambda _1`$, $`\lambda _s`$ and $`_{i=1}^s\lambda _i=0`$, we have
1. $`H^k(M(๐),\stackrel{~}{}_\lambda )=0`$ for $`kn,n1`$,
2. $`H^n(M(๐),\stackrel{~}{}_\lambda )H^{n1}(M(๐),\stackrel{~}{}_\lambda )`$ and $`dimH^n=dimH^{n1}=\left(\genfrac{}{}{0pt}{}{s2}{n1}\right)`$.
3. $`\{e_1e_{i_1}\mathrm{}e_{i_{n1}}:1<i_1<\mathrm{}<i_{n1}<s\}`$ is a basis of $`H^n`$.
4. $`\{\mathrm{\Delta }[e_1:e_{i_1}:\mathrm{}:e_{i_{n1}}]:1<i_1<\mathrm{}<i_{n1}<s\}`$ is a basis of $`H^{n1}`$,
where $`e_j=d\alpha _j/\alpha _j`$ and $`\stackrel{~}{}_\lambda `$ is the rank one local system on $`M(๐)`$ whose monodromy around the hyperplane $`H_j`$ is $`\mathrm{exp}(2\pi \sqrt{1}\lambda _j)`$.
Note that $`H^k(M(๐),\stackrel{~}{}_\lambda )`$ is isomorphic to the twisted de Rham cohomology defined by the one form $`e_\lambda =_{j=1}^s\lambda _je_j`$ (see \[OT2\]) and that $`\omega _\lambda |_U=e_\lambda `$ and $`\eta [i_1,\mathrm{},i_n]|_U=\mathrm{\Delta }[e_{i_1}:\mathrm{}:e_{i_n}]`$. Now suppose that there exists a global ($`n2`$)-form $`\alpha `$ such that $`\eta [i_1,\mathrm{},i_n]=_\lambda (\alpha )`$. Then restricting it to $`U`$, we have $`\mathrm{\Delta }[e_{i_1}:\mathrm{}:e_{i_n}]=\stackrel{~}{}_\lambda (\alpha |_U)`$ where $`\stackrel{~}{}_\lambda =d+e_\lambda `$. However, by the above Lemma, this is a contradiction. Therefore $`\eta [i_1,\mathrm{},i_n]`$ defines a non-vanishing class of degree $`n1`$ for $`1i_1<\mathrm{}<i_ns`$. In a similar fashion, we obtain $`\{\eta [1,i_1,\mathrm{},i_{n1}]:1<i_1<\mathrm{}<i_{n1}<s\}`$ is independent in $`H^{n1}(M,_\lambda )`$.
Assume $`s<m`$ and fix $`m`$. Take $`\eta [m,i_1,\mathrm{},i_n]`$ for $`1i_1<\mathrm{}<i_ns`$. It is easy to see that $`\eta [m,i_1,\mathrm{},i_n]|_U=e_{i_1}\mathrm{}e_{i_n}`$. Therefore, $`\eta [m,i_1,\mathrm{},i_n]`$ defines a non-vanishing class of degree $`n`$. By the same way, using the above Lemma, we have $`\{\eta [m,1,i_1,\mathrm{},i_{n1}]:1<i_1<\mathrm{}<i_{n1}<s\}`$ is independent in $`H^n(M,_\lambda )`$. This completes the proof. โ
If a weight $`\lambda `$ is trivial then $`H^k(M,_\lambda )`$ is isomorphic to the usual de Rham cohomology $`H^k(M)`$ on $`M`$.
###### Corollary 5.
Under the assumption of Theorem 1, we have
$$dimH^k(M)\left(\genfrac{}{}{0pt}{}{s1}{k}\right)\text{ for }1kn1.$$
Moreover, if $`s<m`$ then we have
$$dimH^n(M)\left(\genfrac{}{}{0pt}{}{s1}{n1}\right)\text{ and }dimH^k(M)\left(\genfrac{}{}{0pt}{}{s}{k}\right)\text{ for }1kn1.$$
###### Proof.
In the proof of Theorem 1, since $`๐`$ is the generic arrangement of $`s`$ hyperplanes in $`^n`$ and $`๐๐`$ is in general position, the following is known (\[OT\]).
1. $`dimH^k(M(๐))=\left(\genfrac{}{}{0pt}{}{s}{k}\right)`$ for $`1kn`$ and $`dimH^k(M(d๐))=\left(\genfrac{}{}{0pt}{}{s1}{k}\right)`$ for $`1kn1`$,
2. $`H^k(M(๐))H^k(M(d๐))H^{k1}(M(d๐))`$ for $`1kn1`$, and $`H^n(M(๐))H^{n1}(M(d๐))`$,
3. $`\{e_1e_{i_1}\mathrm{}e_{i_{n1}}:1<i_1<\mathrm{}<i_{n1}s\}`$ is a basis of $`H^n(M(๐))`$.
4. $`\{\mathrm{\Delta }[e_1:e_{i_1}:\mathrm{}:e_{i_k}]:1<i_1<\mathrm{}<i_ks\}\{e_1e_{i_1}\mathrm{}e_{i_{k1}}:1<i_1<\mathrm{}<i_{k1}s\}`$ is a basis of $`H^k(M(๐))`$ for $`1kn1`$.
Since $`\eta [i_1,\mathrm{},i_k]`$ is $`d`$-closed, the same argument leads this corollary. โ
## 4. Generalization
Let $`๐`$ be the arrangement of (affine or projective) hyperplanes. The intersection set $`L(๐)`$ of $`๐`$ is the set of nonempty intersections of elements of $`๐`$. For $`XL(๐)`$, define a central arrangement $`๐_X=\{H๐|XH\}`$. Let $`๐`$ be a central arrangement with center $`_{H๐}H\mathrm{}`$. We call $`๐`$ decomposable if there exist nonempty subarrangements $`๐_1`$ and $`๐_2`$ so that $`๐=๐_1๐_2`$ is a disjoint union and after a linear coordinate change the defining polynomials for $`๐_1`$ and $`๐_2`$ have no common variables. Define $`\mathrm{D}(๐)=\{XL(๐):๐_X\text{ is not decomposable.}\}`$. For a complex weight $`\lambda `$ of $`๐`$ and $`XL(๐)`$, denote $`\lambda _X=_{H๐_X}\lambda _H`$. The construction of a basis of the twisted cohomology for arrangements given by \[FT\] (cf. \[OT2\]), by the same way of proof of Theorem 1, induces the following.
###### Theorem 6.
Let $`1<n<sm`$. Let $`D_1,\mathrm{},D_m`$ be effective divisors on $`^n`$ with same degree and let $`M`$ be the complement of their supports. Suppose (A1) and (A2). Let $`๐`$ be the arrangement of hyperplanes in the dual projective space $`\mathrm{\Lambda }^{}=(^{n1})^{}^{n1}`$ defined by $`D_1,\mathrm{},D_s`$. Let $`\lambda `$ be a non-trivial weight such that $`_{i=1}^s\lambda _i=0`$ and $`\lambda _i=0`$ for $`i=s+1,\mathrm{},m`$. If $`\lambda _X_0`$ for every $`X\mathrm{D}(๐)`$, then we have
$$dimH^{n1}(M,_\lambda )\beta ,$$
and moreover, if $`s<m`$ then we have
$$dimH^n(M,_\lambda )\beta ,$$
where $`\beta `$ is the Euler characteristic $`\chi (M(๐))`$ of $`M(๐)=^{n1}_{H๐}H`$.
###### Remark.
Note that $`\beta `$ is known as the beta invariant of the underlying matroid of $`๐`$. If $`๐`$ is defined over real then $`\beta `$ is the number of bounded chambers in $`^{n1}=^{n1}H`$ for fixed $`H๐`$ (see \[STV, OT2\]).
###### Remark.
Note that $`\lambda _H_0`$ for all hyperplanes in $`๐`$. We can generalize this theorem in the case that there is $`H๐`$ with $`\lambda _H=0`$, by using \[Ka2\].
## 5. Affine case
Let $`V_1^a,\mathrm{},V_m^a`$ be hypersurfaces in the complex affine space $`^n`$ with coordinates $`u=(u_1,\mathrm{},u_n)`$. Denote $`V^a=_{i=1}^mV_i^a`$ and $`M=^nV^a`$. Assume that a polynomial $`f_j(u)`$ of degree $`d_i`$ defines $`V_j^a`$. Let $`\lambda =(\lambda _1,\mathrm{},\lambda _m)`$ be a weight and let $`\omega _\lambda ^a=_{i=1}^m\lambda _idf_i/f_i`$. Then we obtain the twisted de Rham complex $`(\mathrm{\Omega }(V^a),_\lambda ^a)`$ where $`\mathrm{\Omega }(V^a)`$ is the space of rational forms with poles along $`V^a`$ and $`_\lambda ^a=d+\omega _\lambda ^a`$. The Grothendiek-Deligue comparison theorem (\[De\]) asserts
$$H^k(M,_\lambda ^a)H^k(\mathrm{\Omega }(V^a),_\lambda ^a),$$
where $`_\lambda ^a`$ is the rank one local system defined by the flat connection $`_\lambda ^a`$ (cf. \[KN, Ki\]).
Let $`^n`$ be the complex projective space with the infinite hyperplane $`H_{\mathrm{}}`$, which is a compactification of $`^n`$. Let $`[x_0:\mathrm{}:x_n]`$ denote homogeneous coordinates with $`H_{\mathrm{}}=\{x_0=0\}`$. Define the homogeneous polynomial $`F_j(x)=x_0^df_j(x_1/x_0,\mathrm{},x_n/x_0)`$ of degree $`d=\mathrm{max}(d_1,\mathrm{},d_m)`$. Then $`F_j(x)`$ determines the divisor $`D_j`$ of degree $`d`$. If $`d=d_1=\mathrm{}=d_m`$ then the support of $`_{i=1}^mD_i`$ does not contain $`H_{\mathrm{}}`$, otherwise it contains $`H_{\mathrm{}}`$. Note that the weight of $`H_{\mathrm{}}`$ is given by $`_{i=1}^m\lambda _id_j`$. If $`d=d_1=\mathrm{}=d_m`$ and $`_{i=1}^m\lambda _i=0`$ then the weight of $`H_{\mathrm{}}`$ is zero.
###### Corollary 7.
Under the assumption in Theorem 1, if a point $`P`$ in (A2) is not on $`H_{\mathrm{}}`$, then we have
$$dimH^{n1}(M^a,_\lambda ^a)\left(\genfrac{}{}{0pt}{}{s2}{n1}\right)$$
for a non-trivial weight $`\lambda `$ with $`_{i=1}^s\lambda _i=0`$ and $`\lambda _i=0`$ for $`i=s+1,\mathrm{},m`$.
## 6. Examples
### 6.1. $`n=2`$ and $`s=3`$
Let $`D_1`$, $`D_2`$ and $`D_3`$ be irreducible prime divisors on $`^2`$ defined by $`F_1=x_0^2+x_1^22x_2^2`$, $`F_2=x_0^2+2x_1^23x_2^2`$ and $`F_3=2x_0^2+x_1^23x_2^2`$, respectively. They are three generic elements of the pencil of conic curves $`F_{[a:b]}=a(x_0^2x_2^2)+b(x_1^2x_2^2)`$ and they intersect transversally each other at each four intersection points. Let $`M=^2_{i=1}^3D_i`$ and $`M^a=^2_{i=1}^3(D_i^2)`$ where $`^2=^2\{x_2=0\}`$. We have $`H^1(M,_\lambda )0`$ and $`H^1(M^a,_\lambda )0`$ for a non-trivial weight $`\lambda `$ with $`\lambda _1+\lambda _2+\lambda _3=0`$. In the following case, we can get it also.
1. $`F_1=x_0^2+x_1^22x_2^2`$, $`F_2=x_0^2+2x_1^23x_2^2`$ and $`F_3^{}=x_0^2x_1^2`$ (two conics and a set of 2-lines).
2. $`F_1=x_0^2+x_1^22x_2^2`$, $`F_2^{}=x_2^2x_0^2`$ and $`F_3^{}=x_0^2x_1^2`$ (one conics and two sets of 2-lines).
3. $`F_1^{}=x_1^2x_2^2`$, $`F_2^{}=x_2^2x_0^2`$ and $`F_3^{}=x_0^2x_1^2`$ (three sets of 2-lines).
In the last case, it is an arrangement of 6 lines in the Ceva Theorem. In a similar fashion, we get the following. In the degree three case, we get an arrangement of 9 lines in the Pappus Theorem (\[Fa\]). In the degree four case, there are two different arrangements of 12 lines in the Kirkman Theorem and the Steiner Theorem (\[Ka3\]). Those arrangements are 3-nets, whose combinatorial structures are matroids associated to Latin squares (\[LY, Yu2, Ka3\]). In the other hand, the $`B_3`$-arrangement is an example of the case that $`D_i`$โs are not prime. Let $`D_1`$, $`D_2`$ and $`D_3`$ be divisors defined by $`x_2^2(x_0^2x_1^2)`$, $`x_1^2(x_0^2x_2^2)`$ and $`x_0^2(x_1^2x_2^2)`$, respectively. Their divisors can be written by $`D_1=2H_1+H_2+H_3`$, $`D_2=2H_4+H_5+H_6`$ and $`D_3=2H_7+H_8+H_9`$ where $`H_i`$โs are hyperplanes. The arrangement $`๐=\{H_1,\mathrm{},H_9\}`$ is called the $`B_3`$-arrangement. Note that a weight $`\lambda `$ induces the weight $`(2\lambda _1,\lambda _1,\lambda _1,2\lambda _2,\lambda _2,\lambda _2,2\lambda _3,\lambda _3,\lambda _3)`$ of $`๐`$ (cf. \[Fa, Ka3\]).
### 6.2. $`n=2`$ and $`s3`$
We consider the pencil of cubic curves $`F_{[a:b]}=a(x_0^3+x_1^3+x_2^3)+3bx_0x_1x_2`$. A generic element given by $`a0`$ and $`b^31`$ is non-singular. For $`s`$ generic elements $`D_1,\mathrm{},D_s`$, we have $`dimH^1(M,_\lambda )s2`$. Define non-generic elements $`F_1=x_0x_1x_2`$, $`F_2=x_0^3+x_1^3+x_2^33x_0x_1x_2`$, $`F_3=x_0^3+x_1^3+x_2^33\xi x_0x_1x_2`$ and $`F_4=x_0^3+x_1^3+x_2^33\xi ^2x_0x_1x_2`$ where $`\xi =e^{2\pi \sqrt{1}/3}`$. They are four sets of 3-lines and the set $`๐`$ of all lines is called the Hessian configuration, which is the arrangement of 12 lines passing through the nine inflection points of a nonsingular cubic. In this case, we know $`dimH^1(M,_\lambda )=2`$ for a non-trivial weight $`\lambda `$ with $`_{i=1}^4\lambda _i=0`$ (\[Li\]).
### 6.3. $`n=3`$
Let $`F_1=x_0(x_1+x_2+x_3)`$, $`F_2=x_1(x_0+x_2x_3)`$, $`F_3=x_2(x_0x_1+x_3)`$ and $`F_4=x_3(x_0+x_1x_2)`$. Then we can check divisors defined them satisfy the conditions in Theorem 1 and then $`H^2(M,_\lambda )0`$ for a weight $`\lambda `$ with $`_{i=1}^4\lambda _i=0`$. They yield the arrangement of 8 planes, whose underlying matroid is of type $`L_8`$ (\[Ka3\]).
### 6.4. higher dimensional case
Let $`F_i=x_{i1}^dx_i^d`$ for $`i=1,\mathrm{},n`$ and $`F_0=x_n^dx_0^d`$. This support determines the arrangement $`๐`$ of $`(n+1)d`$ hyperplanes in $`^n`$. Note that this is a projective closure of a subarrangement of the monomial arrangement $`๐_{d,d,n+1}`$ (see \[OT, CS\]). Since $`_{k=0}^nF_k=0`$, we have $`H^{n1}(M,_\lambda )0`$ for a weight $`\lambda `$ with $`_{i=0}^n\lambda _i=0`$. The $`n=2`$ case was found in \[CS\]. Note that the underlying matroid of $`๐`$ is a degeneration of the matroid associated to the Latin $`n`$-dimensional hypercube given by the addition table for $`(_d)^n`$ (see \[Ka3\]).
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# Fractal Weyl laws in discrete models of chaotic scattering
## 1. Introduction
The study of resonances, or quasibound states, has a long tradition in theoretical, numerical, and experimental chaotic scattering โ see for instance and references given there. In this paper we discuss the laws for the density of resonances at high energies, or in the semiclassical limit, and the closely related asymptotics of conductance, Fano factors, and โshot noiseโ. Our models are based on a quantization of open baker maps and we focus on fractal Weyl laws for the density of resonances. These laws have origins in the mathematical work on counting resonances .
In ยง2 we present the compact phase space models for chaotic scattering (open bakerโs maps) and their discrete quantizations. The numerical results on counting of quantum resonances showing an agreement with fractal Weyl laws are given in ยง3. In ยง4 we discuss a model which is simpler on the quantum level but more complicated on the classical level (this model can also be interpreted as an alternative quantization of the original bakerโs map, see ยง4.2). In that case we can describe the distribution of resonances very precisely (ยง4.2), showing perfect agreement with the fractal Weyl law. We also find asymptotic expressions for the conductance and the Fano factor (or the โshot noiseโ factor). The fractal Weyl law appears naturally in these asymptotics and an interesting comparison with the random matrix theory is also made (ยง4.3).
To put the fractal Weyl law in perspective we review the usual Weyl law for the density of states in the semiclassical limit. Let $`H(q,p)=p^2+V(q)`$ be a Hamiltonian with a confining potential $`V`$ and let $`E`$ be a nondegenerate energy level,
(1.1)
$$H(q,p)=EH(q,p)0.$$
Assume further that the union of periodic orbits of the Hamilton flow on the surface $`H^1(E)`$ has measure zero with respect to the Liouville measure. Then the spectrum of the quantized Hamiltonian,
(1.2)
$$\widehat{H}=\mathrm{}^2\mathrm{\Delta }+V(q),q^n,$$
near $`E`$ satisfies,
(1.3)
$$\mathrm{\#}\left\{\mathrm{Spec}(\widehat{H})[E\rho \mathrm{},E+\rho \mathrm{}]\right\}=\frac{2\rho \mathrm{}}{(2\pi \mathrm{})^n}\delta \left(H(q,p)E\right)๐q๐p+o(\mathrm{}^{n+1}),$$
see for references to the mathematical literature on this subject.
Suppose now that $`V(q)`$ is not confining. The most extreme case is given by $`V(q)`$ vanishing outside a compact set. An example of such a potential with $`q^2`$ is given in Fig. 1. In that case the eigenvalues are replaced by quantum resonances which can be defined as the poles of the meromorphic continuation of Greenโs function, $`G(z;q^{},q)`$, from $`Imz>0`$ to $`Imz0`$. By Greenโs function we mean the integral kernel of the resolvent:
(1.4)
$$(z\widehat{H})^1u(q^{})=_^nG(z;q^{},q)u(q)๐q,u๐_\mathrm{c}^{\mathrm{}}(^n).$$
We denote the set of resonances by $`\mathrm{Res}(\widehat{H})`$. Near a nondegenerate energy level (1.1) we have the following bound (compare with (1.3) for a closed system) :
(1.5)
$$\mathrm{\#}\left\{\mathrm{Res}(\widehat{H})\left([E\rho \mathrm{},E+\rho \mathrm{}]i[0,\gamma \mathrm{}]\right)\right\}C(\rho ,\gamma )\mathrm{}^{n+1}.$$
When the interaction region is separated from infinity by a barrier, this bound is optimal since resonances are well approximated by eigenvalues of a closed system. In that case the classical trapped set,
(1.6)
$$K_E\stackrel{\mathrm{def}}{=}\{(q,p)H^1(E):\mathrm{\Phi }_H^t(q,p)\to ฬธ\mathrm{},t\pm \mathrm{}\},$$
has a non-empty interior in $`H^1(E)`$, so that its dimension is equal to $`2n1`$.
Suppose now that the classical flow of the Hamiltonian $`H`$ is hyperbolic on $`K_E`$, as is the case for instance in some energy range for the 2-D potential represented in Fig. 1 .
Following the original work of Sjรถstrand , the general upper bound (1.5) is replaced by a bound involving the upper Minkowski dimension of $`K_E`$:
$$dimK_E=2n1sup\{c:\underset{ฯต0}{lim\; sup}ฯต^c\mathrm{vol}(\{\rho \mathrm{\Sigma }_E:d(K_E,\rho )<ฯต\})<\mathrm{}\}.$$
We say that $`K_E`$ is of pure dimension if the supremum is attained. For simplicity we assume that this is the case. Then under the assumption of hyperbolicity of the flow, we have
(1.7)
$$\mathrm{\#}\left\{\mathrm{Res}(\widehat{H})\left([E\rho \mathrm{},E+\rho \mathrm{}]i[0,\gamma \mathrm{}]\right)\right\}C(\rho ,\gamma )\mathrm{}^{\mu _E},2\mu _E+1=dimK_E.$$
This bound is expected to be optimal even though it is not clear what notion of dimension should be used for the lower bounds. The best chance lies in cases in which $`K_E`$ has a particularly nice structure. A class of Hamiltonians for which that happens is given by quotients of hyperbolic space by convex co-compact discrete groups .
A fractal Weyl law for the density of resonances in larger regions are easier to verify and more likely to hold in general:
(1.8)
$$\mathrm{\#}\left\{\mathrm{Res}(\widehat{H})\left([E\delta ,E+\delta ]i[0,\gamma \mathrm{}]\right)\right\}C(\delta ,\gamma )\mathrm{}^{\mu _E1},\delta >0\text{fixed}.$$
The precise meaning of $``$ is left vague in this conjectural statement. The exponent in this relation has been investigated numerically in a variety of settings and the results are encouraging .
### Acknowledgments
The first author thanks Marcos Saraceno for his insights and comments on the various types of quantum bakers. He is also grateful to UC Berkeley for the hospitality in April 2004. Generous support of both authors by the National Science Foundation under the grant DMS-0200732 is also gratefully acknowledged.
## 2. Open baker maps and their quantizations
We consider $`๐^2=[0,1)\times [0,1)`$, the two-torus, as our classical phase space. Classical observables are functions on $`๐^2`$ and the classical dynamics is given in terms of an โopen symplectic mapโ $`B`$, that is a map defined on a subset $`๐๐^2`$, which is invertible and canonical (area and orientation preserving) from $`๐`$ to $`B(๐)`$. The points of $`๐^2๐`$ are interpreted as โfalling in the holeโ, or โescaping to infinityโ
Following a construction performed in , we will be concerned with open versions of the bakerโs map, obtained by restricting the โclosedโ bakerโs map to a subdomain of $`๐^2`$, union of vertical strips. As an example, if we restrict the 3-bakerโs map $`A_3`$:
(2.5)
$$\begin{array}{c}A_3(q,p)\stackrel{\mathrm{def}}{=}(q^{},p^{})=\{\begin{array}{ccc}q^{}=3q,\hfill & p^{}=p/3,\hfill & \text{if}0q<1/3\hfill \\ q^{}=3q1,\hfill & p^{}=(p+1)/3,\hfill & \text{if}1/3q<2/3\hfill \\ q^{}=3q2,\hfill & p^{}=(p+2)/3,\hfill & \text{if}2/3q<1\hfill \end{array}.\end{array}$$
to the domain $`๐_3=๐^2\{1/3q<2/3\}`$, we obtain the open 3-bakerโs map $`B_3`$:
(2.6)
$$(q,p)๐_3,B_3(q,p)=(q^{},p^{})=\{\begin{array}{ccc}q^{}=3q,\hfill & p^{}=p/3,\hfill & \text{if}0q<1/3\hfill \\ q^{}=3q2,\hfill & p^{}=(p+2)/3,\hfill & \text{if}2/3q<1.\hfill \end{array}$$
This open map admits an inverse $`B_3^1`$, which is a canonical map from $`B_3(๐_3)`$ to $`๐_3`$. In this paper we will present numerical results for an open 5-bakerโs map, defined as
(2.7)
$$B_5(q,p)=(q^{},p^{})\stackrel{\mathrm{def}}{=}\{\begin{array}{ccc}q^{}=5q1,\hfill & p^{}=(p+1)/5,\hfill & \text{if}1/5q<2/5\hfill \\ q^{}=5q3,\hfill & p^{}=(p+3)/5,\hfill & \text{if}3/5q<4/5.\hfill \end{array}$$
One can think of $`A_3`$ as model for a Poincarรฉ map for a 2-D closed Hamiltonian system. Removing the domain $`\{1/3q<2/3\}`$ from the torus corresponds to opening the system: the points in this domain will escape through the hole, that is, never come back to the Poincarรฉ section. In the context of mesoscopic quantum dots, such an opening is performed by connecting a lead to the dot, through which electrons are able to escape (see ยง4.3).
For open maps such as $`B=B_3`$, we can define the incoming and outgoing tails, made of points which never escape in the forward (resp. backward) evolution:
$$\begin{array}{cc}\hfill x\mathrm{\Gamma }_{}๐^2& n0,B^n(x)๐.\hfill \\ \hfill x\mathrm{\Gamma }_+& n0,B^n(x)B(๐).\hfill \end{array}$$
In the case of the map (2.6), $`\mathrm{\Gamma }_{}=๐_3\times [0,1)`$, $`\mathrm{\Gamma }_+=[0,1)\times ๐_3`$, where $`๐_3`$ is the standard $`\frac{1}{3}`$-Cantor set on the interval (see Fig. 2).
In analogy with (1.6), we also define the trapped set $`K=\mathrm{\Gamma }_+\mathrm{\Gamma }_{}`$ and, for any point $`xK`$, its stable and unstable manifolds $`W_\pm (x)`$. In the case of the open 3-baker $`B_3`$, we easily check that
$$\mu \stackrel{\mathrm{def}}{=}dim\mathrm{\Gamma }_{}W_+=dim\mathrm{\Gamma }_+W_{}=\frac{1}{2}dimK=\frac{\mathrm{log}2}{\mathrm{log}3}.$$
The quantization of the open map (2.6) is based on the quantization of the โclosedโ bakerโs map $`A_3`$. That, in an outline, is done as follows . To any $`N`$ we associate a space $`_N^N`$ of quantum states on the torus. The components $`\psi _j`$, $`j_N=\{\mathrm{\hspace{0.17em}0},\mathrm{},N1\}`$ of a state $`\psi _N`$ are the amplitudes of $`\psi `$ at the positions $`q=q_j=(j+\frac{1}{2})/N`$, and we will sometimes use Diracโs notation $`\psi _j=q_j|\psi `$. The choice of these โhalf-integers positionsโ is justified by the parity symmetry $`q1q`$ they satisfy , and is further explained in ยง3. The scalar product on $`_N`$ is the standard one on $`^N`$:
(2.8)
$$\psi ,\varphi _N,\varphi |\psi =\underset{j=0}{\overset{N1}{}}\overline{\varphi }_j\psi _j.$$
A classical observable depending on $`q/`$ only, $`f=f(q)`$, is obviously quantized as the multiplication operator
$$\psi _N,[\mathrm{Op}_N(f)\psi ]_j=f\left(\frac{j+1/2}{N}\right)\psi _j.$$
The discrete Fourier transform
(2.9)
$$(๐ข_N)_{j,j^{}}=N^{1/2}e^{2i\pi (j+\frac{1}{2})(j^{}+\frac{1}{2})/N},j,j^{}=0,\mathrm{},N1$$
transforms a โposition vectorโ $`\psi _j=q_j|\psi `$ into the corresponding โmomentum vectorโ $`p_j|\psi =(๐ข_N\psi )_j`$. The momenta are also quantized to values $`p_j=(j+\frac{1}{2})/N`$, $`j=0,\mathrm{},N1`$. Comparing the definition (2.9) with the (standard) Fourier transform on $``$,
$$_{\mathrm{}}u(p)=\frac{1}{\sqrt{2\pi \mathrm{}}}_{}e^{ipq/\mathrm{}}u(q)๐q,$$
we see that the effective Planckโs constant in the discrete model is $`\mathrm{}=(2\pi N)^1`$.
As a result, any observable $`g=g(p)`$ can be quantized as
$$\mathrm{Op}_N(g)\psi =๐ข_N^{}\text{diag}\left(g\left((j+1/2)/N\right)\right)๐ข_N\psi .$$
The Weyl quantization on the torus generalizes this map $`f\mathrm{Op}_N(f)`$ to any classical observable $`f`$, that is any (smooth) function on the torus, in such a way that a real observable $`f`$ is associated with self-adjoint operators, and
$$\frac{i}{\mathrm{}}[\mathrm{Op}_N(f),\mathrm{Op}_N(g)]=\mathrm{Op}_N(\{f,g\})+๐ช(\mathrm{}^2).$$
Let us now consider the following family of unitary operators on $`_N`$, where $`N`$ is taken as a multiple of $`3`$:
(2.10)
$$\widehat{A}_{3,\mathrm{pos}}=A_{3,N}\stackrel{\mathrm{def}}{=}๐ข_N^{}\left(\begin{array}{ccc}๐ข_{N/3}\hfill & 0\hfill & 0\hfill \\ 0\hfill & ๐ข_{N/3}\hfill & 0\hfill \\ 0\hfill & 0\hfill & ๐ข_{N/3}\hfill \end{array}\right).$$
Since $`๐ข_N`$ exchanges position and momentum, the mixed momentum-position representation of $`\widehat{A}_3`$ is given by the matrix
$$\widehat{A}_{3,\mathrm{mom}\mathrm{pos}}=๐ข_NA_{3,N}=\left(\begin{array}{ccc}๐ข_{N/3}& 0& 0\\ 0& ๐ข_{N/3}& 0\\ 0& 0& ๐ข_{N/3}\end{array}\right).$$
In terms of the quantized positions $`q_j`$ and momenta $`p_k`$, the entries of this matrix are given by
$$\left(\widehat{A}_{3,\mathrm{mom}\mathrm{pos}}\right)_{kj}=p_k|\widehat{A}_3|q_j=\frac{1}{\sqrt{2\pi \mathrm{}}}\mathrm{exp}\left(\frac{i}{\mathrm{}}(3q_j\mathrm{})\left(p_k\frac{\mathrm{}}{3}\right)\right),$$
$$\frac{\mathrm{}}{3}q_j<\frac{\mathrm{}+1}{3},\frac{\mathrm{}}{3}p_k<\frac{\mathrm{}+1}{3},\mathrm{}=0,1,2,$$
and zero otherwise. One can then observe that for $`\mathrm{}=0,1,2`$, the function $`S_{\mathrm{}}(p^{},q)=(3q\mathrm{})(p^{}\mathrm{}/3)`$ generates the canonical map $`(q,p)(q^{},p^{})=(3q\mathrm{},p/3+\mathrm{}/3)`$ on the domain $`\{q,p^{}[\mathrm{}/3,(\mathrm{}+1)/3)\}`$, that is, the map $`A_3`$ (2.5) on this domain. The matrix elements $`p_k|\widehat{A}_3|q_j`$ therefore exactly correspond to the Van Vleck semiclassical formula associated with the map $`A_3`$. For this reason (and the unitarity of $`\widehat{A}_3`$), the operator $`\widehat{A}_3`$ was considered a good quantization of $`A_3`$ by Balazs and Voros. A more precise description of the correspondence between $`A_3`$ and $`\widehat{A}_3`$, including the role played by the discontinuities of $`A_3`$, is explained in \[12, ยง4.4\].
To quantize the open baker $`B_3`$ (2.6), we truncate the unitary operator $`\widehat{A}_3`$ using the quantum projector on the domain $`๐`$, $`\mathrm{\Pi }_๐\stackrel{\mathrm{def}}{=}\mathrm{Op}_N(1\mathrm{l}_๐)`$ : in the position basis, we get
(2.11)
$$\widehat{B}_{3,\mathrm{pos}}=B_{3,N}\stackrel{\mathrm{def}}{=}A_{3,N}\mathrm{\Pi }_๐=๐ข_N^{}\left(\begin{array}{ccc}๐ข_{N/3}& 0& 0\\ 0& 0& 0\\ 0& 0& ๐ข_{N/3}\end{array}\right),N3.$$
This subunitary operator is a model for the quantization of a Poincarรฉ map of an open chaotic system . The semiclassical rรฉgime corresponds to the limit $`N\mathrm{}`$. Similarly, the quantum open map associated with the 5-baker $`B_5`$ (2.7) is given by the sequence of matrices:
(2.12)
$$B_{5,N}\stackrel{\mathrm{def}}{=}๐ข_N^{}\left(\begin{array}{ccccc}0& 0& 0& 0& 0\\ 0& ๐ข_{N/5}& 0& 0& 0\\ 0& 0& 0& 0& 0\\ 0& 0& 0& ๐ข_{N/5}& 0\\ 0& 0& 0& 0& 0\end{array}\right),N5.$$
Let us now describe the correspondence between the resonances of a Schrรถdinger operator $`\widehat{H}`$, and the eigenvalues of our subunitary open quantum maps $`B_{3,N}`$ or $`B_{5,N}`$ (denoted generically by $`B_N`$).
Since $`B_N`$ is obtained by truncating the unitary propagator $`A_N`$, it is natural to consider the family of truncated Schrรถdinger propagators $`\chi e^{it\widehat{H}/\mathrm{}}\chi `$, where $`\chi (q)`$ is a cutoff function on some compact set supporting the scatterer. Although the precise eigenvalues of these propagators depend nontrivially on both $`\chi `$ and the time $`t`$, these propagators admit a long-time expansion in terms of the resonances $`z_j`$ of $`\widehat{H}`$ . At an informal level, one may write
$$\chi e^{it\widehat{H}/\mathrm{}}\chi \underset{z_j\mathrm{Res}(\widehat{H})}{}e^{itz_j/\mathrm{}}\widehat{R}_j.$$
On the other hand, the iterated open quantum map $`(B_N)^n`$ can obviously be expanded in terms of the eigenvalues $`\lambda _j`$ of $`B_N`$. For this reason, it makes sense to model the exponentials $`e^{iz_j/\mathrm{}}`$ by the eigenvalues $`\lambda _j`$ of our open quantum map $`B_N`$.
Upon this identification, the boxes in which we count resonances in (1.7), $`[E\rho \mathrm{},E+\rho \mathrm{}]i[0,\gamma \mathrm{}]`$, correspond to the regions
(2.13)
$$๐_{r,\vartheta ,\rho }\stackrel{\mathrm{def}}{=}\{\mathrm{\hspace{0.17em}1}|\lambda |r,|\mathrm{arg}(\lambda e^{i\vartheta })|\rho \},r=\mathrm{exp}(\gamma )(0,1).$$
These analogies induce a conjectural fractal Weyl law for the quantum open bakers (2.6,2.7) which we now describe.
First of all, we consider the partial dimension of the trapped set of the open map $`B`$:
$$\mu =\frac{dimK}{2}=dim(\mathrm{\Gamma }_{}W_+).$$
Then, for any $`r(0,1)`$, there should exist $`C(r)0`$ (a priori, depending on the map $`B`$) such that, in the semiclassical limit, the number of eigenvalues of $`B_N`$ in the sectors (2.13) behaves as
(2.14)
$$\mathrm{\#}\left\{\lambda \mathrm{Spec}(B_N)๐_{r,\vartheta ,\rho }\right\}\frac{\rho }{2\pi }C(r)N^\mu N\mathrm{}.$$
The angular dependence $`\rho /(2\pi )`$ on the RHS means that the distribution of eigenvalues is expected to be asymptotically angular-symmetric.
In , the quantum 2-baker $`A_{2,N}`$ was decomposed into the block $`๐ข_N^{}\left(\begin{array}{cc}๐ข_{N/2}& 0\\ 0& 0\end{array}\right)`$ and the complementary one. The spectral determinant for the unitary map, $`det(1zA_{2,N})`$, was then expanded in terms of these blocks. Although the classical open map associated with each block is quite simple (all points except a fixed one eventually escape), the spectrum of each block was found to be rather complex, and quite different from semiclassical predictions.
A scaling of the type (2.14) was conjectured in for another chaotic map, namely the open kicked rotator. This conjecture was then tested numerically, and a good agreement was observed. The scaling law $`N^\mu `$ was explained heuristically by counting the number of quantum states in an $`\mathrm{}`$-neighbourhood of the incoming tail $`\mathrm{\Gamma }_{}`$ (so that $`\mu `$ was effectively the dimension of $`\mathrm{\Gamma }_{}`$). For the kicked rotator, the fractal exponent $`\mu `$ was not known analytically, and the authors related it to the *mean dwell time* of the dynamics, that is, the average time spent in the cavity before leaving it: $`\mu 1(\lambda \tau _{\mathrm{dwell}})^1`$ ($`\lambda `$ is the mean Lyapounov exponent). For our open bakerโs maps $`B_3`$, the dwell time is $`\tau _{\mathrm{dwell}}=3`$, so the above formula is not valid: $`\mathrm{log}2/\mathrm{log}311/(3\mathrm{log}3)`$. However, the above formula should give a good approximation of $`\mu `$ for a system with a large dwell time (that is, a small opening).
In ยง3 we provide numerical evidence for the validity of (2.14) in the case of the open 5-baker $`B_5`$ (2.7), at least when taking $`N`$ along geometric subsequences. In ยง4.2 we then construct a related quantum model, for which we can prove this Weyl law and calculate $`C(r)`$ explicitly for inverse Planckโs constants of the form $`N=5^k`$, $`k`$.
## 3. Numerical results
We numerically computed the spectra of several open bakerโs maps; in \[12, ยง5\] we showed the numerical results concerning the $`3`$-baker $`B_3`$ (2.6). For a change, we will discuss here the open $`5`$-baker (2.7), quantized in (2.12). For this open map, the partial dimension of the repeller is $`\mu =\mathrm{log}2/\mathrm{log}5=0.4306765\mathrm{}`$ Compared to the 3-baker, this smaller exponent implies that the spectrum of $`\widehat{B}_5`$ is expected to be much sparser than that of $`\widehat{B}_3`$. For this reason, we will represent the spectra using a logarithmic scale (see Fig. 4), and consider regions $`๐_{r,\vartheta ,\rho }`$ for values of $`r`$ ranging from $`r=0.5`$ down to about $`r=0.001`$.
Let us now briefly explain the choice of โhalf-integer quantizationโ for the quantum positions and momenta . The open map $`B_5`$ is symmetric with respect to the parity transformation $`\mathrm{\Pi }(q,p)=(1q,1p)`$: $`\mathrm{\Pi }B_5=B_5\mathrm{\Pi }`$. The choice of quantization is made so that the associated quantum map $`\widehat{B}_5`$ also possess this symmetry, that is, it commutes with the quantum parity operator $`\widehat{\mathrm{\Pi }}`$ defined as $`\widehat{\mathrm{\Pi }}|q_j=|1q_j=|q_{N1j}`$. We can then separately diagonalize the even and odd parts
$$\widehat{B}_{5,\mathrm{ev}}=\widehat{B}_5(1+\widehat{\mathrm{\Pi }})/2\text{ and }\widehat{B}_{5,\mathrm{odd}}=\widehat{B}_5(1\widehat{\mathrm{\Pi }})/2.$$
Both these operators have rank $`N/5`$: together, they give the full nontrivial spectrum of $`\widehat{B}_5`$. We checked that the odd spectrum has the same characteristics as the even one, so we only describe the properties of the latter. It is expected to satisfy the following fractal law (consequence of (2.14)):
(3.1)
$$n(N,r)\stackrel{\mathrm{def}}{=}\mathrm{\#}\{\mathrm{Spec}(B_{5,N,\mathrm{ev}})๐_r\}C(r)(N/5)^{\mathrm{log}2/\mathrm{log}5},๐_r=\{|\lambda |>r\},N\mathrm{}.$$
The simplest set of $`N`$โs to test this fractal Weyl law is given by geometric sequences of the type $`\{N_o\times 5^k,k=0,1,\mathrm{}\}`$: the law (3.1) means that the number of eigenvalues doubles when $`kk+1`$. In table 1 we give some of the numbers $`n(N,r)`$ along the sequence $`N\{5^k\times 20\}`$, for some selected values of $`r`$.
Along each column with $`r0.01`$, the numbers approximately double at each step $`kk+1`$, which seems to confirm the law (3.1). The fact that this law fails for the small radii $`r=0.005,\mathrm{\hspace{0.17em}0.001}`$ may be explained as follows: according to (3.1), when $`N`$ is large the huge majority of the $`N/5`$ eigenvalues of $`B_{5,N,\mathrm{ev}}`$ must be contained within an asymptotically small neighbourhood of the origin; if $`๐_r`$ intersects this neighbourhood, the law (3.1) necessarily fails, since the counting function is proportional to $`N`$ instead of $`N^\mu `$. For the values of $`N`$ listed in the table, this small region seems to be of radius $`0.005`$, explaining the departure from the fractal law in the last two columns.
To further test the validity of the fractal law (3.1), we choose a set of values of $`r`$, and study the $`N`$-dependence of $`n(N,r)`$, for $`N`$ taken along several geometric sequences, generalizing the above table. In Fig. 3, we plot this dependence in logarithmic scale for $`r=0.3`$ (full lines), $`r=0.1`$ (dashed lines), $`r=0.015`$ (dot-dashed lines). Different geometric sequences are represented by a different colors. For almost all pairs $`(N_o,r)`$ the points are almost aligned, and the slope is in very good agreement with the conjectured one $`\mu =\mathrm{log}2/\mathrm{log}5`$. The less convincing data are the ones related to $`r=0.3`$: for this radius, the numbers $`n(N,r)`$ are still quite small, so that fluctuations are much more visible than for the smaller radii. We expect this effect to disappear for larger values of $`N`$.
In Fig. 3 the height of the curves does not only depend on $`r`$, but also on the sequence $`\{N_o\times 5^k\}`$ considered, especially for $`r=0.3`$.
To investigate this apparent contradiction with (3.1) we plot in Fig. 4 the even spectra of $`B_{5,N}`$ along 3 different geometric sequences. These plots suggest that, along a given geometric sequence, the eigenvalue density increases with $`N`$ uniformly with respect to $`\varphi =\mathrm{arg}\lambda `$, but very nonuniformly with respect to $`|\lambda |`$. We see that some regions $`\{r_0<|\lambda |<r_1\}`$ remain empty even for large values of $`N`$. The presence of gaps was already noticeable when comparing the second and third columns of table 1: obviously, for $`N=20\times 5^k`$, there were no eigenvalues in the annulus $`\{0.05<|\lambda |<0.1\}`$, which is confirmed visually in Fig. 4 (bottom). This non-uniform dependence on $`|\lambda |`$ implies that the profile function $`C(r)`$ is nontrivial.
The spectra for the two other geometric sequences also shows the presence of gaps, but the gaps differ from one geometric sequence to the other. This observation also contradicts the law (3.1).
In spite of these problems we nevertheless attempt to compute the profile function $`C(r)`$ appearing in (3.1). Fig. 5 (left) shows $`n(N,r)`$ as functions of $`r(0,1)`$, for $`N`$ along the same three geometric sequences (each one corresponding to a given color/width). We then rescale the vertical coordinate of each curve by the factor $`(N/5)^{\mathrm{log}2/\mathrm{log}5}`$, and plot the rescaled curves in Fig. 5 (right). From far away, these rescaled curves are fairly superposed on each other, which shows that the conjectured scaling (3.1) is approximately correct. Yet, a closer inspection shows that a much better convergence to a single function occurs along each individual geometric sequence. For instance, the curves for $`N=8\times 5^k`$ โpointwiseโ converge to the last one along this sequence ($`N/5=5000`$), which has a plateau on $`\{0.2r0.4\}`$ corresponding to a spectral gap. The curves of the two other sequences seem to converge as well, with plateaux on different intervals.
In the case of the open kicked rotator studied in the rescaled curves $`n(N,r)`$ are more or less superposed, therefore defining a profile function $`C(r)`$. The authors claim that this function corresponds reasonably well with a prediction of random matrix theory . Our results for the 5-bakerโs map contradict this universality: there does not seem to be a global profile function $`C(r)`$, but a family of such functions, which depend on the geometric sequence $`\{N_o\times 5^k\}`$, which could be denoted by $`C(N_o,r)`$. The law (3.1) needs to be adapted by restricting $`N`$ to geometric sequences, which yields the following empirical scaling law:
For any $`N_o>0`$ and $`r(0,1)`$, there exists $`C(N_o,r)0`$ such that, for $`N`$ along $`\{N_o\times 5^k\}`$ and $`k\mathrm{}`$,
(3.2)
$$\mathrm{\#}\left\{\lambda \mathrm{Spec}(B_N)๐_{r,\vartheta ,\rho }\right\}\frac{\rho }{2\pi }C(N_o,r)N^\mu ,$$
where $`๐_{r,\vartheta ,\rho }`$ is given by (2.13). In Fig. 5 the different profile functions are uniformly bounded, $`C_1(r)C(N_o,r)C_2(r)`$, for some envelope functions $`0C_1(r)C_2(r)`$.
This weakening of (2.14) to geometric sequences makes sense for bakerโs maps of the form $`B_3`$, $`B_5`$, which each have a uniform integer expansion factor, leading to number-theoretic properties. In the case of a nonlinear open chaotic map (as the open kicked rotator of ), there is no reason for geometric sequences to play any role, so we expect (2.14) to hold in that case.
## 4. A computable model
Because we are unable to analyze the spectra of the quantum bakers $`B_N`$ rigorously, we introduce simplified models. In the case of the 3-baker, we observe (see Fig. 6, left) that the largest matrix elements are maximal along the โtilted diagonalsโ
(4.1)
$$(n,m)=(3l+ฯต,l+\mathrm{}N/3),\text{with}l\{0,\mathrm{},N/31\},\mathrm{}\{0,2\},ฯต\{0,1,2\}.$$
These โdiagonalsโcorrespond to a discretization of the the map $`B_3`$ projected on the position axis. Away from them, the coefficients do not decrease very fast due to the Gibbs phenomenon (diffraction). The elements on the โdiagonalsโ have moduli $`1/\sqrt{3}+๐ช(1/N)`$ and their phases only depend on $`\mathrm{},ฯต`$ in the above parametrization. Our simplified model is obtained by keeping only the elements on the โdiagonalsโ (see Fig. 6, right), set their moduli to $`1/\sqrt{3}`$ and shift their phases by $`\pi /2`$ (for convenience). Using the parametrization (4.1), we get:
(4.2)
$$(\stackrel{~}{B}_{3,N})_{nm}=\frac{1}{\sqrt{3}}\mathrm{exp}\left(\frac{2i\pi }{3}(ฯต+1/2)(\mathrm{}+1/2)\right).$$
For $`N=9`$ and using $`\omega =e^{2\pi i/3}`$, the matrix reads
$$\stackrel{~}{B}_{3,9}=\frac{\omega ^{1/4}}{\sqrt{3}}\left(\begin{array}{ccccccccc}1\hfill & 0\hfill & 0\hfill & 0\hfill & 0\hfill & 0\hfill & \omega \hfill & 0\hfill & 0\hfill \\ \omega ^{1/2}\hfill & 0\hfill & 0\hfill & 0\hfill & 0\hfill & 0\hfill & \omega ^{1/2}\hfill & 0\hfill & 0\hfill \\ \omega \hfill & 0\hfill & 0\hfill & 0\hfill & 0\hfill & 0\hfill & 1\hfill & 0\hfill & 0\hfill \\ 0\hfill & 1\hfill & 0\hfill & 0\hfill & 0\hfill & 0\hfill & 0\hfill & \omega \hfill & 0\hfill \\ 0\hfill & \omega ^{1/2}\hfill & 0\hfill & 0\hfill & 0\hfill & 0\hfill & 0\hfill & \omega ^{1/2}\hfill & 0\hfill \\ 0\hfill & \omega \hfill & 0\hfill & 0\hfill & 0\hfill & 0\hfill & 0\hfill & 1\hfill & 0\hfill \\ 0\hfill & 0\hfill & 1\hfill & 0\hfill & 0\hfill & 0\hfill & 0\hfill & 0\hfill & \omega \hfill \\ 0\hfill & 0\hfill & \omega ^{1/2}\hfill & 0\hfill & 0\hfill & 0\hfill & 0\hfill & 0\hfill & \omega ^{1/2}\hfill \\ 0\hfill & 0\hfill & \omega \hfill & 0\hfill & 0\hfill & 0\hfill & 0\hfill & 0\hfill & 1\hfill \end{array}\right).$$
The matrix $`\stackrel{~}{B}_{3,N}`$ can obviously not be considered as a โsmall perturbationโ of $`B_{3,N}`$, since we removed many nonnegligible โoff-diagonalโ elements. Actually, by acting with $`\stackrel{~}{B}_{3,N}`$ on Gaussian coherent states, one realizes that these matrices do not quantize the open 3-baker $`B_3`$ (2.6), but rather a more complicated multivalued map $`\stackrel{~}{B}_3`$, built upon $`B_3`$ as follows:
(4.3)
$$(q,p)๐_3,\stackrel{~}{B}_3(q,p)=\underset{j=1}{\overset{1}{}}\{B_3(q,p)+(0,j/3)\}.$$
We refer to \[12, Proposition 6.1\] for a precise statement. As opposed to $`B_3`$, the multivalued map (4.3) is no longer obtained by truncating a canonical transformation, but it comes from three different transformations. $`\stackrel{~}{B}_3`$ can be considered as a model of propagation with ray splitting. Another interpretation is given by considering a Markov process with probabilities $`P(x^{},x)`$ being allocated at each step to the image points $`x^{}`$ of $`x`$. Explicitly, the probabilities take the form
$$P(x^{},x)=f\left(\frac{3p^{}[3q]1/2}{3}\right),f(t)=\left(\frac{\mathrm{sin}(3\pi t)}{3\mathrm{sin}(\pi t)}\right)^2,x=(q,p),x^{}=(q^{},p^{}),$$
so that for each $`x๐_3`$, the sum of the weights associated with the three images of $`x`$ is indeed $`1`$.
Some of the characteristics of the dynamics remain the same as for $`B_3`$. The local dynamics of each branch is the same, and $`\stackrel{~}{B}_3`$ sends all points in $`๐^2๐_3`$ to infinity. One can define incoming and outgoing tails for $`\stackrel{~}{B}_3`$ (see ยง2). As opposed to the case of $`B_3`$, these tails are not symmetrical any more: $`\mathrm{\Gamma }_{}=๐_3\times [0,1)`$, $`\mathrm{\Gamma }_+=๐^2`$. Yet, these formulas are slightly misleading. The second one comes from the fact that any point $`x๐^2`$ has two preimages through $`\stackrel{~}{B}_3`$, namely $`x_0=(q/3,3p),x_2=((q+2)/3,3p)`$, so no point ever escapes to infinity in the past. However, to these preimages are associated the respective weights $`P(x,x_0)`$, $`P(x,x_2)`$, the sum of which is generally $`<1`$: there is thus a loss of probability through $`\stackrel{~}{B}_3^1`$, which is not accounted for by the definition of $`\mathrm{\Gamma }_+`$.
In the next section we will show that the matrices $`\stackrel{~}{B}_{3,N}`$ can nonetheless be interpreted as quantizations of the original open baker $`B_3`$, as long as one switches to a different notion of quantization, derived from a different type of Fourier transform (the Walsh-Fourier transform).
Families of unitary matrices $`\stackrel{~}{A}_{2,N}`$ with a structure similar to $`\stackrel{~}{B}_{3,N}`$ have already been proposed as an alternative quantization of the 2-bakerโs map $`A_2`$ . These matrices can also closely related with the โsemiquantum bakersโ introduced in M. Saraceno, private communication.. In the context of quantum graphs (a recently popular model for quantum chaos), unitary matrices similar with $`\stackrel{~}{A}_{2,N}`$ (but with random phases) occur as โunitary transfer matricesโ associated with binary graphs . In this framework, the matrix $`\stackrel{~}{A}_{2,N}`$ would correspond to a graph with very degenerate bond lengths. In this framework, the matrix $`\stackrel{~}{B}_{3,N}`$ is directly related with a classical transfer matrix defined by $`(_{3,N})_{jj^{}}=\left|(\stackrel{~}{B}_{3,N})_{jj^{}}\right|^2`$, which represents the classical Markov process on the graph. In our case, this transfer matrix is the discretized version of the transfer (Perron-Frobenius) operator associated with the open map $`B_3`$.
### 4.1. The Walsh model interpretation of $`\stackrel{~}{B}_{3,N}`$.
In this section, we represent the matrices $`\stackrel{~}{B}_{3,N}`$ in a way suitable for their spectral analysis. This can be done only in the case where $`N`$ is a power of $`3`$. This representation is connected with the Walsh model of harmonic analysis.
The latter originally appeared in the context of fast signal processing . The major advantage of Walsh harmonic analysis (compared with the usual Fourier analysis) is the possibility to strictly localize wave packets simultaneously in position and in momentum. For our problem, this has the effect to remove the diffraction problems due to the discontinuities of the classical map, which spoil the usual semiclassics .
A recent preprint analyzes some special eigenstates of the โstandardโ quantum 2-baker, using the Walsh-Hadamard transform (which slightly differs from the Walsh transform we give below) as a โfilterโ. We are doing something different here by constructing our simplified model $`\stackrel{~}{B}_{3,N}`$ from the Walsh transform, as $`B_{3,N}`$ was constructed from the discrete Fourier transform (see ยง2).
We first select the expanding coefficient of the bakerโs map, which we denote by $`D`$ (the map (2.6) is associated with $`D=3`$, the map (2.7) with $`D=5`$). Once this is done, we will restrict ourselves to the values of $`N`$ along the geometric sequence $`\{N=D^k,k\}`$. In this case, the Hilbert space can be naturally decomposed as a tensor product of $`k`$ spaces $`^D`$:
(4.4)
$$_N=(^D)_1(^D)_2\mathrm{}(^D)_k.$$
This decomposition appears naturally in the context of quantum computation, where each $`^D`$ represents a โquantum $`D`$-gitโ, that is, a quantum system with $`D`$ levels. Here, we realize this decomposition using the basis of position eigenstates $`|q_j`$ of $`_N`$ (see for the case $`D=2`$). Indeed, each quantum position $`q_j=(j+1/2)/N`$, $`j_{D^k}=\{0,\mathrm{},N1\}`$ is in one-to-one correspondence with a word $`\mathit{ฯต}=ฯต_1ฯต_2\mathrm{}ฯต_k`$ made of symbols ($`D`$-gits) $`ฯต_{\mathrm{}}=ฯต_{\mathrm{}}(j)_D`$:
(4.5)
$$j=\underset{\mathrm{}=1}{\overset{k}{}}ฯต_{\mathrm{}}D^k\mathrm{}.$$
The usual order for $`j_{D^k}`$ corresponds to the lexicographic order for the symbolic words $`\mathit{ฯต}(_D)^k`$. Associating to each $`D`$-git a $`D`$-dimensional vector space $`(^D)_{\mathrm{}}`$ with canonical basis $`\{e_0,e_1,\mathrm{},e_{D1}\}`$, the position eigenstate $`|q_j_N`$ can be decomposed as
(4.6)
$$|q_j=e_{ฯต_1}e_{ฯต_2}\mathrm{}e_{ฯต_k}.$$
This identification realizes the tensor product decomposition (4.4).
The Fourier transforms $`๐ข_N`$ (2.9), and the simpler one without the $`1/2`$ shift,
(4.7)
$$(_N)_{jj^{}}=\frac{e^{2i\pi jj^{}/N}}{\sqrt{N}},j,j^{}_N,N=D^k,$$
are defined by applying the exponential function $`xe^{2i\pi x}`$ to the products
$$\frac{jj^{}}{D^k}=\underset{m=2k}{\overset{k}{}}D^m\stackrel{~}{ฯต}_m(jj^{}),\text{where}\stackrel{~}{ฯต}_m(jj^{})=\underset{\mathrm{}+\mathrm{}^{}=m+k}{}ฯต_{\mathrm{}}(j)ฯต_{\mathrm{}^{}}(j^{}).$$
If we replace in (4.7) the exponential $`e^{2i\pi x}`$ by the piecewise constant function $`e_D(x)=\mathrm{exp}(2i\pi [Dx]/D)`$, and replace each $`\stackrel{~}{ฯต}_m(jj^{})`$ by its value $`ฯต_m(jj^{})`$ modulo $`D`$, we obtain the matrix element
(4.8)
$$\begin{array}{cc}\hfill (๐ฑ_k)_{jj^{}}& \stackrel{\mathrm{def}}{=}D^{k/2}e_D\left(\underset{m=2k}{\overset{k}{}}D^mฯต_m(jj^{})\right)=D^{k/2}\mathrm{exp}\left(\frac{2i\pi }{D}ฯต_1(jj^{})\right)\hfill \\ & =\underset{\mathrm{}=1}{\overset{k}{}}D^{1/2}\mathrm{exp}\left(\frac{2i\pi }{D}ฯต_{\mathrm{}}(j)ฯต_{k+1\mathrm{}}(j^{})\right).\hfill \end{array}$$
The matrix $`๐ฑ_k`$ defines the Walsh transform in dimension $`D^k`$.
Because we have used the โhalf-integerโ Fourier transform (2.9) to define our quantum bakerโs map, we will need a slightly different version of Walsh transform, namely
$$(๐ฒ_k)_{jj^{}}\stackrel{\mathrm{def}}{=}\underset{\mathrm{}=1}{\overset{k}{}}D^{1/2}\mathrm{exp}\left(\frac{2i\pi }{D}(ฯต_{\mathrm{}}(j)+1/2)(ฯต_{k+1\mathrm{}}(j^{})+1/2)\right).$$
Both $`๐ฑ_k`$ and $`๐ฒ_k`$ preserve the tensor product structure (4.4): for any $`v_1,\mathrm{},v_n^D`$,
(4.9)
$$๐ฑ_k(v_1\mathrm{}v_k)=_Dv_k\mathrm{}_Dv_1,๐ฒ_k(v_1\mathrm{}v_k)=๐ข_Dv_k\mathrm{}๐ข_Dv_1.$$
These expressions show that $`๐ฑ_k`$ and $`๐ฒ_k`$ are unitary.
Specializing the computations to $`D=3`$, we are now in position to define the toy model $`\stackrel{~}{B}_{3,N}`$ (in the case $`N=3^k`$) as the โWalsh-quantizationโ of the 3-baker (2.6) (as opposed to the โstandardโ quantization of the multivalued map $`\stackrel{~}{B}_3`$ (4.3)). Indeed, one can check that the matrix (4.2) can be expressed as
(4.10)
$$\stackrel{~}{B}_{3,N}=๐ฒ_k^{}\left(\begin{array}{ccc}๐ฒ_{k1}\hfill & 0\hfill & 0\hfill \\ 0\hfill & 0\hfill & 0\hfill \\ 0\hfill & 0\hfill & ๐ฒ_{k1}\hfill \end{array}\right).$$
This formula is clearly the Walsh analogue of the definition (2.11) of the โstandardโ quantum open baker $`B_{3,N}`$. From this definition and (4.9), we see the action of $`B_{3,N}`$ on tensor products:
(4.11)
$$\stackrel{~}{B}_{3,N}(v_1\mathrm{}v_k)=v_2\mathrm{}v_k๐ข_3^{}\pi _{0,2}v_1,v_j^3,$$
where $`\pi _{0,2}`$ is the orthogonal projector (in $`^3`$) on $`e_0e_2`$.
### 4.2. Distribution of resonances
Using (4.11) we can explicitly describe the spectrum of $`\stackrel{~}{B}_{3,N}`$ for $`N=3^k`$. The computation is identical with \[12, Section 6.2\], so we only give the results. The generalized kernel of $`\stackrel{~}{B}_{3,N}`$ is spanned by the position states $`|q_j`$ such that $`ฯต_{\mathrm{}}(j)=1`$ for at least one index $`1\mathrm{}k`$. This corresponds to positions $`q_j`$ โfarโ from the Cantor set $`๐_3`$, so that the classical points $`(q_j,p)`$ are sent to infinity at a time $`nk`$. This kernel has dimension $`3^k2^k=NN^{\mathrm{log}2/\mathrm{log}3}`$.
The nonzero eigenvalues of $`\stackrel{~}{B}_{3,N}`$ are given by the set (see Fig. 7)
$$\{\lambda _+\}\{\lambda _{}\}\underset{\mathrm{}=0}{\overset{k1}{}}\underset{p=1}{\overset{k1}{}}\{e^{2i\pi \mathrm{}/k}\lambda _+^{1p/k}\lambda _{}^{p/k}\},\text{where}\lambda _+=1,\lambda _{}=\frac{i}{\sqrt{3}}.$$
For each $`p\{1,\mathrm{},k1\}`$, the $`k`$ eigenvalues of modulus $`|\lambda _{}|^{p/k}=3^{p/2k}`$ asymptotically have the same degeneracy $`\left(\genfrac{}{}{0pt}{}{k}{p}\right)/k`$ as $`k\mathrm{}`$ (semiclassical limit), which shows that their distribution is circular-symmetric. Taking these multiplicities into account, we obtain the following Weyl law for the eigenvalues of $`\stackrel{~}{B}_{3,N}`$ inside a region (2.13), along the sequence $`N\{3^k\}`$, $`k\mathrm{}`$:
(4.14)
$$\begin{array}{c}\mathrm{\#}\left\{\mathrm{Spec}(\stackrel{~}{B}_{3,N})๐_{r,\vartheta ,\rho }\right\}=\frac{\rho }{2\pi }N^\mu (C(1,r)+o(1))\\ \mu =dim(\mathrm{\Gamma }_{}W_+)=\frac{\mathrm{log}2}{\mathrm{log}3},C(1,r)=1\mathrm{l}_{(0,3^{1/4}]}(r).\end{array}$$
The values $`\lambda _{}`$, $`\lambda _+`$ in (4.14) are the nonzero eigenvalues of the matrix $`๐ข_3^{}\pi _{0,2}`$ appearing in (4.11). We used the notation $`C(1,r)`$ for the profile function to be consistent with our notations in (3.2), that is, to emphasize that this estimate is valid only along the sequence $`N\left\{\mathrm{\hspace{0.17em}1}\times 3^k\right\}`$.
We notice that the spectrum of the classical transfer matrix $`_{3,N}`$ defined at the end of ยง4 is drastically different: this matrix admits one simple nontrivial eigenvalue $`\lambda =2/3`$ (interpreted as the classical escape rate), the rest of the spectrum lying in the generalized kernel. Therefore, the features of the quantum spectrum is intimately related with the oscillatory phases of $`\stackrel{~}{B}_{3,N}`$ along the โdiagonalsโ.
### 4.3. Conductance and Shot Noise
In this section, we consider an open bakerโs map as a model of quantum transport through a โchaotic quantum dotโ, that is a 2-dimensional cavity connected to the outside world through a certain number of โleadsโ carrying the current; each lead is connected to the cavity along a segment $`L_j`$ of the boundary (see Fig. 8), and the connection is assumed to be โperfectโ: a particule inside the cavity which hits the boundary along $`qL_j`$ is completely evacuated to the lead. Therefore, the phase space domain $`L_j\times [0,1)`$ above this segment is a part of the โholeโ, in the terminology of ยง2, whereas the remaining set $`I=[0,1)(L_j)`$ represents the boundary of the quantum dot, which lifts to the phase space domain $`๐=I\times [0,1)`$.
In the previous sections we have studied the open quantum map obtained by projecting a unitary quantum dynamics (called generically $`U_N`$) onto a subdomain $`๐`$ of the phase space: resonances were defined as the eigenvalues of $`U_N\mathrm{\Pi }_๐`$. These resonances are supposed to represent the metastable quantum states inside the open quantum dot, after it has been opened. In the present section, we want to study another aspect of the open system, namely the โtransportโ through the dot, using the formalism of . We will focus on the case where the opening $`L`$ splits into two segments $`L=L_1L_2`$, and we study the transmission matrix from the lead $`L_1`$ to the lead $`L_2`$.
Once we are given, on one side, the quantum map $`U_N`$ associated with the closed dynamics inside the โcavityโ, on the other side, the projectors on the leads $`\mathrm{\Pi }_{L_i}`$ and on the โinteriorโ $`\mathrm{\Pi }_I=\mathrm{\Pi }_๐`$, the transmission matrix (from $`L_1`$ to $`L_2`$) is defined as the block
(4.15)
$$t(\vartheta )=\underset{n1}{}e^{in\vartheta }\mathrm{\Pi }_{L_2}U_N(\mathrm{\Pi }_IU_N)^{n1}\mathrm{\Pi }_{L_1}.$$
The parameter $`\vartheta [0,2\pi )`$ is the โquasi-energyโ of the particles. According to Landauerโs theory of coherent transport, each eigenvalue $`T_i(\vartheta )`$ of the matrix $`t(\vartheta )t^{}(\vartheta )`$ corresponds to a โtransmission channelโ. The dimensionless conductance of the system is then given by
(4.16)
$$g(\vartheta )=\mathrm{tr}\left(t(\vartheta )t^{}(\vartheta )\right).$$
A transmission channel is โclassicalโ if the eigenvalue $`T_i`$ is very close to unity (perfect transmission) or close to zero (perfect reflection). The intermediate values characterize โnonclassical channelsโ (governed by strong interference effects). The number number of the latter can be estimated by the noise power
(4.17)
$$P(\vartheta )=\mathrm{tr}\left(t(\vartheta )t^{}(\vartheta )\left(Idt(\vartheta )t^{}(\vartheta )\right)\right),$$
or equivalently the Fano factor, $`F=P/g`$. It is sometimes necessary to perform an ensemble averaging over $`\vartheta `$ to obtain significant results . However, for the model we study here, these quantities will depend very little on $`\vartheta `$.
The closed quantum dot will be modeled by the following quantum map: we consider the 4-bakerโs map $`A_4`$ and quantize it using the Walsh transform $`๐ฑ_k`$ (4.8) with $`D=4`$. In dimension $`N=4^k`$, our unitary propagator is therefore
$$U_N=\stackrel{~}{A}_{4,N}=๐ฑ_k^{}\left(\begin{array}{cccc}๐ฑ_{k1}\hfill & 0\hfill & 0\hfill & 0\hfill \\ 0\hfill & ๐ฑ_{k1}\hfill & 0\hfill & 0\hfill \\ 0\hfill & 0\hfill & ๐ฑ_{k1}\hfill & 0\hfill \\ 0\hfill & 0\hfill & 0\hfill & ๐ฑ_{k1}\hfill \end{array}\right).$$
We attach the leads on the intervals $`L_1=[0,1/4]`$ and $`L_2=[3/4,1]`$: this way, the projectors $`\mathrm{\Pi }_{L_i}`$ as well as the projector $`\mathrm{\Pi }_I=Id\mathrm{\Pi }_{L_1}\mathrm{\Pi }_{L_2}`$ can be represented as tensor product operators:
$$\mathrm{\Pi }_{L_1}=\pi _0Id_4Id_4\mathrm{},\mathrm{\Pi }_{L_2}=\pi _3Id_4\mathrm{},\mathrm{\Pi }_I=\pi _IId_4\mathrm{}.$$
Here $`\pi _i`$ is the orthogonal projector on the basis state $`e_i`$ of $`^4`$, and $`\pi _I=\pi _1\pi _2`$. This tensor action, together with the action of $`\stackrel{~}{A}_{4,N}`$ (analogous to (4.11)), allow us to compute all quantities explicitly.
The spectrum of the โinsideโ propagator for this model, $`\stackrel{~}{B}_{4,N}=\stackrel{~}{A}_{4,N}\mathrm{\Pi }_I`$, satisfies a fractal Weyl law of the type (4.14) along the sequence $`N=4^k`$, with exponent $`\mu =\mathrm{log}2/\mathrm{log}4=1/2`$, and profile $`C(1,r)=1\mathrm{l}_{[0,2^{3/4}]}(r)`$.
For this model and in the semiclassical limit $`k\mathrm{}`$, we could compute the dimensionless conductance (4.16). The computation \[12, ยง7.2\] requires to control the time evolution up to $`n=Ck`$ for some $`1<C<2`$: this is of the order of the Ehrenfest time $`\tau _E=k`$ for the system. For any $`\vartheta `$ we obtain
(4.18)
$$g(N=4^k,\vartheta )=\frac{4^{k1}}{2}(1+o(2^k))=\frac{N/4}{2}(1+o(1)),k\mathrm{}.$$
Here $`N/4`$ is the number of transmission channels from $`L_1`$ to $`L_2`$, that is the rank of the matrix $`t(\vartheta )`$. We see that, as could be expected, approximately one half of the scattering channels get transmitted from one lead to the other, the other half being reflected back.
Asymptotics for the shot noise (4.17) (which counts the โnonclassicalโ transmission channelsโ) are more interesting and again independent of $`\vartheta `$:
(4.19)
$$P(N=4^k,\vartheta )=2^{k1}\left(\frac{11}{80}+๐ช(e^{Ck})\right)=\frac{11}{80}(N/4)^\mu (1+o(1)),k\mathrm{}.$$
Here $`\mu =1/2`$ is the dimension appearing in the fractal Weyl law for the resonances. A similar fractal law for the shot noise had been observed in in the case of the quantum kicked rotator; the power law $`N^\mu `$ for the number of nonclassical channels was explained there through a study of the dynamics up to the Ehrenfest time.
The constant $`11/80`$ in (4.19) gives the average โshot noiseโ per nonclassical transmission channel. This number is close to the random matrix theory prediction for this quantity, namely $`1/8`$ . The precise number $`11/80`$ certainly depends on which bakerโs map one starts from, and which quantization one uses. For instance, we did not check whether the โhalf-integerโ Walsh quantization of the 4-baker leads to the same prefactor, but we expect the result to be close to it. It would be interesting to actually check the full distribution for the transmission eigenvalues $`T_i`$, and compare it with the prediction of random matrix theory .
The near agreement with random matrix theory is in contrast with the fact that the semiclassical resonance spectrum of the propagator $`\stackrel{~}{B}_{4,N}`$ inside the dot is very different from that of a random subunitary matrix. Somehow, the matrix $`t(\vartheta )`$, obtained by summing iterates of $`\stackrel{~}{B}_{4,N}`$, has acquired some โrandomnessโ, as far as the distribution of its singular values is concerned.
The transport properties of chaotic cavities has also been studied within the framework of quantum graphs. The shot noise (4.17) could be semiclassically estimated in the case of a โstar graphโ, by summing over transmitting trajectories on the graph (they studied the case of โsmall openingsโ). The authors show that one needs to take into account the โaction correlationsโ between different trajectories, in order to reproduce the random matrix result. As mentioned before, the matrix $`\stackrel{~}{A}_{4,N}`$ can be interpreted as the unitary transfer matrix for a different type of graph , with bonds having degenerate lengths. Somehow, our use of the tensor product structure implicitly takes into account the action correlations for this particular graph.
## 5. Conclusions
Quantum open bakerโs maps provide a simple and elegant model for the study of quantum resonances of open chaotic systems. The numerical investigation of these models is easily accessible and, as shown in ยง3, gives a good agreement with the fractal Weyl law on โsmall energy scalesโ, which is (1.7) in the case of Hamiltonian flows. Only larger energy scales (1.8) were considered previously. It would be interesting to investigate the spectrum of the model operator (2.11) for higher values of $`N\mathrm{}^1`$. The naรฏve numerical approach we took (full diagonalization of the matrices $`B_N`$) only allowed to reach values $`N5000`$. It would make more sense to use an algorithm allowing us to extract only the largest eigenvalues (which are the ones we are interested in), instead of the full spectrum.
By modifying the standard quantum bakerโs map, in a way which still fits in the framework of quantization of chaotic dynamics, we obtained a model for which the fractal Weyl law (4.14) can be rigorously proven. Since the spectrum of this model is explicitely computable (and forms a lattice), it is forcibly nongeneric. However, the explicit computation of other physical quantities associated with our model, namely the conductance and the โshot noiseโ, shows more generic properties. The fractal Weyl law is also present in the calculation of the โshot noiseโ, and the prefactor is (unexpectedly) close to random matrix predictions.
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# Examples of bosonic de Finetti states over finite dimensional Hilbert spaces
## 1 Introduction
According to the Quantum de Finetti Theorem , locally normal infinite-particle states with Bose-Einstein symmetry can be represented as mixtures of infinite tensor powers of vector states. This note presents examples of infinite-particle states with Bose-Einstein symmetry that arise as limits of Gibbs ensembles on finite dimensional spaces, and displays their de Finetti representations.
The central example is as follows. If the single-particle Hilbert space $``$ is finite dimensional, the projector onto the symmetric subspace of the $`n`$-particle space can be normalized, and this defines the infinite-temperature ensemble for $`n`$ bosons with single-particle space $``$. For each fixed $`m`$, the $`m`$-particle reduced density operators under the $`n`$-boson infinite-temperature ensembles converge, as $`n`$ tends to infinity, to the density operator describing the $`m`$-particle statistics under a certain bosonic infinite-particle state $`\omega _0`$. The infinite-particle state $`\omega _0`$ has a de Finetti representation as a mixture of infinite tensor powers of vector states $`P_\mathrm{v}`$, where $`\mathrm{v}`$ is a unit vector and $`P_\mathrm{v}=|\mathrm{v}\mathrm{v}|`$ denotes the projector onto the span of $`\mathrm{v}`$. In the de Finetti mixture for $`\omega _0`$, the weight of the tensor power state $`P_{\mathrm{v}}^{}{}_{}{}^{\mathrm{}}`$ is the probability density for $`\mathrm{v}\stackrel{}{=}^{d+1}`$ to equal
$$\mathrm{v}(p,\theta )=(e^{i\theta _0}\sqrt{p_0},e^{i\theta _1}\sqrt{p_1},\mathrm{},e^{i\theta _d}\sqrt{p_d})$$
(1)
when $`p=(p_0,p_1,\mathrm{},p_d)`$ is sampled uniformly from the $`d`$-dimensional simplex $`\mathrm{\Delta }_d`$ and the phase angles $`\theta _i`$ in $`\theta =(\theta _0,\theta _1,\mathrm{},\theta _d)`$ are each sampled uniformly from $`[0,2\pi )`$, independently of one another and of $`p`$. Thus the infinite-particle state $`\omega _0`$ corresponds to the uniform probability measure on $`\mathrm{\Delta }_d\times [0,2\pi )^{d+1}`$.
Similar limits are obtained for finite temperature Gibbs ensembles, provided the temperature is scaled properly. Suppose $`H`$ is a Hermitian operator on the single-particle space $`\stackrel{}{=}^{d+1}`$ and $`\mathrm{\Gamma }_n(\beta )`$ denotes the Gibbs canonical ensemble for $`n`$ noninteracting bosons with single-particle Hamiltonian $`H`$ at inverse temperature $`\beta `$. Then, as $`n`$ tends to infinity, the $`m`$-particle reduced density operators under $`\mathrm{\Gamma }_n(\beta /n)`$ converge to the $`m`$-particle density of a certain bosonic infinite-particle state $`\omega _\beta `$. The infinite-particle state $`\omega _\beta `$ is an average of states $`P_{\mathrm{v}}^{}{}_{}{}^{\mathrm{}}`$ with respect to the probability density on $`\mathrm{\Delta }_d\times [0,2\pi )^{d+1}`$ that minimizes the โfree energyโ
$$\underset{[0,2\pi )^{d+1}}{}_{\mathrm{\Delta }_d}\mathrm{v},H\mathrm{v}f(p,\theta )๐p๐\theta +\frac{1}{\beta }\underset{[0,2\pi )^{d+1}}{}_{\mathrm{\Delta }_d}f(p,\theta )\mathrm{ln}f(p,\theta )๐p๐\theta ,$$
where $`\mathrm{v}=\mathrm{v}(p,\theta )`$ is as in (1). We obtain similar results for bosons with โmean fieldโ interactions, but again we must scale temperature in proportion to the number of particles. This stands in contrast to the analogous mean field limits for distinguishable particles, which are obtained without any peculiar scaling of temperature .
The physical relevance of these facts is limited. On the one hand, they concern limits of canonical ensembles, which are appropriate when the number of bosons is fixed, and therefore not appropriate for massless bosons (e.g., photons). On the other hand, massive bosons inhabit infinite dimensional Hilbert spaces, so to speak, whereas our results concern finite dimensional Hilbert spaces. However, the sort of ensemble we study is appropriate for (noninteracting) systems of $`n`$ material bosons in thermal equilibrium, in case it is known that every one of these bosons is trapped in a potential well of depth $`E`$. The statistical state of that system would be a conditional Gibbs ensemble, supported on the finite dimensional Hilbert space spanned by the symmetrized products of trapped (bound) states. Only noninteracting systems of trapped bosons are considered, because the conditional Gibbs ensemble only makes sense if the Hamiltonian of the system commutes with the observable that every particle is trapped.
Our results are presented in Section 3, after a quick review of the Quantum de Finetti Theorem in the next section.
## 2 The Quantum de Finetti Theorem
Let $``$ be Hilbert space (which we will call the single-particle Hilbert space) and let $`^n`$ denote the $`n`$-fold tensor power of $``$ (the $`n`$-particle Hilbert space). When $`\pi `$ denotes a permutation of $`\{1,2,\mathrm{},n\}`$, let $`U_\pi `$ denote the unitary โpermutationโ operator on $`^n`$ defined by
$$U_\pi (x_1x_2\mathrm{}x_n)=x_{\pi (1)}x_{\pi (2)}\mathrm{}x_{\pi (n)}.$$
For each $`n`$ let $`D_n`$ be a density operator on the $`n`$-particle Hilbert space $`^n`$, the $`n`$-fold tensor power of $``$. We want the density operators $`D_n`$ to be symmetric, and we assume
* (A) for all $`n`$, the density operator $`D_n`$ commutes with any permutation operator $`U_\pi `$ on $`^n`$ .
We are especially interested here in systems of bosons, for which
* (B) for all $`n`$, $`D_nU_\pi =D_n`$ for any permutation operator $`U_\pi `$ on $`^n`$.
Condition (B) is stronger than (A). We also want the sequence $`\{D_n\}`$ of density operators to be consistent with respect to โsubsamplingโ in the sense that
* (C) for all $`m<n`$, $`D_{n:m}=D_m`$,
where $`D_{n:m}`$ denotes the $`m^{th}`$ partial trace of $`D_n`$, i.e., the operator such that
$$\mathrm{Tr}(D_{n:m}A)=\mathrm{Tr}(D_n(AI\stackrel{nmtimes}{\mathrm{}}I))$$
for all $`A(^m)`$.
The structure of sequences $`\{D_n\}`$ of density operators satisfying (C) and (A) or (B) is given by the quantum analogue of the de Finetti Theorem of probability theory . Let $`\rho `$ be a density operator on $``$. A sequence $`\{D_n\}`$ of density operators of the form
$`D_1`$ $`=`$ $`\rho `$
$`D_2`$ $`=`$ $`\rho \rho `$
$`D_3`$ $`=`$ $`\rho \rho \rho ,\text{et cetera}`$ (2)
always satisfies (A) and (C), but it satisfies (B) and (C) if and only if $`\rho `$ is a pure state, i.e., a rank one projector on $``$. Roughly speaking, any sequence of density operators satisfying (A) and (C) is uniquely representable as a mixture of sequences of the form (2). That is, if $`\{D_n\}`$ satisfies (A) and (C) then there exists a unique probability measure $`\mu `$ supported on the single-particle density operators such that
$$D_n=\rho ^n\mu (d\rho )$$
(3)
for all $`n`$. Furthermore, if $`\{D_n\}`$ satisfies (B) and (C), then the measure $`\mu (d\rho )`$ in the integral representation (3) is even supported on the set of vector states $`\rho =P_\psi `$. This paraphrases the propositions of , ignoring the technical details; we now restate the results with more care.
For $`mn`$, let $`j_{mn}`$ denote the \*-isomorphic embedding
$$j_{mn}(B)=BI^{nm}$$
of $`(^m)`$ into $`(^n)`$. The system of C\* algebras $`(^n)`$ and isomorphic injections $`j_{mn}`$ has an inductive limit $`๐`$. The inductive or direct limit in the category of C\* algebras may be constructed as in \[3, Proposition 11.4.1\]. The inductive limit $`๐`$ is unique up to isomorphism, and for each $`n`$ there is a \*-isomorphism $`i_n`$ from $`(^n)`$ into $`๐`$ such that $`i_nj_{mn}=i_m`$ for all $`mn`$ and the union of the images $`i_n((^n))`$ is dense in $`๐`$. A sequence $`\{D_n\}`$ of density operators satisfying the conditions (C) can be used to define a continuous positive linear functional $`\omega `$ on $`๐`$ by
$$\omega (i_n(B))=\mathrm{Tr}(D_nB)B(^n).$$
(4)
This is well-defined thanks to the consistency conditions (C) and the density of $`i_n((^n))`$ in $`๐`$. In particular, $`\omega (e)=1`$, where $`e`$ is the identity element of the C\* algebra $`๐`$. If $`\{D_n\}`$ satisfies (A) as well as (C) then $`\omega `$ is symmetric in the sense that
$$\omega (i_n(U_\pi BU_\pi ^{}))=\omega (i_n(B))$$
(5)
for all $`n`$, all $`B(^n)`$, and all $`\pi \mathrm{\Pi }_n`$, the set of permutations of $`\{1,2,\mathrm{},n\}`$. The set of all โsymmetric statesโ on $`๐`$, i.e., the set
$$\mathrm{SS}=\left\{\omega ๐^{}\right|\omega (e)=1\text{and}\omega (x^{}x)0x๐\text{and}\omega \text{satisfies}(\text{5})\},$$
is a convex subset of the Banach dual $`๐^{}`$ of $`๐`$, and it is compact with respect to the weak\* topology. Let $`\mathrm{SS}_1`$ denote the space of single-particle states, i.e., the set
$$\mathrm{SS}_1=\left\{\rho ()^{}\right|\rho (I)=1\text{and}\omega (A^{}A)0A()\}$$
endowed with the relative weak\* topology it inherits as a subset of the Banach dual $`()^{}`$ of $`()`$ . It was first shown in that each $`\omega \mathrm{SS}`$ has a unique representation as an integral of product states
$$\omega =_{\mathrm{SS}_1}\rho \rho \rho \mathrm{}\mathrm{}\mu (d\rho )=_{\mathrm{SS}_1}\rho ^{\mathrm{}}\mu (d\rho ),$$
(6)
where $`\mu `$ is a probability measure on the $`\sigma `$-algebra $`_1`$ generated by the intersections with $`\mathrm{SS}_1`$ of weak\* open sets in $`()^{}`$. We sketch a proof of this, following reference : First, the extreme points of $`\mathrm{SS}`$ are identified as the product states $`\rho ^{\mathrm{}}`$. Thus, the set of extreme points is the image of the compact space $`\mathrm{SS}_1`$ under the continuous injection $`\rho \rho ^{\mathrm{}}`$, and it follows that the extreme set is closed in $`\mathrm{SS}`$. The existence of an integral representation (6) is then a consequence of the Krein-Milman Theorem, and its uniqueness is shown in by a direct argument.
It is further shown in that the measure $`\mu (d\rho )`$ appearing in the integral representation (6) of $`\omega `$ is supported on the measurable subset of normal states on $`()`$ if $`\omega `$ is determined, as in formula (4) above, by sequences of density operators satisfying (A) and (C). If, in addition, the sequence of density operators defining $`\omega `$ satisfies (B), then the measure $`\mu (d\rho )`$ is even supported on the vector states $`\rho (A)=\psi ,A\psi `$ with $`\psi =1`$.
## 3 Examples of bosonic de Finetti states
In this section we exhibit some sequences $`\{D_n\}`$ satisfying (B) and (C) that are obtained from natural statistical ensembles. In all of these examples, the single-particle Hilbert space $``$ is finite dimensional. After introducing the notation, we will state all of our results before proceeding to their proofs.
Let $`=^{d+1}`$ and let $`^{(n)}`$ denote the subspace of symmetric vectors in $`^n`$. Let $`\mathrm{\Sigma }_n`$ denote the symmetrizing projector
$$\mathrm{\Sigma }_n=\frac{1}{n!}\underset{\pi \mathrm{\Pi }_n}{}U_\pi $$
(7)
from $`^n`$ onto $`^{(n)}`$. We now introduce notation for the occupation number basis of $`^{(n)}`$ relative to a fixed orthonormal (ordered) basis $`\{e_j\}`$ of $``$. Let $`๐ง=(n_0,n_1,\mathrm{},n_d)`$ be an ordered $`d+1`$-tuple of nonnegative integers (occupation numbers) and let $`\mathrm{\#}๐ง`$ denote $`n_j`$. We use the notation
$$\left(\genfrac{}{}{0pt}{}{n}{๐ง}\right)=n!/\underset{i=0}{\overset{d}{}}n_i!$$
for multinomial coefficients. The vector
$$\mathrm{\Psi }_๐ง=\sqrt{\left(\genfrac{}{}{0pt}{}{n}{๐ง}\right)}\mathrm{\Sigma }_n(e_0^{n_0}e_1^{n_1}\mathrm{}e_d^{n_d})$$
(8)
is a unit vector in $`^{(n)}`$, and the set of vectors $`\{\mathrm{\Psi }_๐ง|\mathrm{\#}๐ง=n\}`$ is an orthonormal basis of $`^{(n)}`$. Let $`P_๐ง`$ denote the rank-one projector onto the span of $`\mathrm{\Psi }_๐ง`$:
$$P_๐ง\mathrm{\Phi }=\mathrm{\Psi }_๐ง,\mathrm{\Phi }\mathrm{\Psi }_๐ง.$$
(9)
We begin by considering the โuniformly mixedโ density operators supported on $`^{(n)}`$:
###### Proposition 1
Let $`\mathrm{\Sigma }_n`$ denote the symmetrizing projector (7). For each $`m`$,
$$S_m\underset{n\mathrm{}}{lim}\frac{1}{\mathrm{Tr}\mathrm{\Sigma }_n}\mathrm{\Sigma }_{n:m}=\underset{๐ฆ:\mathrm{\#}๐ฆ=m}{}\left\{\left(\genfrac{}{}{0pt}{}{m}{๐ฆ}\right)_{\mathrm{\Delta }_d}\underset{i=0}{\overset{d}{}}p_i^{m_i}\lambda _d(dp)\right\}P_๐ฆ,$$
(10)
where $`\lambda _d(dp)`$ denotes normalized Lebesgue measure on the $`d`$-dimensional simplex
$$\mathrm{\Delta }_d=\left\{p=(p_0,p_1,\mathrm{},p_d)^{d+1}\right|0p_ii=1,2,\mathrm{},d\text{and}\underset{i=0}{\overset{d}{}}p_i=1\}.$$
The sequence $`\{S_m\}`$ satisfies (B) and (C) of Section 2. By the Quantum de Finetti Theorem, there exists a measure $`\mu `$ supported on the pure states on $`^{d+1}`$ such that
$$S_m=P^m\mu (dP)$$
for all $`m`$. This measure can be described as follows. Define the map
$$\mathrm{v}:\mathrm{\Delta }_d\times [0,2\pi )^{d+1}^{d+1}$$
by
$$\mathrm{v}(p_0,p_1,\mathrm{},p_d,\theta _0,\theta _1,\mathrm{},\theta _d)=\underset{j=0}{\overset{d}{}}e^{i\theta _j}\sqrt{p_j}e_j$$
(11)
where $`\{e_i\}`$ is the standard basis of $`^{d+1}`$. The map $`\mathrm{v}`$ is many-one onto the set of unit vectors in $`^{d+1}`$. The probability measure $`\mu (dP)`$ is the one induced via $`\mathrm{v}`$ from the uniform measure
$$\lambda (dp)\sigma (d\theta )\lambda (dp)\frac{d\theta _0}{2\pi }\frac{d\theta _1}{2\pi }\mathrm{}\frac{d\theta _d}{2\pi }$$
on $`\mathrm{\Delta }_d\times [0,2\pi )^{d+1}`$. In other words,
###### Proposition 2
The density operator (10) equals
$$_{\mathrm{\Delta }_d}_{[0,2\pi )^{d+1}}\left(P_{\mathrm{v}(p,\theta )}\stackrel{\text{m times}}{\mathrm{}}P_{\mathrm{v}(p,\theta )}\right)\sigma (d\theta )\lambda _d(dp).$$
(12)
Next we consider Gibbs ensembles for noninteracting systems of bosons. Let
$$H_n=\underset{i=1}{\overset{n}{}}T_i$$
(13)
be the Hamiltonian for $`n`$ noninteracting bosons with single-particle space $`=^{d+1}`$. Let $`\{e_j\}`$ be an orthonormal basis of $``$ consisting of eigenvectors of the single-particle operator $`T`$, so that $`Te_j=ฯต_je_j`$. The Gibbs density operator for the $`n`$ boson system is
$$\mathrm{\Gamma }_n(\beta )=\frac{1}{Z_{n,\beta }}\underset{๐ง:\mathrm{\#}๐ง=n}{}\underset{i=0}{\overset{d}{}}e^{\beta n_iฯต_i}P_๐ง\mathrm{with}Z_{n,\beta }=\underset{๐ง:\mathrm{\#}๐ง=n}{}\underset{i=0}{\overset{d}{}}e^{\beta n_iฯต_i}.$$
(14)
An interesting limit is attained if temperature is scaled in proportion to $`n`$ as $`n\mathrm{}`$. If the temperature is not scaled as $`n\mathrm{}`$ then a sort of Bose-Einstein condensation is attained in the limit.
###### Proposition 3
Let $`H_n`$ be the noninteracting Hamiltonian (13) and let $`\mathrm{\Gamma }_n(\beta )`$ denote the Gibbs density (14). Let $`\{e_j\}`$ be an orthonormal basis of $``$ consisting of eigenvectors of the single-particle operator $`T`$, so that $`Te_j=ฯต_je_j`$.
* (i) For each $`m`$, the limit $`\underset{n\mathrm{}}{lim}\mathrm{\Gamma }_{n:m}(\beta /n)`$ exists and equals
$$\underset{๐ฆ:\mathrm{\#}๐ฆ=m}{}\left\{\left(\genfrac{}{}{0pt}{}{m}{๐ฆ}\right)_{\mathrm{\Delta }_d}\underset{i=0}{\overset{d}{}}p_i^{m_i}Z_\beta ^1\underset{i=0}{\overset{d}{}}e^{\beta ฯต_ip_i}\lambda _d(dp)\right\}P_๐ฆ$$
with $`Z_\beta ^1=_{\mathrm{\Delta }_d}_{i=0}^d\mathrm{exp}(\beta ฯต_ip_i)\lambda _d(dp)`$.
* (ii) If $`ฯต_0<ฯต_1\mathrm{}ฯต_d`$, then for each $`m`$,
$$\underset{n\mathrm{}}{lim}\mathrm{\Gamma }_{n:m}(\beta )=P_{(m,0,\mathrm{},0)}=P_{e_0}^{}{}_{}{}^{m}.$$
Finally, we consider systems with two-particle interactions in the โmean fieldโ scaling. Let $`V`$ be a Hamiltonian operator on $``$ such that $`V(xy)=V(yx)`$ for all $`x,y`$. For $`n>2`$, define the Hamiltonian
$$H_n=\underset{i=1}{\overset{n}{}}T_i+\frac{1}{n1}\underset{1i<jn}{}V_{ij},$$
(15)
where $`V_{ij}`$ denotes the operator obtained by applying $`V`$ to the $`i^{th}`$ and $`j^{th}`$ factors of $`^n`$. For any $`n`$ and any $`\beta `$, the $`n`$-particle Gibbs density at inverse temperature $`\beta `$ for the Hamiltonian (15) is
$$\mathrm{\Gamma }_n(\beta )=\frac{1}{\mathrm{Tr}(e^{\beta H_n}\mathrm{\Sigma }_n)}e^{\beta H_n}\mathrm{\Sigma }_n.$$
(16)
###### Proposition 4
Let $`\mathrm{\Gamma }_n(\beta )`$ denote the Gibbs density (16). For each $`m`$, the limit
$$G_m=\underset{n\mathrm{}}{lim}\left\{\mathrm{\Gamma }_n(\beta /n)\right\}_{:m}$$
exists and defines a density operator on $`(^{d+1})^n`$. The de Finetti representation of $`G_m`$ is
$$\frac{1}{Z_\beta }_{\mathrm{\Delta }_d}_{[0,2\pi )^{d+1}}\stackrel{\text{m times}}{P_\mathrm{v}\mathrm{}P_\mathrm{v}}e^{\beta \{\mathrm{Tr}(TP_\mathrm{v})+\mathrm{Tr}(V(P_\mathrm{v}P_\mathrm{v}))/2\}}\sigma (d\theta )\lambda _d(dp)$$
with $`\mathrm{v}=\mathrm{v}(p,\theta )`$ as in (11) and
$$Z_\beta =_{\mathrm{\Delta }_d}_{[0,2\pi )^{d+1}}e^{\beta \{\mathrm{Tr}(TP_\mathrm{v})+\mathrm{Tr}(V(P_\mathrm{v}P_\mathrm{v}))/2\}}\sigma (d\theta )\lambda _d(dp).$$
The rest of this section is devoted to proving the above propositions.
Recall the definition (9) of the projectors $`P_๐ง`$. For each $`n`$, let $`\rho _n`$ be an $`n`$-particle density
$$\rho _n=\underset{๐ง:\mathrm{\#}๐ง=n}{}w_n(๐ง)P_๐ง,$$
where $`w_n`$ is a probability measure on the set $`\{๐ง|\mathrm{\#}๐ง=n\}`$. Each probability measure $`w_n`$ can be associated with the discrete probability measure
$$\omega _n=\underset{๐ง:\mathrm{\#}๐ง=n}{}w_n(๐ง)\delta (p\frac{1}{n}๐ง)$$
on the $`d`$-dimensional simplex $`\mathrm{\Delta }_d`$. It may be verified that
$$P_{๐ง:m}=\left(\genfrac{}{}{0pt}{}{n}{m}\right)^1\underset{๐ฆ:\mathrm{\#}๐ฆ=m}{}\underset{i=0}{\overset{d}{}}\left(\genfrac{}{}{0pt}{}{n_i}{m_i}\right)P_๐ฆ$$
(this equals $`0`$ if any $`m_i>n_i`$ for any $`i`$), and therefore
$`\rho _{n:m}`$ $`=`$ $`\left({\displaystyle \genfrac{}{}{0pt}{}{n}{m}}\right)^1{\displaystyle \underset{๐ฆ:\mathrm{\#}๐ฆ=m}{}}\left[{\displaystyle \underset{๐ง:\mathrm{\#}๐ง=n}{}}w_n(๐ง){\displaystyle \underset{i=0}{\overset{d}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{n_i}{m_i}}\right)\right]P_๐ฆ`$
$`=`$ $`{\displaystyle \underset{๐ฆ:\mathrm{\#}๐ฆ=m}{}}\left({\displaystyle \genfrac{}{}{0pt}{}{m}{๐ฆ}}\right)\left[{\displaystyle \underset{๐ง:\mathrm{\#}๐ง=n}{}}w_n(๐ง){\displaystyle \frac{_{i=0}^d\frac{n_i}{n}(\frac{n_i}{n}\frac{1}{n})\mathrm{}(\frac{n_i}{n}\frac{m_i1}{n})}{1(1\frac{1}{n})(1\frac{2}{n})\mathrm{}(1\frac{m1}{n})}}\right]P_๐ฆ.`$
The coefficient of $`P_๐ฆ`$ in (LABEL:rewriteMe) may be written
$$\left(\genfrac{}{}{0pt}{}{m}{๐ฆ}\right)_{\mathrm{\Delta }_d}f_n(p)\omega _n(dp),$$
where
$$f_n(p)=1\mathrm{l}_{\{p_i>(m_i1)/ni\}}(p)\frac{_{i=0}^dp_i(p_i\frac{1}{n})\mathrm{}(p_i\frac{m_i1}{n})}{1(1\frac{1}{n})(1\frac{2}{n})\mathrm{}(1\frac{m1}{n})}.$$
The functions $`f_n(p)`$ converge uniformly to $`_{i=0}^dp_i^{m_i}`$ on $`\mathrm{\Delta }_d`$. Therefore, if $`\omega _n`$ converges weakly to some probability measure $`\omega (dp)`$ on $`\mathrm{\Delta }_d`$, then
$$\underset{n\mathrm{}}{lim}\rho _{n:m}=\underset{๐ฆ:\mathrm{\#}๐ฆ=m}{}\left(\genfrac{}{}{0pt}{}{m}{๐ฆ}\right)_{\mathrm{\Delta }_d}\underset{i=0}{\overset{d}{}}p_i^{m_i}\omega (dp)P_๐ฆ.$$
(18)
The probability measures on $`\mathrm{\Delta }_d`$ corresponding to the Gibbs density operators (14) for noninteracting bosons are
$$\omega _n=Z_{n,\beta }^1\underset{๐ง:\mathrm{\#}๐ง=n}{}\underset{i=0}{\overset{d}{}}e^{\beta n_iฯต_i}\delta (p\frac{1}{n}๐ง).$$
(19)
If all of the eigenvalues of $`T`$ are equal, then the measures (19) converge weakly to $`\lambda _d(dp)`$, the uniform probability measure on the simplex, but if $`ฯต_0`$ is strictly smaller than all of the other eigenvalues of $`T`$, then the measures (19) converge weakly to $`\delta (p(1,0,\mathrm{},0))`$, a point-mass at the lowest energy vertex of the simplex. This convergence implies Propositions 1 and assertion (ii) of Proposition 3 by formula (18). On the other hand, the probability measures corresponding to the Gibbs density operators $`\mathrm{\Gamma }_n(\beta /n)`$ for noninteracting bosons are
$$\omega _n=Z_{n,\beta }^1\underset{๐ง:\mathrm{\#}๐ง=n}{}\underset{i=0}{\overset{d}{}}e^{\beta ฯต_in_i/n}\delta (p\frac{1}{n}๐ง),$$
and these converge weakly to
$$Z_\beta ^1\underset{i=0}{\overset{d}{}}e^{\beta ฯต_ip_i}\lambda _d(dp)$$
with $`Z_\beta =_{\mathrm{\Delta }_d}_{i=0}^d\mathrm{exp}(\beta ฯต_ip_i)\lambda _d(dp)`$. This proves assertion (i) of Proposition 3.
To prove Proposition 2, we will show that (12) and (10) are equal. Define the rank-one operators $`Q_{jk}(x)=e_k,xe_j`$. From (11),
$$P_{\mathrm{v}(p,\theta )}=\underset{j,k=0}{\overset{d}{}}e^{i(\theta _j\theta _k)}\sqrt{p_jp_k}Q_{jk}$$
and therefore $`P_{\mathrm{v}(p,\theta )}^{}{}_{}{}^{m}`$ equals
$$\underset{j_1,\mathrm{},j_m=0}{\overset{d}{}}\underset{k_1,\mathrm{},k_m=0}{\overset{d}{}}\underset{r=0}{\overset{d}{}}\sqrt{p_{j_r}p_{k_r}}e^{i(\theta _{j_r}\theta _{k_r})}Q_{j_1k_1}Q_{j_2k_2}\mathrm{}Q_{j_mk_m}.$$
(20)
For $`i=0,1,\mathrm{},d`$, let $`N_i:\{0,1,\mathrm{},d\}^m`$ be defined by
$$N_i(x_1,x_2,\mathrm{},x_m)=\mathrm{\#}\{k\{1,2,\mathrm{},m\}:x_k=i\}$$
and define
$$N(x_1,x_2,\mathrm{},x_m)=(N_0(x_1,x_2,\mathrm{},x_m),\mathrm{},N_d(x_1,x_2,\mathrm{},x_m)).$$
If $`N(j_1,\mathrm{},j_m)=N(k_1,\mathrm{},k_m)`$ then
$$_{[0,2\pi )^{d+1}}\underset{r=0}{\overset{d}{}}e^{i(\theta _{j_r}\theta _{k_r})}\sigma (d\theta )=1,$$
but otherwise it equals $`0`$. Thus, from (20),
$`{\displaystyle _{\mathrm{\Delta }_d}}{\displaystyle _{[0,2\pi )^{d+1}}}P_{\mathrm{v}(p,\theta )}^{}{}_{}{}^{m}\sigma (d\theta )\lambda _d(dp)`$
$`=`$ $`{\displaystyle \underset{๐ฆ:\mathrm{\#}๐ฆ=m}{}}{\displaystyle _{\mathrm{\Delta }_d}}{\displaystyle \underset{i=0}{\overset{m}{}}}p_i^{m_i}\lambda _d(dp){\displaystyle \underset{\stackrel{j_1,\mathrm{},j_m:}{N(j_1,\mathrm{},j_m)=๐ฆ}}{}}{\displaystyle \underset{\stackrel{k_1,\mathrm{},k_m:}{N(k_1,\mathrm{},k_m)=๐ฆ}}{}}Q_{j_1k_1}\mathrm{}Q_{j_mk_m}`$
$`=`$ $`{\displaystyle \underset{๐ฆ:\mathrm{\#}๐ฆ=m}{}}{\displaystyle _{\mathrm{\Delta }_d}}{\displaystyle \underset{i=0}{\overset{m}{}}}p_i^{m_i}\lambda _d(dp)\left({\displaystyle \genfrac{}{}{0pt}{}{m}{๐ฆ}}\right)P_๐ฆ`$
by the definition (9) of $`P_๐ฆ`$. This proves Proposition 2.
Finally, we derive Proposition 4 from Proposition 2. Define $`W=TI+IT+V`$. Then the Hamiltonian (15) can be written
$$H_n=\frac{1}{n1}\underset{1i<jn}{}W_{ij}.$$
We claim that
$$\underset{n\mathrm{}}{lim}\frac{n^j}{\mathrm{Tr}\mathrm{\Sigma }_n}\{(H_n)^j\mathrm{\Sigma }_n\}_{:m}=2^j\left\{W_{m+1,m+2}W_{m+3,m+4}\mathrm{}W_{m+2j1,m+2j}S_{m+2j}\right\}_{:m}$$
(21)
for each $`j,m`$. This is so because $`(H_n)^j`$ contains $`\left(\genfrac{}{}{0pt}{}{n}{2}\right)^j`$ terms of the form $`(n1)^jW_{a_1b_1}W_{a_2b_2}\mathrm{}W_{a_jb_j}`$, and, when $`n`$ is large, the majority of these terms are such that the indices $`a_1,b_1,\mathrm{},a_j,b_j`$ are all distinct and greater than $`m`$. The sum of the remaining terms in $`(H_n)^j`$ is $`o(n^j)`$ and does not contribute to the limit (21). By the symmetry of $`\mathrm{\Sigma }_n`$,
$`\left\{W_{a_1b_1}W_{a_2b_2}\mathrm{}W_{a_jb_j}\mathrm{\Sigma }_n\right\}_{:m}`$ $`=`$ $`\left\{W_{m+1,m+2}\mathrm{}W_{m+2j1,m+2j}\mathrm{\Sigma }_n\right\}_{:m}`$
$`=`$ $`\left\{W_{m+1,m+2}\mathrm{}W_{m+2j1,m+2j}\mathrm{\Sigma }_{n:m+2j}\right\}_{:m}`$
if $`a_1,b_1,\mathrm{},a_j,b_j`$ are all distinct and greater than $`m`$. There are asymptotically $`n^{2j}/2`$ such terms, so (21) follows from Proposition 1.
Now, to prove Proposition 4, expand
$$\frac{1}{\mathrm{Tr}\mathrm{\Sigma }_n}e^{\beta n^1H_n}\mathrm{\Sigma }_n=\underset{j=0}{\overset{\mathrm{}}{}}\frac{1}{j!}(\beta )^jn^j(H_n)^j\frac{1}{\mathrm{Tr}\mathrm{\Sigma }_n}\mathrm{\Sigma }_n$$
and take the $`m^{th}`$ partial trace:
$$\frac{1}{\mathrm{Tr}\mathrm{\Sigma }_n}\left\{e^{\beta n^1H_n}\mathrm{\Sigma }_n\right\}_{:m}=\underset{j=0}{\overset{\mathrm{}}{}}\frac{1}{j!}(\beta )^jn^j\left\{(H_n)^j\frac{1}{\mathrm{Tr}\mathrm{\Sigma }_n}\mathrm{\Sigma }_n\right\}_{:m}.$$
(22)
The $`j^{th}`$ term of the series in (22) converges to
$$(1)^j\frac{1}{j!}\left(\frac{\beta }{2}\right)^j\left\{W_{m+1,m+2}W_{m+3,m+4}\mathrm{}W_{m+2j1,m+2j}S_{m+2j}\right\}_{:m}$$
as $`n\mathrm{}`$ by (21) and is bounded by $`\frac{1}{j!}\beta ^jW^j`$ uniformly in $`n`$. Since the series in (22) are majorized by the convergent series $`_j\frac{1}{j!}\beta ^jW^j`$ and converge term-by-term as $`n\mathrm{}`$, it follows that
$$\underset{n\mathrm{}}{lim}\frac{1}{\mathrm{Tr}\mathrm{\Sigma }_n}\left\{e^{\beta n^1H_n}\mathrm{\Sigma }_n\right\}_{:m}=\underset{j=0}{\overset{\mathrm{}}{}}(1)^j\frac{1}{j!}\left(\frac{\beta }{2}\right)^j\left\{W_{m+1,m+2}\mathrm{}W_{m+2j1,m+2j}S_{m+2j}\right\}_{:m}.$$
(23)
Substituting the integral representations (12) for $`S_{m+2j}`$ into (23) yields
$`\underset{n\mathrm{}}{lim}{\displaystyle \frac{1}{\mathrm{Tr}\mathrm{\Sigma }_n}}\left\{e^{\beta n^1H_n}\mathrm{\Sigma }_n\right\}_{:m}`$
$`=`$ $`{\displaystyle \underset{j=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{j!}}\left({\displaystyle \frac{\beta }{2}}\right)^j{\displaystyle _{\mathrm{\Delta }_d}}{\displaystyle _{[0,2\pi )^{d+1}}}\left[\mathrm{Tr}\left(WP_{\mathrm{v}(p,\theta )}^{}{}_{}{}^{2}\right)\right]^jP_{\mathrm{v}(p,\theta )}^{}{}_{}{}^{m}\sigma (d\theta )\lambda _d(dp)`$
$`=`$ $`{\displaystyle _{\mathrm{\Delta }_d}}{\displaystyle _{[0,2\pi )^{d+1}}}e^{\beta {\scriptscriptstyle \frac{1}{2}}\mathrm{Tr}(WP_{\mathrm{v}(p,\theta )}P_{\mathrm{v}(p,\theta )})}P_{\mathrm{v}(p,\theta )}^{}{}_{}{}^{m}\sigma (d\theta )\lambda _d(dp).`$
Proposition 4 follows from the preceding equation and the definition (16) of $`\mathrm{\Gamma }_n(\beta )`$.
## 4 Acknowledgments
I would like to thank Lucien Le Cam for listening very patiently to some of this story and for his kind encouragement. This work was supported by the Austrian START project โNonlinear Schrรถdinger and quantum Boltzmann equationsโ of Norbert J. Mauser (contract Y-137-Tec).
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# Methods
## Methods
### Spectro-astrometry
Conceptually the principles of spectro-astrometry are easy to understand. The profile of a star is smeared by atmospheric turbulence to appear gaussian (at least to a first approximation) rather than point-like. Whereas the width of the profile is determined by the so-called seeing, how accurately we can determine the centroid of emission is, in theory for fixed seeing, limited only by the strength of the observed signal to noise ratio. Increasing the total number of detected photons increases the positional, or astrometric, accuracy, so that, in principle, milliarcsecond precision is possible with very large ground based telescopes<sup>16-18</sup>.
Consider now a long-slit spectrum of a close binary system consisting of two virtually identical stars. We will assume that the slit is orientated along the same position angle as the binary. (We note that strictly this is not necessary: it is only necessary that the slit is not orthogonal). If the separation of the binary is considerably less than the seeing, the profile of the system in the spatial direction will consist of a single gaussian-that is, the system is unresolved and the centroid of emission will lie exactly between the two components. Now suppose that one of the two stars differs slightly from the other in being a strong H$`\alpha `$ emitter; in such a case, the emission centroid will shift towards that star in the spectrum at the position of the H$`\alpha `$ line. In this way it is possible to resolve certain types of binaries with separations well within the seeing limit<sup>19</sup>. In the case of a jet (pure emission line region) plus star (continuum source), one can go further and interpolate the continuum across a line, thereby allowing its contribution to be removed. It is then possible to measure separately the spatial centroid of the pure emission line region and determine its offset with respect to the continuum, that is, the parent star. Moreover, as the line can be emitted over a range of wavelengths, owing to the Doppler effect, it may also be feasible to recover spatio-kinematic information. For example, if the jet is bipolar, that is, it has oppositely directed blue- and redshifted flows from the source; the emission centroid of the red and blue wings of the line will be displaced to opposite sides of the continuum centre.
The detailed method by which we measure offsets can briefly be described as follows. First, the centroid of the continuum emission in the spatial direction is determined using a one-dimensional gaussian fit. The line of such centroids, in the dispersion direction but excluding any region where emission lines are present, is then fitted with a second-order polynomial, over a range of typically 200-300 . In this way, instrumental curvature and tilting, with a characteristic frequency many times larger than the width of any line, is determined. The fit, to the centre of the continuum, is then subtracted from the actual measured centroids, leaving residuals that are evenly scattered about the abscissa (that is, the fit defines the zero offset line). The continuum data points shown in Fig. 1 are thus the residuals. Finally, the two-dimensional fit to the continuum, broadened to take account of the point spread function, is subtracted from the emission lines. Any emission line offsets are then measured.
The accuracy (in arcseconds) of the method is set by the error in the centroid of the gaussian fit, which depends on the seeing and the number of detected photons, N. Formally, the error is given by Seeing/\[2(2 ln 2)1/2N1/2\], assuming that photon noise is the only source of noise. N, of course, is a function of the binning and the spatial sampling (pixel width). This explains why, for example, we can achieve a higher spectro-astrometric accuracy with a bright line, such as H$`\alpha `$ than a weak one, for example, the \[SII\]$`\lambda `$6731 line. In some cases, it is necessary to bin up a weak line in the dispersion direction, as we have done to varying degrees for the \[O I\] doublet and the \[SII\]$`\lambda `$6731 line, to achieve sufficient signal to noise ratio. Note that we sometimes use different binning factors for the continuum, in comparison with the line, so as to achieve a similar signal to noise ratio in both. This allows us to have comparable offset errors in both components, and to define the common 1$`\sigma `$ error lines shown in Fig. 1. As can be seen from the plots, the typical limiting offset that we can measure in the spatial direction (3$`\sigma `$) is around 30 mas. This corresponds to 4.5 au at the distance to the $`\rho `$ Ophiuchi cloud.
### Echelle spectroscopy
The high resolution spectra of $`\rho `$ Oph 102 were taken with the UV-visual Echelle Spectrograph (UVES) on the European Southern Observatoryโs 8 m Kueyen Telescope, one of the telescopes in the Very Large Telescope (VLT) suite, in May 2003. A total of three 45 min exposures of the target were made, together with a series of flats and biases as well as an observation of an arc lamp for wavelength calibration. The slit was orientated north-south and had a width of 1 โณwhile the seeing was 0$`\stackrel{}{.}`$65. The central wavelength was set at 580 nm, giving a spectral range of 450-680 nm. Only the red part of the spectrum from 580 to 680 nm, however, was analysed. The pixel scale was 0$`\stackrel{}{.}`$182 and the spectral resolution R=40,000. The data were reduced using standard Image Reduction and Analysis Facility (IRAF) routines.
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2. Kรถnigl, A. & Pudritz, R. E. Disk winds and the accretion-outflow connection. Protostars and Planets IV, University of Arizona Press. 759-787 (2000).
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8. Comerรณn, F., Fernรกndez, M., Baraffe, I., Neuhรคuser, R., & Kaas, A. A. New low-mass members of the Lupus 3 dark cloud: Further indications of pre-main-sequence evolution strongly affected by accretion. Astron. Astrophys. 406, 1001-1017 (2003).
9. Fernรกndez, M. & Comerรณn, F. Intense accretion and mass loss of a very low mass young stellar object. Astron. Astrophys. 380, 264-276 (2001).
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13. Hirth, G. A., Mundt, R., & Solf, J. Spatial and kinematic properties of the forbidden emission line region of T Tauri stars. Astron. Astrophys. Suppl. 126, 437-469 (1997).
14. Masciadri, E., & Raga, A. C. Looking for outflows from brown dwarfs. Astrophys. J., 615, 850-854 (2004).
15. Muzerolle, J., Calvet, N., & Hartmann, L. Emission-line diagnostics of T Tauri magnetospheric accretion. II. Improved model tests and insights into accretion physics. Astrophys. J. 550, 944-961 (2001).
16. Takami, M., Bailey, J., Gledhill, T. M., Chrysostomou, A., & Hough, J. H. Circumstellar structure of RU Lupi down to AU scales. Mon. Not. R. Astron. Soc. 323 177-187 (2001).
17. Whelan, E. T., Ray, T. P., & Davis, C. J. Paschen beta emission as a tracer of outflow activity from T-Tauri stars, as compared to optical forbidden emission. Astron. Astrophys. 417 247-261 (2004).
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Supplementary Information is linked to the online version of the paper at www.nature.com/nature but requests for supplementary materials can also be addressed to ewhelan@cp.dias.ie
Acknowledgements This work was supported in part by Science Foundation Ireland and the JETSET Marie Curie reserach training network.
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# Pressure effects on the electronic properties and superconductivity in the ๐ฝ-pyrochlore oxides: ๐ดOs2O6 (๐ด = Na, K, Rb, Cs)
## I Introduction
During the last decade or so, the search and interest in the superconductivity of non-Cu based oxides has extended from an effort to understand the pairing mechanism in the cuprates to a broader search of superconductivity in materials in which electron correlations are thought to play a determining role. To this end, researchers try to exploit a diverse range of factors, from the orbital degrees of freedom to the crystal structure of the material. A particularly interesting example is the recently discovered family of superconducting Os-oxide $`\beta `$-pyrochlores AOs<sub>2</sub>O<sub>6</sub>, with $`A`$ = K, Rb, and Cs,yonezawa04a ; yonezawa04b ; yonezawa04c which have superconducting transition temperatures $`T_c=9.6`$ K, 6.3 K, and 3.3 K, respectively. The pyrochlore structure is a network of corner sharing tetrahedra and is a geometrically frustrated spin system if the ions bear a localized magnetic moment interacting antiferromagnetically with their nearest neighbors. Further, the Os ionsโlocated at the tetrahedra verticesโpossess a formal oxidation state of $`5.5+`$ ($`5d^{2.5}`$). As pointed out by Hiroi and co-workers,hiroi04 compared to other transition metal pyrochlore oxides, this places these materials between Cd<sub>2</sub>Re<sub>2</sub>O<sub>7</sub> (Re$`{}_{}{}^{5+}:5d^2`$), which is a good conductor at low temperatureshanawa01 ; jin01 (becoming superconductor at $`1`$ K) and Cd<sub>2</sub>Os<sub>2</sub>O<sub>7</sub> (Os$`{}_{}{}^{5+}:5d^3`$), which is an insulator at low temperatures and exhibits antiferromagnetic ordering.mandrus01
While the several experimental resultsyonezawa04a ; yonezawa04b ; yonezawa04c ; bruehwiler04 ; magishi04 ; arai04 ; koda04 ; khasanov04 ; muramatsu04 reported during the past year provide key information, the pairing mechanism in these materials is still under debate. These experiments indicate similarities but also differences among the compounds with different alkali metal, tending to single out KOs<sub>2</sub>O<sub>6</sub>. Let us mention briefly some of the observations that should be taken into account by any proposed pairing mechanism. Firstly, the very change in $`T_c`$ upon substitution of the alkali atom ($`A`$) may be initially counterintuitive. Indeed, the negative chemical pressure leads to an increase of the lattice constant with the ionic radius of $`A`$, so that one may expect $`T_c`$ to increase because band narrowing should lead to an increase of the density of states (DOS) at the Fermi level ($`E_\mathrm{F}`$). As shown by the reported $`T_c`$โs above, however, the opposite occurs. In line with these findings, the change in $`T_c`$ under applied hydrostatic pressure is found to be positive initially in all these materials.khasanov04 ; muramatsu04 This has been interpreted as suggesting that the pairing mechanism in these materials is unconventional, or non-BCS-like.yonezawa04c Of interest is the fact that in KOs<sub>2</sub>O<sub>6</sub> the increase of $`T_c`$ with pressure reaches a maximum at 0.56 GPa and tends to vanish gradually at higher pressures.muramatsu04 ; comment0
Further observations are that the temperature dependence of the resistivity shows an unusual concave behavior at low temperature in the case of KOs<sub>2</sub>O<sub>6</sub>,yonezawa04a while a $`T^2`$ behavior is observedyonezawa04b just above $`T_c`$ in the case of RbOs<sub>2</sub>O<sub>6</sub> and on a larger temperature intervalyonezawa04c in the case of CsOs<sub>2</sub>O<sub>6</sub>. Also, nuclear magnetic resonance (NMR) experiments reveal a weak temperature dependence of the Knight shift of both the <sup>39</sup>K and <sup>87</sup>Rb nuclei in the corresponding pyrochlores.magishi04 ; arai04 In the normal state, however, the relaxation rate divided by temperature ($`1/TT_1`$) follows the Korringa relaxation in the case of Rb,magishi04 or deviates weakly from it,arai04 but deviates more strongly from this behavior in the case of K.arai04 In both cases, however, there appears to be evidence for antiferromagnetic spin fluctuations. Finally, regarding the superconducting gap, Magishi and co-workers find that the relaxation rate in the superconducting state suggests an anisotropic but nodeless gapmagishi04 in RbOs<sub>2</sub>O<sub>6</sub>; at the same time, Koda and collaborators interpret their muon spin rotation ($`\mu `$SR) study of the magnetic penetration depth in KOs<sub>2</sub>O<sub>6</sub> as pointing to an anisotropic gap with nodes.koda04
We report here on a first-principles study of the electronic structure and superconducting parameters of the compounds $`A`$Os<sub>2</sub>O<sub>6</sub> ($`A`$ = Na, K, Rb, and Cs) and on the effects of hydrostatic pressure. We find that the main traits of the electronic structure reported previouslysaniz04 in the case of KOs<sub>2</sub>O<sub>6</sub> are common to all these materials, with relatively small qualitative and quantitative changes. The differences stem essentially from the energy level with respect to $`E_\mathrm{F}`$ of the van Hove singularity (vHS) with momentum $`๐ค`$ near the center of the $`\mathrm{\Gamma }`$-$`L`$ line. In particular, the density of states at the Fermi energy, $`N(E_\mathrm{F})`$, tends to increase with the size of $`A`$ because the vHS is pushed closer to $`E_\mathrm{F}`$. The effect of applied hydrostatic pressure is to push the vHS away from $`E_\mathrm{F}`$. This is very clearly reflected by the increase or suppression of a constriction between the two $`\mathrm{\Gamma }`$-point centered Fermi surface shells, basically due to a bending of the outer shell that depends on the proximity of the vHS to $`E_\mathrm{F}`$.
We further estimate $`T_c`$ with the well-known McMillan-Allen-Dynes expression,bennemann72 with the electron-phonon coupling constant calculated within the crude rigid muffin-tin approximation (RMTA).gaspari72 ; pettifor77 We also calculate the Stoner susceptibility enhancement parameter and estimate the electron-spin coupling constant within the Doniach-Engelsberg approximation.doniach66 This allows us to show that spin fluctuations contribute importantly to the effective electron mass, significantly reducing $`T_c`$. Despite the approximations implicit in these calculations, we find, remarkably, that the calculated $`T_c`$ follows rather well the trends observed in experiment, both upon substitution of the alkali metal and under hydrostatic pressure. Our results, thus, bring further support to the electron-phonon coupling description of these superconductors.bruehwiler04 ; magishi04 ; khasanov04 ; saniz04
Section II is devoted to the methodology of our calculations as well as to structural properties; in section III, we present and discuss our results in relation to the experimental findings mentioned above.
## II Methodology and structural aspects
We use the highly precise full-potential linearized augmented plane-wave (FLAPW)wimmer81 implementation of the density functional approach to the electronic structure and properties of crystalline solids. We make our self-consistent calculations within the Perdew, Burke, and Ernzerhof generalized gradient approximation (GGA)perdew97 to the exchange-correlation potential and include the spin-orbit coupling (SOC) term in the Hamiltonian. Angular momenta up to $`l=8`$ are used for both the charge density in the muffin-tins and the wave functions. The irreducible part of the Brillouin zone is sampled with a uniform mesh of 120 k-points. The Os $`5p`$ and K $`3p`$ states are treated as valence electrons.
The $`\beta `$-pyrochlores crystallize in a cubic structure with space group $`Fd\overline{3}m`$. There are 18 atoms in the unit cell: two $`A`$ atoms ($`8b`$), four Os atoms ($`16c`$), and twelve O atoms ($`48f`$). The Os atoms are octahedrally coordinated by six O atoms. An internal parameter, $`x`$, fixes the positions of the latter and thereby also determines the degree of rhombohedral distortion of the octahedra enclosing the Os atoms. In all the cases considered, we determine the lattice constant, $`a`$, as well as $`x`$, by minimizing the total energy and ensuring the total force on the O atoms is less than $`10^4`$ a.u. The muffin-tin radii used are 2.2 a.u. for the Os ions and 1.3 a.u. for the O ions. The corresponding values for the Na, K, Rb, and Cs ions are, respectively, 2.6, 2.8, 3.0, and 3.1 a.u.
The values of the calculated structural parameters are reported in Table 1, including the smaller Os-O-Os angle defining the main rectangular cross section of the octahedra and the O-Os-O angle characterizing the staggered Os-O chains on the underlying the pyrochlore lattice. The calculated ($`T=0`$ K) lattice constants differ from the (room temperature) experimental results of Hiroi and co-workers (see Ref. muramatsu04, ) by $`+1.95`$% for KOs<sub>2</sub>O<sub>6</sub>, $`+1.98`$% for RbOs<sub>2</sub>O<sub>6</sub>, and $`+2.00`$% for CsOs<sub>2</sub>O<sub>6</sub>. Although the compound with Na has not been synthesizedโprobably because of its small sizeโwe have included it in our study to better identify the trends followed by the different properties upon $`A`$ substitution. With respect to the other parameters in Table 1, unfortunately the only experimental values reported are those of Brรผhwiler et al.bruehwiler04 for RbOs<sub>2</sub>O<sub>6</sub>. In this case, the difference of the calculated Os-O bond length from experiment is $`+1.8\%`$, while the difference for the internal parameter, and the O-Os-O and Os-O-Os angles are, notably, all below 0.1%. From Table 1, it is clear that the angle in the Os-O chains increases with decreasing $`T_c`$. Assuming a phonon mediating pairing, it was suggestedbruehwiler04 that this angle plays a role in determining $`T_c`$. Future studies of the phonon spectra of these materials should allow one to verify this interesting point. In Table 1, we also give the calculated bulk moduli, which are necessary to calculate volume changes under pressure.
## III Electronic structure and properties
### III.1 Band structure and density of states
As we reported previously in the case of KOs<sub>2</sub>O<sub>6</sub>,saniz04 we find that in all the $`\beta `$-pyrochlores the band structure around $`E_\mathrm{F}`$ is given by a manifold of twelve bands arising mainly from Os $`5d`$ states and O $`2p`$ states. The dispersion of the bands is generally the same for the different compounds, but there is an important difference near $`E_\mathrm{F}`$, which is that the vHS near the center of the $`\mathrm{\Gamma }`$-$`L`$ line moves up closer to the Fermi level as the size of $`A`$ increases. This is clearly illustrated in Fig. 1, where we compare the energy bands of the different compounds for k-points along the $`\mathrm{\Gamma }`$-$`L`$. The consequence of this for DOS is also intriguing. Indeed, although the peak due to the vHS tends to decrease for the cases with a larger $`A`$ ion, $`N(E_\mathrm{F})`$ increases because the vHS is closer to $`E_\mathrm{F}`$. This is clear from Fig. 2, where we show a close-up of the total DOS around $`E_\mathrm{F}`$ for the different compounds. In Table 2, where we list the total DOS at $`E_\mathrm{F}`$, as well as the muffin-tin sphere projected DOS for O and Os. For reference, we also indicate the values of the bare band Sommerfeld coefficient and of the band Pauli paramagnetic susceptibility. Comparing with measurements on powder samples, the specific heat mass-enhancements, $`\gamma _{\mathrm{exp}}/\gamma _\mathrm{b}`$, appear to be 3.3 for KOs<sub>2</sub>O<sub>6</sub> ($`\gamma _{\mathrm{exp}}=19`$ mJ/K<sup>2</sup> mol Os)hiroi04 and 3.7 for RbOs<sub>2</sub>O<sub>6</sub> ($`\gamma _{\mathrm{exp}}=44`$ mJ/K<sup>2</sup> mol Os).bruehwiler05 If the results reported by Muramatsu et al. are used for RbOs<sub>2</sub>O<sub>6</sub> and CsOs<sub>2</sub>O<sub>6</sub> (both with $`\gamma 20`$ mJ/K<sup>2</sup> mol Os)muramatsu04 one finds mass enhancements of 3.4 and 3.2, respectively. We note, however, that very recently Hiroi and co-workers reportedhiroi05 specific heat measurements on a single crystal sample of KOs<sub>2</sub>O<sub>6</sub>, which, making an estimate similar to the powder sample case, yields a surprising $`\gamma _{\mathrm{exp}}`$=64.8 mJ/K<sup>2</sup> mol Os. This results in an unusually large $`\gamma _{\mathrm{exp}}/\gamma _\mathrm{b}11.4`$.comment1
The effect of hydrostatic pressure is basically to push the eigenvalues around $`E_\mathrm{F}`$ downward. We illustrate this in the case of RbOs<sub>2</sub>O<sub>6</sub> in Fig. 3, where we show a close look at the bands around $`E_\mathrm{F}`$ both for a sample under zero pressure and a sample under simulated pressure such that the change in the lattice constant is $`\mathrm{\Delta }a/a=2`$%. While this corresponds to a relatively large pressure, it shows clearly the effect on the vHS, pushing it away from $`E_\mathrm{F}`$. This naturally leads to a decrease of $`N(E_\mathrm{F})`$. As a further, quantitative illustration, in Table 3 we list $`N(E_\mathrm{F})`$ for various pressures in the case of KOs<sub>2</sub>O<sub>6</sub>.
As one may expect, the above effects are reflected in the topology of the Fermi surface. This is of general interest because of its direct relation to electronic properties and the possible effect of the vHS on quasi-particle lifetimes. The Fermi surface consists of two closed electron-like sheets centered at the $`\mathrm{\Gamma }`$ point, and a third hole-like sheet giving rise to a tubular network. These Fermi surface sheets are shown in Fig. 4, plotted for clarity in the reciprocal unit cell. The tubular network actually does not present any major difference among the compounds considered; as a typical example we show the case of CsOs<sub>2</sub>O<sub>6</sub> in Fig. 4(a). In contrast, the closed shells exhibit a clear difference near the midpoint of the $`\mathrm{\Gamma }`$-$`L`$ line, where the vHS is located. Indeed, in the case of NaOs<sub>2</sub>O<sub>6</sub>, in Fig. 4(b), the two shells show no narrowing of the distance at this point, while the narrowing is obvious in the case of CsOs<sub>2</sub>O<sub>6</sub>, as shown in Fig. 4(c). As discussed above, this is due to the closeness of the vHS to $`E_\mathrm{F}`$ in the latter case. The topology of the Fermi surface is also relevant to the superconducting gap. In relation to this, we note that the multi-band character of the Fermi surface and the different symmetry of its sheets may be of significance to some of the experimental results on the $`\beta `$-pyrochlores. As pointed out above, Koda and co-workers interpret their results on the linear field dependence of the penetration depth, $`\lambda `$, as suggesting a non-conventional pairing mechanism in KOs<sub>2</sub>O<sub>6</sub>,koda04 possibly mediated by magnetic fluctuations. However, in MgB<sub>2</sub>, which is a phonon-mediated superconductor, $`\lambda `$ also exhibits such a linear dependence on the applied field.ohishi03 In the case of MgB<sub>2</sub> this arises because of its two-gap nature, which in turn is due to the particular character of its Fermi surface sheets, which we find akin to the present case to some extent.
### III.2 Superconducting parameters
In the following, we examine the ability of a phonon-mediated pairing scenario to account for the experimental evidence, and, more particularly, the effects of alkali metal substitution and of pressure. To this end, we estimate the electron-phonon coupling constant $`\lambda _{\mathrm{ep}}`$ within the McMillan-Hopfield framework,mcmillan68 ; hopfield69 and the crude rigid-muffin-tin approximation.gaspari72 ; pettifor77 The spherically averaged Hopfield parameter can be writtenskriver85
$$\eta _i=\underset{l}{}M_{i_{l,l+1}}^2\frac{2l+2}{(2l+1)(2l+3)}\left[\frac{N_{i_l}(E_\mathrm{F})N_{i_{l+1}}(E_\mathrm{F})}{N(E_\mathrm{F})}\right],$$
(1)
where $`N_{i_l}(E_\mathrm{F})`$ is the $`l`$-angular momentum DOS projected on the muffin-tin sphere of atom $`i`$; $`M_{i_{l,l+1}}=\varphi _{i_l}\varphi _{i_{l+1}}[(D_{i_l}l)(D_{i_{l+1}}+l+2)+(E_\mathrm{F}V_i)R_i^2]`$ is an electron-phonon matrix element, in terms of the logarithmic derivatives ($`D_{i_l}`$) and the partial wave amplitudes ($`\varphi _{i_l}`$), both evaluated at $`E_\mathrm{F}`$ and at the muffin-tin radius ($`R_i`$); $`V_i`$ is the one-electron potential at $`R_i`$. Then we have $`\lambda _{\mathrm{ep}}=_i\eta _i/\overline{M}\omega _i^2`$, where $`\overline{M}`$ is the average mass.comment2 As is common, the average phonon frequency is estimated in terms of the Debye temperature as $`\omega ^2^{1/2}=0.69\mathrm{\Theta }_\mathrm{D}`$.
Regarding the Debye temperatures, there is unfortunately no experimental information for KOs<sub>2</sub>O<sub>6</sub> and CsOs<sub>2</sub>O<sub>6</sub>. In the case of RbOs<sub>2</sub>O<sub>6</sub>, recent measurements suggest that the specific heat does not follow the usual $`T^3`$ Debye model,comment3 so that it is not clear what effective $`\mathrm{\Theta }_\mathrm{D}`$ is appropriate in our case. With this caveat, in our estimates we use the value $`\mathrm{\Theta }_\mathrm{D}=285`$ K at $`T=6`$ K obtained if a parametrization of the heat capacity in terms of a $`T`$ dependent $`\mathrm{\Theta }_\mathrm{D}`$ is enforced.comment3 In Table 4, we present our results for the $`\eta _i`$ and $`\lambda _{\mathrm{ep}}`$ for the different compounds.comment4
The critical temperature is subsequently calculated with the McMillan-Allen-Dynes equationbennemann72 ; allen75
$$T_c=\frac{\omega ^2^{1/2}}{1.2}\mathrm{exp}\left[\frac{1.04(1+\lambda _{\mathrm{ep}}+\mu _{\mathrm{sp}})}{\lambda _{\mathrm{ep}}(\mu ^{}+\mu _{\mathrm{sp}})(1+0.62\lambda _{\mathrm{ep}})}\right],$$
(2)
where $`\mu ^{}`$ is the Coulomb pseudopotential and $`\mu _{\mathrm{sp}}`$ is an effective electron-spin coupling constant. Note that, lacking the knowledge of the Eliashberg function $`\alpha ^2F(\omega )`$, in lieu of the logarithmic average $`\omega _{\mathrm{ln}}`$ we take the average phonon frequency, which we estimate in terms of $`\mathrm{\Theta }_\mathrm{D}`$, as indicated above.comment5 The Coulomb pseudopotential can be estimated through $`\mu ^{}=0.26n(E_\mathrm{F})/[1+n(E_\mathrm{F})]`$, where $`n(E_\mathrm{F})`$ is the DOS at $`E_\mathrm{F}`$ per eV and atom.bennemann72 To estimate $`\mu _{\mathrm{sp}}`$ we use the expression derived by Doniach and Engelsberg $`\mu _{\mathrm{sp}}3IN(E_\mathrm{F})\mathrm{ln}\{1+0.03IN(E_\mathrm{F})/[1IN(E_\mathrm{F})]\}.`$doniach66 ; comment6 The Stoner parameter in this expression, $`I`$, is calculated following the band formulation of Gunnarssongunnarsson76 and Brooks et al.brooks87 within spin-density-functional theory. More specifically, we have $`I=_in_iI_i,`$ where $`n_i`$ is the number of atoms of type $`i`$, and $`I_i`$ the atomic Stoner parameter written as $`I_i=\widehat{n}_{i_{ll^{}}}J_{i_{ll^{}}}`$. Here $`\widehat{n}_{i_{ll^{}}}=N_{i_l}(E_\mathrm{F})N_{i_l^{}}(E_\mathrm{F})/N_i^2(E_\mathrm{F})`$ and $`J_{i_{ll^{}}}=๐r|K(r)|\varphi _{i_l}^2(r)\varphi _{i_l^{}}^2(r)`$ with, again, the partial wave amplitudes $`\varphi _{i_l}`$ calculated at $`E_\mathrm{F}`$. The exchange-correlation kernel $`K`$ used is the one given by Gunnarsson.gunnarsson76 In our case, the alkali atom contribution is completely negligible and only the diagonal $`l=1`$ term in O and the diagonal $`l=2`$ in Os contribute because of the dominance of the respective partial DOS at $`E_\mathrm{F}`$ (see, e.g., Ref. brooks87, ). We note that our calculations are done taking spin-orbit coupling into account.
Our results for the superconducting parameters, as well as for $`I`$ and the Stoner enhancement factor $`S=1/[1IN(E_\mathrm{F})/2]`$, are given in Table 5. We note first that our calculated $`T_c`$ (last column) follows well the experimental trend, although the range of the experimental $`T_c`$โs is noticeably larger. Clearly, however, a more refined calculation of $`\lambda _{\mathrm{ep}}`$ based on, e.g., the frequency dependence of the Eliashberg function $`\alpha ^2F`$, can easily account for the difference of a fraction to a few K between our results and experiment seen in Table 5. In this regard, Kuneลก and co-workerskunes04 have found that the alkali ions in these materials possess a varying degree of anharmonicity in the potential, which would add to the difference in their $`T_c`$โs. Secondly, we note the significant role of spin fluctuations. Indeed, with $`\mu _{\mathrm{sp}}=0`$ the predicted $`T_c`$ (given as $`T_c^{}`$) is $``$72% (KOs<sub>2</sub>O<sub>6</sub>) to $``$144% (CsOs<sub>2</sub>O<sub>6</sub>) higher. Thus, although the calculated Stoner enhancement factors ($`2<S<3`$) indicate that these systems are not close to a ferromagnetic instability, they are sufficient to produce a significant electron-spin coupling.
Finally, we have calculated the superconducting parameters of KOs<sub>2</sub>O<sub>6</sub> (again with $`\mathrm{\Theta }_\mathrm{D}=285`$ K) under the simulated effect of pressure, to study the initial change of $`T_c`$ with pressure. We have considered pressures up to 1.166 GPa, which is of the order of the pressures used in the experimental report by Muramatsu et al.muramatsu04 Our results are given in Table 6. As in experiment,khasanov04 ; muramatsu04 we see that $`T_c`$ increases with pressure, although $`\lambda _{\mathrm{ep}}`$ decreases. The main reason is that both $`\mu ^{}`$ and $`\mu _{\mathrm{sp}}`$ also decrease, and more importantly the latter than the former. Hence, the initial increase of $`T_c`$ with pressure appears to be due mainly to a decrease of spin fluctuations, driven by the decrease of $`N(E_\mathrm{F})`$ (the Stoner parameter $`I`$ is almost unchanged). Comparing our results with the ratio $`T_c/T_c^0=1.04`$ found at 0.56 GPa by Muramatsu and collaborators,muramatsu04 we see that the change in $`T_c`$ is of the same order and is essentially accounted for. Indeed, if the main cause were to be phononic, i.e., a change in $`\mathrm{\Theta }_\mathrm{D}`$, the latter would have to decrease with pressure, contrary to any likelihood.comment7 Again, we surmise that a more refined calculation can readily account for the somewhat steeper initial increase of $`T_c`$ observed. In their experimental study of the change of $`T_c`$ with pressure in the case of RbOs<sub>2</sub>O<sub>6</sub>, Khasanov and co-workers had already concluded that the observed positive slope at low pressures must arise mainly from electronic contributions, as opposed to phononic ones, although the mechanism behind the effect was still an open question.khasanov04
We note that the very recent report by Muramatsu and collaboratorsmuramatsu05 shows that in the present family of superconductors, after reaching a maximum value, with increasing pressure $`T_c`$ will gradually decrease, falling below its ambient pressure value, until superconductivity is eventually suppressed. To make predictions of $`T_c`$ at those high pressures, however, it would be necessary to have a minimum information on the phonon spectra and how they are affected by pressure. We do not know what would be a sensible value of $`\mathrm{\Theta }_\mathrm{D}`$ to make such estimates of $`T_c`$ within our approach. However, we can try to understand what happens as follows. If the Grรผneisen parameter for KOs<sub>2</sub>O<sub>6</sub> is around 1.8, which is a rough average for most substances,gschneider64 then for a pressure $`P=3.5`$ GPa (corresponding to $`\mathrm{\Delta }V/V3`$%), $`\mathrm{\Theta }_\mathrm{D}301`$ K. The calculated superconducting parameters ($`\mu ^{}=0.085`$, $`\mu _{\mathrm{sp}}=0.048`$, and $`\lambda _{\mathrm{ep}}=0.738`$) then would lead to $`T_c/T_c^00.88`$ (against $`0.66`$ in experiment). This result suggests that the decrease of $`T_c`$ at higher pressures could be understood as reflecting the fact that in that regime the phononic properties become dominant.
It is remarkable, given the approximate nature of our estimate of $`\lambda _{\mathrm{ep}}`$ and $`T_c`$, that our results account rather well for the behavior of $`T_c`$ under substitution of the alkali metal and under applied pressure. We believe this brings strong support for the phonon-mediated pairing scenario. To understand more fully the properties of the $`\beta `$-pyrochlore Os oxides, however, further investigation will be required, experimentally and theoretically. For instance, the different nuclear spin-lattice relaxation rates of the alkali ions in KOs<sub>2</sub>O<sub>6</sub> and RbOs<sub>2</sub>O<sub>6</sub> remain to be clarified. This could be partly due to the rapid variation of the DOS close to $`E_\mathrm{F}`$. Indeed, Fig. 2 shows that $`N(E_\mathrm{F})`$ can change by as much as 50% within an energy range of $`\pm `$25 meV (more so the heavier the alkali ion). A further contribution may come from the changing character of the alkali atom DOS at $`E_\mathrm{F}`$. We find that its $`s`$ character falls from 73% in NaOs<sub>2</sub>O<sub>6</sub>, to 11% in CsOs<sub>2</sub>O<sub>6</sub>, passing by 31% in KOs<sub>2</sub>O<sub>6</sub> and 23% in RbOs<sub>2</sub>O<sub>6</sub> (at the same time its $`p`$ character rises from 16% in NaOs<sub>2</sub>O<sub>6</sub>, to 55% in CsOs<sub>2</sub>O<sub>6</sub>, passing by 31% in KOs<sub>2</sub>O<sub>6</sub> and 38% in RbOs<sub>2</sub>O<sub>6</sub>).
Furthermore, the origin of the unusual behavior of the resistivity in KOs<sub>2</sub>O<sub>6</sub> at low $`T`$ is also not understood, nor is the rather large specific heat mass enhancements $`\gamma _{\mathrm{exp}}/\gamma _\mathrm{b}`$ in all these materials. The electron-phonon and electron-spin coupling constants obtained above are clearly insufficient to account for the observed enhancements of 3โ4. It is possible that the vHS near $`E_\mathrm{F}`$ and the nesting exhibited by the Fermi surfacecomment8 contribute to both the resistivity and the specific heat. The observed enhancements suggest, in our view, that the unusual non-Debye behavior at low temperature of the specific heat mentioned above is a generic property.
###### Acknowledgements.
We thank Lin-Hui Ye, S. H. Rhim, J. B. Ketterson, and W. Halperin for helpful discussions. We are also grateful to B. Barbiellini and G. Grimvall for their suggestions and to B. Battlogg and M. Brรผhwiler for sharing data with us prior to publication. This work was supported by the Department of Energy (under grant No. DE-FG02-88ER 45372/A021 and a computer time grant at the National Energy Research Scientific Computing Center).
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# Characteristic Functions for Ergodic Tuples
## 0. Introduction
If $`Z=_{i=1}^dA_iA_i^{}`$ is a normal, unital, ergodic, completely positive map on $`B()`$, the bounded linear operators on a complex separable Hilbert space, and if there is a (necessarily unique) invariant vector state for $`Z`$, then we also say that $`\underset{ยฏ}{A}=(A_1,\mathrm{},A_d)`$ is a coisometric, ergodic row contraction with a one-dimensional invariant subspace for the adjoints. Precise definitions are given below. This is the main setting to be investigated in this paper.
In Section 1 we give a concise review of a result on the dilations of $`Z`$ obtained by R. Gohm in \[Go04\] in a chapter called โCocycles and Coboundariesโ. There exists a conjugacy between a homomorphic dilation of $`Z`$ and a tensor shift, and we emphasize an explicit infinite product formula that can be obtained for the intertwining unitary. \[Go04\] may also be consulted for connections of this topic to a scattering theory for noncommutative Markov chains by B. Kรผmmerer and H. Maassen (cf. \[KM00\]) and more general for the relevance of this setting in applications.
In this work we are concerned with its relevance in operator theory and correspondingly in Section 2 we shift our attention to the row contraction $`\underset{ยฏ}{A}=(A_1,\mathrm{},A_d)`$. Our starting point has been the observation that the intertwining unitary mentioned above has many similarities with the notion of characteristic function occurring in the theory of functional models of contractions, as initiated by B. Sz.-Nagy and C. Foias (cf. \[NF70, FF90\]). In fact, the center of our work is the commuting diagram 3.3 in Section 3, which connects the results in \[Go04\] mentioned above with the theory of minimal isometric dilations of row contractions by G. Popescu (cf. \[Po89a\]) and shows that the intertwining unitary determines a multi-analytic inner function, in the sense introduced by G. Popescu in \[Po89c, Po95\]. We call this inner function the extended characteristic function of the tuple $`\underset{ยฏ}{A}`$, see Definition 3.3.
Section 4 is concerned with an explicit computation of this inner function. In Section 5 we show that it is an extension of the characteristic function of the $``$-stable part $`\stackrel{}{\underset{ยฏ}{A}}`$ of $`\underset{ยฏ}{A}`$, the latter in the sense of Popescuโs generalization of the Sz.-Nagy-Foias theory to row contractions (cf. \[Po89b\]). This explains why we call our inner function an extended characteristic function. The row contraction $`\underset{ยฏ}{A}`$ is a one-dimensional extension of the $``$-stable row contraction $`\stackrel{}{\underset{ยฏ}{A}}`$, and in our analysis we separate the new part of the characteristic function from the part already given by Popescu.
G. Popescu has shown in \[Po89b\] that for completely non-coisometric tuples, in particular for $``$-stable ones, his characteristic function is a complete invariant for unitary equivalence. In Section 6 we prove that our extended characteristic function does the same for the tuples $`\underset{ยฏ}{A}`$ described above. In this sense it is characteristic. This is remarkable because the strength of Popescuโs definition lies in the completely non-coisometric situation while we always deal with a coisometric tuple $`\underset{ยฏ}{A}`$. The extended characteristic function also does not depend on the choice of the decomposition $`_{i=1}^dA_iA_i^{}`$ of the completely positive map $`Z`$ and hence also characterizes $`Z`$ up to conjugacy. We think that together with its nice properties established earlier this clearly indicates that the extended characteristic function is a valuable tool for classifying and investigating such tuples respectively such completely positive maps.
Section 7 contains a worked example for the constructions in this paper.
## 1. Weak Markov dilations and conjugacy
In this section we give a brief and condensed review of results in \[Go04\], Chapter 2, which will be used in the following and which, as described in the introduction, motivated the investigations documented in this paper. We also introduce notation.
A theory of weak Markov dilations has been developed in \[BP94\]. For a (single) normal unital completely positive map $`Z:B()B()`$, where $`B()`$ consists of the bounded linear operators on a (complex, separable) Hilbert space, it asks for a normal unital $`{}_{}{}^{}`$endomorphism $`\widehat{J}:B(\widehat{})B(\widehat{})`$, where $`\widehat{}`$ is a Hilbert space containing $``$, such that for all $`n`$ and all $`xB()`$
$$Z^n(x)=p_{}\widehat{J}^n(xp_{})|_{}.$$
Here $`p_{}`$ is the orthogonal projection onto $``$. There are many ways to construct $`\widehat{J}`$. In \[Go04\], 2.3, we gave a construction analogous to the idea of โcoupling to a shiftโ used in \[Kรผ85\] for describing quantum Markov processes. This gives rise to a number of interesting problems which remain hidden in other constructions.
We proceed in two steps. First note that there is a Kraus decomposition $`Z(x)=_{i=1}^da_ixa_i^{}`$ with $`(a_i)_{i=1}^dB()`$. Here $`d=\mathrm{}`$ is allowed in which case the sum should be interpreted as a limit in the strong operator topology. Let $`๐ซ`$ be a $`d`$-dimensional Hilbert space with orthonormal basis $`\{ฯต_1,\mathrm{},ฯต_d\}`$, further $`๐ฆ`$ another Hilbert space with a distinguished unit vector $`\mathrm{\Omega }_๐ฆ๐ฆ`$. We identify $``$ with $`\mathrm{\Omega }_๐ฆ๐ฆ`$ and again denote by $`p_{}`$ the orthogonal projection onto $``$. For $`๐ฆ`$ large enough there exists an isometry
$$u:๐ซ๐ฆ\text{s.t.}p_{}u(hฯต_i)=a_i(h),$$
for all $`h,i=1,\mathrm{},d`$, or equivalently,
$$u^{}(h\mathrm{\Omega }_๐ฆ)=\underset{i=1}{\overset{d}{}}a_i^{}(h)ฯต_i.$$
Explicitly, one may take $`๐ฆ=^{d+1}`$ (resp. infinite-dimensional) and identify
$$๐ฆ(\mathrm{\Omega }_๐ฆ)\underset{1}{\overset{d}{}}\underset{1}{\overset{d}{}}.$$
Then, using isometries $`u_1,\mathrm{},u_d:_1^d`$ with orthogonal ranges and such that $`a_i=p_{}u_i`$ for all $`i`$ (for example, such isometries are explicitly constructed in Popescuโs formula for isometric dilations, cf. \[Po89a\] or equation 3.2 in Section 3), we can define
$$u(hฯต_i):=u_i(h)$$
for all $`h,i=1,\mathrm{},d`$ and check that $`u`$ has the desired properties. Now we define a $`{}_{}{}^{}`$homomorphism
$`J:B()`$ $``$ $`B(๐ฆ),`$
$`x`$ $``$ $`u(x\mathrm{๐}_๐ซ)u^{}.`$
It satisfies
$$p_{}J(x)(h\mathrm{\Omega }_๐ฆ)=p_{}u(x\mathrm{๐})u^{}(h\mathrm{\Omega }_๐ฆ)$$
$$=p_{}u(x\mathrm{๐})\left(\underset{i=1}{\overset{d}{}}a_i^{}(h)ฯต_i\right)=\underset{i=1}{\overset{d}{}}a_ixa_i^{}(h)=Z(x)(h),$$
which means that $`J`$ is a kind of first order dilation for $`Z`$.
For the second step we write $`\stackrel{~}{๐ฆ}:=_1^{\mathrm{}}๐ฆ`$ for an infinite tensor product of Hilbert spaces along the sequence $`(\mathrm{\Omega }_๐ฆ)`$ of unit vectors in the copies of $`๐ฆ`$. We have a distinguished unit vector $`\mathrm{\Omega }_{\stackrel{~}{๐ฆ}}`$ and a (kind of) tensor shift
$$R:B(\stackrel{~}{๐ฆ})B(๐ซ\stackrel{~}{๐ฆ}),\stackrel{~}{y}\mathrm{๐}_๐ซ\stackrel{~}{y}.$$
Finally $`\stackrel{~}{}:=\stackrel{~}{๐ฆ}`$ and we define a normal $`{}_{}{}^{}`$endomorphism
$`\stackrel{~}{J}:B(\stackrel{~}{})`$ $``$ $`B(\stackrel{~}{}),`$
$`B()B(\stackrel{~}{๐ฆ})x\stackrel{~}{y}`$ $``$ $`J(x)\stackrel{~}{y}B(๐ฆ)B(\stackrel{~}{๐ฆ}).`$
Here we used von Neumann tensor products and (on the right hand side) a shift identification $`๐ฆ\stackrel{~}{๐ฆ}\stackrel{~}{๐ฆ}`$. We can also write $`\stackrel{~}{J}`$ in the form
$$\stackrel{~}{J}()=u(Id_{}R)()u^{},$$
where $`u`$ is identified with $`u\mathrm{๐}_{\stackrel{~}{๐ฆ}}`$. The natural embedding $`\mathrm{\Omega }_{\stackrel{~}{๐ฆ}}\stackrel{~}{}`$ leads to the restriction $`\widehat{J}:=\stackrel{~}{J}|_\widehat{}`$ with $`\widehat{}:=\overline{\text{span}}_{n0}\stackrel{~}{J}^n(p_{})(\stackrel{~}{})`$, which can be checked to be a normal unital -endomorphism satisfying all the properties of a weak Markov dilation for $`Z`$ described above. See \[Go04\], 2.3.
A Kraus decomposition of $`\widehat{J}`$ can be written as
$$\widehat{J}(x)=\underset{i=1}{\overset{d}{}}t_ixt_i^{},$$
where $`t_iB(\widehat{})`$ is obtained by linear extension of $`\stackrel{~}{๐ฆ}h\stackrel{~}{k}u_i(h)\stackrel{~}{k}=u(hฯต_i)\stackrel{~}{k}(๐ฆ)\stackrel{~}{๐ฆ}\stackrel{~}{๐ฆ}`$. Because $`\widehat{J}`$ is a normal unital $`{}_{}{}^{}`$endomorphism the $`(t_i)_{i=1}^d`$ generate a representation of the Cuntz algebra $`๐ช_d`$ on $`\widehat{}`$ which we called a coupling representation in \[Go04\], 2.4. Note that the tuple $`(t_1,\mathrm{},t_d)`$ is an isometric dilation of the tuple $`(a_1,\mathrm{},a_d)`$, i.e., the $`t_i`$ are isometries with orthogonal ranges and $`p_{}t_i^n|_{}=a_i^n`$ for all $`i=1,\mathrm{},d`$ and $`n`$.
The following multi-index notation will be used frequently in this work. Let $`\mathrm{\Lambda }`$ denote the set $`\{1,2,\mathrm{},d\}.`$ For operator tuples $`(a_1,\mathrm{},a_d),`$ given $`\alpha =(\alpha _1,\mathrm{},\alpha _m)`$ in $`\mathrm{\Lambda }^m`$, $`a_\alpha `$ will stand for the operator $`a_{\alpha _1}a_{\alpha _2}\mathrm{}a_{\alpha _m}`$, $`|\alpha |:=m`$. Further $`\stackrel{~}{\mathrm{\Lambda }}:=_{n=0}^{\mathrm{}}\mathrm{\Lambda }^n`$, where $`\mathrm{\Lambda }^0:=\{0\}`$ and $`a_0`$ is the identity operator. If we write $`a_\alpha ^{}`$ this always means $`(a_\alpha )^{}=a_{\alpha _m}^{}\mathrm{}a_{\alpha _1}^{}`$.
Back to our isometric dilation, it can be checked that
$$\overline{\text{span}}\{t_\alpha h:h,\alpha \stackrel{~}{\mathrm{\Lambda }}\}=\widehat{},$$
which means that we have a minimal isometric dilation, cf. \[Po89a\] or the beginning of Section 3. For more details on the construction above see \[Go04\], 2.3 and 2.4.
Assume now that there is an invariant vector state for $`Z:B()B()`$ given by a unit vector $`\mathrm{\Omega }_{}`$. Equivalent: There is a unit vector $`\mathrm{\Omega }_๐ซ=_{i=1}^d\overline{\omega }_iฯต_i๐ซ`$ such that $`u(\mathrm{\Omega }_{}\mathrm{\Omega }_๐ซ)=\mathrm{\Omega }_{}\mathrm{\Omega }_๐ฆ`$. Also equivalent: For $`i=1,\mathrm{},d`$ we have $`a_i^{}\mathrm{\Omega }_{}=\overline{\omega }_i\mathrm{\Omega }_{}`$. Here $`\omega _i`$ with $`_{i=1}^d|\omega _i|^2=1`$ and we used complex conjugation to get nice formulas later. See \[Go04\], A.5.1, for a proof of the equivalences.
On $`\stackrel{~}{๐ซ}:=_1^{\mathrm{}}๐ซ`$ along the unit vectors $`(\mathrm{\Omega }_๐ซ)`$ in the copies of $`๐ซ`$ we have a tensor shift
$$S:B(\stackrel{~}{๐ซ})B(\stackrel{~}{๐ซ}),\stackrel{~}{y}\mathrm{๐}_๐ซ\stackrel{~}{y}.$$
Its Kraus decomposition is $`S(\stackrel{~}{y})=_{i=1}^ds_i\stackrel{~}{y}s_i^{}`$ with $`s_iB(\stackrel{~}{๐ซ})`$ and $`s_i(\stackrel{~}{k})=ฯต_i\stackrel{~}{k}`$ for $`\stackrel{~}{k}\stackrel{~}{๐ซ}`$ and $`i=1,\mathrm{},d`$. In \[Go04\], 2.5, we obtained an interesting description of the situation when the dilation $`\widehat{J}`$ is conjugate to the shift endomorphism $`S`$. This result will be further analyzed in this paper. We give a version suitable for our present needs but the reader should have no problems to obtain a proof of the following from \[Go04\], 2.5.
###### Theorem 1.1.
Let $`Z:B()B()`$ be a normal unital completely positive map with an invariant vector state $`\mathrm{\Omega }_{},\mathrm{\Omega }_{}`$. Notation as introduced above, $`d2`$. The following assertions are equivalent:
* $`Z`$ is ergodic, i.e., the fixed point space of $`Z`$ consists of multiples of the identity.
* The vector state $`\mathrm{\Omega }_{},\mathrm{\Omega }_{}`$ is absorbing for $`Z`$, i.e., if $`n\mathrm{}`$ then $`\varphi (Z^n(x))\mathrm{\Omega }_{},x\mathrm{\Omega }_{}`$ for all normal states $`\varphi `$ and all $`xB()`$. (In particular, the invariant vector state is unique.)
* $`\widehat{J}`$ and $`S`$ are conjugate, i.e., there exists a unitary $`๐ฐ:\widehat{}\stackrel{~}{๐ซ}`$ such that
$$\widehat{J}(\widehat{x})=๐ฐ^{}S(๐ฐ\widehat{x}๐ฐ^{})๐ฐ.$$
* The $`๐ช_d`$representations corresponding to $`\widehat{J}`$ and $`S`$ are unitarily equivalent, i.e.,
$$๐ฐt_i=s_i๐ฐ\text{for}i=1,\mathrm{},d.$$
An explicit formula can be given for an intertwining unitary as occurring in (c) and (d). If any of the assertions above is valid then the following limit exists strongly,
$$\stackrel{~}{๐ฐ}=\underset{n\mathrm{}}{lim}u_{0n}^{}\mathrm{}u_{01}^{}:\stackrel{~}{๐ฆ}\stackrel{~}{๐ซ},$$
where we used a leg notation, i.e., $`u_{0n}=(Id_{}R)^{n1}(u)`$. In other words $`u_{0n}`$ is $`u`$ acting on $``$ and on the $`n`$th copy of $`๐ซ`$. Further $`\stackrel{~}{๐ฐ}`$ is a partial isometry with initial space $`\widehat{}`$ and final space $`\stackrel{~}{๐ซ}\mathrm{\Omega }_{}\stackrel{~}{๐ซ}\stackrel{~}{๐ซ}`$ and we can define $`๐ฐ`$ as the corresponding restriction of $`\stackrel{~}{๐ฐ}`$.
To illustrate the product formula for $`๐ฐ`$, which will be our main interest in this work, we use it to derive (d).
$$๐ฐt_i(h\stackrel{~}{k})=๐ฐ\left[u(hฯต_i)\stackrel{~}{k}\right]=\underset{n\mathrm{}}{lim}u_{0n}^{}\mathrm{}u_{01}^{}u_{01}(hฯต_i\stackrel{~}{k})$$
$$=\underset{n\mathrm{}}{lim}u_{0n}^{}\mathrm{}u_{02}^{}(hฯต_i\stackrel{~}{k})=s_i๐ฐ(h\stackrel{~}{k}).$$
Let us finally note that Theorem 1.1 is related to the conjugacy results in \[Pow88\] and \[BJP96\]. Compare also Proposition 2.4.
## 2. Ergodic coisometric row contractions
In the previous section we considered a map $`Z:B()B()`$ given by $`Z(x)=_{i=1}^dA_ixA_i^{}`$, where $`(A_i)_{i=1}^dB()`$. We can think of $`(A_i)_{i=1}^d`$ as a $`d`$-tuple $`\underset{ยฏ}{A}=(A_1,\mathrm{},A_d)`$ or (with the same notation) as a linear map
$$\underset{ยฏ}{A}=(A_1,\mathrm{},A_d):\underset{i=1}{\overset{d}{}}.$$
(Concentrating now on the tuple we have changed to capital letters $`A`$. We will sometimes return to lower case letters $`a`$ when we want to emphasize that we are in the (tensor product) setting of Section 1.) We have the following dictionary.
$`Z(\mathrm{๐})\mathrm{๐}`$ $``$ $`{\displaystyle \underset{i=1}{\overset{d}{}}}A_iA_i^{}\mathrm{๐}`$
$``$ $`\underset{ยฏ}{A}\text{is a contraction}`$
$`Z(\mathrm{๐})=\mathrm{๐}`$ $``$ $`{\displaystyle \underset{i=1}{\overset{d}{}}}A_iA_i^{}=\mathrm{๐}`$
$`\left(Z\text{is called unital}\right)`$ $`\left(\underset{ยฏ}{A}\text{is called coisometric}\right)`$
$`\mathrm{\Omega }_{},\mathrm{\Omega }_{}=\mathrm{\Omega }_{},Z()\mathrm{\Omega }_{}`$ $``$ $`A_i^{}\mathrm{\Omega }_{}=\overline{\omega }_i\mathrm{\Omega }_{},\omega _i,{\displaystyle \underset{i=1}{\overset{d}{}}}|\omega _i|^2=1`$
$`\left(\text{ invariant vector state}\right)`$ $`\left(\text{ common eigenvector for adjoints}\right)`$
$`Z\text{ergodic}`$ $``$ $`\{A_i,A_i^{}\}^{}=\mathbf{\hspace{0.17em}1}`$
$`\left(\text{trivial fixed point space}\right)`$ $`\left(\text{trivial commutant}\right)`$
The converse of the implication at the end of the dictionary is not valid. This is related to the fact that the fixed point space of a completely positive map is not always an algebra. Compare the detailed discussion of this phenomenon in \[BJKW00\].
By a slight abuse of language we call the tuple (or row contraction) $`\underset{ยฏ}{A}=(A_1,\mathrm{},A_d)`$ ergodic if the corresponding map $`Z`$ is ergodic. With this terminology we can interpret Theorem 1.1 as a result about ergodic coisometric row contractions $`\underset{ยฏ}{A}`$ with a common eigenvector $`\mathrm{\Omega }_{}`$ for the adjoints $`A_i^{}`$. This will be examined starting with Section 3. To represent these objects more explicitly let us write $`\stackrel{}{}:=\mathrm{\Omega }_{}`$. With respect to the decomposition $`=\mathrm{\Omega }_{}\stackrel{}{}`$ we get $`2\times 2`$ block matrices
$$A_i\left(\begin{array}{cc}\omega _i& 0\\ |\mathrm{}_i& \AA _i\end{array}\right),A_i^{}\left(\begin{array}{cc}\overline{\omega }_i& \mathrm{}_i|\\ 0& \AA _i^{}\end{array}\right).$$
(2.1)
Here $`\AA _iB(\stackrel{}{})`$ and $`\mathrm{}_i\stackrel{}{}`$. For the off-diagonal terms we used a Dirac notation that should be clear without further comments.
Note that the case $`d=1`$ is rather uninteresting in this setting because if $`A`$ is a coisometry with block matrix $`\left(\begin{array}{cc}\omega & 0\\ |\mathrm{}& \AA \end{array}\right)`$ then because
$$\left(\begin{array}{cc}1& 0\\ 0& \mathrm{๐}\end{array}\right)=AA^{}\left(\begin{array}{cc}|\omega |^2& \omega \mathrm{}|\\ \overline{\omega }|\mathrm{}& |\mathrm{}\mathrm{}|+\AA \AA ^{}\end{array}\right)$$
we always have $`\mathrm{}=0`$. But for $`d2`$ there are many interesting examples arising from unital ergodic completely positive maps with invariant vector states. See Section 1 and also Section 7 for an explicit example. We always assume $`d2`$.
###### Proposition 2.1.
A coisometric row contraction $`\underset{ยฏ}{A}=(A_1,\mathrm{},A_d)`$ is ergodic with common eigenvector $`\mathrm{\Omega }_{}`$ for the adjoints $`A_1^{},\mathrm{},A_d^{}`$ if and only if $`\stackrel{}{}`$ is invariant for $`A_1,\mathrm{},A_d`$ and the restricted row contraction $`(\AA _1,\mathrm{},\AA _d)`$ on $`\stackrel{}{}`$ is $``$-stable, i.e., for all $`h\stackrel{}{}`$
$$\underset{n\mathrm{}}{lim}\underset{|\alpha |=n}{}\AA _\alpha ^{}h^2=0.$$
Here we used the multi-index notation introduced in Section 1. Note that $``$-stable tuples are also called pure, we prefer the terminology from \[FF90\].
###### Proof.
It is clear that $`\mathrm{\Omega }_{}`$ is a common eigenvector for the adjoints if and only if $`\stackrel{}{}`$ is invariant for $`A_1,\mathrm{},A_d`$. Let $`Z()=_{i=1}^dA_iA_i^{}`$ be the associated completely positive map. With $`q:=\mathrm{๐}|\mathrm{\Omega }_{}\mathrm{\Omega }_{}|`$, the orthogonal projection onto $`\stackrel{}{}`$, and by using $`qA_iq=A_iq\AA _i`$ for all $`i`$, we get
$$Z^n(q)=\underset{|\alpha |=n}{}A_\alpha qA_\alpha ^{}=\underset{|\alpha |=n}{}\AA _\alpha \AA _\alpha ^{}$$
and thus for all $`h\stackrel{}{}`$
$$\underset{|\alpha |=n}{}\AA _\alpha ^{}h^2=h,Z^n(q)h.$$
Now it is well known that ergodicity of $`Z`$ is equivalent to $`Z^n(q)0`$ for $`n\mathrm{}`$ in the weak operator topology. See \[GKL06\], Prop. 3.2. This completes the proof. $`\mathrm{}`$
###### Remark 2.2.
Given a coisometric row contraction $`\underset{ยฏ}{a}=(a_1,\mathrm{},a_d)`$ we also have the isometry $`u:๐ซ๐ฆ`$ from Section 1. We introduce the linear map $`a:๐ซB(),ka_k`$ defined by
$$a_k^{}(h)k:=(\mathrm{๐}_{}|kk|)u^{}(h\mathrm{\Omega }_๐ฆ).$$
Compare \[Go04\], A.3.3. In particular $`a_i=a_{ฯต_i}`$ for $`i=1,\mathrm{},d`$, where $`\{ฯต_1,\mathrm{},ฯต_d\}`$ is the orthonormal basis of $`๐ซ`$ used in the definition of $`u`$. Arvesonโs metric operator spaces, cf. \[Ar03\], give a conceptual foundation for basis transformations in the operator space linearly spanned by the $`a_i`$. Similarly, in our formalism a unitary in $`B(๐ซ)`$ transforms $`\underset{ยฏ}{a}=(a_1,\mathrm{},a_d)`$ into another tuple $`\underset{ยฏ}{a}^{}=(a_1^{},\mathrm{},a_d^{})`$. If $`\mathrm{\Omega }_{}`$ is a common eigenvector for the adjoints $`a_i^{}`$ then $`\mathrm{\Omega }_{}`$ is also a common eigenvector for the adjoints $`(a_i^{})^{}`$ but of course the eigenvalues are transformed to another tuple $`\underset{ยฏ}{\omega }^{}=(\omega _1^{},\mathrm{},\omega _d^{})`$. We should consider the tuples $`\underset{ยฏ}{a}`$ and $`\underset{ยฏ}{a}^{}`$ to be essentially the same. This also means that the complex numbers $`\omega _i`$ are not particularly important and they should not play a role in classification. They just reflect a certain choice of orthonormal basis in the relevant metric operator space. Independent of basis transformations is the vector $`\mathrm{\Omega }_๐ซ=_{i=1}^d\overline{\omega }_iฯต_i๐ซ`$ satisfying $`u(\mathrm{\Omega }_{}\mathrm{\Omega }_๐ซ)=\mathrm{\Omega }_{}\mathrm{\Omega }_๐ฆ`$ (see Section 1) and the operator $`a_{\mathrm{\Omega }_๐ซ}=_{i=1}^d\overline{\omega }_ia_i`$.
For later use we show
###### Proposition 2.3.
Let $`\underset{ยฏ}{A}=(A_1,\mathrm{},A_d)`$ be an ergodic coisometric row contraction such that $`A_i^{}\mathrm{\Omega }_{}=\overline{\omega }_i\mathrm{\Omega }_{}`$ for all $`i`$, further $`A_{\mathrm{\Omega }_๐ซ}:=_{i=1}^d\overline{\omega }_iA_i`$. Then for $`n\mathrm{}`$ in the strong operator topology
$$(A_{\mathrm{\Omega }_๐ซ}^{})^n|\mathrm{\Omega }_{}\mathrm{\Omega }_{}|.$$
###### Proof.
We use the setting of Section 1 to be able to apply Theorem 1.1. From $`u^{}(h\mathrm{\Omega }_๐ฆ)=_{i=1}^da_i^{}(h)ฯต_i`$ we obtain
$$u^{}(h\mathrm{\Omega }_๐ฆ)=a_{\mathrm{\Omega }_๐ซ}^{}(h)\mathrm{\Omega }_๐ซh^{}$$
with $`h^{}\mathrm{\Omega }_๐ซ^{}`$. Assume that $`h\stackrel{}{}`$. Because $`u^{}`$ is isometric on $`\mathrm{\Omega }_๐ฆ`$ we conclude that
$$u^{}(\mathrm{\Omega }_{}\mathrm{\Omega }_๐ฆ)=\mathrm{\Omega }_{}\mathrm{\Omega }_๐ซu^{}(h\mathrm{\Omega }_๐ฆ)$$
(2.2)
and thus also $`a_{\mathrm{\Omega }_๐ซ}^{}(h)\stackrel{}{}`$. In other words,
$$a_{\mathrm{\Omega }_๐ซ}^{}(\stackrel{}{})\stackrel{}{}.$$
Let $`q_n`$ be the orthogonal projection from $`_1^n๐ซ`$ onto $`\mathrm{\Omega }_{}_1^n๐ซ`$. From Theorem 1.1 it follows that
$$(\mathrm{๐}q_n)u_{0n}^{}\mathrm{}u_{01}^{}(h\underset{1}{\overset{n}{}}\mathrm{\Omega }_๐ฆ)0(n\mathrm{}).$$
On the other hand, by iterating the formula from the beginning,
$$u_{0n}^{}\mathrm{}u_{01}^{}(h\underset{1}{\overset{n}{}}\mathrm{\Omega }_๐ฆ)=\left((a_{\mathrm{\Omega }_๐ซ}^{})^n(h)\underset{1}{\overset{n}{}}\mathrm{\Omega }_๐ซ\right)h^{}$$
with $`h^{}(_1^n\mathrm{\Omega }_๐ซ)^{}`$. It follows that also
$$(\mathrm{๐}q_n)\left((a_{\mathrm{\Omega }_๐ซ}^{})^n(h)\underset{1}{\overset{n}{}}\mathrm{\Omega }_๐ซ\right)0.$$
But from $`a_{\mathrm{\Omega }_๐ซ}^{}(\stackrel{}{})\stackrel{}{}`$ we have $`q_n\left((a_{\mathrm{\Omega }_๐ซ}^{})^n(h)_1^n\mathrm{\Omega }_๐ซ\right)=0`$ for all $`n`$. We conclude that $`(a_{\mathrm{\Omega }_๐ซ}^{})^n(h)0`$ for $`n\mathrm{}`$. Further
$$a_{\mathrm{\Omega }_๐ซ}^{}\mathrm{\Omega }_{}=\underset{i=1}{\overset{d}{}}\omega _ia_i^{}\mathrm{\Omega }_{}=\underset{i=1}{\overset{d}{}}\omega _i\overline{\omega }_i\mathrm{\Omega }_{}=\mathrm{\Omega }_{},$$
and the proposition is proved. $`\mathrm{}`$
The following proposition summarizes some well known properties of minimal isometric dilations and associated Cuntz algebra representations.
###### Proposition 2.4.
Suppose $`\underset{ยฏ}{A}`$ is a coisometric tuple on $``$ and $`\underset{ยฏ}{V}`$ is its minimal isometric dilation. Assume $`\mathrm{\Omega }_{}`$ is a distinguished unit vector in $``$ and $`\underset{ยฏ}{\omega }=(\omega _1,\mathrm{},\omega _d)^d,_i|\omega _i|^2=1`$. Then the following are equivalent.
1. $`\underset{ยฏ}{A}`$ is ergodic and $`A_i^{}\mathrm{\Omega }_{}=\overline{\omega }_i\mathrm{\Omega }_{}`$ for all $`i`$.
2. $`\underset{ยฏ}{V}`$ is ergodic and $`V_i^{}\mathrm{\Omega }_{}=\overline{\omega }_i\mathrm{\Omega }_{}`$ for all $`i`$.
3. $`V_i^{}\mathrm{\Omega }_{}=\overline{\omega }_i\mathrm{\Omega }_{}`$ and $`\underset{ยฏ}{V}`$ generates the GNS-representation of the Cuntz algebra $`๐ช_d=C^{}\{g_1,\mathrm{},g_d\}`$ ($`g_i`$ its abstract generators) with respect to the Cuntz state which maps
$$g_\alpha g_\beta ^{}\omega _\alpha \overline{\omega }_\beta ,\alpha ,\beta \stackrel{~}{\mathrm{\Lambda }}.$$
Cuntz states are pure and the corresponding GNS-representations are irreducible.
This Proposition clearly follows from Theorem 5.1 of \[BJKW00\], Theorem 3.3 and Theorem 4.1 of \[BJP96\]. Note that in Theorem 1.1(d) we already saw a concrete version of the corresponding Cuntz algebra representation.
## 3. A new characteristic function
First we recall some more details of the theory of minimal isometric dilations for row contractions (cf. \[Po89a\]) and introduce further notation.
The full Fock space over $`^d`$ ($`d2`$) denoted by $`\mathrm{\Gamma }(^d)`$ is
$$\mathrm{\Gamma }(^d):=^d(^d)^^2\mathrm{}(^d)^^m\mathrm{}.$$
$`10\mathrm{}`$ is called the vacuum vector. Let $`\{e_1,\mathrm{},e_d\}`$ be the standard orthonormal basis of $`^d`$. Recall that we include $`d=\mathrm{}`$ in which case $`^d`$ stands for a complex separable Hilbert space of infinite dimension. For $`\alpha \stackrel{~}{\mathrm{\Lambda }}`$, $`e_\alpha `$ will denote the vector $`e_{\alpha _1}e_{\alpha _2}\mathrm{}e_{\alpha _m}`$ in the full Fock space $`\mathrm{\Gamma }(^d)`$ and $`e_0`$ will denote the vacuum vector. Then the (left) creation operators $`L_i`$ on $`\mathrm{\Gamma }(^d)`$ are defined by
$$L_ix=e_ix$$
for $`1id`$ and $`x\mathrm{\Gamma }(^d).`$ The row contraction $`\underset{ยฏ}{L}=(L_1,\mathrm{},L_d)`$ consists of isometries with orthogonal ranges.
Let $`\underset{ยฏ}{T}=(T_1,\mathrm{},T_d)`$ be a row contraction on a Hilbert space $``$. Treating $`\underset{ยฏ}{T}`$ as a row operator from $`_{i=1}^d`$ to $`,`$ define $`D_{}:=(\mathrm{๐}\underset{ยฏ}{T}\underset{ยฏ}{T}^{})^{\frac{1}{2}}:`$ and $`D:=(\mathrm{๐}\underset{ยฏ}{T}^{}\underset{ยฏ}{T})^{\frac{1}{2}}:_{i=1}^d_{i=1}^d`$. This implies that
$$D_{}=(\mathrm{๐}\underset{i=1}{\overset{d}{}}T_iT_i^{})^{\frac{1}{2}},D=(\delta _{ij}\mathrm{๐}T_i^{}T_j)_{d\times d}^{\frac{1}{2}}.$$
(3.1)
Observe that $`\underset{ยฏ}{T}D^2=D_{}^2\underset{ยฏ}{T}`$ and hence $`\underset{ยฏ}{T}D=D_{}\underset{ยฏ}{T}.`$ Let $`๐:=\text{Range }D`$ and $`๐_{}:=\text{Range }D_{}.`$ Popescu in \[Po89a\] gave the following explicit presentation of the minimal isometric dilation of $`\underset{ยฏ}{T}`$ by $`\underset{ยฏ}{V}`$ on $`(\mathrm{\Gamma }(^d)๐)`$,
$$V_i(h\underset{\alpha \stackrel{~}{\mathrm{\Lambda }}}{}e_\alpha d_\alpha )=T_ih[e_0D_ih+e_i\underset{\alpha \stackrel{~}{\mathrm{\Lambda }}}{}e_\alpha d_\alpha ]$$
(3.2)
for $`h`$ and $`d_\alpha ๐.`$ Here $`D_ih:=D(0,\mathrm{},0,h,0,\mathrm{},0)`$ and $`h`$ is embedded at the $`i^{th}`$ component.
In other words, the $`V_i`$ are isometries with orthogonal ranges such that $`T_i^{}=V_i^{}|_{}`$ for $`i=1,\mathrm{},d`$ and the spaces $`V_\alpha `$ with $`\alpha \stackrel{~}{\mathrm{\Lambda }}`$ together span the Hilbert space on which the $`V_i`$ are defined. It is an important fact, which we shall use repeatedly, that such minimal isometric dilations are unique up to unitary equivalence (cf. \[Po89a\]).
Now, as in Section 2, let $`\underset{ยฏ}{A}=(A_1,\mathrm{},A_d),A_iB()`$, be an ergodic coisometric tuple with $`A_i^{}\mathrm{\Omega }_{}=\overline{\omega }_i\mathrm{\Omega }_{}`$ for some unit vector $`\mathrm{\Omega }_{}`$ and some $`\underset{ยฏ}{\omega }^d,_i|\omega _i|^2=1`$. Let $`\underset{ยฏ}{V}=(V_1,\mathrm{},V_d)`$ be the minimal isometric dilation of $`\underset{ยฏ}{A}`$ given by Popescuโs construction (see equation 3.2) on $`\left(\mathrm{\Gamma }(^d)๐_A\right)`$. Because $`A_i^{}=V_i^{}|_{}`$ we also have $`V_i^{}\mathrm{\Omega }_{}=\overline{\omega }_i\mathrm{\Omega }_{}`$ and because $`\underset{ยฏ}{V}`$ generates an irreducible $`๐ช_d`$representation (Proposition 2.4), we see that $`\underset{ยฏ}{V}`$ is also a minimal isometric dilation of $`\underset{ยฏ}{\omega }:^d`$. In fact, we can think of $`\underset{ยฏ}{\omega }`$ as the most elementary example of a tuple with all the properties stated for $`\underset{ยฏ}{A}`$. Let $`\underset{ยฏ}{\overset{~}{V}}=(\stackrel{~}{V}_1,\mathrm{},\stackrel{~}{V}_d)`$ be the minimal isometric dilation of $`\underset{ยฏ}{\omega }`$ given by Popescuโs construction on $`(\mathrm{\Gamma }(^d)๐_\omega )`$.
Because $`\underset{ยฏ}{A}`$ is coisometric it follows from equation 3.1 that $`D`$ is in fact a projection and hence $`D=(\delta _{ij}\mathrm{๐}A_i^{}A_j)_{d\times d}.`$ We infer that $`D(A_1^{},\mathrm{},A_d^{})^T=0`$, where $`T`$ stands for transpose. Applied to $`\underset{ยฏ}{\omega }`$ instead of $`\underset{ยฏ}{A}`$ this shows that $`D_\omega =(\mathrm{๐}|\overline{\underset{ยฏ}{\omega }}\overline{\underset{ยฏ}{\omega }}|)`$ and
$$๐_\omega (\overline{\omega }_1,\mathrm{},\overline{\omega }_d)^T=^d,$$
where $`\overline{\underset{ยฏ}{\omega }}=(\overline{\omega }_1,\mathrm{},\overline{\omega }_d)`$.
###### Remark 3.1.
Because $`\mathrm{\Omega }_{}`$ is cyclic for $`\{V_\alpha ,\alpha \stackrel{~}{\mathrm{\Lambda }}\}`$ we have
$$\overline{\text{span}}\{A_\alpha \mathrm{\Omega }_{}:\alpha \stackrel{~}{\mathrm{\Lambda }}\}=\overline{\text{span}}\{p_{}V_\alpha \mathrm{\Omega }_{}:\alpha \stackrel{~}{\mathrm{\Lambda }}\}=.$$
Using the notation from equation 2.1 this further implies that
$$\overline{\text{span}}\{\AA _\alpha l_i:\alpha \stackrel{~}{\mathrm{\Lambda }},1id\}=\stackrel{}{}.$$
As minimal isometric dilations of the tuple $`\underset{ยฏ}{\omega }`$ are unique up to unitary equivalence, there exists a unitary
$$W:(\mathrm{\Gamma }(^d)๐_A)(\mathrm{\Gamma }(^d)๐_\omega ),$$
such that $`WV_i=\stackrel{~}{V}_iW`$ for all $`i.`$
After showing the existence of $`W`$ we now proceed to compute $`W`$ explicitly. For $`\underset{ยฏ}{A}`$, by using Popescuโs construction, we have its minimal isometric dilation $`\underset{ยฏ}{V}`$ on $`(\mathrm{\Gamma }(^d)๐_A).`$ Another way of constructing a minimal isometric dilation $`\underset{ยฏ}{t}`$ of $`\underset{ยฏ}{a}`$ was demonstrated in Section 1 on the space $`\widehat{}`$ (obtained by restricting to the minimal subspace of $`\stackrel{~}{๐ฆ}`$ with respect to $`\underset{ยฏ}{t}`$). Identifying $`\underset{ยฏ}{A}`$ and $`\underset{ยฏ}{a}`$ on the Hilbert space $``$ there is a unitary $`\mathrm{\Gamma }_A:\widehat{}(\mathrm{\Gamma }(^d)๐_A)`$ which is the identity on $``$ and satisfies $`V_i\mathrm{\Gamma }_A=\mathrm{\Gamma }_At_i`$.
By Theorem 1.1(d) the tuple $`\underset{ยฏ}{s}`$ on $`\stackrel{~}{๐ซ}`$ arising from the tensor shift is unitarily equivalent to $`\underset{ยฏ}{t}`$ (resp. $`\underset{ยฏ}{V}`$), explicitly $`๐ฐt_i=s_i๐ฐ`$ for all $`i`$. An alternative viewpoint on the existence of $`๐ฐ`$ is to note that $`\underset{ยฏ}{s}`$ is a minimal isometric dilation of $`\underset{ยฏ}{\omega }.`$ In fact, $`s_i^{}\mathrm{\Omega }_{\stackrel{~}{๐ซ}}=ฯต_i,\mathrm{\Omega }_๐ซ\mathrm{\Omega }_{\stackrel{~}{๐ซ}}=\overline{\omega }_i\mathrm{\Omega }_{\stackrel{~}{๐ซ}}`$ for all $`i`$. Hence there is also a unitary $`\mathrm{\Gamma }_\omega :\stackrel{~}{๐ซ}(\mathrm{\Gamma }(^d)๐_\omega )`$ with $`\mathrm{\Gamma }_\omega \mathrm{\Omega }_{\stackrel{~}{๐ซ}}=1`$ which satisfies $`\stackrel{~}{V}_i\mathrm{\Gamma }_\omega =\mathrm{\Gamma }_\omega s_i`$.
###### Remark 3.2.
It is possible to describe $`\mathrm{\Gamma }_\omega `$ in an explicit way and in doing so to construct an interesting and natural (unitary) identification of $`_1^{\mathrm{}}^d`$ and $`(\mathrm{\Gamma }(^d)^{d1})`$. In fact, recall (from Section 1) that $`\stackrel{~}{๐ซ}=_1^{\mathrm{}}๐ซ`$ and the space $`๐ซ`$ is nothing but a $`d`$-dimensional Hilbert space. Hence we can identify
$$^d๐ซ=\stackrel{}{๐ซ}\mathrm{\Omega }_๐ซ๐_\omega \overline{\underset{ยฏ}{\omega }}^T^{d1}$$
In this identification the orthonormal basis $`(ฯต_i)_{i=1}^d`$ of $`๐ซ`$ goes to the canonical basis $`(e_i)_{i=1}^d`$ of $`^d`$, in particular the vector $`\mathrm{\Omega }_๐ซ=_i\overline{\omega }_iฯต_i`$ goes to $`\overline{\underset{ยฏ}{\omega }}^T=(\overline{\omega }_1,\mathrm{},\overline{\omega }_d)^T`$ and we have $`\stackrel{}{๐ซ}๐_\omega `$. Then we can write
$`\mathrm{\Gamma }_\omega :\mathrm{\Omega }_{\stackrel{~}{๐ซ}}`$ $``$ $`1,`$
$`k\mathrm{\Omega }_{\stackrel{~}{๐ซ}}`$ $``$ $`e_0k`$
$`ฯต_\alpha k\mathrm{\Omega }_{\stackrel{~}{๐ซ}}`$ $``$ $`e_\alpha k,`$
where $`k\stackrel{}{๐ซ},\alpha \stackrel{~}{\mathrm{\Lambda }},ฯต_\alpha =ฯต_{\alpha _1}\mathrm{}ฯต_{\alpha _n}_1^n๐ซ`$ (the first $`n`$ copies of $`๐ซ`$ in the infinite tensor product $`\stackrel{~}{๐ซ}`$), $`e_\alpha =e_{\alpha _1}\mathrm{}e_{\alpha _n}\mathrm{\Gamma }(^d)`$ as usual. It is easily checked that $`\mathrm{\Gamma }_\omega `$ given in this way indeed satisfies the equation $`\stackrel{~}{V}_i\mathrm{\Gamma }_\omega =\mathrm{\Gamma }_\omega s_i`$ (for all $`i`$), which may thus be seen as the abstract characterization of this unitary map (together with $`\mathrm{\Gamma }_\omega \mathrm{\Omega }_{\stackrel{~}{๐ซ}}=1`$).
Summarizing, for $`i=1,\mathrm{},d`$
$`V_i\mathrm{\Gamma }_A=\mathrm{\Gamma }_At_i,๐ฐt_i=s_i๐ฐ,\stackrel{~}{V}_i\mathrm{\Gamma }_\omega =\mathrm{\Gamma }_\omega s_i`$
and we have the commuting diagram
(3.3)
From the diagram we get
$$W=\mathrm{\Gamma }_\omega ๐ฐ\mathrm{\Gamma }_A^1.$$
Combined with the equations above this yields $`WV_i=\stackrel{~}{V}_iW`$ and we see that $`W`$ is nothing but the dilations-intertwining map which we have already introduced earlier. Hence $`๐ฐ`$ and $`W`$ are essentially the same thing and for the study of certain problems it may be helpful to switch from one picture to the other.
In the following we analyze $`W`$ to arrive at an interpretation as a new kind of characteristic function. First we have an isometric embedding
$$\widehat{C}:=W|_{}:(\mathrm{\Gamma }(^d)๐_\omega ).$$
(3.4)
Note that $`\widehat{C}\mathrm{\Omega }_{}=W\mathrm{\Omega }_{}=1`$. The remaining part is an isometry
$$M_{\widehat{\mathrm{\Theta }}}:=W|_{\mathrm{\Gamma }(^d)๐_A}:\mathrm{\Gamma }(^d)๐_A\mathrm{\Gamma }(^d)๐_\omega .$$
(3.5)
From equation 3.2 we get for all $`i`$
$$V_i|_{\mathrm{\Gamma }(^d)๐_A}=(L_i\mathrm{๐}_{๐_A}),$$
$$\stackrel{~}{V}_i|_{\mathrm{\Gamma }(^d)๐_\omega }=(L_i\mathrm{๐}_{๐_\omega }),$$
and we conclude that
$$M_{\widehat{\mathrm{\Theta }}}(L_i\mathrm{๐}_{๐_A})=(L_i\mathrm{๐}_{๐_\omega })M_{\widehat{\mathrm{\Theta }}},1id.$$
(3.6)
In other words, $`M_{\widehat{\mathrm{\Theta }}}`$ is a multi-analytic inner function in the sense of \[Po89c, Po95\]. It is determined by its symbol
$$\widehat{\theta }:=W|_{e_0๐_A}:๐_A\mathrm{\Gamma }(^d)๐_\omega ,$$
(3.7)
where we have identified $`e_0๐_A`$ and $`๐_A`$. In other words, we think of the symbol $`\widehat{\theta }`$ as an isometric embedding of $`๐_A`$ into $`\mathrm{\Gamma }(^d)๐_\omega `$.
###### Definition 3.3.
We call $`M_{\widehat{\mathrm{\Theta }}}`$ (or $`\widehat{\theta }`$) the extended characteristic function of the row contraction $`\underset{ยฏ}{A}`$,
See Sections 5 and 6 for more explanation and justification of this terminology.
## 4. Explicit computation of the extended characteristic function
To express the extended characteristic function more explicitly in terms of the tuple $`\underset{ยฏ}{A}`$ we start by defining
$`\widehat{D}_{}:\stackrel{}{}=\mathrm{\Omega }_{}`$ $``$ $`\stackrel{}{๐ซ}=๐ซ\mathrm{\Omega }_๐ซ๐_\omega ,`$ (4.1)
$$h\left(\mathrm{\Omega }_{}|\mathrm{๐}_๐ซ\right)u^{}(h\mathrm{\Omega }_๐ฆ),$$
where $`u:๐ซ๐ฆ`$ is the isometry introduced in Section 1. That indeed the range of $`\widehat{D}_{}`$ is contained in $`\stackrel{}{๐ซ}`$ follows from equation 2.2, i.e., $`u^{}(h\mathrm{\Omega }_๐ฆ)\mathrm{\Omega }_{}\mathrm{\Omega }_๐ซ`$ for $`h\stackrel{}{}`$. With notations from equation 2.1 we can get a more concrete formula.
###### Lemma 4.1.
For all $`h\stackrel{}{}`$ we have $`\widehat{D}_{}(h)=_{i=1}^d\mathrm{}_i,hฯต_i`$.
###### Proof.
$`\left(\mathrm{\Omega }_{}|\mathrm{๐}_๐ซ\right)u^{}(h\mathrm{\Omega }_๐ฆ)=_{i=1}^d\mathrm{\Omega }_{},a_i^{}hฯต_i=_{i=1}^d\mathrm{}_i,hฯต_i.`$ $`\mathrm{}`$
###### Proposition 4.2.
The map $`\widehat{C}:(\mathrm{\Gamma }(^d)๐_\omega )`$ from equation 3.4 is given explicitly by $`\widehat{C}\mathrm{\Omega }_{}=1`$ and for $`h\stackrel{}{}`$ by
$$\widehat{C}h=\underset{\alpha \stackrel{~}{\mathrm{\Lambda }}}{}e_\alpha \widehat{D}_{}\AA _\alpha ^{}h.$$
###### Proof.
As $`W\mathrm{\Omega }_{}=1`$ also $`\widehat{C}\mathrm{\Omega }_{}=1`$. Assume $`h\stackrel{}{}`$. Then
$`u_{01}(h\mathrm{\Omega }_{\stackrel{~}{๐ฆ}})`$ $`=`$ $`{\displaystyle \underset{i}{}}a_i^{}hฯต_i\mathrm{\Omega }_{\stackrel{~}{๐ฆ}}`$
$`=`$ $`{\displaystyle \underset{i}{}}\mathrm{}_i,h\mathrm{\Omega }_{}ฯต_i\mathrm{\Omega }_{\stackrel{~}{๐ฆ}}+{\displaystyle \underset{i}{}}\underset{i}{\overset{}{\stackrel{}{a}}}hฯต_i\mathrm{\Omega }_{\stackrel{~}{๐ฆ}}.`$
Because $`u^{}(\mathrm{\Omega }_{}\mathrm{\Omega }_๐ฆ)=\mathrm{\Omega }_{}\mathrm{\Omega }_๐ซ`$ we obtain (with Lemma 4.1) for the first part
$`\underset{n\mathrm{}}{lim}u_{0n}^{}\mathrm{}u_{02}^{}({\displaystyle \underset{i}{}}\mathrm{}_i,h\mathrm{\Omega }_{}ฯต_i\mathrm{\Omega }_{\stackrel{~}{๐ฆ}})`$
$`=`$ $`{\displaystyle \underset{i}{}}\mathrm{}_i,h\mathrm{\Omega }_{}ฯต_i\mathrm{\Omega }_{\stackrel{~}{๐ซ}}=\mathrm{\Omega }_{}\widehat{D}_{}h\mathrm{\Omega }_{\stackrel{~}{๐ซ}}\widehat{D}_{}h\mathrm{\Omega }_{\stackrel{~}{๐ซ}}\stackrel{~}{๐ซ}.`$
Using the product formula from Theorem 1.1 and iterating the argument above we get
$$\widehat{C}(h)=Wh=\mathrm{\Gamma }_\omega ๐ฐ\mathrm{\Gamma }_A^1(h)$$
$$=\mathrm{\Gamma }_\omega (\widehat{D}_{}h\mathrm{\Omega }_{\stackrel{~}{๐ซ}})+\mathrm{\Gamma }_\omega \underset{n\mathrm{}}{lim}u_{0n}^{}\mathrm{}u_{02}^{}\underset{i}{}\underset{i}{\overset{}{\stackrel{}{a}}}hฯต_i\mathrm{\Omega }_{\stackrel{~}{๐ฆ}}$$
$$=e_0\widehat{D}_{}h+\mathrm{\Gamma }_\omega \underset{n\mathrm{}}{lim}u_{0n}^{}\mathrm{}u_{03}^{}\underset{j,i}{}(\mathrm{}_j,\underset{i}{\overset{}{\stackrel{}{a}}}h\mathrm{\Omega }_{}+\underset{j}{\overset{}{\stackrel{}{a}}}\underset{i}{\overset{}{\stackrel{}{a}}}h)ฯต_iฯต_j\mathrm{\Omega }_{\stackrel{~}{๐ฆ}}$$
$$=e_0\widehat{D}_{}h+\underset{i=1}{\overset{d}{}}e_i\widehat{D}_{}\underset{i}{\overset{}{\stackrel{}{a}}}h+\mathrm{\Gamma }_\omega \underset{n\mathrm{}}{lim}u_{0n}^{}\mathrm{}u_{03}^{}\underset{j,i}{}\underset{j}{\overset{}{\stackrel{}{a}}}\underset{i}{\overset{}{\stackrel{}{a}}}hฯต_iฯต_j\mathrm{\Omega }_{\stackrel{~}{K}}$$
$$=\mathrm{}$$
$$=\underset{|\alpha |<m}{}e_\alpha \widehat{D}_{}\underset{\alpha }{\overset{}{\stackrel{}{a}}}h+\mathrm{\Gamma }_\omega \underset{n\mathrm{}}{lim}u_{0n}^{}\mathrm{}u_{0,m+1}^{}\underset{|\alpha |=m}{}\underset{\alpha }{\overset{}{\stackrel{}{a}}}hฯต_\alpha \mathrm{\Omega }_{\stackrel{~}{๐ฆ}}.$$
From Proposition 2.1 we have $`_{|\alpha |=m}\underset{\alpha }{\overset{}{\stackrel{}{a}}}h^20`$ for $`m\mathrm{}`$ and we conclude that the last term converges to $`0`$. It follows that the series converges and this proves Proposition 4.2. $`\mathrm{}`$
###### Remark 4.3.
Another way to prove Proposition 4.2 for $`h\stackrel{}{}`$ consists in repeatedly applying the formula
$$u^{}(h\mathrm{\Omega }_๐ฆ)=a_{\mathrm{\Omega }_๐ซ}^{}h\mathrm{\Omega }_๐ซ+h^{},h^{}\stackrel{}{๐ซ}$$
to the $`u_{0n}^{}(h\mathrm{\Omega }_๐ฆ)`$ and then using $`(a_{\mathrm{\Omega }_๐ซ}^{})^nh0`$, see Proposition 2.3. This gives some insight how the infinite product in Theorem 1.1 transforms into the infinite sum in Proposition 4.2.
Now we present an explicit computation of the extended characteristic function. One way of writing $`๐_A`$ is
$$๐_A=\overline{\text{span}}\{(V_iA_i)h:i\mathrm{\Lambda },h\}.$$
Let $`d_h^i:=(V_iA_i)h.`$ Then
$$\widehat{\theta }d_h^i=W(V_iA_i)h=\stackrel{~}{V}_i\widehat{C}h\widehat{C}A_ih.$$
Case I: Take $`h=\mathrm{\Omega }_{}.`$
$$\stackrel{~}{V}_i\widehat{C}\mathrm{\Omega }_{}=\stackrel{~}{V}_i1=\omega _i[e_0(\mathrm{๐}|\overline{\underset{ยฏ}{\omega }}\overline{\underset{ยฏ}{\omega }}|)ฯต_i],$$
$$\widehat{C}A_i\mathrm{\Omega }_{}=\omega _i\underset{\alpha }{}e_\alpha \widehat{D}_{}\AA _\alpha ^{}l_i$$
and thus
$`\widehat{\theta }d_\mathrm{\Omega }_{}^i`$ $`=`$ $`e_0[(\mathrm{๐}|\overline{\underset{ยฏ}{\omega }}\overline{\underset{ยฏ}{\omega }}|)ฯต_i\widehat{D}_{}l_i]{\displaystyle \underset{|\alpha |1}{}}e_\alpha \widehat{D}_{}\AA _\alpha ^{}l_i`$
$`=`$ $`e_0[ฯต_i{\displaystyle \underset{j}{}}\overline{\omega }_j\omega _iฯต_j{\displaystyle \underset{j}{}}l_j,l_iฯต_j]{\displaystyle \underset{|\alpha |1}{}}e_\alpha {\displaystyle \underset{j}{}}l_j,\AA _\alpha ^{}l_iฯต_j`$
$`=`$ $`e_0[ฯต_i{\displaystyle \underset{j}{}}(\overline{\omega }_j\omega _i+l_j,l_i)ฯต_j]{\displaystyle \underset{|\alpha |1}{}}e_\alpha {\displaystyle \underset{j}{}}\AA _\alpha l_j,l_iฯต_j`$
$$=e_0[ฯต_i\underset{j}{}A_j\mathrm{\Omega }_{},A_i\mathrm{\Omega }_{}ฯต_j]\underset{|\alpha |1}{}e_\alpha \underset{j}{}\AA _\alpha l_j,l_iฯต_j.$$
(4.2)
Case II: Now let $`h\stackrel{}{}`$. With $`i\mathrm{\Lambda }`$
$$\stackrel{~}{V}_i\widehat{C}h=(L_i\mathrm{๐})\widehat{C}h=\underset{\alpha }{}e_ie_\alpha \widehat{D}_{}\AA _\alpha ^{}h,$$
$$\widehat{C}A_ih=\underset{\beta }{}e_\beta \widehat{D}_{}\AA _\beta ^{}\AA _ih.$$
Finally
$$\widehat{\theta }d_h^i=\underset{\alpha }{}e_ie_\alpha \widehat{D}_{}\AA _\alpha ^{}h\underset{\beta }{}e_\beta \widehat{D}_{}\AA _\beta ^{}\AA _ih$$
$$=e_0\widehat{D}_{}\AA _ih+e_i\underset{\alpha }{}e_\alpha \widehat{D}_{}\AA _\alpha ^{}(\mathrm{๐}\AA _i^{}\AA _i)h+\underset{ji}{}e_j\underset{\alpha }{}e_\alpha \widehat{D}_{}\AA _\alpha ^{}(\AA _j^{}\AA _i)h$$
$$=e_0\widehat{D}_{}\AA _ih+\underset{j=1}{\overset{d}{}}e_j\underset{\alpha \stackrel{~}{\mathrm{\Lambda }}}{}e_\alpha \widehat{D}_{}\AA _\alpha ^{}(\delta _{ji}\mathrm{๐}\AA _j^{}\AA _i)h.$$
(4.3)
## 5. Case II is the characteristic function of $`\stackrel{}{\underset{ยฏ}{A}}`$
In this section we show that case II in the previous section can be identified with the characteristic function of the $``$-stable tuple $`\stackrel{}{\underset{ยฏ}{A}}`$, in the sense introduced by Popescu in \[Po89b\]. This is the reason why we have called $`\widehat{\theta }`$ an extended characteristic function. All information about $`\underset{ยฏ}{A}`$ beyond $`\stackrel{}{\underset{ยฏ}{A}}`$ must be contained in case I.
First recall the theory of characteristic functions for row contractions, as developed by G. Popescu in \[Po89b\], generalizing the theory of B. Sz.-Nagy and C. Foias (cf. \[NF70\]) for single contractions. We only need the results about a $``$-stable tuple $`\stackrel{}{\underset{ยฏ}{A}}=(\AA _1,\mathrm{},\AA _d)`$ on $`\stackrel{}{}`$. In this case, with $`\underset{}{\overset{}{D}}=(\mathrm{๐}\stackrel{}{\underset{ยฏ}{A}}\stackrel{}{\underset{ยฏ}{A}^{}})^{\frac{1}{2}}:\stackrel{}{}\stackrel{}{}`$ and $`\underset{}{\overset{}{๐}}`$ its range, the map
$$\stackrel{}{C}:\stackrel{}{}\mathrm{\Gamma }(^d)\underset{}{\overset{}{๐}}$$
(5.1)
$$h\underset{\alpha \stackrel{~}{\mathrm{\Lambda }}}{}e_\alpha \underset{}{\overset{}{D}}\AA _\alpha ^{}h$$
is an isometry (Popescuโs Poisson kernel). If, as usual, $`\stackrel{}{D}=(\mathrm{๐}\stackrel{}{\underset{ยฏ}{A}^{}}\stackrel{}{\underset{ยฏ}{A}})^{\frac{1}{2}}:_1^d\stackrel{}{}_1^d\stackrel{}{}`$, with $`\stackrel{}{๐}`$ its range, and if $`P_j`$ is the projection onto the $`j`$-th component, then the characteristic function $`\theta _\AA `$ of $`\stackrel{}{\underset{ยฏ}{A}}`$ can be defined as
$$\theta _\AA :\stackrel{}{๐}\mathrm{\Gamma }(^d)\underset{}{\overset{}{๐}}$$
(5.2)
$$fe_0\underset{j=1}{\overset{d}{}}\AA _jP_jf+\underset{j=1}{\overset{d}{}}e_j\underset{\alpha \stackrel{~}{\mathrm{\Lambda }}}{}e_\alpha \underset{}{\overset{}{D}}\AA _\alpha ^{}P_j\stackrel{}{D}f.$$
See \[Po89b\] for details, in particular for the important result that $`\theta _\AA `$ characterizes the $``$-stable tuple $`\stackrel{}{\underset{ยฏ}{A}}`$ up to unitary equivalence.
Now consider again the tuple $`\underset{ยฏ}{A}`$ of the previous section, with extended characteristic function $`\widehat{\theta }`$. From equation 2.1
$$A_i\left(\begin{array}{cc}\omega _i& 0\\ |\mathrm{}_i& \AA _i\end{array}\right),A_i^{}\left(\begin{array}{cc}\overline{\omega }_i& \mathrm{}_i|\\ 0& \AA _i^{}\end{array}\right)$$
and hence
$$A_iA_i^{}=\left(\begin{array}{cc}|\overline{\omega }_i|^2& \overline{\omega }_il_i|\\ |\overline{\omega }_il_i& |l_il_i|+\AA _i\AA _i^{}\end{array}\right).$$
Recall that $`D_{}^2=\mathrm{๐}_iA_iA_i^{}`$ which is $`0`$ as $`\underset{ยฏ}{A}`$ is coisometric. Thus $`_i\overline{\omega }_il_i=0`$ and $`\mathrm{๐}_i\AA _i\AA _i^{}=_i|l_il_i|.`$ The first equation means that $`A_{\mathrm{\Omega }_๐ซ}^{}(\stackrel{}{})\stackrel{}{}`$ and that
$$\widehat{๐}_{}h,\mathrm{\Omega }_๐ซ=\underset{i}{}\mathrm{}_i,hฯต_i,\underset{j}{}\overline{\omega }_jฯต_j=\underset{i}{}\overline{\omega }_i\mathrm{}_i,h=0,$$
which we already know (see 4.1).
The second equation yields
$$\underset{}{\overset{2}{\stackrel{}{D}}}=\mathrm{๐}\underset{i}{}\AA _i\AA _i^{}=\underset{i}{}|l_il_i|.$$
###### Lemma 5.1.
There exists an isometry $`\gamma :\underset{}{\overset{}{๐}}\stackrel{}{๐ซ}๐_\omega `$ defined for $`h\stackrel{}{}`$ as
$$\underset{}{\overset{}{D}}h\underset{i}{}l_i,hฯต_i=\widehat{D}_{}h.$$
###### Proof.
Take $`h\stackrel{}{}.`$ By Lemma 4.1 we have $`\widehat{D}_{}(h)=_{i=1}^d\mathrm{}_i,hฯต_i`$. Now we can compute
$$\widehat{D}_{}h^2=\underset{i}{}l_i,hฯต_i,\underset{j}{}l_j,hฯต_j=\underset{i}{}h,l_il_i,h=h,\underset{}{\overset{2}{\stackrel{}{D}}}h=\underset{}{\overset{}{D}}h^2.$$
Hence $`\gamma :\underset{}{\overset{}{D}}h\widehat{D}_{}h`$ is isometric. $`\mathrm{}`$
###### Theorem 5.2.
Let $`\underset{ยฏ}{A}=(A_1,\mathrm{},A_d),A_iB()`$, be an ergodic coisometric tuple with $`A_i^{}\mathrm{\Omega }_{}=\overline{\omega }_i\mathrm{\Omega }_{}`$ for some unit vector $`\mathrm{\Omega }_{}`$ and some $`\underset{ยฏ}{\omega }^d,_i|\omega _i|^2=1`$. Let $`\widehat{\theta }`$ be the extended characteristic function of $`\underset{ยฏ}{A}`$ and let $`\theta _\AA `$ be the characteristic function of the ($``$-stable) tuple $`\underset{ยฏ}{\AA }`$. For $`h\stackrel{}{}`$
$`\gamma \underset{}{\overset{}{D}}h`$ $`=`$ $`\widehat{D}_{}h,`$
$`(\mathrm{๐}\gamma )\stackrel{}{C}h`$ $`=`$ $`\widehat{C}h,`$
$$(\mathrm{๐}\gamma )\theta _\AA d_h^i=\widehat{\theta }d_h^i\text{for}i\mathrm{\Lambda }.$$
In other words, the part of $`\widehat{\theta }`$ described by case II in the previous section is equivalent to $`\theta _\AA `$.
###### Proof.
We only have to use Lemma 5.1 and compare Proposition 4.2 and equation 5.1 as well as equations 4.3 and 5.2. For the latter note that $`d_h^i=\stackrel{}{D}(0,\mathrm{},0,h,0\mathrm{},0)`$, where $`h`$ is embedded at the $`i`$-th position. Hence
$$\gamma \underset{j}{}\AA _jP_jd_h^i=\gamma \underset{j}{}\AA _jP_j\stackrel{}{D}(0,\mathrm{},0,h,0\mathrm{},0)=\gamma \stackrel{}{\underset{ยฏ}{A}}\stackrel{}{D}(0,\mathrm{},0,h,0\mathrm{},0)$$
$$=\gamma \underset{}{\overset{}{D}}\stackrel{}{\underset{ยฏ}{A}}(0,\mathrm{},0,h,0\mathrm{},0)=\widehat{D}_{}\AA _ih$$
and also
$$P_j\stackrel{}{D}d_h^i=P_j\stackrel{2}{\stackrel{}{D}}(0,\mathrm{},0,h,0\mathrm{},0)=(\delta _{ji}\mathrm{๐}\AA _j^{}\AA _i)h.$$
$`\mathrm{}`$
Of course, Theorem 5.2 explains why we have called $`\widehat{\theta }`$ an extended characteristic function.
## 6. The extended characteristic function is a complete unitary invariant
In this section we prove that the extended characteristic function is a complete invariant with respect to unitary equivalence for the row contractions investigated in this paper. Suppose that $`\underset{ยฏ}{A}=(A_1,\mathrm{},A_d)`$ and $`\underset{ยฏ}{B}=(B_1,\mathrm{},B_d)`$ are ergodic and coisometric row contractions on Hilbert spaces $`_A`$ and $`_B`$ such that $`A_i^{}\mathrm{\Omega }_A=\overline{\omega }_i\mathrm{\Omega }_A`$ and $`B_i^{}\mathrm{\Omega }_B=\overline{\omega }_i\mathrm{\Omega }_B`$ for $`i=1,\mathrm{},d`$, where $`\mathrm{\Omega }_A_A`$ and $`\mathrm{\Omega }_B_B`$ are unit vectors and $`\underset{ยฏ}{\omega }=(\omega _1,\mathrm{},\omega _d)`$ is a tuple of complex numbers. Recall from Remark 2.2 that it is no serious restriction of generality to assume that it is the same tuple of complex numbers in both cases because this can always be achieved by a transformation with a unitary $`d\times d`$matrix (with scalar entries). We will use all the notations introduced earlier with subscripts $`A`$ or $`B`$.
Let us say that the extended characteristic functions $`\widehat{\theta }_A`$ and $`\widehat{\theta }_B`$ are equivalent if there exists a unitary $`V:๐_A๐_B`$ such that $`\widehat{\theta }_A=\widehat{\theta }_BV`$. Note that the ranges of $`\widehat{\theta }_A`$ and $`\widehat{\theta }_B`$ are both contained in $`\mathrm{\Gamma }(^d)๐_\omega `$ and thus this definition makes sense. Let us further say that $`\underset{ยฏ}{A}`$ and $`\underset{ยฏ}{B}`$ are unitarily equivalent if there exists a unitary $`U:_A_B`$ such that $`UA_i=B_iU`$ for $`i=1,\mathrm{},d`$. By ergodicity the unit eigenvector $`\mathrm{\Omega }_A`$ (resp. $`\mathrm{\Omega }_B`$) is determined up to an unimodular constant (see Theorem 1.1(b)) and hence in the case of unitary equivalence we can always modify $`U`$ to satisfy additionally $`U\mathrm{\Omega }_A=\mathrm{\Omega }_B`$.
###### Theorem 6.1.
The extended characteristic functions $`\widehat{\theta }_A`$ and $`\widehat{\theta }_B`$ are equivalent if and only if $`\underset{ยฏ}{A}`$ and $`\underset{ยฏ}{B}`$ are unitarily equivalent.
###### Proof.
If $`\underset{ยฏ}{A}`$ and $`\underset{ยฏ}{B}`$ are unitarily equivalent then all constructions differ only by naming and it follows that $`\widehat{\theta }_A`$ and $`\widehat{\theta }_B`$ are equivalent. Conversely, assume that there is a unitary $`V:๐_A๐_B`$ such that $`\widehat{\theta }_A=\widehat{\theta }_BV`$. Now from the commuting diagram 3.3 and the definitions following it
$`W_B_B`$ $`=`$ $`\left(\mathrm{\Gamma }(^d)๐_\omega \right)M_{\widehat{\mathrm{\Theta }}_B}\left(\mathrm{\Gamma }(^d)๐_B\right)`$
$`=`$ $`\left(\mathrm{\Gamma }(^d)๐_\omega \right)M_{\widehat{\mathrm{\Theta }}_B}\left(\mathrm{\Gamma }(^d)V๐_A\right)`$
$`=`$ $`\left(\mathrm{\Gamma }(^d)๐_\omega \right)M_{\widehat{\mathrm{\Theta }}_A}\left(\mathrm{\Gamma }(^d)๐_A\right)`$
$`=`$ $`W_A_A,`$
where we used equation 3.6, i.e., $`M_{\widehat{\mathrm{\Theta }}}(L_i\mathrm{๐}_๐)=(L_i\mathrm{๐}_{๐_\omega })M_{\widehat{\mathrm{\Theta }}},1id`$, to deduce $`M_{\widehat{\mathrm{\Theta }}_A}=M_{\widehat{\mathrm{\Theta }}_B}(\mathrm{๐}V)`$ from $`\widehat{\theta }_A=\widehat{\theta }_BV`$. Now we define the unitary $`U`$ by
$$U:=W_B^1W_A|__A:_A_B.$$
Because $`W_A\mathrm{\Omega }_A=1=W_B\mathrm{\Omega }_B`$ we have $`U\mathrm{\Omega }_A=\mathrm{\Omega }_B`$. Further for all $`i=1,\mathrm{},d`$ and $`h_A`$,
$$UA_ih=W_B^1W_AA_ih=W_B^1W_AP__AV_i^Ah=P__BW_B^1W_AV_i^Ah$$
$$=P__BW_B^1\stackrel{~}{V}_iW_Ah=P__BV_i^BW_B^1W_Ah=B_iUh,$$
i.e., $`\underset{ยฏ}{A}`$ and $`\underset{ยฏ}{B}`$ are unitarily equivalent. $`\mathrm{}`$
###### Remark 6.2.
An analogous result for completely non-coisometric tuples has been shown by G. Popescu in \[Po89b\], Theorem 5.4.
Note further that if we change $`\underset{ยฏ}{A}=(A_1,\mathrm{},A_d)`$ into $`\underset{ยฏ}{A}^{}=(A_1^{},\mathrm{},A_d^{})`$ by applying a unitary $`d\times d`$matrix with scalar entries (as described in Remark 2.2), then $`\widehat{\theta }_A=\widehat{\theta }_A^{}`$. In fact, this follows immediately from the definition of $`W`$ as an intertwiner in Section 3, from which it is evident that $`W`$ does not change if we take the same linear combinations on the left and on the right. This does not contradict Theorem 6.1 because $`\underset{ยฏ}{\omega }`$ and $`\underset{ยฏ}{\omega }`$ are now different tuples of eigenvalues and Theorem 6.1 is only applicable when the same tuple of eigenvalues is used for $`\underset{ยฏ}{A}`$ and $`\underset{ยฏ}{B}`$.
For another interpretation, let $`Z`$ be a normal, unital, ergodic, completely positive map with an invariant vector state $`\mathrm{\Omega }_A,\mathrm{\Omega }_A`$. If we consider two minimal Kraus decompositions of $`Z`$, i.e.,
$$Z=\underset{i=1}{\overset{d}{}}A_iA_i^{}=\underset{i=1}{\overset{d}{}}A_i^{}(A_i^{})^{},$$
with $`d`$ minimal, then the tuples $`\underset{ยฏ}{A}=(A_1,\mathrm{},A_d)`$ into $`\underset{ยฏ}{A}^{}=(A_1^{},\mathrm{},A_d^{})`$ are related in the way considered above (see for example \[Go04\], A.2). It follows that $`\widehat{\theta }_A=\widehat{\theta }_A^{}`$ does not depend on the decomposition but can be associated to $`Z`$ itself. Hence we have the following reformulation of Theorem 6.1.
###### Corollary 6.3.
Let $`Z_1,Z_2`$ be normal, unital, ergodic, completely positive maps on $`B(_1),B(_2)`$ with invariant vector states $`\mathrm{\Omega }_1,\mathrm{\Omega }_1`$ and $`\mathrm{\Omega }_2,\mathrm{\Omega }_2`$. Then the associated extended characteristic functions $`\widehat{\theta }_1`$ and $`\widehat{\theta }_2`$ are equivalent if and only if $`Z_1`$ and $`Z_2`$ are conjugate, i.e., there exists a unitary $`U:_1_2`$ such that
$$Z_1(x)=U^{}Z_2(UxU^{})U\text{for all}xB(_1).$$
## 7. Example
The following example illustrates some of the constructions in this paper.
Consider $`=^3`$ and
$$A_1=\frac{1}{\sqrt{2}}\left(\begin{array}{ccc}0& 0& 0\\ 1& 0& 0\\ 0& 1& 1\end{array}\right),A_2=\frac{1}{\sqrt{2}}\left(\begin{array}{ccc}1& 1& 0\\ 0& 0& 1\\ 0& 0& 0\end{array}\right).$$
Then $`_{i=1}^2A_iA_i^{}=\mathrm{๐}.`$ Take the unital completely positive map $`Z:M_3M_3`$ by $`Z(x)=_{i=1}^2A_ixA_i^{}.`$ It is shown in Section 5 of \[GKL06\] (and not difficult to verify directly) that this map is ergodic. We will use the same notations here as in previous sections. Observe that the vector $`\mathrm{\Omega }_{}:=\frac{1}{\sqrt{3}}(1,1,1)^T`$ gives an invariant vector state for $`Z`$ as
$$\mathrm{\Omega }_{},Z(x)\mathrm{\Omega }_{}=\mathrm{\Omega }_{},x\mathrm{\Omega }_{}=\frac{1}{3}\underset{i,j=1}{\overset{3}{}}x_{ij}.$$
$`A_i^{}\mathrm{\Omega }_{}=\frac{1}{\sqrt{2}}\mathrm{\Omega }_{}`$ and hence $`\underset{ยฏ}{\omega }=\frac{1}{\sqrt{2}}(1,1).`$ The orthogonal complement $`\stackrel{}{}`$ of $`\mathrm{\Omega }_{}`$ in $`^3`$ and the orthogonal projection $`Q`$ onto $`\stackrel{}{}`$ are given by
$$\stackrel{}{}=\{\left(\begin{array}{c}k_1\\ k_2\\ (k_1+k_2)\end{array}\right):k_1,k_2\},Q=\frac{1}{3}\left(\begin{array}{ccc}2& 1& 1\\ 1& 2& 1\\ 1& 1& 2\end{array}\right).$$
From this we get for $`\AA _i=QA_iQ=A_iQ`$
$$\AA _1=\frac{1}{3\sqrt{2}}\left(\begin{array}{ccc}0& 0& 0\\ 2& 1& 1\\ 2& 1& 1\end{array}\right),\AA _2=\frac{1}{3\sqrt{2}}\left(\begin{array}{ccc}1& 1& 2\\ 1& 1& 2\\ 0& 0& 0\end{array}\right).$$
We notice that the tuple $`\stackrel{}{\underset{ยฏ}{A}}=(\AA _1,\AA _2)`$ is $``$-stable as (by induction)
$$\underset{|\alpha |=n}{}\AA _\alpha \AA _\alpha ^{}=\frac{1}{3\times 2^{n1}}\left(\begin{array}{ccc}1& 1& 0\\ 1& 2& 1\\ 0& 1& 1\end{array}\right)0(n\mathrm{}).$$
Here $`๐ซ=^2`$ and $`\stackrel{}{๐ซ}:=๐ซ\mathrm{\Omega }_๐ซ`$ with $`\mathrm{\Omega }_๐ซ=\frac{1}{\sqrt{2}}(1,1)^T`$. Easy calculation shows that $`\widehat{D}_{}:\stackrel{}{}\stackrel{}{๐ซ}`$ is given by
$$\left(\begin{array}{c}k_1\\ k_2\\ (k_1+k_2)\end{array}\right)\frac{1}{\sqrt{6}}(2k_1+k_2)\left(\begin{array}{c}1\\ 1\end{array}\right).$$
Moreover $`\underset{}{\overset{}{D}}=\frac{1}{\sqrt{6}}\left(\begin{array}{ccc}1& 0& 1\\ 0& 0& 0\\ 1& 0& 1\end{array}\right).`$ There exists an isometry $`\gamma :\underset{}{\overset{}{๐}}\stackrel{}{๐ซ}`$ such that $`\left(\begin{array}{c}1\\ 0\\ 1\end{array}\right)\left(\begin{array}{c}1\\ 1\end{array}\right)`$ and $`\gamma (\underset{}{\overset{}{D}}h)=\widehat{D}_{}h`$ for $`h\stackrel{}{}`$.
The map $`\widehat{C}:\mathrm{\Gamma }(^d)๐_\omega `$ is given by $`\widehat{C}(\mathrm{\Omega }_{})=1`$ and for $`h\stackrel{}{}`$ by
$`\widehat{C}\left(\begin{array}{c}k_1\\ k_2\\ (k_1+k_2)\end{array}\right)`$ $`=`$ $`e_0{\displaystyle \frac{(2k_1+k_2)}{\sqrt{6}}}\left(\begin{array}{c}1\\ 1\end{array}\right)+{\displaystyle \underset{\alpha ,\alpha _1=1}{}}e_\alpha ({\displaystyle \frac{1}{\sqrt{2}}})^{|\alpha |}{\displaystyle \frac{(k_1+2k_2)}{\sqrt{6}}}`$
$`\times \left(\begin{array}{c}1\\ 1\end{array}\right)+{\displaystyle \underset{\alpha ,\alpha _1=2}{}}e_\alpha ({\displaystyle \frac{1}{\sqrt{2}}})^{|\alpha |}{\displaystyle \frac{(k_1k_2)}{\sqrt{6}}}\left(\begin{array}{c}1\\ 1\end{array}\right)`$
where the summations are taken over all $`0\alpha \stackrel{~}{\mathrm{\Lambda }}`$ such that $`\alpha _i\alpha _{i+1}`$ for all $`1i|\alpha |`$ and fixing $`\alpha _1`$ to $`1`$ or $`2`$ as indicated. This simplification occurs because $`\AA _i^2=0`$ for $`i=1,2`$. All the summations below in this section are also of the same kind.
Now using the equations 4.2 and 4.3 for $`\widehat{\theta }_A:๐_A\mathrm{\Gamma }(^d)๐_\omega `$ and simplifying we get
$`\widehat{\theta }_Ad_\mathrm{\Omega }_{}^1`$ $`=`$ $`e_0{\displaystyle \frac{1}{6}}\left(\begin{array}{c}1\\ 1\end{array}\right)+{\displaystyle \underset{\alpha ,\alpha _1=1}{}}e_\alpha ({\displaystyle \frac{1}{\sqrt{2}}})^{|\alpha |}{\displaystyle \frac{1}{6}}\left(\begin{array}{c}1\\ 1\end{array}\right)`$
$`+{\displaystyle \underset{\alpha ,\alpha _1=2}{}}e_\alpha ({\displaystyle \frac{1}{\sqrt{2}}})^{|\alpha |}{\displaystyle \frac{1}{6}}\left(\begin{array}{c}1\\ 1\end{array}\right)=\widehat{\theta }_Ad_\mathrm{\Omega }_{}^2,`$
and for $`h\stackrel{}{},`$
$`\widehat{\theta }_Ad_h^1`$ $`=`$ $`e_0{\displaystyle \frac{k_1}{2\sqrt{3}}}\left(\begin{array}{c}1\\ 1\end{array}\right)+e_1{\displaystyle \frac{(k_1+k_2)}{\sqrt{6}}}\left(\begin{array}{c}1\\ 1\end{array}\right)+{\displaystyle \underset{\alpha ,\alpha _1=1}{}}e_1e_\alpha `$
$`({\displaystyle \frac{1}{\sqrt{2}}})^{|\alpha |}{\displaystyle \frac{(k_1+2k_2)}{\sqrt{6}}}\left(\begin{array}{c}1\\ 1\end{array}\right){\displaystyle \underset{\alpha ,\alpha _1=2}{}}e_1e_\alpha ({\displaystyle \frac{1}{\sqrt{2}}})^{|\alpha |}{\displaystyle \frac{k_2}{\sqrt{6}}}\left(\begin{array}{c}1\\ 1\end{array}\right)`$
$`+{\displaystyle \underset{\alpha ,\alpha _1=2}{}}e_\alpha ({\displaystyle \frac{1}{\sqrt{2}}})^{|\alpha |}{\displaystyle \frac{k_1}{2\sqrt{3}}}\left(\begin{array}{c}1\\ 1\end{array}\right),`$
$`\widehat{\theta }_Ad_h^2`$ $`=`$ $`e_0{\displaystyle \frac{(k_1+k_2)}{2\sqrt{3}}}\left(\begin{array}{c}1\\ 1\end{array}\right)+{\displaystyle \underset{\alpha ,\alpha _1=1}{}}e_\alpha ({\displaystyle \frac{1}{\sqrt{2}}})^{|\alpha |}{\displaystyle \frac{(k_1+k_2)}{2\sqrt{3}}}\left(\begin{array}{c}1\\ 1\end{array}\right)`$
$`+e_2{\displaystyle \frac{k_1}{\sqrt{6}}}\left(\begin{array}{c}1\\ 1\end{array}\right)+{\displaystyle \underset{\alpha ,\alpha _1=1}{}}e_2e_\alpha ({\displaystyle \frac{1}{\sqrt{2}}})^{|\alpha |}{\displaystyle \frac{k_2}{\sqrt{6}}}\left(\begin{array}{c}1\\ 1\end{array}\right)`$
$`+{\displaystyle \underset{\alpha ,\alpha _1=2}{}}e_2e_\alpha ({\displaystyle \frac{1}{\sqrt{2}}})^{|\alpha |}{\displaystyle \frac{(k_1k_2)}{\sqrt{6}}}\left(\begin{array}{c}1\\ 1\end{array}\right).`$
Form this we can easily obtain $`\stackrel{}{C}`$ and $`\theta _\AA `$ for $`h\stackrel{}{}`$ by using the following relations from Theorem 5.2,
$$(\mathrm{๐}\gamma )\stackrel{}{C}h=\widehat{C}h,$$
$$(\mathrm{๐}\gamma )\theta _\AA d_h^i=\widehat{\theta }_Ad_h^i.$$
Further
$$l_1=A_1\mathrm{\Omega }_{}\frac{1}{\sqrt{2}}\mathrm{\Omega }_{}=\frac{1}{\sqrt{6}}\left(\begin{array}{c}1\\ 0\\ 1\end{array}\right),l_2=A_2\mathrm{\Omega }_{}\frac{1}{\sqrt{2}}\mathrm{\Omega }_{}=\frac{1}{\sqrt{6}}\left(\begin{array}{c}1\\ 0\\ 1\end{array}\right),$$
$`\AA _1l_1=\frac{1}{2\sqrt{3}}\left(\begin{array}{c}0\\ 1\\ 1\end{array}\right)`$ and clearly $`\stackrel{}{}=\overline{\text{span}}\{\AA _\alpha l_i:i=1,2\text{ and }\alpha \stackrel{~}{\mathrm{\Lambda }}\},`$ as already observed in Remark 3.1.
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# Summary of theoretical contributions
## 1 Theoretical overview
We have witnessed an exciting conference with an excellent program, heated scientific discussions and lots of new data and theoretical ideas. Our topics ranged from astrophysics to field theory, from heavy-ion reaction phenomenology to big-bang cosmology.
My task, i.e. to review all in all about 30 theory talks, is combined in this paper with a cross-disciplinary analysis of experiments, which verify - pardon, falsify the theoretical conjectures in many cases - there is an old saying that a theory can never be verified: even if lots of data support the theory, at some point the theory will always go astrayโฆ
Let me rearrange the order of the theory talks on the topics of our meeting<sup>1</sup><sup>1</sup>1Instead of giving explicit references I refer the reader to the electronic proceedings available on the Web:
* Equation of State
* Collective Dynamics
* Jets : Production and Quenching
* Results from $`p+p`$, $`p(d)+A`$ and $`A+A`$ collisions
* Signatures of Quark Gluon Plasma
* QCD at Finite Temperature and Density
* Multiparticle production, fluctuations and correlations
* Cosmological Implications of the QCD Phase Transition
* QCD Phenomenology
* Low $`x`$ behaviour of QCD
* Strangeness and heavy flavor production
The common interest is given by the titel of the conference: โPhysics and astrophysics of the quark gluon plasma<sup>2</sup><sup>2</sup>2Somewhat in this summary, also pardon the fact that some of the many interesting items have not been taken up here, because I did not witness the first days of the conference.
* Astrophysics
* Lattice
* Colored Glass
* Fluctuations & DCCs
* J/Psi & EM Probes
* Strangeness
* Transport Theory
* Hydro & Jets
John Ellis gave a beautiful survey of the common issues in both heavy-ion physics and the big bang cosmology: We do in both cases study a very fast expansion of dense/hot strongly interacting matter, and do have the task to reconcile whatever happened in the first few nanoseconds of the big bang from the sparse debris found nowadays. The connection to the matter-anti matter asymmetry problems is particularly exciting for future topical studies at the LHC. This is quite analogous to the transient 6-8 fm/c $`22.5\times 10^{23}s`$ timescale of the collision processes at RHIC.
The intense astrophysics discussions between Bombaci and Banyopadyhyay about the possible occurence of massive strange quark stars (SQS), the transition of neutron stars to strange hyperon-, hybrid- and quark stars, and the relation to the gamma ray bursts (quark-deconfinement nova-model) has been of particular interest - this transition is predicted to yield a radius-collapse of several kilometers.
The first observation of the โdouble delight pulsarโ psr-j0737-3039 will enable us to pin down the mass-radius curves by the spin-orbit effect with high precision. D. Bandyopadyhyay showed that soft equations of state (EOS) are ruled out by EXO 0748-676. The connection of conjectured different color superconducting phases to the cooling curves of SQS have been pointed out in the paper by Mishra and Mishra.
The lattice-QCD (lQCD) discussions between Gavai and Laermann centered about the questions on the order of the phase transition and on the speed of sound. Laermann stated that there is no indication for criticality, while Gavai and friends showed that the critical endpoint is at $`T=0.95T_c`$, $`\mu _B/T=1.11.3`$, i.e. less than half of the $`\mu _B`$=400 MeV values given by the Swansea-Bielefeld and Wuppertal-Budapest collaborations (cf. Fig 1). Gavai showed also that all the way up to $`T=2T_c`$, the speed of sound is much less, $`c_s^2=0.15`$ at $`T=1.1T_c`$, than that of a noninteracting ultrarelativistic (massless) gas, $`c_s^2=1/3`$.
In the colored glass condensate section, Venugopalan explored the demise of the structure function, in particular how the dipole and higher multipole operators may turn out to be the more relevant observables at high energies. Adding valence quark contributions, Kovchegov showed a quite satisfactory agreement of the Color-Glass-Condensate (CGC)-model to the observed rapidity dependence of the $`p_T`$-distributions. McLerran iterated the theme of the Color Glass Condensate as THE Medium: Pomerons, Odderons, Reggeons as Quasiparticle excitations of the CGC - does this mean that the CGC is the initial phase for the QGP? Is the strong Quark-Gluon Plasma (sQGP) really the CGC? Is rapid โthermalizationโ due to the CGC? Does flow arise largely from the CGC? Well, definitely LHC is THE CGC machine โ according to McLerran.
Fluctuations and Disordered Chiral Condensates (DCCโs) were discussed by Koch, Csรถrgรถ, Chandrasekar and Randrup, among others. $`K/\pi `$ fluctuations increase towards lower beam energy with a significant enhancement over the hadronic cascade model UrQMD (cf. Fig. 2)! On the other hand, $`p/\pi `$ fluctuations are negative โ this indicates a strong contribution from resonance decays, as was shown by Koch in comparing NA49-data to UrQMD results.
Dileptons, J/Psi, and photons have been discussed by Lee, Mustafa and Koch (among others). Large corrections on the QCD NLO Quarkonium- Gluon/hadron dissociation cross section have been reported even for the Ypsilon system, especially near threshold. The thermal width of the $`J/\mathrm{\Psi }`$ should be $``$ 1 GeV at T=600 MeV according to Leeโs estimates.
Strangegeness and equilibration has been the main topic of Rafelski, Cleymans, Braun-Munzinger and Bleicher. The structure in the $`K/\pi `$ ratios reported by NA49 near $`\sqrt{(}s)`$ = 8 GeV is not reproduced by any model (cf. Fig. 3), but Peter Braun-Munzinger notes: the natural smearing is 3 GeV near that energy - how can the โhornโ then be so steep? Hadron-string models work well globally, as Bleicher reports, but these models do NOT give MULTI-STRANGE BARYONs! Is the alternative a four parameter nonequilibrium thermal model, with $`T,\mu ,\mathrm{\Gamma }_\mu ,\mathrm{\Gamma }_s`$, by Rafelski et al.?
The extreme density/temperature dependence of the characteristic equilibration time, $`\tau _{eq.}T^{60}`$, was pointed out by Braun-Munzinger, which implies that all particles freeze out at about the hadronization time. According to Braun-Munzinger this might be due to Carsten Greinerโs conjecture of Hagedorn states as intermediate doorway states.
Deeply bound $`\overline{p}`$ and $`K^{}`$ states as gateway to cold and dense matter were discussed by Walter Greiner: $`\overline{p}`$โs โ due to $`G`$-parity in the strong interactions โ and $`K^{}`$ can suppress repulsive vector fields, thus predicting discrete bound states with binding energies of several 100th MeV and 20 fm/c life times . Formation of such cold and highly dense nuclear system at densities $`\rho 35\rho _0`$ will be studied in dense $`\overline{p}`$ \- nuclear systems at FAIR (GSI)and the $`K`$-nucleus collisions at J-PARC.
Jacak, Shuryak, Heinz and Chauduri discussed applications of hydrodynamics to RHIC-collisions. The reasons why hydro does reasonably well fit both, radial and elliptic flow for a large number of hadron species (cf. Fig. 4), is still not fully settled. The question of early thermalization and the unsatisfactory rapidity distributions from ideal hydrodynamics remain open.
Bass showed in his talk, however, that the recombination/quark coalescence models (cf. Fig. 5) can help analyze the participant scaling and even the charm flow. However, as Bleicher showed, even the hadron/string model UrQMD may exhibit โrecombinationโ and participant scaling.
Jacak showed the PHENIX jet-pair distributions, which clearly give a novel signal to the away-side jet suppression (cf. Fig. 6 for STAR results), i.e. the recent topic of Mach-cones induced by stopped jets in the quark-gluon liquid . This is most important as an observable, because it links the parton dynamics and collective flow and the jet tomography to the measurement of the speed of sound in the medium - be it a weakly or strongly coupled plasma: the opening angle of the Mach-shock-wakes directly gives the speed of sound in the medium, which is linked to both, the appearance of vector potentials and the parton/constituent mass parameters.
## 2 Interlude on Mach shocks
Sideward peaks around the away-side jet have been predicted recently as a signature of Mach shock waves created by stopping partonic jets propagating through a QGP formed in an ultrarelativistic heavyโion collision. Analogous Mach shock waves were studied long ago for heavy-ion induced Mach shocks travelling through cold hadronic matter as well as in nuclear Fermi liquids . It has been argued that Machโlike motions of quarkโgluon matter can appear via the excitation of collective plasmon waves by the moving color charge associated with the leading jet .
Pointโlike perturbations (a small body, a hadron or parton etc.) moving with a supersonic speed in the spatially homogeneous ideal fluid produce the Mach region of the perturbed matter . In the fluid rest frame (FRF) the Mach region has a conical shape (cf. Fig. 7) with an opening angle with respect to the direction of particle propagation given by<sup>3</sup><sup>3</sup>3 Here and below the quantities in the FRF are marked by a tilde.
$$\stackrel{~}{\theta }_M=\mathrm{sin}^1\left(\frac{c_s}{\stackrel{~}{v}}\right),$$
(1)
where $`c_s`$ denotes the sound velocity of the unperturbed (upstream) fluid and $`\stackrel{~}{๐}`$ is the particle velocity with respect to the fluid. In the FRF, trajectories of fluid elements (perpendicular to the surface of the Mach cone) are inclined at the angle $`\mathrm{\Delta }\theta =\pi /2\stackrel{~}{\theta }_M`$ with respect to $`\stackrel{~}{๐}`$ . Strictly speaking, formula (1) is applicable only for weak, soundโlike perturbations and certainly not valid for spaceโtime regions close to a leading particle. Nevertheless, it suffices for a qualitative analysis of flow effects. Following Refs. one can estimate the angle of preferential emission of secondaries associated with a fast jet in the QGP. Substituting $`\stackrel{~}{v}=1,c_s=1/\sqrt{3}`$ into Eq. (1) gives the value $`\mathrm{\Delta }\theta `$ 0.96 rad = 61<sup>o</sup>แนชhis agrees well with positions of maxima of the awayโside twoโparticle distributions observed by the STAR Collaboration (cf. 6) in central Au+Au collisions at RHIC energies (cf. also B. Jacakโs talk).
Let us consider the case when the awayโside jet propagates with velocity $`๐`$ parallel to the matter flow velocity $`๐`$ . Assuming that $`๐`$ does not change with space and time, and performing the Lorentz boost to the FRF, one sees that a weak Mach shock has a conical shape with the axis along $`๐`$ . In this reference frame, the shock front angle $`\stackrel{~}{\theta }_M`$ is given by (1). Transformation from the FRF to the c.m. frame (CMF) shows that the Mach region remains conical, but the Mach angle becomes smaller in the CMF:
$$\mathrm{tan}\theta _M=\frac{1}{\gamma _u}\mathrm{tan}\stackrel{~}{\theta }_M,$$
(2)
where $`\gamma _u(1u^2)^{1/2}`$ is the Lorentz factor corresponding to the flow velocity u . The resulting expression for the Mach angle in the CMF is
$$\theta _M=\mathrm{tan}^1\left(c_s\sqrt{\frac{1u^2}{\stackrel{~}{v}^2c_s^2}}\right),$$
(3)
where
$$\stackrel{~}{v}=\frac{vu}{1vu},$$
(4)
and upper (lower) sign corresponds to the jetโs motion in (or opposite to) the direction of collective flow. For ultrarelativistic jets ($`v1`$) one can take $`\stackrel{~}{v}1`$ which leads to a simpler expression
$$\theta _M\mathrm{tan}^1\left(\frac{c_s\gamma _s}{\gamma _u}\right)=\mathrm{sin}^1\left(c_s\sqrt{\frac{1u^2}{1u^2c_s^2}}\right),$$
(5)
where $`\gamma _s=(1c_s^2)^{1/2}`$ . According to (5), in the ultrarelativistic limit $`\theta _M`$ does not depend on the direction of flow with respect to the jet. The Mach cone becomes more narrow as compared to jet propagation in static matter. This narrowing effect has a purely relativistic origin. Indeed, the difference between $`\theta _M`$ from (5) and the Mach angle in absence of flow ($`\underset{u0}{lim}\theta _M=\mathrm{sin}^1c_s`$) is of second order in the collective velocity $`u`$ . The Mach angle calculated from (5) is shown in Fig. 8 (from ) as a function of $`u`$ for different sound velocities $`c_s`$ . Following Ref. , the value $`c_s^2=1/5`$ is identified with the hadronic matter and $`c_s^2=1/3`$ with ideal QGP composed of massless quarks and gluons. The value $`c_s^2=2/3`$ may be chosen to represent a strongly coupled QGP . We see that precise measurements will provide valuable information on the properties of the quark-gluon liquid .
## 3 Future directions
I propose future correlation measurements which can yield spectroscopic information on the plasma:
1. Measure the sound velocity of the expanding plasma by the emission pattern of the plasma particles traveling sideways with respect to the jet axis: The dispersive wave generated by the wake of the jet in the plasma yields preferential emission to an angle (relative to the jet axis) which is given by the ratio of the leading jet particlesโ velocity, devided by the sound velocity in the hot dense plasma rest frame. The speed of sound for a non-interacting gas of relativistic massless plasma particles is $`c_s\frac{c}{\sqrt{3}}57\%c`$, while for a plasma with strong vector interactions, $`c_sc`$, since strong shocks can yield larger speeds. They are also related โ unlike the linearized sound waves โ to strong matter flow with high flow velocities $`v_f`$ approaching the speed of light relative to the expanding medium. Hence, the emission angle measurement can yield information of the interactions in the plasma.
2. The NA49 collaboration has observed the collapse of both, $`v_1`$\- and $`v_2`$-collective flow of protons (cf. Fig. 9), in Pb+Pb collisions at 40 A$``$GeV, which presents first evidence for a first order phase transition in baryon-rich dense matter. It should be possible to study the nature of this transition and the properties of the expected chirally restored and deconfined phase both at the forward fragmentation region at RHIC, with upgraded and/or second generation detectors, and at the new GSI facility FAIR.
3. A critical discussion of the use of collective flow as a barometer for the equation of state (EoS) of hot dense matter at RHIC showed that hadronic rescattering models can explain $`<30\%`$ of the observed elliptic flow $`v_2`$ for $`p_T>2`$ GeV/c . I interpret this as evidence for the production of superdense matter at RHIC with initial pressure way above hadronic pressure, $`p>1`$ GeV/fm<sup>3</sup>.
4. The fluctuations in the flow, $`v_1`$ and $`v_2`$, should be measured. Ideal Hydrodynamics predicts that they are larger than 50 % due to initial state fluctuations. The QGP coefficient of viscosity may be determined experimentally from the fluctuations observed and proof the conjecture of Ref. .
5. The connection of $`v_2`$ to jet suppression has proven experimentally that the collective flow is not faked by minijet fragmentation and theoretically that the away-side jet suppression can only partially ($`<`$ 50%) be due to pre-hadronic or hadronic rescattering (cf. Fig. 10).
6. I propose upgrades and second generation experiments at RHIC, which inspect the first order phase transition in the fragmentation region, i.e. at $`\mu _B200400`$ MeV ($`y35`$), where the collapse of the proton flow โ analogous to the 40 A$``$GeV data โ should be seen.
Let me finally express my birthday greetings to Bikash and thank him and his crew for decades of exciting physics conjectures, his strong involvement into our field and courage to built up such a great school of young successful scientists in India, which are highly competitive in the whole world.
## Acknowledgments
This work is partially supported by GSI, BMFT, DFG, and DAAD. Let me thank Elena Bratkovskaya, Leonid Satarov and Igor Mishustin (from FIAS) for their contributions to this Summary.
## References
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# Optical Constraints on an X-ray Transient Source in M31
## 1 Introduction
Many low-mass X-ray binaries (LMXBs) that undergo bright, transient X-ray bursts have been shown to harbor compact objects with masses $`>`$3 M (see McClintock & Remillard, 2004, and references therein). These sources are some of the most securely identified black holes known. Such sources are therefore of great interest for future studies of black hole accretion disk physics and general relativity.
The Chandra X-ray observatory is well-suited to searching for similar X-ray sources in nearby galaxies. In particular, the bulge of M31 can be entirely searched for bright, transient X-ray events with a single 5 ks observation. Monitoring of M31 by Chandra has shown that such transient sources appear in the M31 bulge about once each month (Williams et al., 2004). These efforts have been successful in finding dozens of transient sources (Kong et al., 2002; Di Stefano et al., 2004; Williams et al., 2004).
Recently, this monitoring effort has been combined with follow-up $`HST`$ observations, and optical counterparts for some transient X-ray sources have been found, placing new constraints on the physical properties of these potential black hole binary systems (Williams et al., 2004, 2005a, 2005b). With a sufficient number of optical counterparts, the orbital period distribution of these likely black hole binaries can be determined. This distribution is a fundamental observable parameter that must be matched by any model of binary stellar evolution.
Here, we report a bright outburst from an X-ray source in M31. $`HST`$ imaging prior to, during, and after the outburst reveals no optical sources that exhibited strong variability during the X-ray outburst within the location uncertainty, suggesting that the optical counterpart of the X-ray nova (XRN) was fainter than $`B=25.5`$. Section 2 discusses the details of the data, including the reduction and analysis techniques used. Section 3 provides the results of the analysis, and ยง 4 explains the implications of these results. Finally, ยง 5 summarizes our conclusions.
## 2 Data Reduction and Analysis
### 2.1 X-ray
We obtained observations of the M31 bulge with the Chandra ACIS-I camera on 2004-January-31 (ObsID 4681), 2004-May-23 (ObsID 4682), and 2004-July-17 (ObsID 4719). Observations 4681 and 4682 were Guaranteed Time Observations supplied by S. Murray. All of the data were obtained in โalternating readout modeโ which reduces event pileup for bright sources but lowers the effective exposure by $`20\%`$. Observation 4681 was observed centered on R.A.=00:42:44.4, Dec.=41:16:08.3 with a roll angle of 305.6 degrees for an effective exposure time of 4.09 ks. Observation 4682 was centered on the same coordinates with a roll angle of 80.0 degrees for an effective exposure time of 3.95 ks, and observation 4719 was observed centered on R.A.=00:42:44.3, Dec.=41:16:08.4 with a roll angle of 116.8 degrees for an effective exposure time of 4.12 ks.
These observations were all reduced using the software package CIAO v3.1 with the CALDB v2.28. We created exposure maps for the images using the task merge\_all,<sup>1</sup><sup>1</sup>1http://cxc.harvard.edu/ciao/download/scripts/merge\_all.tar and we found and measured positions and fluxes of the sources in the image using the CIAO task wavdetect.<sup>2</sup><sup>2</sup>2http://cxc.harvard.edu/ciao3.0/download/doc/detect\_html\_manual/Manual.html Each data set detected sources down to (0.3โ10 keV) fluxes of $``$6$`\times `$10<sup>-6</sup> photons cm<sup>-2</sup> s<sup>-1</sup> or 0.3โ10 keV luminosities of $``$10<sup>36</sup> erg s<sup>-1</sup> for a typical X-ray binary system in M31.
We aligned the coordinate system of the X-ray images with the optical images of the Local Group Survey (LGS; Massey et al., 2001). These images have an assigned J2000 (FK5) world coordinate system accurate to $``$0.25<sup>โฒโฒ</sup>, and they provided the standard coordinate system to which we aligned all of our data for this project. The positions of X-ray sources with known globular cluster counterparts were aligned with the globular cluster centers in the LGS images using the IRAF<sup>3</sup><sup>3</sup>3IRAF is distributed by the National Optical Astronomy Observatory, which is operated by the Association of Universities for Research in Astronomy, Inc., under cooperative agreement with the National Science Foundation. task ccmap, allowing adjustments of the pixel scale as well as rotation and shifts in $`X`$ and $`Y`$. The alignment had errors of 0.073<sup>โฒโฒ</sup> in R.A. and 0.034<sup>โฒโฒ</sup> in Dec. as determined by ccmap.
In addition, we measured the position errors for the X-ray sources using the IRAF task imcentroid, which projects the counts from the source onto each axis and calculates the error in the position by dividing the standard deviation of the pixel positions of all of the source counts by the square root of the number of counts. Because the pixels in Chandra images are aligned with north up and east to the left, the X position error was taken to be the R.A. error, and the Y position error was taken to be the Dec. error.
We cross-correlated the X-ray source positions against all previously published X-ray catalogs and the Simbad<sup>4</sup><sup>4</sup>4http://simbad.u-strasbg.fr/ database to look for any new, bright X-ray source likely to be an X-ray nova (XRN). Herein we focus on one bright source in particular which was detected only in observation 4682 at R.A.=00:42:33.428, Dec.=+41:17:03.37. This source was detected once by ROSAT in 1990 (\[PFJ93\] 31; Primini et al., 1993), but it has not been detected since. We name this source CXOM31 J004233.4+411703, following the naming convention described in Kong et al. (2002). We also give this source the short name, r2-71, based on the source position in M31 using the description given in Williams et al. (2004). The source is 2.0 west and 0.9 north of the M31 nucleus. Upper-limits to the X-ray flux at this position in observations 4681 and 4719 were measured by determining the flux necessary to produce a detection 4-$`\sigma `$ above the background flux.
We extracted the X-ray spectrum of r2-71 from the detection in observation 4682, which contained 198 counts, using the CIAO task psextract.<sup>5</sup><sup>5</sup>5http://cxc.harvard.edu/ciao/ahelp/psextract.html We binned the spectrum in energy so that each bin contained $`\stackrel{>}{}`$10 counts to allow for standard $`\chi ^2`$ statistics. We then fit the spectrum using CIAO 3.1/Sherpa. Two spectral models were fitted to the spectrum: a power-law with absorption and a disk blackbody model with absorption. The spectrum did not contain sufficient information to provide useful constraints on the foreground absorption. We therefore fixed the absorption to the typical Galactic foreground value (6$`\times `$10<sup>20</sup> cm<sup>-2</sup>). This value provided acceptable fits for both model types. Results are discussed in ยง 3.
### 2.2 Optical
We obtained three sets of HST ACS data. The first ACS image was taken 2004-January-23 in pursuit of another transient source. This observation was pointed at R.A.=00:42:43.86, Dec.=41:16:30.1 with an orientation of 55 deg, fortuitously containing the location of r2-71 in the northern corner of the image even though r2-71 had not been detected in any Chandra observations. The second observation was taken 2004-June-14. Intended to search for the optical counterpart of r2-71, this observation was pointed at R.A.=00:42:33.39, Dec.=41:17:43.4 with an orientation of 260 deg. The slightly northern pointing was done to avoid a dangerously bright star from landing on the CCD. The third observation was taken on 2004-August-15, after the X-ray source had faded. This observation was pointed at R.A.=00:42:41.5, Dec.=41:17:00.0 with an orientation of 220 deg in order to pursue another transient source in the field.
All three observations were taken using the standard ACS box 4-point dither pattern to recover the highest possible spatial resolution. All exposures were taken through the F435W filter. The total exposure times were 2200 seconds for each data set. We aligned and drizzled each set of 4 images into high-resolution (0.025<sup>โฒโฒ</sup> pixel<sup>-1</sup>) images using the PyRAF<sup>6</sup><sup>6</sup>6PyRAF is a product of the Space Telescope Science Institute, which is operated by AURA for NASA. task multidrizzle,<sup>7</sup><sup>7</sup>7multidrizzle is a product of the Space Telescope Science Institute, which is operated by AURA for NASA. http://stsdas.stsci.edu/pydrizzle/multidrizzle which has been optimized to process ACS imaging data. The task removes cosmic ray events and geometric distortions, and it combines the dithered frames together into one final photometric image.
We aligned the $`HST`$ images to the LGS coordinate system with ccmap using stars common to both data sets. The resulting alignment had rms errors of 0.04<sup>โฒโฒ</sup> (less than 1 ACS pixel). The consistency of this alignment can be seen by the excellent agreement between the resulting coordinate systems of the three $`HST`$ images, each independently aligned with the LGS coordinate system, shown in Figure 1.
We processed the relevant sections of the final images with DAOPHOT II and ALLSTAR (Stetson et al., 1990) to obtain photometry for the resolved stars within the error circle of the X-ray transient. This region is extremely crowded with stars at only 2.2 from the M31 nucleus. The faintest source visible in the error circles in the images shown in Figure 1 is $`B=26.7`$; the brightest is $`B=24.7`$ (on the northeast rim of the circle).
We subtracted the images taken before and after the X-ray nova from the 2004-June-14 observation, during which the X-ray source was likely to be active. Prior to subtraction, the images were transformed to have pixels aligned in North-up, East-left orientation. These transformations were performed with the IRAF tasks geomap and geotran. The relevant sections of the subtracted images are shown in Figure 2.
## 3 Results
### 3.1 X-ray
The three X-ray images from our $`Chandra`$ observations are shown in Figure 1. These images clearly show r2-71 detected only in the 2004-May-23 observation. This detection had a 0.3โ10 keV flux of (1.7$`\pm `$0.1$`)\times `$10<sup>-4</sup> photons cm<sup>-2</sup> s<sup>-1</sup>. The previous observation and following observation had 4$`\sigma `$ upper limits of $`<`$1.4$`\times `$10<sup>-5</sup> photons cm<sup>-2</sup> s<sup>-1</sup>. The second upper limit shows that the source decayed by a factor of at least 12 in 55 days. Therefore the $`e`$-folding decay time of r2-71 was $``$22 days.
The non-detection of this source in the survey of Kong et al. (2002) provides a 0.3โ7 keV flux upper limit of $`<`$8$`\times `$10<sup>-7</sup> photons cm<sup>-2</sup> s<sup>-1</sup>, showing that the source changed in flux by more than a factor of 100. Therefore r2-71 is certainly a transient X-ray source in M31.
The errors in the centroid determination of the X-ray source were 0.08<sup>โฒโฒ</sup> in R.A. and 0.10<sup>โฒโฒ</sup> in Declination. We added these errors in quadrature to the errors in the alignment of the X-ray and optical images (0.073<sup>โฒโฒ</sup> in R.A. and 0.034<sup>โฒโฒ</sup> in Declination) to obtain the final (1$`\sigma `$) position errors of 0.11<sup>โฒโฒ</sup> and 0.11<sup>โฒโฒ</sup> in R.A. and Declination respectively. These errors resulted in the 2$`\sigma `$ error circle shown in Figure 1.
The X-ray spectrum of r2-71 was well-fitted by both the power-law and the disk-blackbody models. Fortunately, both of these fits give the same measurement for the absorption-corrected 0.3โ7 keV flux and the corresponding 0.3โ7 keV luminosity. The best-fitting power-law has a slope of 1.5$`\pm `$0.1 with $`\chi ^2/\nu `$=17.87/17 (probability = 0.40). This fit yields an unabsorbed 0.3โ7 keV flux of ($`5.4\pm 0.5`$)$`\times `$10<sup>-13</sup> erg cm<sup>-2</sup> s<sup>-1</sup>. The best fitting disk blackbody has an inner disk temperature of kT=1.4$`\pm `$0.2 keV, an inner disk radius of $`(6\pm 2)/cos^{1/2}(i)`$ km with a $`\chi ^2/\nu =16.38/17`$ (probability=0.50). The resulting unabsorbed 0.3โ7 keV flux is ($`4.9\pm 2.2`$)$`\times `$10<sup>-13</sup> erg cm<sup>-2</sup> s<sup>-1</sup>.
Assuming a distance to M31 of 780 kpc (Williams, 2003), the results are both consistent with the X-ray luminosity of (3.9$`\pm `$0.4)$`\times `$10<sup>37</sup> erg s<sup>-1</sup> obtained from the power-law fit, similar to the luminosity seen in the 1990 ROSAT data by Primini et al. (1993). The spectrum of r2-71 is in the normal range of M31 X-ray transients as measured by Williams et al. (in preparation), with hardness ratios of $`(MS)/(M+S)=0.59\pm 0.11`$ and $`(HS)/(H+S)=0.52\pm 0.11`$, where S, M, and H represent the number of counts detected in the energy ranges 0.3โ1.0, 1.0โ2.0, and 2.0โ7.0 keV, respectively. Gaussian errors were measured for the background-subtracted number of counts in each energy bin with the CIAO task dmextract.
### 3.2 Optical
Analysis of the optical data initially provided a few possible variable stars within the r2-71 error circle. We scrutinized each possibility through aperture photometry, completeness tests, and difference imaging. The results show no strong detection of optical variability and provide an upper-limit to the $`B`$ magnitude of any highly variable counterpart to r2-71.
The region of interest for all three of our optical observations is shown in the images in Figure 1. The DAOPHOT II output revealed one star in the error circle that was significantly brighter in the 2004-June-14 observation, when the X-ray source was most likely active. The star at R.A.=00:42:33.414, decl.=41:17:03.47, was measured to have $`B=25.50\pm 0.06`$ in the 2004-June-14 observation. DAOPHOT II failed to find this star in the 2004-January-23 observation. Because inspection of the 2004-January-23 image reveals a source at this location, this non-detection was likely due to the effects of the bright neighboring star 0.05<sup>โฒโฒ</sup> to the southeast. Aperture photometry of the location in the 2004-January-23 observation with a 0.075<sup>โฒโฒ</sup> radius aperture yields $`B=26.0\pm 0.1`$. DAOPHOT II measured the star to be $`B=25.98\pm 0.09`$ in the 2004-August-15 observation. Therefore according to the standard errors an increase in brightness during the 2004-June-14 at the 4$`\sigma `$ confidence level occurred; however, the standard errors do not take into account the uncertainty introduced by the close brighter neighbor. Any added uncertainty due to crowding would decrease the significance of this brightness increase.
In addition to this suspicious counterpart candidate, there were 2 fainter stars detected by DAOPHOT II in the error circle of the 2004-June-14 observation that were not detected in the other 2 observations, when the X-ray source was not active. These stars had $`B`$ magnitudes of 26.6$`\pm `$0.2 and 26.7$`\pm `$0.2 in the 2004-June-14 observation. These stars may not have been detected in the other observations because crowding issues caused our photometry to be incomplete at these faint magnitudes.
We determined the completeness of the optical data in the area of r2-71 by comparing the DAOPHOT II output from the 2004-June-14 observation to those of the 2004-August-15 observation. Figure 3 shows two histograms. The solid histogram shows the percentage of stars detected by the DAOPHOT analysis in the 2004-June-14 observation within 3<sup>โฒโฒ</sup> of the center of the error circle but not detected by the same analysis in the 2004-August-15 observation. The dotted histogram shows the percentage of stars detected by the DAOPHOT analysis in the 2004-August-15 observation within 3<sup>โฒโฒ</sup> of the center of the error circle but not detected by the same analysis in the 2004-June-14 observation. The number of lost stars begins to increase at $`B=25.5`$, suggesting that the completeness of the data begins to decrease at that magnitude, most likely due to the crowding in this dense region of M31. This result is consistent with the failure of DAOPHOT II to detect the $`B=26.0`$ star in the error circle in the 2004-January-23 observation, even though there appears to be emission at that location in the image in Figure 1. Therefore all of the variable candidates in the r2-71 error circle are attributable to crowding and completeness issues, as they are all fainter than $`B=25.5`$.
The lack of any strong variability detection within the r2-71 error circle is confirmed with the difference images shown in Figure 2. The most variable location in the difference images is marked with arrows; this variability is not statistically significant. The DAOPHOT analysis did not measure a brightness increase for this star in the 2004-June-14 data. As a second check for a possible brightness increase, we performed aperture photometry of the location in all three observations with a 0.1<sup>โฒโฒ</sup> radius aperture. The location had $`B=25.07\pm 0.07`$ in the 2004-June-14 observation and $`B=25.23\pm 0.08`$ in the other observations ($`\mathrm{\Delta }B=0.16\pm 0.11`$), showing variability of only 1.5$`\sigma `$. We discounted this low-significance peak in the difference image as a counterpart candidate for r2-71.
Concisely, no variable star inside the error circle of r2-71 was found other than suspicious DAOPHOT detections affected by crowding and completeness. Therefore, we were unable to identify the optical counterpart to r2-71; however, our completeness results suggest that any highly variable counterpart must have had $`B25.5`$ during the 2004-June-14 $`HST`$ observation.
## 4 Discussion
### 4.1 Duty Cycle
Our search of the literature found one previous detection of r2-71 in 1990. The detection was in only one ROSAT observation (source 31 in Primini et al., 1993), and the lack of detections of outbursts of this source before or since (Trinchieri & Fabbiano, 1991; Kong et al., 2002; Williams et al., 2004) helps to constrain the duty cycle of this system. M31 was observed by Einstein in the summers and winters of 1979 and 1980 (Trinchieri & Fabbiano, 1991) and by ROSAT in the summers of 1990 (Primini et al., 1993), 1991 (Supper et al., 1997), 1992 (Supper et al., 2001), 1994, and 1995. We searched the ROSAT HRI images of the M31 bulge taken in the summers of 1994 and 1995 and found no detections of r2-71. It was therefore detected only once in about 9 months of monitoring spanning 7 years before Chandra. These observations constrain the duty cycle of the source to be $`\stackrel{<}{}`$0.1, since a larger duty cycle would have allowed more than one detection in these early data sets.
Now, in 2004, r2-71 has been seen for the first time in about 5 years of Chandra monitoring, going back to late 1999 (Williams et al., 2004; Kong et al., 2002). Assuming the 1990 outburst also lasted $``$1 month, the source has been active for at least 2 months of a 14 year timespan, providing a lower limit on its duty cycle of $`>`$0.01. Adding the $``$8 months per year for 5 years of Chandra monitoring to the 9 months of monitoring before Chandra implies 2 months of activity in $``$49 months of monitoring, or a duty cycle of $``$0.04. Therefore, our best estimate of the duty cycle of r2-71 is $``$0.04, and it can be reliably constrained to the range 0.01โ0.1.
### 4.2 Orbital Period
The HST data provide indirect evidence that r2-71 is an LMXB. This preliminary classification allows us to predict the range in which the orbital period of the system will fall assuming the system is similar to Galactic LMXB transient systems.
The $`HST`$ data set rules against the possibility that the X-ray source is an HMXB. Even the brightest star in the error circle of r2-71 has $`B=24.7`$, which implies M$`{}_{B}{}^{}=0.2`$ (assuming $`mM=24.47`$ and $`A_B=0.4`$). Even so, this $`B`$-band luminosity is fainter than massive O and B type stars. Furthermore, this brightest star is not considered a counterpart candidate because it did not show significant variability. Because r2-71 is not an HMXB, we continue under the assumption that it is an LMXB.
Van Paradijs & McClintock (1994) identified an empirical correlation between the optical/X-ray luminosity ratios of LMXBs in outburst and their orbital periods. Their model assumes that the optical emission arises from X-ray heating of the accretion disk. Larger disks form in longer period systems and glow brighter in the optical than smaller disks. In this model, the faint upper-limit on the optical brightness of r2-71 would suggest that it is a small accretion disk system with a short orbital period. Counter-examples to the correlation exist, like XTE J1118+480 (Williams et al., 2005a); however, such counter-examples usually stand out as odd in other ways. For example, XTE J1118+480 was fainter and harder than typical XRNe. Therefore, for the purposes of the prediction, we assume that r2-71 is an XRN similar to the many Galactic LMXB transient events that fit the correlation well. A few of these are described in detail below.
We checked the applicability of the correlation to the specific case of our data because our X-ray and optical data are not precisely contemporaneous. If the errors in the correlation and optical luminosity are taken into account, our investigations show that the correlation provides reliable period range predictions for both โclassicalโ and more recently discovered Galactic XRNe, even if the optical luminosity is measured 3 weeks after the X-ray luminosity.
First, we checked the application of the correlation to โclassicalโ Galactic XRNe, those with smooth exponential decays (e.g. A0620-00, Nova Mus, GRO 0422+32, etc.). These types of events have optical decay timescales that average $``$2.2 times longer than their X-ray decay timescales (Chen et al., 1997). For example, A0620-00 has an $`e`$-folding optical decay time of $``$75 days and an X-ray decay time of $``$25 days (Esin et al., 2000), so that its optical flux decreases by 25% in 3 weeks, for a 0.3 mag change. Such an effect is small compared to the large dynamic range of the correlation, which covers 8 optical magnitudes and 3 orders of magnitude in X-ray luminosity. The X-ray flux of A0620 in outburst was $``$50 Crab (Esin et al., 2000), which translates to $``$8$`\times `$10<sup>37</sup> erg s<sup>-1</sup> at the appropriate distance (1.05$`\pm `$0.40 kpc; Shahbaz et al., 1994). Applying the 75 day optical decay time to the peak optical magnitude ($`V=11.2`$; Liu et al., 2001), A0620 was $`V=11.5`$ three weeks after peak. Assuming an extinction of $`A_V=1.2`$ (Liu et al., 2001), M<sub>V</sub>=0.2$`\pm `$0.9 at that time. The orbital period prediction from these values is 0.8$`{}_{0.7}{}^{}{}_{}{}^{+7.6}`$ days, which is consistent with the known period of 0.3 days (Liu et al., 2001). Similar results are seen when the correlation is applied to other classical systems, including, for example, 4U 1543-47 (Williams et al., 2005a). The predicted period range is therefore reliable for classical XRNe even if the optical observation is 3 weeks after the X-ray observation.
In addition, the correlation even holds for several more recent X-ray transient sources that have exhibited complex light curves, such as GRO J1655-40 and XTE 1550-564 (Williams et al., 2005a). We tested the effects of the 3 week interval between the X-ray and optical observations of r2-71 on our period prediction using the complicated optical and X-ray lightcurves of the recent XRN XTE J1550-564. Inspection of the lightcurves of Jain et al. (2001) suggests that if we observed XTE J1550-564 $``$8 days after its X-ray peak, when its X-ray luminosity was $``$4$`\times `$10<sup>37</sup> erg s<sup>-1</sup> (for a distance of 5.3$`\pm `$2.3 kpc; Orosz et al., 2002), and then observed the location in the optical 3 weeks later, we would have seen the counterpart at $`V19`$. Applying an extinction of $`A_V=4.75`$ (Orosz et al., 2002) implies M$`{}_{V}{}^{}=0.6_{0.8}^{+1.2}`$. Putting these numbers into the empirical correlation provides a period prediction of $`0.7_{0.6}^{+4.8}`$ days. If we were fortunate enough to catch the brightest X-ray flare, with a flux of 1.6$`\times `$10<sup>-7</sup> erg cm<sup>-2</sup> s<sup>-1</sup> (Sobczak et al., 2000), the source luminosity would have been L$`{}_{X}{}^{}5\times `$10<sup>38</sup> erg s<sup>-1</sup>. In this case, we would have measured $`V17.5`$ three weeks later (Jain et al., 1999), 0.9 mag fainter than the peak of $`V=16.6`$ (Liu et al., 2001). The later optical measurement would have yielded M$`{}_{V}{}^{}=0.9_{0.8}^{+1.2}`$ and a period prediction of $`1.1_{1.0}^{+13.3}`$ days. The actual period of XTE J1550-564 is 1.55 days (Orosz et al., 2002), within the predicted range.
Succinctly, the HST data show that r2-71 is not an HMXB, and therefore may be an LMXB. The van Paradijs & McClintock (1994) correlation provides reliable orbital period range predictions for such objects even when the observations are separated by 3 weeks and the relation is applied to a complex transient lightcurve. We therefore apply the correlation to our measurements of r2-71 under the assumption that, as in the above Galactic examples, the errors in absolute $`V`$ magnitude and in the correlation are sufficient to account for complications in the lightcurve and the 3-week gap between X-ray and optical observations. We note that these predictions rely on the assumption that r2-71 behaves in a similar way to many Galactic XRNe.
Our $`B`$-band brightness upper-limit of $`B25.5`$, from our completeness results, can be converted to a $`V`$-band luminosity by assuming the same foreground extinction we assumed for the X-ray spectral fit and converting to optical extinction using the relation of Predehl & Schmitt (1995). Assuming $`mM=24.47`$, M$`{}_{B}{}^{}0.6`$. Then using the mean $`BV`$ colors of Galactic LMXBs in the Liu et al. (2001) catalog (-0.09 +/- 0.14), M$`{}_{V}{}^{}0.5`$. Placing this upper-limit on the optical luminosity and our 0.3โ7 keV X-ray luminosity of 3.9 $`\times 10^{37}`$ erg s<sup>-1</sup> into the van Paradijs & McClintock (1994) correlation, including their quoted errors, we obtain a prediction for the period of the LMXB system r2-71 of $`P1.6`$ days.
## 5 Conclusions
We have constrained the X-ray and optical properties of a repeating X-ray transient source in the M31 bulge, which we have named CXOM31 J004233.4+411703 or r2-71. This source has undergone at least two X-ray outbursts brighter than 10<sup>37</sup> erg s<sup>-1</sup> in the past two decades. Previous X-ray observations reveal that the source has varied by at least a factor of 100 in X-ray luminosity, and our Chandra monitoring program shows that the outburst in May of 2004 had an $`e`$-folding decay time of less than a month. The observed activity of the source from 1979 to the present suggests that it has a duty cycle of 0.04$`{}_{0.03}{}^{}{}_{}{}^{+0.06}`$.
Optical observations of the location of r2-71 with $`HST`$ ACS before, during, and after the X-ray outburst show no clear optical counterpart to this transient X-ray event in the M31 bulge. The stellar content of the region rules out the presence of an HMXB transient system at the location of r2-71. We therefore assume r2-71 is an LMXB system. No reliable variability was detected in the r2-71 error circle, so that we did not detect the optical counterpart of the XRN. A difference image of the region confirms the lack of significant variability. The optical data therefore place an upper-limit on the $`B`$-band brightness of the outburst of $`B25.5`$. The corresponding upper-limit on the $`V`$-band luminosity along with the X-ray luminosity measured from the Chandra spectrum provide a prediction of $``$1.6 days for the orbital period of the LMXB system.
Support for this work was provided by NASA through grant number GO-9087 from the Space Telescope Science Institute and through grant number GO-3103X from the Chandra X-Ray Center. MRG acknowledges support from NASA LTSA grant NAG5-10889. SSM acknowledges the support of the HRC contract NAS8-03060. JEM acknowledges the support of NASA grant NNG0-5GB31G.
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# CO in HI Self-Absorbed Clouds in Perseus
## 1 Introduction
The study of cold (T $`<100`$ K) atomic gas, a major component of the interstellar medium in the Galaxy, is a difficult problem. While HI emission lines can be used to easily map the distribution of the atomic gas, it is often difficult to separate the warm and cold components. On the other hand, direct observations of cold HI gas can be obtained from HI lines seen in absorption against warm, background HI emission (called HI self-absorption or HISA). Studies of HISA in molecular clouds have shown that the HISA is often well mixed with molecular gas (Jackson et al. 2002, Li & Goldsmith 2003).
Using data from the Canadian Galactic Plane Survey (CGPS; Taylor 1999, Taylor et al. 2003), Gibson et al. (2000; 2005a; 2005b) have revealed HISA in unprecedented detail over a square degree in Perseus. In Gibson et al. (2000; hereafter G2000) two regions in the Perseus arm were found to be of particular interest. The first, labeled the $`globule`$, is a compact (unresolved in the 1 main beam of the CGPS), dark (absorption line center contrast $`>`$ 42 K; Gibson et al. 2005b) region with narrow ($`\mathrm{\Delta }V_{FWHM}=2.5`$ km s<sup>-1</sup>) HISA substructure. IRAS 60 $`\mu `$m dust continuum emission was observed in the globule, but there was no detected <sup>12</sup>CO J=1$``$0 emission (Heyer et al. 1998). The second, labeled the $`complex`$, is part of a larger HI region. It also shows a deep HISA feature (line center contrast $``$ 28 K; Gibson et al. 2005b) and 60 $`\mu `$m emission. However, unlike the globule, it has detected <sup>12</sup>CO J=1$``$0 emission (T$`{}_{R}{}^{}`$ 1.12 K; Heyer et al. 1998). Figure 1 shows the general region of study, with both the globule and complex labeled. Unfortunately, while HISA lines can be readily identified throughout the Galaxy, it is difficult to extract physical parameters (density, temperature, column density etc.) from the observations.
Using a new technique, G2000 were able to set limits on mass, temperature, optical depth, density and column density in the globule and the complex. However, the range of these limits was quite large, mainly due to the unknown molecular gas fraction in the two regions. In the globule, no <sup>12</sup>CO J=1$``$0 was detected to the limiting sensitivity of the Outer Galaxy Survey of Heyer et al. (1998) (hereafter OGS). In the complex, <sup>12</sup>CO J=1$``$0 was detected but, with only one molecular transition, it is difficult to determine the total <sup>12</sup>CO column density. Therefore, in this paper, we present observations of <sup>12</sup>CO J=1$``$0 and J=2$``$1 in both the globule and the complex in order to determine the molecular gas content. Our goals are to complement the HISA observations presented in G2000, and to better constrain the physical properties of the gas in these two regions.
## 2 Observations
Our observations of the J=1$``$0 transition of <sup>12</sup>CO ($`\nu `$ = 115.3 GHz) and <sup>13</sup>CO ($`\nu `$ = 110.2 GHz) were obtained at the Five Colleges Radio Astronomy Observatory (FCRAO) in 2000 February and May. These <sup>12</sup>CO data were obtained in order to either detect the J=1$``$0 or, at least, achieve a better sensitivity than the OGS data. The <sup>13</sup>CO data were obtained to help determine whether the multiple spectral features seen in some of the <sup>12</sup>CO observations were due to multiple line-of-sight clouds or self-absorption in the <sup>12</sup>CO. The FCRAO 14 m telescope has a full width half maximum (FWHM) beamsize of 45.5<sup>โฒโฒ</sup> at 115 GHz, and main beam efficiency ($`\eta _{mb}`$) of 42%. The velocity resolution (channel width) of these observations is 0.252 km s<sup>-1</sup>.
The J=2$``$1 transition of <sup>12</sup>CO ($`\nu `$ = 230.6 GHz) was observed at the Caltech Submillimeter Observatory (CSO) in 2001 September. The CSO is a 10.4 m telescope on Mauna Kea, Hawaii, with a FWHM beamsize of 32<sup>โฒโฒ</sup>, and $`\eta _{mb}`$ of 69% (at 230 GHz). The velocity resolution of these observations was 0.063 km s<sup>-1</sup>.
In the globule, we observed a single point at the central position in the J=1$``$0 transition and an 8 point map with full CSO (32<sup>โฒโฒ</sup>) beam spacings in the J=2$``$1 transition. A ninth position at a half beam spacing ($`\mathrm{\Delta }\alpha =0,\mathrm{\Delta }\delta =16^{\prime \prime }`$) was also observed at the CSO. In the complex, we observed a 9 point strip in both <sup>12</sup>CO transitions, as well as the J=1$``$0 transition of <sup>13</sup>CO. The strip starts at the central position (labeled as Position 1 in Figure 1), where <sup>12</sup>CO J=1$``$0 emission had been detected in the OGS, and moves outwards at constant Galactic longitude but increasing Galactic latitude (from $`b`$ = +0.91 to $`+0.97^{}`$ in steps of 25<sup>โฒโฒ</sup>) to where <sup>12</sup>CO J=1$``$0 had not been previously detected. The spacing of the complex observations were approximately equal to FCRAO half beam spacings. The observed positions in both the globule and the complex are shown in Figure 1.
## 3 Results
In Sections 3.1 and 3.2 that follow, we first provide details of the observed spectra and then show how the results are used to obtain the molecular gas column densities.
### 3.1 Description of the Spectra
Figure 2 shows a postage stamp map of the <sup>12</sup>CO J=2$``$1 transition (solid line) and the HISA (dotted line) in the globule, with the single half-beam spaced ($`\mathrm{\Delta }\delta `$ = -16<sup>โฒโฒ</sup>) observation shown in the middle panel for clarity. Gaussian profiles were fit to the spectra, the parameters of which are listed in Table 1 under the heading 32<sup>โฒโฒ</sup> Resolution. Positions with no detected emission are listed with their $`1\sigma `$ rms noise limits.
To facilitate comparison of the J=1$``$0 and J=2$``$1 transitions, the J=2$``$1 observations were convolved to match the 45<sup>โฒโฒ</sup> beam of the FCRAO. The far right panel in Figure 2 shows the convolved <sup>12</sup>CO J=2$``$1 spectrum (solid line) along with the FCRAO J=1$``$0 spectrum (dashed line) at the central position of the globule. The <sup>12</sup>CO J=2$``$1 emission is clearly visible in the convolved spectrum, but below the peak-to-peak noise level of the J=1$``$0 observations. This shows that molecular gas is present, but that even our deeper <sup>12</sup>CO J=1$``$0 observations did not reach the sensitivity limit required to detect it. A Gaussian profile was fit to the J=2$``$1 transition, the parameters of which are listed in Table 1 under the heading Convolved to 45<sup>โฒโฒ</sup> resolution. Since there was no signal detected in the J=1$``$0 transition, the 1$`\sigma `$ rms noise was used as an upper limit on signal strength.
Figure 2 clearly shows that, while the HISA component of the globule lies at $`V_{LSR}`$ = -41 km s<sup>-1</sup>, the <sup>12</sup>CO emission is centered at $`V_{LSR}`$ = -45 km s<sup>-1</sup>. This suggests that these two emission lines are tracing separate clouds along the line of sight. This suggestion is supported by recent observations of the region with an increased signal to noise ratio in the <sup>12</sup>CO J=1$``$0 transition (Brunt, 2005) which reveal a number of weak <sup>12</sup>CO knots in the vicinity of the globule at a $`V_{LSR}`$ of -45 km s<sup>-1</sup>. To draw comparisons with the atomic globule at $`V_{LSR}`$ -41 km s<sup>-1</sup> we use the 1 $`\sigma `$ rms noise limit for both the J=1$``$0 and J=2$``$1 transitions ($`T_{mb}`$ = 0.58 and 0.06 K respectively) as an upper limit on line strength. As an approximation to a line width, we use the FWHM of the HISA lines ($`\mathrm{\Delta }`$V<sub>FWHM</sub> = 2.5 km s<sup>-1</sup>). This is a valid assumption if the atomic and molecular gas components are mixed and the line widths are dominated by turbulence. This assumption is supported by the similarities between the HISA and <sup>12</sup>CO line widths seen in the complex. These line parameters are used in Section 3.2.1 to constrain the molecular gas content of the globule.
The bottom half of each panel in Figure 3 shows the HISA as a function of position in the complex. The top halves show the <sup>12</sup>CO J=2$``$1 (solid line), <sup>12</sup>CO J=1$``$0 (dashed line), and <sup>13</sup>CO J=1$``$0 (dotted line) at the same positions. Note that Figure 3 is a single strip of observations (of increasing Galactic latitude) and is only presented as what appears to be a nine point map for readability. Figure 3 shows interesting changes in the spectral line profiles progressing from the cloud core (Position 1) to the edge (Position 9). The first two positions show emission peaks at $``$ -40 km s<sup>-1</sup> with blueshifted shoulders (or secondary peaks) at $``$ -41 km s <sup>-1</sup>. The third position is singly peaked, while the fourth through sixth positions have emission peaks at -41 km s<sup>-1</sup> with redshifted shoulders (or secondary peaks) at -40 km s<sup>-1</sup>. The emission drops off rapidly after Position 6, with Positions 7 through 9 showing no significant signal. <sup>13</sup>CO and <sup>12</sup>CO observations show similar line profiles, which is suggestive of emission from multiple clouds (as opposed to self absorption). We will test this assumption in the next section. Gaussian profiles were fit to the spectra in the complex and are listed in Table 2. Note that Positions 1 and 2 were fit with 2 separate Gaussians.
Figure 4 shows the result of smoothing the <sup>12</sup>CO J=2$``$1 in the complex to the velocity resolution of the HI observations (solid lines of <sup>12</sup>CO J=2$``$1 are overlayed on the dotted HISA profiles). This spectral resampling blends the multiple <sup>12</sup>CO components (seen best in Positions 1 and 2) into one component. While there is general agreement between the HISA and the smoothed <sup>12</sup>CO emission lines, the profiles do not match exactly. This could be due to a number of different factors: a greater spatial extent for the HISA than for the <sup>12</sup>CO, different beamsizes of the observations, or errors in the background subtraction of the HISA features. Positions 7 though 9 are not shown due to the lack of molecular emission at those positions.
### 3.2 Molecular Gas Column Densities
There are a number of techniques for solving the equations of radiative transfer and detailed balance in molecular clouds. However, to determine the physical parameters of the molecular components of the globule and the complex, we used a Large Velocity Gradient (LVG; i.e. Goldsmith et al. 1983) model which offers a rudimentary approach to determining the bulk properties of a region based on the integrated intensities of the observed lines. These models employed the collision rates of Flower & Launay (1985) and Shinke et al. (1985).
In this section, we describe the method by which we determined the molecular gas column densities from our CO observations. The results of this analysis will be used in Section 4 to constrain the physical properties of the HISA clouds, as well as the atomic gas fraction.
#### 3.2.1 The Globule
To determine the molecular gas content in the HISA globule itself, we are interested in <sup>12</sup>CO emission at V<sub>LSR</sub> = -41 km s<sup>-1</sup>. However, since there is no detectable signal at this velocity we set an upper limit to the <sup>12</sup>CO column density by using the 1$`\sigma `$ rms noise of the <sup>12</sup>CO J=2$``$1 transition ($`T_{mb}`$ = 0.06 K). In lieu of any actual CO detection in the globule, we used a range of parameters chosen to cover the likely range of physical conditions. Thus, for the globule, we used kinetic temperatures from 8 K to 50 K and fixed the density at $`10^2`$, $`10^3`$, and $`10^4`$ cm<sup>-3</sup> to check for differences between sub-thermal and LTE excitation of the <sup>12</sup>CO lines. For each temperature-density combination, we used our LVG code to calculate the <sup>12</sup>CO J=2$``$1 brightness for a series of 50 different column densities. The column densities ranged from $`5\times 10^{14}5\times 10^{17}`$ cm<sup>-2</sup> and were logarithmically-spaced. Thus, by comparing our 1$`\sigma `$ noise limit to the LVG models, we were able to set an upper limit to the <sup>12</sup>CO column density for each temperature-density combination. The maximum <sup>12</sup>CO column density, corresponding to the lowest temperature-density combination (8K and 100 cm<sup>-3</sup> respectively), is $`9.5\times 10^{15}`$ cm<sup>-2</sup>.
At such small column densities, <sup>12</sup>CO self-shielding is very inefficient and the H<sub>2</sub>/<sup>12</sup>CO abundance ratios can be strongly affected. However, using the photodissociation region (PDR) models of van Dishoeck & Black (1988), and assuming that the strength of the UV field is that of the average interstellar radiation field (i.e. G<sub>o</sub> = 1), we can estimate the total H<sub>2</sub> column density in the globule. For our maximum <sup>12</sup>CO column density of $`9.5\times 10^{15}`$ cm<sup>-2</sup>, the van Dishoeck & Black curves provide an H<sub>2</sub> column density of N(H<sub>2</sub>) $`<1.1\times 10^{21}`$ cm<sup>-2</sup>.
#### 3.2.2 The Complex
For the complex, we only modeled the spectrally smoothed data since they are most directly comparable to the HISA spectra. In our models we used kinetic temperatures of 12 K to 50 K (at temperatures lower than 12K we were unable to find a LVG fits to our <sup>12</sup>CO data). At each temperature we created a $`50\times 50`$ logarithmically-spaced grid of LVG models in density-column density parameter space. The densities ranged from $`10^25\times 10^4`$ cm<sup>-3</sup>, and the <sup>12</sup>CO column densities from $`5\times 10^{15}5\times 10^{18}`$ cm<sup>-2</sup>. The observed line intensities were fit to the grid of LVG models using a $`\chi ^2`$ minimization routine to find the density$``$column density combination that best fit the observations at each temperature. Again, comparing our <sup>12</sup>CO column densities to the models of van Dishoeck & Black (1988), we also determined the H<sub>2</sub> column densities. The results of this LVG analysis will be presented in Section 4.1.
To test our assumption that the multiple velocity components seen in the complex were due to separate clouds rather than self-absorption in a single cloud, we also ran an extensive series of Accelerated Monte Carlo (AMC; Hogerheijde & van der Tak, 2000) models. AMC models have the advantages of producing model spectra based on the input physical parameters of the cloud (i.e. temperature, density and velocity gradient), and a greater available range of parameter space. While LVG models incorporate the simplifying assumption that the concentric model shells are radiatively decoupled, this is not the case for the more robust AMC models.
We first attempted to model the cloud as a single entity with a variety of density, temperature, abundance, and velocity gradients. Some models, in which the abundance ratio was varied with a power law on the order of 0.5 came close to matching the <sup>12</sup>CO spectra, but no single model (for any position) was able to match the <sup>12</sup>CO and <sup>13</sup>CO line profiles and ratios, despite attempts at varying the <sup>12</sup>CO/<sup>13</sup>CO abundance ratio. In addition, none of these models were capable of matching the shift in the red-blue asymmetry discussed previously. In all, we ran several thousand models. The total parameter space covered is shown in Table 3. Column 1 shows the parameters that were varied and column 2 shows the range of central values used for each parameter. Models were run as powerlaw functions of the form $`x(r)=x_o\left(\frac{r_o}{r}\right)^\alpha `$ where $`x_o`$ is the value of the parameter (abundance, temperature, density) at some characteristic radius ($`r_o`$), $`\alpha `$ is the power law index , and $`r`$ is the radius. $`\alpha `$ was allowed to vary between -2 and +2, except for variations in the density for which $`\alpha >0`$ would be unphysical. We also tried modeling each individual position as a separate and distinct โcloudletโ, with its own unique temperature, density and velocity gradients. Again, none of the models were able to reproduce the spectra across the entire complex. Therefore, we are reasonably confident that the multiple velocity components seen in the spectra are due to separate line of sight clouds.
## 4 Discussion
### 4.1 Atomic and Molecular Gas Components in HISA Clouds
Gibson et al. (2000) presented a preliminary view of HISA features seen in the CGPS, finding a number of small scale HISA features that had not previously been seen in single-dish surveys. In more recent papers (Gibson et al. 2005a & 2005b) present a much more detailed view of HISA features through a re-analysis of the CGPS data, using automated algorithms to locate HISA features in the HI data cube instead of the โby-eyeโ identification used in the earlier paper. In addition, Gibson et al (2005b) contains a revision of the HISA data presented in G2000. In the earlier paper the HISA amplitudes were in error; the correct HISA amplitudes are lower than those listed in G2000. In this section we use these improved HISA observations and the technique presented in Gibson et al. (2000) to determine the physical properties of HISA gas.
The technique presented by G2000 allows one to constrain the physical properties of the HISA gas in the complex and the globule by placing limits on the optical depth and spin temperature of the HISA features. The optical depth of the HISA can be determined through:
$$\tau _{_{HISA}}=\frac{CN_{tot}f}{T_K\mathrm{\Delta }v}$$
(1)
Where $`C`$ is combination of constants for the HI spin-flip transition = 5.2$`\times 10^{19}`$ (Dickey & Lockman 1990), $`\mathrm{\Delta }v`$ is the line width, $`T_K`$ is the kinetic temperature of the gas which, under the assumptions of G2000 should be the same as $`T_s`$, the Hi spin temperature, and $`N_{tot}`$ is the total column density (i.e. $`N_{atomic}+N_{molecular}`$).
This analysis was able to find families of solutions based upon allowed ranges of the unknown atomic gas fraction ($`f=`$N(HI)/\[N(HI) + N(H<sub>2</sub>)\]) and $`p`$, the fraction of HI emission originating behind the HISA cloud. Figure 6 of G2000 shows their results. Even if we make the reasonable assumption that $`p`$ is close to unity, indicating that most of the cold HI is in the foreground as would be needed for strong HI self-absorption, the unknown atomic gas fraction still limits our ability to determine the amount and temperature of the HISA gas. For example, Figure 6 of G2000 shows that in the complex, if the HISA is warm ($`T_k>40`$ K) then $`f1`$ and the HISA opacity ($`\tau _{HISA}`$) is approximately 1 to 2, whereas if the HISA is cold ($`T_k<10`$K) then $`f<0.1`$ and $`\tau _{HISA}`$ is constrained to$`0.40.5`$.
With our LVG analysis of the <sup>12</sup>CO observations, we are able to better constrain the HISA properties ($`f`$, $`T`$, and $`\tau `$) in the globule and complex, by finding the overlap between the original parameter space of G2000 and the parameter space defined by our <sup>12</sup>CO observations. For the HISA gas, this was done according to the procedure detailed in G2000 but using the corrected data as given in Gibson et al. 2005b. For our <sup>12</sup>CO observations, we took the derived H<sub>2</sub> column density for each position (Section 3.2) and, using the above equation, calculated $`\tau _{HISA}`$ for a range of atomic gas fractions ($`f`$) from 0.01 to 1.
Figure 5 shows an example for the Globule in which we plot $`\tau _{HISA}`$ vs $`T_{spin}`$. For this figure we have assumed that n(H<sub>2</sub>) $`=10^2`$ cm<sup>-3</sup> in our LVG calculation. The dotted curves show the values of $`p`$, the thin, tilted strips (starting at the upper-left corner of the plot) show $`\tau _{HISA}`$ vs $`T_{spin}`$ as determined from the Gibson et al. (2000) analysis of the HISA features for a range of assumed atomic gas fraction ($`f`$) , and the wider, tilted strips (starting in the bottom-right corner of the plot) show $`\tau _{HISA}`$ vs $`T_{spin}`$ as determined from our <sup>12</sup>CO analysis for a range of atomic gas fractions. The bold โsharkfinโ-shaped region shows the union of the HISA and <sup>12</sup>CO solutions for all values of $`f`$, where each individual intersection is for a particular $`f`$ only. As can be seen from Figure 5, the inclusion of the <sup>12</sup>CO data significantly constrains the allowed range of physical conditions in the HISA gas. The โsharkfinโ corresponds to a solution of 8 K $`<T_{spin}<11`$ K and $`0.5<\tau _{HISA}<5.8`$. If we make the further assumption that $`p>0.8`$ (i.e. that most of the Hi emission is in the background as would be required to produce strong HISA features), then $`0.5<\tau _{HISA}<0.8`$.
In addition, since the intersection of HISA and <sup>12</sup>CO solution also corresponds to an intersection of $`f`$, we can constrain the atomic gas fraction, and find that $`f=0.020.06`$. A similar analysis was done for assumed densities of $`10^2`$ and $`10^3`$ cm<sup>-3</sup> in the globule, and for each position in the complex. The results are given in Table 4. Note that no solutions were found for the Complex - Positions 1 & 2. This is probably due to the presence of the strong second spectral feature. Since the LVG analysis was done on the spectrally smoothed data (to match the spectral resolution of the HISA observations) the two line-of-sight clouds are treated as a single component (see Figure 4). The atomic gas fraction in Positions 3 through 6 were found to be $`f=0.020.1`$. The combination of the HISA and molecular observations also allow us to better constrain the molecular gas density and column density since we now have a better constraint on the kinetic temperature. The results are given in Table 5.
Given the small atomic gas fractions, the complex and the globule seem to be predominantly molecular in composition. While this seems obvious for the central positions in the complex where the <sup>12</sup>CO lines are relatively strong, it is less obvious for Position 6, where the <sup>12</sup>CO lines are weak. Nevertheless our analysis suggests that $`>`$ 90% of the gas is molecular. It is even more puzzling in the globule where the <sup>12</sup>CO lines are undetectable and yet our limits on the <sup>12</sup>CO column density suggest that up to 95% to 98% of the gas could be molecular. Thus it is possible that the globule and the edge of the complex are regions where the gas has a large molecular component that is not well-traced by <sup>12</sup>CO. There are, however, possible alternative explanations. Gibson et al (2000) derived the spin temperature of the HISA features via the equation
$$T_s=\sqrt{\frac{Pf_nC\mathrm{\Delta }s}{k\tau _o\mathrm{\Delta }V}}$$
(2)
Two main assumptions that go into producing equation (2) are: a) the volume density is equal to the column density divided by the path length through the cloud ($`\mathrm{\Delta }s`$) and, b) the volume density is related to the spin temperature through the ideal gas law. Therefore, if either the gas pressure $`P`$ is below the canonical value of $`P/k=4000`$ K cm<sup>-3</sup>, or the HISA features are thinner along the line-of-sight than assumed by G2000, then the atomic gas fraction would have to be larger than predicted to produce the same spin temperature. However, to produce gas which is primarily atomic (i.e. $`f1`$) would require either a large drop in gas pressure or significantly foreshortened clouds. While there is no evidence to support the hypothesis of significantly lower than average gas pressures, there is evidence to support the notion that observationally identified clouds may be preferentially elongated along the line of sight rather than foreshortened (e.g. Heiles 1997) due to observational selection effects.
We can estimate the total mass contained in the 45<sup>โฒโฒ</sup> beam of our FCRAO <sup>12</sup>CO $`J=12`$ observations from the solutions given above and using a distance to Perseus of 2 kpc. The upper limit for the total gas mass in the globule (M(HI) + M(H<sub>2</sub>)) is $`1.33M_{}`$ for the given range of temperatures, H<sub>2</sub> column densities, and atomic gas fractions. This range is between 7 and 60 times lower than the Jeans masses calculated for the globule for the given ranges of temperature and density. A similar calculation for the complex finds that each 45<sup>โฒโฒ</sup> beam contains $`34M_{}`$ of material. This is again 7 to 60 times lower than the Jeans mass implying that neither the globule nor the individual positions in the complex are gravitationally bound.
### 4.2 The Nature of the HISA Clouds
So what are these HISA clouds? They appear to be non-gravitationally bound regions of cold, primarily molecular, gas that is not detected in <sup>12</sup>CO. Such regions are not unusual in high-latitude cirrus clouds and have been traced via excess IR emission (e.g. Heiles, Reach, & Koo 1988; Reach, Koo & Heiles 1994). However, most of the cirrus clouds in which <sup>12</sup>CO does not seem to trace the total amount of molecular gas are relatively warm (T$`{}_{K}{}^{}>30`$K), low density (n(H<sub>2</sub>) $`<100`$ cm<sup>-3</sup>) regions in which the HI does not appear to be self absorbed. In contrast, the temperatures (10 K - 25 K) and densities (n(H<sub>2</sub>) $`1001200`$ cm<sup>-3</sup>) in our HISA clouds are similar to those seen in molecular clouds. Thus, the HISA gas has temperatures and densities that bridge the gap between the ambient atomic ISM and the colder, denser molecular medium.
Could these HISA clouds be sites where atomic gas is condensing into the molecular phase required for star formation? Although critical to the evolution of matter in the Galaxy, molecular condensation is poorly understood. Using a model that assumes that molecular clouds form from atomic gas after the passage of shock waves, Bergin et al. (2004) find that the molecular cloud formation timescale is not controlled by the formation rate of H<sub>2</sub> on grains but, rather, by the shielding of molecules from the UV radiation. While the H<sub>2</sub> can self-shield quite efficiently, <sup>12</sup>CO formation requires shielding of the interstellar radiation by dust grains. If the total A<sub>V</sub> is greater than $`0.7`$ then there is enough material present to effectively shield the <sup>12</sup>CO.
Using the 60 and 100 $`\mu `$m IRAS HIRES data, we have estimated the dust temperature and the amount of visual extinction in the globule and at each of our observed positions in the complex. Following the procedure outlined in Wood et al. (1994), we calculated the dust temperature from the ratio of the 60 and 100 $`\mu `$m fluxes, assuming that the dust is optically thin, that the 60 and 100 $`\mu `$m beam solid angles are roughly equivalent, and that the dust emissivity spectral index is 1.5. We find that the dust temperatures are all $`<30`$K. If the dust is optically thin, knowing the temperature allows us to determine the 100 $`\mu `$m dust opacity from the ratio of the observed 100 $`\mu `$m flux to the Planck function. The dust visual extinction was then calculated from the relationship between A<sub>V</sub> and the 100 $`\mu `$m dust opacity provided by Wood et al. (1994), which is a functional fit to the data given in Jarrett et al. (1989). In the complex, the dust visual extinction ranges from A<sub>V</sub> = 3.8 at Position 1 to A<sub>V</sub> = 2.2 at Position 6. For Positions 7 through 9, A<sub>V</sub> drops below 2, as it does in the globule where A<sub>V</sub> = 1.2. Note that these are the visual extinctions through the entire cloud. The edge-to-center visual extinctions, which are directly comparable to the Bergin et al. (2004) models, are half these values. Thus, these results are consistent with the scenario given by Bergin et al (2004) in which the absence of detectable <sup>12</sup>CO in the globule, and in Positions 7 - 9 of the complex, is due to the limited UV shielding provided by the dust in these regions. However, there could still be a considerable amount of molecular gas in Positions 7 - 9 and in the globule, since H<sub>2</sub> can form at considerably earlier times and lower column densities than <sup>12</sup>CO.
Since A<sub>V</sub> in the globule and at the edge of the complex is close to the critical value needed to shield <sup>12</sup>CO from the interstellar radiation field, it is possible that these regions are in the process of forming <sup>12</sup>CO. Thus, over time, <sup>12</sup>CO lines may eventually become observable as the <sup>12</sup>CO abundance continues to increase. Bergin et al. (2004) predict the observed intensities of various transitions as a function of time (as clouds evolve from atomic to molecular). In their model, a cloud with a <sup>12</sup>CO $`J=21`$ line strength less that 0.06 K (our 1$`\sigma `$ detection limit) would be less than $`10^7`$ years old. This โageโ is consistent with the minimum transit time between spiral arms in the outer Galaxy ($`10^7`$ years; Heyer & Tereby 1998).
An alternative scenario is that instead of being molecular clouds in the process of formation, the HISA clouds could represent transient events, or even the dispersal of molecular clouds into the atomic medium. While we cannot rule out these possibilities, the formation scenario seems more likely since the Complex and the Globule both seem to be correlated with a region slightly downstream of the spiral shock wave in the Perseus arm (Gibson et al 2005; Roberts 1972) where the gas is densest. This is precisely the type of region where we would expect H<sub>2</sub> to be condensing from HI.
## 5 Conclusions
Using <sup>12</sup>CO $`J=21`$ observations from the CSO, and <sup>12</sup>CO and <sup>13</sup>CO $`J=10`$ observations from the FCRAO, we have determined the molecular gas content in two regions of HI Self-Absorption (HISA) in Perseus. In the globule we observed a small 8-point map at the CSO to match the 45<sup>โฒโฒ</sup> resolution, single-point $`J=10`$ observation taken with the FCRAO. In the complex we observed a nine point strip from the center of the cloud to the edge of the cloud in increasing galactic latitude. No <sup>12</sup>CO $`J=21`$ emission was detected in the globule to a 1$`\sigma `$ rms limit of 0.06 K. <sup>12</sup>CO was detected in Positions 1 through 7 in the complex but fell below the 1$`\sigma `$ rms noise limit in Position 8 and 9. Positions 1 & 2 were found to contain two spectral features which we determined to be due to separate line-of-sight clouds rather than to <sup>12</sup>CO self absorption.
Using both Large Velocity Gradient and Monte Carlo radiative transfer codes, we were able to determine the molecular gas content in the globule and complex. In the globule N(<sup>12</sup>CO) $`<6.0\times 10^{15}`$ cm<sup>-2</sup>, implying that N(H<sub>2</sub>) $`<9.9\times 10^{20}`$ cm<sup>-2</sup>. In the complex we found that the H<sub>2</sub> column densities ranged from $`1.22.2\times 10^{21}`$ cm<sup>-2</sup>.
By comparing the HISA and <sup>12</sup>CO observations we are able to constrain the physical conditions and atomic gas fraction ($`f`$) of these two regions. In the globule, 8 K $`<T_{spin}<22`$ K and $`0.02<f<0.2`$ depending on whether the (unknown) gas density is $`10^2`$, $`10^3`$, or $`10^4`$ cm<sup>-3</sup>. In the complex, 12 K $`<T_{spin}<24`$ K, $`0.02<f<0.05`$, and the gas density is constrained ($`100<n<1200`$ cm<sup>-3</sup>. These results imply that the gas in the HISA clouds is colder and denser than that usually associated with the atomic ISM and, indeed, is similar to that seen in molecular clouds. The small atomic gas fractions also imply that there is a significant molecular component in HISA clouds, even when little or no <sup>12</sup>CO is detected. The level of <sup>12</sup>CO detected and the visual extinction due to dust is consistent with the idea that these HISA clouds are undergoing a transition from the atomic to molecular phase.
The authors would like to thank Floris van der Tak for his help with the Monte Carlo Models, and the Natural Sciences and Engineering Research Council of Canada for their financial support. The Five College Radio Astronomy Observatory is operated with the permission of the Metropolitan District Commission, Commonwealth of Massachusetts, and with the support of the National Science Foundation under grant AST 01-00793. The CSO is funded under a grant from the National Science Foundation.
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# The non-classical scissors mode of a vortex lattice in a Bose-Einstein condensate
## I Model and a numerical simulation
We consider a weakly interacting quasi two-dimensional (2D) Bose-Einstein condensate in an anisotropic harmonic trap rotating at the angular frequency $`\mathrm{\Omega }`$ quasi2D . In the rotating frame, the trapping potential is static and given by
$$U(x,y)=\frac{1}{2}m\omega ^2\left[(1ฯต)x^2+(1+ฯต)y^2\right]$$
(1)
where $`m`$ is the mass of an atom, $`ฯต>0`$ is the trap anisotropy and $`\omega `$ is the mean oscillation frequency of the atoms. From now on, we shall always remain in the rotating frame.
The condensate is initially in a stationary state of the Gross-Pitaevskii equation:
$$\left[\frac{\mathrm{}^2}{2m}\mathrm{\Delta }+g|\psi |^2+U\mathrm{\Omega }L_z\mu \right]\psi =0$$
(2)
where $`g`$ is the 2D coupling constant describing the atomic interactions Dum , $`\mu `$ is the chemical potential, $`\psi `$ is the condensate field normalized to the number of particles $`N`$ and $`L_z=xp_yyp_x`$ is the angular momentum operator along $`z`$. In this paper, we shall concentrate on the case that the rotation frequency $`\mathrm{\Omega }`$ is large enough so that several vortices are present in the field $`\psi `$, forming a regular array.
The standard procedure to excite the scissors mode is to rotate abruptly the trapping potential by a small angle $`\theta `$ and to keep it stationary afterwards. This is theoretically equivalent to abruptly rotating the field $`\psi `$ by the angle $`\theta `$ while keeping the trap unperturbed:
$$\psi (t=0^+)=e^{i\theta L_z/\mathrm{}}\psi (t=0^{}).$$
(3)
The subsequent evolution of $`\psi `$ is given by the time dependent Gross-Pitaevskii equation.
We have solved this time dependent Gross-Pitaevskii equation numerically, with the FFT splitting techniques detailed in Dum , the initial state being obtained by the conjugate gradient method Modugno . The results for a 4-vortex and a 10-vortex configurations are shown in figure 1 for $`ฯต=0.025`$, for an initial rotation of the trapping potential by an angle of 10 degrees: we see a small amplitude slow oscillation of the vortex lattice as a whole around the new axis of the trap, accompanied by weak internal oscillations of the lattice and of the condensate. This suggests that, in this low $`ฯต`$ limit, the excitation procedure mainly excites the scissors mode, and that this scissors mode has a very low frequency, much lower than the trap frequency $`\omega `$.
## II Classical hydrodynamics
In the approximation of a coarse-grained vorticity, a condensate with many vortices is described by a density profile $`\rho (๐ซ)`$ and a velocity field $`๐ฏ(๐ซ)`$ solving the classical hydrodynamics equations:
$`\left(_t+๐ฏ\mathrm{๐ ๐ซ๐๐}\right)๐ฏ`$ $`=`$ $`{\displaystyle \frac{1}{m}}\mathrm{๐ ๐ซ๐๐}(U+\rho g)2๐\times ๐ฏ๐\times (๐\times ๐ซ)`$ (4)
$`_t\rho +\mathrm{div}\rho ๐ฏ`$ $`=`$ $`0.`$ (5)
The first equation is Eulerโs equation in the rotating frame, including the Coriolis and centrifugal forces. The second one is the continuity equation.
The velocity field $`๐ฏ`$ is the velocity field of the fluid in the rotating frame. In a stationary state, we set $`๐ฏ=\mathrm{๐}`$ which amounts to assuming the solid body velocity field $`๐\times ๐ซ`$ in the lab frame. The corresponding stationary density profile is given by a quadratic ansatz:
$$\rho _0(x,y)=\frac{m}{g}(\mu a_1x^2a_2y^2).$$
(6)
From Eq.(4) one then finds
$$a_{1,2}=\frac{1}{2}\left[\omega ^2(1ฯต)\mathrm{\Omega }^2\right].$$
(7)
What happens after the rotation of the density profile by a small angle? To answer this question analytically, we linearize the hydrodynamics equations around the steady state, setting $`\rho =\rho _0+\delta \rho `$ and $`๐ฏ=\delta ๐ฏ`$:
$`_t\delta ๐ฏ`$ $`=`$ $`{\displaystyle \frac{1}{m}}\mathrm{๐ ๐ซ๐๐}\delta \rho g2๐\times \delta ๐ฏ`$ (8)
$`_t\delta \rho +\mathrm{div}\rho _0\delta ๐ฏ`$ $`=`$ $`0.`$ (9)
At time $`t=0^+`$, $`\delta ๐ฏ=0`$ and, to first order in the rotation angle $`\theta `$, $`\delta \rho =2\theta mxy(a_1a_2)/g`$. The subsequent evolution is given by the time dependent polynomial ansatz:
$`\delta \rho `$ $`=`$ $`{\displaystyle \frac{m}{g}}\left[c(t)๐ซ\delta A(t)๐ซ\right]`$ (10)
$`\delta ๐ฏ`$ $`=`$ $`\delta B(t)๐ซ`$ (11)
where $`\delta A(t)`$ is a time-dependent 2$`\times `$2 symmetric matrix and $`\delta B(t)`$ is a time-dependent 2$`\times `$2 general matrix. These matrices evolve according to
$`\delta \dot{A}`$ $`=`$ $`A\mathrm{Tr}\delta B(t)A\delta B(t)\delta B^T(t)A`$ (12)
$`\delta \dot{B}`$ $`=`$ $`2\delta A(t)+2\mathrm{\Omega }\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right)\delta B(t)`$ (13)
where we have set $`A=\left(\begin{array}{cc}a_1& 0\\ 0& a_2\end{array}\right)`$. The constant term evolves as $`\dot{c}+\mu \mathrm{Tr}\delta B=0`$.
First, one may look for eigenfrequencies of the system solved by $`\delta A`$ and $`\delta B`$. For $`ฯต=0`$, this may be done analytically Cozzini : one finds one mode of zero frequency, and 6 modes of non-zero frequencies:
$$\nu =\pm 2\omega ,\pm \left[(2\omega ^2\mathrm{\Omega }^2)^{1/2}+\mathrm{\Omega }\right],\pm \left[(2\omega ^2\mathrm{\Omega }^2)^{1/2}\mathrm{\Omega }\right].$$
(14)
At this stage, the presence of a zero energy mode for vanishing anisotropy looks promising to explain the low frequency of the scissors mode. At weak but non-zero $`ฯต`$, however, a numerical diagonalization of the resulting 7$`\times `$7 matrix shows that the zero frequency is unchanged, whereas the others change very little. Analytically, one then easily finds the zero-frequency mode for arbitrary $`ฯต`$:
$$\delta A=\mathrm{\Omega }A\delta B=\left(\begin{array}{cc}0& a_2\\ a_1& 0\end{array}\right)$$
(15)
Why is there such a zero-frequency mode in the classical hydrodynamics? We have found that this is because there a continuous branch of stationary solutions of the classical hydrodynamics equations parametrized by two real numbers $`\alpha `$ and $`\beta `$:
$`๐ฏ`$ $`=`$ $`\left(\begin{array}{c}\alpha y\\ \beta x\end{array}\right)`$ (16)
$`g\rho /m`$ $`=`$ $`\mu a_+x^2a_{}y^2`$ (17)
where
$`a_+`$ $`=`$ $`a_1\mathrm{\Omega }\beta +{\displaystyle \frac{1}{2}}\alpha \beta `$ (18)
$`a_{}`$ $`=`$ $`a_2+\mathrm{\Omega }\alpha +{\displaystyle \frac{1}{2}}\alpha \beta .`$ (19)
These real numbers are not independent since they have to satisfy
$$\alpha \beta (\alpha +\beta )+2\alpha a_1+2\beta a_2=0.$$
(20)
This is a second degree equation for $`\beta `$ at a given $`\alpha `$ so it can be solved explicitly, giving rise to two branches. One of them contains the $`๐ฏ=\mathrm{๐}`$ stationary solution as a particular case, with $`\alpha =\beta =0`$; it terminates in a point where $`a_+=a_{}=0`$. Each stationary solution on this branch has a zero-frequency mode.
So what happens after the scissors mode excitation, in the classical hydrodynamics approximation? We have numerically integrated the linearized equations Eqs.(12,13) with the initial conditions specified above these equations (see figure 2). We find a scissors mode oscillation, however at a large frequency close to the $`ฯต=0`$ prediction $`\nu =\left[(2\omega ^2\mathrm{\Omega }^2)^{1/2}\mathrm{\Omega }\right]`$, in disagreement with the simulations of the previous section note\_pour\_carlos .
## III Rotational symmetry breaking: from Goldstone mode to scissors mode
Having failed to find the scissors mode of section I using the classical hydrodynamics approximation we now discard it and turn to a quantum treatment of the problem. Consider a stationary solution of the Gross-Pitaevskii equation with a vortex lattice present, in the case of an isotropic trap, $`ฯต=0`$. This solution clearly breaks the rotational symmetry SO(2) of the Hamiltonian. Since this is a continuous symmetry group, Goldstoneโs theorem guarantees us the existence of a degree of freedom behaving as a massive boson Blaizot . In particular, this implies the existence of a zero energy eigenmode of the condensate, corresponding to the rotation of the wave function in real space. The โmassโ of the Goldstone boson can be shown to be
$$M=_\mathrm{\Omega }\psi ^{}L_z\psi .$$
(21)
The variable $`\theta `$ conjugate to the Goldstone momentum $`P`$ has the physical meaning of being a rotation angle of the lattice with respect to a reference direction. Note that a quantum spreading of this angular variable takes place when time proceeds, as a consequence of the Goldstone Hamiltonian $`P^2/2M`$.
Now, for a non-zero value of the trap anisotropy, there is no SO(2) symmetry any longer in the Hamiltonian, so that the Goldstone mode is turned into a regular mode of the condensate, behaving as a harmonic oscillator with a finite frequency. This means that the variable $`\theta `$ now experiences a potential. An intuitive estimate of this potential is obtained by taking the stationary solution, rotating it by an angle $`\theta `$ and calculating its Gross-Pitaevskii energy in the presence of a trap anisotropy:
$$=\frac{P^2}{2M}+V(\theta )$$
(22)
where we keep only the $`\theta `$ dependent part of the energy:
$$V(\theta )=\frac{1}{2}m\omega ^2ฯต\left[x^2y^2(1\mathrm{cos}2\theta )+2xy\mathrm{sin}2\theta \right]$$
(23)
where $`\mathrm{}=\psi ^{}\mathrm{}\psi `$. We note in passing that the fact that $`\psi `$ is a local minimum of energy imposes that $`V(\theta )0`$ for all rotation angle $`\theta `$ so that
$`xy`$ $`=`$ $`0`$ (24)
$`x^2y^2`$ $``$ $`0.`$ (25)
Quadratising $`V(\theta )`$ around $`\theta =0`$ leads to the prediction that the angle $`\theta `$ oscillates when initially its value is different from zero. The resulting motion of the condensate is an oscillation of the lattice as a whole, i. e. a scissors mode with an angular frequency:
$$\nu _{\mathrm{scissors}}^2ฯต\frac{2m\omega ^2x^2y^2}{_\mathrm{\Omega }L_z}.$$
(26)
This results in a $`\sqrt{ฯต}`$ dependence of the mode frequency in the case when $`x^2y^2`$ does not vanish in the low $`ฯต`$ limit.
It is now clear why classical hydrodynamics does not exhibit this low energy mode: since it does not break SO(2) symmetry, the corresponding Goldstone mode is absent and so the low energy scissors mode does not appear. We note that, in cases where SO(2) symmetry breaking would occur in the hydrodynamics equations, the scissors mode frequency would be predicted to tend to zero, even if the motion of the particles is not quantized; this was indeed shown by an explicit solution of the superfluid hydrodynamic equation for a rotating vortex-free condensate exhibiting rotational symmetry breaking and the $`\sqrt{ฯต}`$ dependence of the scissors mode frequency was seen experimentally Jean\_et\_Sandro .
## IV Analytical results from Bogoliubov theory
### IV.1 Technicalities of Bogoliubov theory
A systematic way of calculating the eigenfrequencies of the condensate is simply to linearize the time dependent Gross-Pitaevskii equation around a stationary solution, taking as unknown the vector $`(\delta \psi ,\delta \psi ^{})`$ where $`\delta \psi `$ is the deviation of the field from its stationary value. One then faces the following operator,
$$_{\mathrm{GP}}=\left(\begin{array}{cc}\frac{\mathrm{}^2}{2m}\mathrm{\Delta }+U+2g|\psi |^2\mathrm{\Omega }L_z\mu & g\psi ^2\\ g\psi ^2& \left[\frac{\mathrm{}^2}{2m}\mathrm{\Delta }+U+2g|\psi |^2\mathrm{\Omega }L_z\mu \right]^{}\end{array}\right)$$
(27)
where the complex conjugate of an operator $`X`$ is defined as $`๐ซ|X^{}|๐ซ^{}=๐ซ|X|๐ซ^{}^{}`$. The scissors mode we are looking for is an eigenstate of this operator, and the corresponding eigenvalue over $`\mathrm{}`$ gives the frequency of the mode.
At this stage, a technical problem appears, coming from the fact that the operator $`_{\mathrm{GP}}`$ is not diagonalizable CastinDum . This can be circumvented by using the number-conserving Bogoliubov theory, in which the eigenmodes of the condensate are the eigenvectors of the following operator:
$$=\left(\begin{array}{cc}Q& 0\\ 0& Q^{}\end{array}\right)_{\mathrm{GP}}\left(\begin{array}{cc}Q& 0\\ 0& Q^{}\end{array}\right)$$
(28)
where $`Q`$ projects orthogonally to the condensate wave function
$$\varphi =\frac{\psi }{N^{1/2}},$$
(29)
and $`Q^{}`$ therefore projects orthogonally to $`\varphi ^{}`$. We shall assume here that the operator $``$ is diagonalizable for non-zero values of the trap anisotropy $`ฯต`$. We shall see in the next subsection that it is in general not diagonalizable for $`ฯต=0`$ when several vortices are present.
One then introduces the eigenmodes of $``$, written as $`(u_k,v_k)`$, such that $`|u_k|^2|v_k|^2=1`$; to each of these eigenmodes of eigenvalue $`ฯต_k`$ corresponds an eigenmode of $``$ of eigenvalue $`ฯต_k^{}`$, given by $`(v_k^{},u_k^{})`$ Houches . We shall assume here that the condensate wave function is a dynamically stable solution of the stationary Gross-Pitaevskii equation, so that all the $`ฯต_k`$ are real. We shall also assume that the condensate wave function is a local minimum of the Gross-Pitaevskii energy functional (condition of thermodynamical metastability) so that all the $`ฯต_k`$ are positive.
### IV.2 An upper bound for the lowest energy Bogoliubov mode
We can then use the following result. For a given deviation $`\delta \psi ^{}`$ of the condensate field from its stationary value, we form the vectors
$$\stackrel{}{e}_q^q\left(\begin{array}{c}Q\delta \psi \\ Q^{}\delta \psi ^{}\end{array}\right)$$
(30)
where $`q`$ is an integer. As shown in the appendix A we then have the following inequality:
$$\mathrm{min}_kฯต_k\left(\frac{\stackrel{}{e}_0,\stackrel{}{e}_1}{\stackrel{}{e}_0,\stackrel{}{e}_1}\right)^{1/2}$$
(31)
where the scalar product $`,`$ is of signature $`(1,1)`$:
$$\left(\begin{array}{c}u\\ v\end{array}\right),\left(\begin{array}{c}u^{}\\ v^{}\end{array}\right)=u^{}u^{}v^{}v^{}.$$
(32)
We apply this inequality, taking for $`\delta \psi `$ the deviation originating from the rotation of $`\psi `$ by an infinitesimal angle: expanding Eq.(3) to first order in the rotation angle, we put
$$\stackrel{}{e}_0=\left(\begin{array}{c}iQL_z\psi \\ iQ^{}L_z^{}\psi ^{}\end{array}\right).$$
(33)
After lengthy calculations detailed in the appendix A, we obtain the upper bound
$$\mathrm{min}_kฯต_k^2\frac{\mathrm{}^2\varphi |(_\theta ^2U)|\varphi }{_\mathrm{\Omega }(\varphi |L_z|\varphi )}$$
(34)
where $`\theta `$ is the angle of polar coordinates in the $`xy`$ plane.
In the limit of a small $`ฯต`$, we assume that the scissors mode is the lowest energy Bogoliubov mode, as motivated in section III. Eq.(34) then gives an upper bound to the scissors mode frequency. An explicit calculation of $`_\theta ^2U`$ shows that the upper bound Eq.(34) coincides with the intuitive estimate Eq.(26):
$$_\theta ^2U=2m\omega ^2ฯต\varphi |x^2y^2|\varphi .$$
(35)
This upper bound suggests two possibilities in the low $`ฯต`$ limit. In the first one, $`\varphi |x^2y^2|\varphi `$ tends to a non-zero value, which implies that the scissors mode frequency vanishes at most as $`ฯต^{1/2}`$. In the second possibility, $`\varphi |x^2y^2|\varphi `$ tends to zero; under the reasonable assumption of a $`\varphi |x^2y^2|\varphi `$ vanishing linearly with $`ฯต`$, this implies a scissors mode frequency vanishing at most as $`ฯต`$. This second case we term โdegenerateโ.
We point out that the degenerate case contains all the cases where the stationary wave function $`\psi `$ at vanishing trap anisotropy has a discrete rotational symmetry with an angle $`\gamma `$ different from $`\pi `$. To show this, we introduce the coordinates rotated by an angle $`\gamma `$,
$`x^{}`$ $`=`$ $`x\mathrm{cos}\gamma y\mathrm{sin}\gamma `$ (36)
$`y^{}`$ $`=`$ $`x\mathrm{sin}\gamma +y\mathrm{cos}\gamma .`$ (37)
The symmetry of $`\psi `$ implies that $`\varphi |x^2y^2|\varphi =\varphi |x^2y^2|\varphi `$ and $`\varphi |xy|\varphi =\varphi |x^{}y^{}|\varphi `$. By expansion of these identities, and assuming $`1\mathrm{cos}2\gamma 0`$, one gets $`\varphi |x^2y^2|\varphi =0`$.
An important question is to estimate the relative importance of the modes other than the scissors mode excited by the sudden rotation of the condensate. As shown in the appendix A, under the assumption that the scissors mode is the only mode of vanishing frequency when $`ฯต0`$, and anticipating some results of the next subsection, the weight of the initial excitation $`\stackrel{}{e}_0`$ on the non-scissors modes behaves as
$$\stackrel{}{e}_0^{\mathrm{non}\mathrm{scissors}}^2=O(ฯต^2)$$
(38)
both in the degenerate and the non-degenerate cases.
### IV.3 Perturbation theory in $`ฯต`$
We now treat the anisotropic part of the trapping potential as a perturbation, $`\delta U=m\omega ^2ฯต(y^2x^2)/2`$. The Bogoliubov operator $``$, considered as a function of $`ฯต`$, can be written as
$$_ฯต=_0+\delta $$
(39)
where $`\delta `$ is a perturbation. Note that the explicit expression of $`\delta `$ cannot be given easily, as it involves not only $`\delta U`$ but also the effect of the first order change of the condensate wave function entering the mean field terms of $``$; but we shall not need such an explicit expression.
Apart from eigenmodes of non-zero eigenfrequency, the operator $`_0`$ has a normal zero energy mode characterized by the vector $`\stackrel{}{e}_n`$ and an anomalous mode characterized by the vector $`\stackrel{}{e}_a`$, in accordance with Goldstoneโs theorem. They are given by
$`\stackrel{}{e}_n`$ $`=`$ $`\left(\begin{array}{c}Q_0L_z\psi _0\\ Q_0^{}L_z^{}\psi _0^{}\end{array}\right)`$ (40)
$`\stackrel{}{e}_a`$ $`=`$ $`\left(\begin{array}{c}Q_0_\mathrm{\Omega }\psi _0\\ Q_0^{}_\mathrm{\Omega }\psi _0^{}\end{array}\right)`$ (41)
where $`\psi _0`$ is the stationary condensate field for $`ฯต=0`$ and $`Q_0`$ is the projector on the $`ฯต=0`$ condensate wavefunction $`\varphi _0`$. As shown in the appendix A, one has indeed $`_0\stackrel{}{e}_n=0`$ and $`_0\stackrel{}{e}_a=\stackrel{}{e}_n`$. Within the subspace generated by $`\stackrel{}{e}_n`$ and $`\stackrel{}{e}_a`$, the operator $`_0`$ has therefore the Jordan canonical form:
$$_0|_{\mathrm{subspace}}=\left(\begin{array}{cc}0& 1\\ 0& 0\end{array}\right).$$
(42)
Note that the normal and the anomalous vectors are, up to a global factor, adjoint vectors for the modified scalar product Eq.(32), in the sense that $`\stackrel{}{e}_n,\stackrel{}{e}_a=M_0`$ where $`M_0`$ is given in Eq.(21) for $`ฯต=0`$.
As usual in first order perturbation theory, one takes the restriction of the perturbation $`\delta `$ to the subspace and one then diagonalizes it:
$$|_{\mathrm{subspace}}=\left(\begin{array}{cc}\delta _{nn}& 1+\delta _{na}\\ \delta _{an}& \delta _{aa}\end{array}\right).$$
(43)
The notation $`\delta _{an}`$ means than one takes the component of $`\delta \stackrel{}{e}_n`$ onto the vector $`\stackrel{}{e}_a`$. Since the adjoint vector of $`\stackrel{}{e}_a`$ is $`\stackrel{}{e}_n/M_0`$ one has
$$\delta _{an}=M_0^1\stackrel{}{e}_n,\delta \stackrel{}{e}_n.$$
(44)
A similar notation is used for $`\delta _{nn}`$ and $`\delta _{aa}`$. Forming the characteristic polynomial of the 2$`\times `$2 matrix Eq.(43) and taking all the $`\delta `$โs to be $`O(ฯต)`$, one realizes that to leading order in $`ฯต`$, the eigenvalues are $`\pm \delta _{an}^{1/2}`$, that is they scale as $`ฯต^{1/2}`$ if $`\delta _{an}0`$.
Finally, we use the exact identity $`\stackrel{}{e_0},\stackrel{}{e_0}=N\varphi |(_\theta ^2U)|\varphi `$, proved in the appendix A; noting that $`\stackrel{}{e}_n`$ and $`\stackrel{}{e}_0`$ differ by terms of first order in $`ฯต`$, expanding $``$ in $`_0+\delta `$ and using the fact that $`\stackrel{}{e},_0\stackrel{}{e}_n=\stackrel{}{e}_n,_0\stackrel{}{e}=0`$ whatever the vector $`\stackrel{}{e}`$, we get for the scissors mode angular frequency
$$\nu _{\mathrm{scissors}}^2=\delta _{an}+O(ฯต^2)=\frac{\varphi _0|(_\theta ^2U)|\varphi _0}{_\mathrm{\Omega }\varphi _0|L_z|\varphi _0}+O(ฯต^2).$$
(45)
When $`x^2y^2_00`$, this coincides with the upper bound Eq.(34) to leading order in $`ฯต`$: the upper bound is then saturated for low $`ฯต`$.
### IV.4 Analytic expressions in the Thomas-Fermi limit for the non-degenerate case
In the Thomas-Fermi limit ($`\mu \mathrm{}\omega `$) there exist asymptotic functionals giving the energy of the vortex lattice as a function of the vortex positions in the isotropic Dum and anisotropic Aftalion1 cases. Here we shall use these functionals to evaluate the frequency of the scissors mode in Eq.(45).
From the general expression (Eq.(2.12) in Aftalion1 ) we take the simplifying assumption that the distances of the vortex cores to the trap center are much smaller that the Thomas-Fermi radius of the condensate, to obtain
$`\mathrm{\Delta }E`$ $`=`$ $`\mathrm{\Delta }E_0+{\displaystyle \frac{\eta ^2m\omega ^2}{2}}\left({\displaystyle \underset{i}{}}(1ฯต)x_i^2+(1+ฯต)y_i^2\right)\left(2\pi |\mathrm{log}\eta |+2\sqrt{2\pi }{\displaystyle \frac{\mathrm{\Omega }}{\eta \omega }}\right)`$ (46)
$`\sqrt{2\pi }\mathrm{}\omega \eta {\displaystyle \underset{ij}{}}\mathrm{log}(d_{ij}/l)+O(ฯต^2)`$
where $`\mathrm{\Delta }E`$ is the energy difference per atom between the configurations with and without vortices, $`\eta =\mathrm{}/\sqrt{2Nmg}=\mathrm{}\omega /(\mu \sqrt{2\pi })1`$ (with $`\mu `$ being the Thomas-Fermi chemical potential of the non-rotating gas), $`(x_i,y_i)`$ are the coordinates of the i<sup>th</sup> vortex core, $`d_{ij}`$ is the distance between the vortex cores $`i`$ and $`j`$ and $`l`$ is a length scale on the order of the Thomas-Fermi radius and independent of $`ฯต`$ tech2 . $`\mathrm{\Delta }E_0`$ is a quantity independent of $`ฯต`$ and of the vortex positions: therefore we shall not need its explicit expression. It is assumed that $`\mathrm{\Omega }\omega `$, and that
$$\mathrm{\Omega }>\mathrm{\Omega }_m=\left(\frac{\pi }{2}\right)^{1/2}\eta \omega |\mathrm{log}\eta |$$
(47)
i.e. the rotation frequency is large enough to ensure that each vortex experiences an effective trapping potential close to the trap center. An important property of the simplified energy functional is that the positions of the vortices minimizing Eq.(46) are universal quantities not depending on $`\mathrm{\Omega }`$ Aftalion1 when they are rescaled by the length $`\lambda `$ o\_vs\_om :
$$\lambda ^2=\frac{\mathrm{}}{2m(\mathrm{\Omega }\mathrm{\Omega }_m)}.$$
(48)
Now we calculate $`x^2y^2`$ using the Hellman-Feynman theorem:
$$\frac{d}{dฯต}E=\frac{1}{2}m\omega ^2(x^2y^2)|\varphi |^2$$
(49)
where $`E`$ is the energy per atom and $`\varphi `$ is the condensate wave function. We have also that
$$\frac{d}{dฯต}E(ฯต=0)=\frac{d}{dฯต}\mathrm{\Delta }E(ฯต=0)$$
(50)
since the vortex-free configuration is rotationally symmetric for $`ฯต=0`$. We thus obtain really\_technical
$$\underset{ฯต0}{lim}\varphi |(x^2y^2)|\varphi =2\sqrt{2\pi }\eta \frac{\mathrm{\Omega }\mathrm{\Omega }_m}{\omega }\underset{i}{}x_i^2y_i^2.$$
(51)
Using this formula we study vortex lattices with up to ten vortices. By considering first the case $`ฯต=0`$, we find that they are all degenerate except the case of two, nine and ten vortices, see figure 3a for $`n_v=9`$ and figure 3b for $`n_v=10`$. Interestingly the $`ฯต=0`$ 10-vortex configuration differs from the one of the full numerical simulation of figure 1. By minimizing the simplified energy functional for a non-zero $`ฯต`$, we find that the 10-vortex configuration of figure 1 ceases indeed to be a local minimum of energy when $`ฯต`$ is smaller than 0.0182: this prevents us from applying Eq.(45) of perturbation theory to the calculation of the scissors mode frequency of the numerical simulations. The 10-vortex configuration of figure 3b ceases to be a local minimum of energy when $`ฯต`$ is larger than 0.0173. The 10-vortex configuration minimizing the energy for $`0.0173<ฯต<0.0182`$ is a distorted configuration which breaks both the $`x`$ and $`y`$ reflection symmetries of the energy functional, but not the parity invariance, see figure 3c.
The last point is to calculate the derivative of the mean angular momentum with respect to $`\mathrm{\Omega }`$. In the Thomas-Fermi limit, for $`ฯต=0`$, and to first order in the squared distance of the vortex cores from the trap center, the angular momentum is given by Dum ; Aftalion1 :
$$\underset{ฯต0}{lim}L_z\mathrm{}n_v\sqrt{2\pi }m\omega \eta \underset{i}{}x_i^2+y_i^2,$$
(52)
where $`n_v`$ is the number of vortices. Using the rescaling by $`\lambda `$ to isolate the $`\mathrm{\Omega }`$ dependence we get
$$\frac{d}{d\mathrm{\Omega }}L_z(ฯต=0)=\sqrt{2\pi }m\omega \eta \frac{d\lambda ^2/d\mathrm{\Omega }}{\lambda ^2}\underset{i}{}x_i^2+y_i^2$$
(53)
where $`d\lambda ^2/d\mathrm{\Omega }=\lambda ^2/(\mathrm{\Omega }\mathrm{\Omega }_m)<0`$. In this way, $`\nu _{\mathrm{scissors}}`$ for low $`ฯต`$ can be calculated from Eq.(45) in the non-degenerate case analytically in terms of the equilibrium positions of the vortices:
$$\nu _{\mathrm{scissors}}^2=4ฯต(\mathrm{\Omega }\mathrm{\Omega }_m)^2\frac{_ix_i^2y_i^2}{_ix_i^2+y_i^2}.$$
(54)
### IV.5 Numerical results in the degenerate case
In the degenerate case we must go to higher order in perturbation theory to get the leading order prediction for the frequency of the scissors mode. A simpler alternative is to calculate numerically the frequency of this mode by iterating the operator $`^1`$ Brachet starting with the initial guess $`\stackrel{}{e}_0`$ defined in Eq.(33) tech . The corresponding numerical results are presented in figure 4 in the non-degenerate case of two vortices, and in figure 5 in the degenerate case of three vortices. In the non-degenerate case, $`\nu _{\mathrm{scissors}}^2`$ is found to scale linearly with $`ฯต`$ for low $`ฯต`$, as expected, and the corresponding slope is is very good agreement with the prediction Eq.(45); the agreement with the asymptotic formula Eq.(54) is poor, which could be expected since the parameters are not deeply enough in the Thomas-Fermi regime o\_vs\_om . In the degenerate case with three vortices, we find numerically that $`\nu _{\mathrm{scissors}}`$ scales as $`ฯต^{3/2}`$ for low $`ฯต`$, which is indeed compatible with Eq.(34) which leads to a scissor frequency upper bound scaling as $`ฯต`$. The fact that a strictly higher exponent ($`3/2>1`$) is obtained shows that some cancellation happens in the next order of perturbation theory, may be due to the threefold symmetry of the vortex configuration.
## V Conclusion
In this paper we have studied the problem of a non-classical scissors mode of a condensate containing a vortex lattice. In 2D simulations of the Gross-Pitaevskii equation we showed that, when such a condensate experiences a sudden rotation by a small angle of the anisotropic harmonic trapping potential, the orientation of its vortex lattice will undergo very low frequency oscillations of the scissors mode type. Motivated by these numerical results we searched for this mode using the well-known classical hydrodynamics approximation where the condensate density is a smooth function showing no sign of the presence of vortices and where their effect on the velocity field is taken into account through a coarse-grained vorticity. We showed that this approximation does not contain a low energy scissors mode. We then were able to explain this discrepancy using the Gross-Pitaevskii equation that treats quantum mechanically the motion of the particles. In this case, the density profile of the vortex lattice breaks the rotational symmetry which naturally gives rise to a Goldstone mode in the limit of an isotropic trap and, for a finite anisotropy it becomes a low energy scissors mode. The existence of a scissors mode at a low value of the mode frequency is therefore a direct consequence of the discrete nature of the vortices in a condensate, which is itself a consequence of the quantization of the motion of the particles.
We obtained quantitative predictions for the mode frequency using two separate methods. First we calculated an upper bound on the frequency of the mode using an inequality involving the Bogoliubov energy spectrum. Second, using perturbation theory in $`ฯต`$, the anisotropy of the trapping harmonic potential, we showed that this inequality becomes an equality to leading order in $`ฯต`$ when the expectation value $`x^2y^2_0`$ taken in the unperturbed state does not vanish; in this case the frequency of the scissors mode tends to zero as $`ฯต^{1/2}`$, and we gave an analytic prediction for the coefficient in front of $`ฯต^{1/2}`$ in the Thomas-Fermi limit. However, in the cases where the expectation value does vanish (which we termed โdegenerateโ), the frequency will be at most linear in $`ฯต`$. We have illustrated this using a three-vortex lattice where the frequency vanishes as $`ฯต^{3/2}`$. Also, in the general case, we have shown that the relative weight of the non-scissors modes excited by a sudden infinitesimal rotation of the trap tends to zero as $`O(ฯต^2)`$ so that the excitation procedure produces a pure scissors mode in the low trap anisotropy limit.
Finally we point out that the existence of the non-classical scissors mode does not rely on the fact that we are dealing with a Bose-Einstein condensate per se. In particular, we expect, based on the very general arguments of section III, that this mode will equally be present in a Fermi superfluid containing a vortex lattice and that its frequency will tend to zero as $`\sqrt{ฯต}`$ in the non-degenerate case just like in its bosonic counterpart.
Laboratoire Kastler Brossel is a research unit of Ecole normale supรฉrieure and Universitรฉ Paris 6, associated to CNRS. This work is part of the research program on quantum gases of the Stichting voor Fundamenteel Onderzoek der Materie (FOM), which is financially supported by the Nederlandse Organisatie voor Wetenschappelijk Onderzoek (NWO).
## Appendix A Derivation of the upper bound on the Bogoliubov energies
In this appendix, we prove several results used to control analytically the frequency and the excitation weight of the scissors mode.
We first demonstrate the inequality Eq.(31). The key assumption is that the the condensate wavefunction is a local minimum of energy and does not break any continuous symmetry (which implies here a non-zero trap anisotropy $`ฯต0`$). Then the Bogoliubov operator $``$ of the number-conserving Bogoliubov theory is generically expected to be diagonalizable, with all eigenvalues being real and non-zero. We then expand $`\stackrel{}{e}_0`$ on the eigenmodes of $``$:
$$\stackrel{}{e}_0=\underset{k_+}{}b_k\left(\begin{array}{c}u_k\\ v_k\end{array}\right)+b_k^{}\left(\begin{array}{c}v_k^{}\\ u_k^{}\end{array}\right)$$
(55)
where $`_+`$ is the set of modes normalized as $`|u_k|^2|v_k|^2=1`$, and which have here strictly positive energies $`ฯต_k>0`$ since the condensate wavefunction is a local minimum of energy Houches . The $`b_k`$โs are complex numbers; the amplitudes on the modes of the $`_{}`$ family (of energies $`ฯต_k`$) are simply $`b_k^{}`$ since $`\stackrel{}{e}_0`$ is of the form $`(f,f^{})`$. Using the fact that for the modified scalar product Eq.(32), different eigenmodes are orthogonal, and each eigenmode has a โnorm squaredโ equal to $`+1`$ for $`_+`$ and $`1`$ for $`_{}`$, we get:
$$\frac{\stackrel{}{e}_0,\stackrel{}{e}_1}{\stackrel{}{e}_0,\stackrel{}{e}_1}=\frac{_{k_+}2ฯต_k|b_k|^2}{_{k_+}2ฯต_k^1|b_k|^2}.$$
(56)
This is simply the expectation value of $`ฯต_k^2`$ with the positive weights $`ฯต_k^1|b_k|^2`$. Hence Eq.(31).
Next, starting from $`\stackrel{}{e}_0`$ given by Eq.(33), we have to calculate $`\stackrel{}{e}_1`$ and $`\stackrel{}{e}_1`$ to obtain Eq.(34).
To calculate $`\stackrel{}{e}_1`$, we take the derivative of the Gross-Pitaevskii equation Eq.(2) with respect to the rotation frequency $`\mathrm{\Omega }`$, which leads to
$$_{\mathrm{GP}}\left(\begin{array}{c}_\mathrm{\Omega }\psi \\ _\mathrm{\Omega }\psi ^{}\end{array}\right)=(_\mathrm{\Omega }\mu )\left(\begin{array}{c}\psi \\ \psi ^{}\end{array}\right)+\left(\begin{array}{c}L_z\psi \\ L_z^{}\psi ^{}\end{array}\right).$$
(57)
We have to get $``$ instead of $`_{\mathrm{GP}}`$: since $`\psi `$ is normalized to the number of particles $`N`$, which does not depend on $`\mathrm{\Omega }`$, we have
$$\psi ^{}_\mathrm{\Omega }\psi =i\gamma $$
(58)
where $`\gamma `$ is real, so that $`_\mathrm{\Omega }\psi =i\gamma \psi /N+Q(_\mathrm{\Omega }\psi )`$ where $`Q`$ projects orthogonally to $`\psi `$. Using the fact that $`(\psi ,\psi ^{})`$ is in the kernel of $`_{\mathrm{GP}}`$, we obtain:
$$\stackrel{}{e}_1=\left(\begin{array}{c}iQ(_\mathrm{\Omega }\psi )\\ iQ^{}(_\mathrm{\Omega }\psi ^{})\end{array}\right)$$
(59)
which allows to get
$$\stackrel{}{e}_0,\stackrel{}{e}_1=_\mathrm{\Omega }\psi ^{}L_z\psi .$$
(60)
In passing, we note that the fact that all $`ฯต_k`$ are positive implies that $`\stackrel{}{e}_0,\stackrel{}{e}_1>0`$ so that the mean angular momentum is an increasing function of $`\mathrm{\Omega }`$, a standard thermodynamic stability constraint for a system in contact with a reservoir of angular momentum Rokhsar .
We now proceed with the calculation of $`\stackrel{}{e}_1\stackrel{}{e}_0`$. Since $`L_z`$ is hermitian,
$$L_z\psi =Q(L_z\psi )\alpha \psi $$
(61)
where $`\alpha `$ is real. Since $`(\psi ,\psi ^{})`$ is in the kernel of $`_{\mathrm{GP}}`$, we conclude that
$$_{\mathrm{GP}}\stackrel{}{e}_0=i[_{\mathrm{GP}},\left(\begin{array}{cc}L_z& 0\\ 0& L_z^{}\end{array}\right)]\left(\begin{array}{c}\psi \\ \psi ^{}\end{array}\right)$$
(62)
where $`[,]`$ stands for the commutator. One calculates the commutator and then computes its action on $`(\psi ,\psi ^{})`$: various simplifications occur so that
$$_{\mathrm{GP}}\stackrel{}{e}_0=\left(\begin{array}{c}\mathrm{}(_\theta U)\psi \\ \mathrm{}(_\theta U)\psi ^{}\end{array}\right)$$
(63)
where $`\theta `$ is the angle of polar coordinates in the $`xy`$ plane and where we used $`L_z=i\mathrm{}_\theta `$. $`\stackrel{}{e}_1`$ results from this expression by projection. We finally obtain
$$\stackrel{}{e}_0,\stackrel{}{e}_1=\mathrm{}^2\psi ^{}(_\theta ^2U)\psi .$$
(64)
In what concerns the scissors mode frequency, the last point is to justify the identities $`_0\stackrel{}{e}_n=0`$ and $`_0\stackrel{}{e}_a=\stackrel{}{e}_n`$, where $`_0`$ is the $`ฯต=0`$ limit of $``$, $`\stackrel{}{e}_n`$ is defined in Eq.(40) and $`\stackrel{}{e}_a`$ is defined in Eq.(41). First, one notes that, within a global factor $`i`$, $`\stackrel{}{e}_n`$ is the limit of $`\stackrel{}{e}_0`$ for $`ฯต0`$. Then taking the limit $`ฯต0`$ in Eq.(63), one immediately gets $`_0\stackrel{}{e}_n=0`$ since $`_\theta U`$ vanishes when the trapping potential $`U`$ is rotationally symmetric. Second, one notes from Eq.(59) that, within a global factor $`i`$, $`\stackrel{}{e}_a`$ is the zero $`ฯต`$ limit of $`\stackrel{}{e}_1`$. Then taking the zero $`ฯต`$ limit of the identity $`\stackrel{}{e}_1=\stackrel{}{e}_0`$ leads to $`_0\stackrel{}{e}_a=\stackrel{}{e}_n`$.
Finally we control the weight of the modes other than the scissors mode that are excited by the sudden infinitesimal rotation of the condensate:
$$\stackrel{}{e}_0^{\mathrm{non}\mathrm{scissors}}^2\underset{k_+\{k_s\}}{}b_k\left(\begin{array}{c}u_k\\ v_k\end{array}\right)+b_k^{}\left(\begin{array}{c}v_k^{}\\ u_k^{}\end{array}\right)^2$$
(65)
where $`\mathrm{}`$ is the usual $`_2`$ norm and $`k_s`$ is the index of the scissors mode. The basic assumption is that the scissors mode is the only one with vanishing frequency in the $`ฯต0`$ limit. We rewrite Eq.(64) as
$$\underset{k_+}{}2ฯต_k|b_k|^2=2m\omega ^2ฯต(x^2y^2)|\psi |^2.$$
(66)
In the degenerate case the right-hand side is $`O(ฯต^2)`$ so that each (positive) term $`ฯต_k|b_k|^2`$ is $`O(ฯต^2)`$. For the normal modes, both $`ฯต_k^1`$ and the mode functions $`u_k,v_k`$ have a finite limit for $`ฯต0`$, which proves $`\stackrel{}{e}_0^{\mathrm{non}\mathrm{scissors}}^2=O(ฯต^2)`$.
In the non-degenerate case the same reasoning leads to a weight being $`O(ฯต)`$. A better estimate can be obtained from
$$\stackrel{}{e}_0,\stackrel{}{e}_1ฯต_{k_s}^2\stackrel{}{e}_0,\stackrel{}{e}_1=\underset{k_+\{k_s\}}{}2\left(ฯต_kฯต_{k_s}^2/ฯต_k\right)|b_k|^2.$$
(67)
Then using the explicit expression for the scalar products in the left-hand side, see Eq.(60) and Eq.(64), and using the result Eq.(45) of the perturbative expansion for $`\nu _{\mathrm{scissors}}=ฯต_{k_s}/\mathrm{}`$, one realizes that the left-hand side of Eq.(67) is $`O(ฯต^2)`$, which leads to a weight on non-scissors modes being $`O(ฯต^2)`$ as in the degenerate case.
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# Acknowledgment
## Acknowledgment
The authors wish to thank I. P. Neupane for useful comments.
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# On the extrinsic topology of Lagrangian submanifolds
## 1. Introduction
The fact that a manifold $`L`$ admits an embedding into a symplectic manifold as a Lagrangian submanifold yields restrictions on the topology of $`L`$. This intrinsic topology of Lagrangian submanifolds has been studied quite extensively, cf. the recent survey .
In this paper, we investigate the *extrinsic topology* of Lagrangian submanifolds: We address the question how a Lagrangian submanifold lies homologically in the ambient symplectic manifold. Among other things this turns into a new Lagrangian intersection mechanism.
In the basic theorem 3.1, we provide for a monotone closed Lagrangian submanifold $`L`$ in a closed symplectic manifold $`M`$ a Floer-theoretic representation of the homomorphisms $`\iota _k:\mathrm{H}_k(L;/2)\mathrm{H}_k(M;/2)`$ for degrees $`k>dimL+1N_L`$, where $`N_L`$ denotes the *minimal Maslov number* of $`L`$ (see definition 2.2). In the second basic theorem 3.3 we remove the restriction on the degree of $`\iota _k`$ for Hamiltonian functions having sufficiently small Hofer norm. If the Lagrangian submanifold $`L`$ is displaceable, theorem 3.1 implies a vanishing result for the homomorphisms $`\iota _k:\mathrm{H}_k(L;/2)\mathrm{H}_k(M;/2)`$:
###### Theorem.
Let $`(M,\omega )`$ be a monotone closed symplectic manifold and $`LM`$ a monotone closed Lagrangian submanifold with $`N_L2`$. If the Lagrangian submanifold $`L`$ is (Hamiltonianly) displaceable, then the homomorphism $`\iota _k`$ vanishes for degrees $`k>dimL+1N_L`$.
###### Corollary.
For any displaceable monotone Lagrangian submanifold $`L`$ with $`N_L2`$,
$$[L]=0\mathrm{H}_m(M;/2).$$
Applying the basic theorems 3.1 and 3.3 to spectral capacities leads to a new mechanism for producing Lagrangian intersection results. Namely, in a symplectic manifold with finite spectral capacity all monotone Lagrangian submanifolds $`L`$ of minimal Maslov number $`N_L>dimL+1`$ intersect each other, cf. theorem 5.5. An instance of this new mechanism is:
###### Corollary.
Any two simply connected Lagrangian submanifolds in $`\mathrm{P}^n\times \mathrm{P}^n`$ intersect.
An example of a simply connected Lagrangian submanifold of $`\mathrm{P}^n\times \mathrm{P}^n`$ is the anti-diagonal $`\overline{\mathrm{\Delta }}=\{(\overline{z},z)z\mathrm{P}^n\}`$.
Organization of the paper. In section 2 we recall the construction of Floer homology for semi-positive symplectic manifolds and the Piunikhin-Salamon-Schwarz isomorphism. Furthermore, we give the definition of the action selector and of spectral capacities via the action filtration on Floer homology. Theorems 3.1 and 3.3 are stated in section 3 and proved in section 4. In section 5 we derive various applications, which are divided into two subsections, on the extrinsic topology and on Lagrangian intersections.
Acknowledgements. The results of this paper are by-products of my Ph.D. thesis and are partly contained therein. I would like to express my gratitude to my advisor Matthias Schwarz.
I was financially supported by the DFG through the Graduiertenkolleg โAnalysis, Geometrie und ihre Verbindung zu den Naturwissenschaftenโ at the University of Leipzig and the Schwerpunktprogramm โGlobale Differentialgeometrieโ, grant SCHW 892/2-1.
## 2. Floer homology and spectral capacities
Despite the fact that we consider here Floer homology only for monotone symplectic manifolds we recall briefly and without proofs the construction of Floer homology for semi-positive symplectic manifolds with the help of Novikov rings. Even though Novikov rings can be avoided when dealing with monotone symplectic manifold, the statements are much clearer when using Novikov rings. All details can be found in .
###### Definition 2.1.
A closed symplectic manifold $`(M^{2m},\omega )`$ is called semi-positive if
$$\omega (A)>0,c_1(A)3mc_1(A)0$$
(SP)
for all $`A\pi _2(M)`$, where $`c_1`$ is the first Chern class of $`(M,\omega )`$.
Most symplectic manifolds treated here are monotone, i.e. they satisfy $`\omega |_{\pi _2(M)}=\lambda c_1|_{\pi _2(M)}`$ for some constant $`\lambda >0`$. Monotone symplectic manifolds are clearly semi-positive.
###### Definition 2.2.
Let $`(M,\omega )`$ be a symplectic manifold and $`LM`$ a Lagrangian submanifold. The minimal area of a non-constant holomorphic sphere in $`M`$ is
$$A_M=inf\left\{_{S^2}s^{}\omega \right|s\text{ is a non-constant holomorphic sphere in }M\}.$$
The minimal area of a non-constant holomorphic disk on $`L`$ is
$$A_L=inf\left\{_{D^2}d^{}\omega \right|d\text{ is a non-constant holomorphic disk with boundary on }L\}.$$
The infimum over the empty set is by definition $`+\mathrm{}`$. Furthermore, the minimal Chern number $`N_M`$ of $`M`$ is the positive generator of the image $`c_1(\pi _2(M))`$ of the first Chern class $`c_1`$, and the minimal Maslov number $`N_L`$ of $`L`$ is the positive generator of the image $`\mu _{\mathrm{Maslov}}(\pi _2(M,L))`$ of the Maslov index $`\mu _{\mathrm{Maslov}}`$. In case that $`c_1`$ or $`\mu _{\mathrm{Maslov}}`$ vanish we set $`N_M=+\mathrm{}`$ or $`N_L=+\mathrm{}`$.
###### Remark 2.3.
The Maslov index equals $`2c_1`$ on spheres and is congruent modulo 2 to the first Stiefel-Whitney class $`w_1(L)`$ of $`L`$ evaluated on the boundary loop of a disk. Since $`w_1(L)=0`$ if $`L`$ is orientable, the Maslov index is an even number for orientable Lagrangian submanifolds. Due to the fact that we set $`N_L=+\mathrm{}`$ in case $`\mu _{\mathrm{Maslov}}0`$ we conclude $`N_L2`$ for all orientable Lagrangian submanifolds.
The Novikov ring. Let $`\mathrm{\Gamma }`$ be the finitely generated Abelian group
$$\mathrm{\Gamma }:=\frac{\pi _2(M)}{\mathrm{ker}c_1\mathrm{ker}\omega }.$$
(2.1)
A grading on $`\mathrm{\Gamma }`$ is defined via $`\mathrm{deg}:\mathrm{\Gamma }`$, $`\mathrm{deg}(A):=2c_1(A)`$. We set $`\mathrm{\Gamma }_k:=\mathrm{deg}^1(k)`$ and define the Novikov ring associated to $`\mathrm{\Gamma }`$ by
$$\mathrm{\Lambda }(\mathrm{\Gamma }):=_k\mathrm{\Lambda }_k(\mathrm{\Gamma }),\text{ where}$$
$$\mathrm{\Lambda }_k(\mathrm{\Gamma }):=\left\{\underset{A\mathrm{\Gamma }_k}{}\lambda _Ae^A\right|\mathrm{\#}\{A\mathrm{\Gamma }_k|\lambda _A0,\omega (A)c\}<\mathrm{}c\},$$
(2.2)
with coefficients $`\lambda _A`$ in $`/2`$. The ring structure is given by
$`\mathrm{\Lambda }_k(\mathrm{\Gamma })\times \mathrm{\Lambda }_l(\mathrm{\Gamma })`$ $`\mathrm{\Lambda }_{k+l}(\mathrm{\Gamma }),`$ (2.3)
$`\left({\displaystyle \underset{A\mathrm{\Gamma }_k}{}}\lambda _Ae^A\right)\left({\displaystyle \underset{B\mathrm{\Gamma }_l}{}}\mu _Be^B\right)`$ $`{\displaystyle \underset{A}{}}{\displaystyle \underset{B}{}}(\lambda _A\mu _B)e^{A+B}`$
$`={\displaystyle \underset{C}{}}\left({\displaystyle \underset{A}{}}\lambda _A\mu _{CA}\right)e^C.`$
$`\mathrm{\Lambda }_0(\mathrm{\Gamma })`$ is a field and, in particular, $`\mathrm{\Lambda }(\mathrm{\Gamma })`$ is a vector space over $`\mathrm{\Lambda }_0(\mathrm{\Gamma })`$.
###### Example (symplectically aspherical case).
If the symplectic manifold $`(M,\omega )`$ is symplectically aspherical, i.e. $`\omega |_{\pi _2(M)}=0`$ and $`c_1|_{\pi _2(M)}=0`$, then
$$\mathrm{\Gamma }=\mathrm{\Gamma }_0=\{0\}\text{and}\mathrm{\Lambda }(\mathrm{\Gamma })=\mathrm{\Lambda }_0(\mathrm{\Gamma })=/2.$$
(2.4)
###### Example (monotone case).
If $`(M,\omega )`$ is monotone, i.e. $`\lambda >0`$ such that $`\omega |_{\pi _2(M)}=\lambda c_1|_{\pi _2(M)}`$, then $`\mathrm{\Gamma }_0=\{0\}`$ and $`\mathrm{\Gamma }`$. This implies that
$$\mathrm{\Lambda }_0(\mathrm{\Gamma })=/2\text{and}\mathrm{\Lambda }(\mathrm{\Gamma })\left(/2\right)[[q]][q^1],$$
(2.5)
i.e. $`\mathrm{\Lambda }(\mathrm{\Gamma })`$ is isomorphic to the ring of Laurent series with coefficients in $`/2`$.
There are various types of Novikov rings for which Floer homology and quantum cohomology can be defined, cf. \[11, chapter 11.1\]. We choose here the field $`\mathrm{\Lambda }_0(\mathrm{\Gamma })`$.
Floer homology over $`\mathrm{\Lambda }_0(\mathrm{\Gamma })`$. The group $`\mathrm{\Gamma }`$ gives rise to a covering $`\mathrm{\Gamma }\stackrel{~}{M}M`$ of the space $`M`$ of *contractible loops* in $`M`$. The elements of $`\stackrel{~}{M}`$ are represented by equivalence classes $`\overline{x}[x,d_x]`$, where $`xM`$ and $`d_x`$ is an extension of the contractible loop $`x`$ to the unit disk $`D^2`$. Two pairs $`(x,d_x)`$ and $`(y,d_y)`$ are equivalent if
$$(x,d_x)(y,d_y)x=y,c_1(d_x\mathrm{\#}\overline{d}_y)=0,\omega (d_x\mathrm{\#}\overline{d}_y)=0,$$
(2.6)
where $`d_x\mathrm{\#}\overline{d}_y`$ is the sphere obtained by glueing the two disks along their common boundary $`x=y`$. The group $`\mathrm{\Gamma }`$ acts on $`\stackrel{~}{M}`$ by concatenating the disk $`d_x`$ with the sphere $`A`$
$$\overline{x}\mathrm{\#}A[x,d_x]\mathrm{\#}A:=[x,A\mathrm{\#}d_x]\text{for all }A\mathrm{\Gamma }\text{ and }\overline{x}\stackrel{~}{M}.$$
(2.7)
For a Hamiltonian function $`H:S^1\times M`$ we define the action functional $`๐_H`$ on $`\stackrel{~}{M}`$ by
$$๐_H([x,d_x]):=_{D^2}d_x^{}\omega _0^1H(t,x(t))๐t$$
(2.8)
and the Hamiltonian vector field $`X_H`$ associated to $`H`$ by $`\omega (X_H,)=dH`$. The time-1-map $`\mathrm{\Phi }_H`$ of the flow induced by the Hamiltonian vector field $`X_H`$ is called a Hamiltonian diffeomorphism. The set of critical points of $`๐_H`$ is
$$\stackrel{~}{๐ซ}(H)=\{[x,d_x]\stackrel{~}{M}\dot{x}(t)=X_H(t,x(t))\}$$
(2.9)
and is graded by the Conley-Zehnder index $`\mu _{\mathrm{CZ}}:\stackrel{~}{๐ซ}(H)`$, $`[x,d_x]\mu _{\mathrm{CZ}}([d_x])`$. We abbreviate $`\stackrel{~}{๐ซ}_k(H):=\mu _{\mathrm{CZ}}^1(k)`$. The covering $`\stackrel{~}{M}`$ is chosen in such a way that the action functional $`๐_H`$ and the Conley-Zehnder index $`\mu _{\mathrm{CZ}}`$ are $``$ resp. $``$-valued.
The action functional and the grading behave as follows under the action of the group $`\mathrm{\Gamma }`$
$$๐_H(\overline{x}\mathrm{\#}A)=๐_H(\overline{x})+\omega (A)\text{and}\mu _{\mathrm{CZ}}(\overline{x}\mathrm{\#}A)=\mu _{\mathrm{CZ}}(\overline{x})+\mathrm{deg}(A).$$
(2.10)
For a non-degenerate Hamiltonian function $`H`$, i.e. if the graph of $`\mathrm{\Phi }_H`$ is transverse to the diagonal $`\mathrm{\Delta }M\times M`$, we define the Floer chain groups by
$$\mathrm{CF}_k(H):=\left\{\underset{\overline{x}\stackrel{~}{๐ซ}_k(H)}{}a_{\overline{x}}\overline{x}\right|\mathrm{\#}\{\overline{x}|a_{\overline{x}}0,๐_H(\overline{x})c\}<\mathrm{}c\}.$$
(2.11)
The Floer chain groups become finite-dimensional vector spaces over the field $`\mathrm{\Lambda }_0(\mathrm{\Gamma })`$ if we set
$`\left({\displaystyle \underset{A\mathrm{\Gamma }_0}{}}\lambda _Ae^A\right)\left({\displaystyle \underset{\overline{x}}{}}a_{\overline{x}}\overline{x}\right)`$ $`:={\displaystyle \underset{\overline{x}}{}}{\displaystyle \underset{A}{}}\lambda _Aa_{\overline{x}}\overline{x}\mathrm{\#}(A)`$ (2.12)
$`={\displaystyle \underset{\overline{x}}{}}\left({\displaystyle \underset{A}{}}\lambda _Aa_{\overline{x}\mathrm{\#}A}\right)\overline{x},`$
(cf. ). In particular, $`๐_H\left(e^A\overline{x}\right)=๐_H(\overline{x}\mathrm{\#}(A))=๐_H(\overline{x})\omega (A)`$, i.e. the Novikov action of $`A`$ decreases the value of the action functional by $`\omega (A)`$.
The boundary operator is defined by counting connecting Floer cylinders. We define
$$(\overline{x},\overline{y};J,H):=\left\{uC^{\mathrm{}}(\times S^1,M)\right|\begin{array}{cc}& _su+J(u)\left(_tuX_H(t,u)\right)=0\hfill \\ & u(\mathrm{})=x,u(+\mathrm{})=y\hfill \\ & d_x\mathrm{\#}ud_y\hfill \end{array}\},$$
(2.13)
where $`d_x\mathrm{\#}ud_y`$ is in the Novikov sense, namely $`c_1(d_x\mathrm{\#}u\mathrm{\#}\overline{d}_y)=0=\omega (d_x\mathrm{\#}u\mathrm{\#}\overline{d}_y)`$, and where $`J`$ is a compatible almost complex structure on $`M`$.
###### Theorem (Floer-Hofer-Salamon ).
For generic choices of $`H`$ and $`J`$ the moduli space $`(\overline{x},\overline{y};J,H)`$ is a finite-dimensional manifold of dimension
$$dim(\overline{x},\overline{y};J,H)=\mu _{\mathrm{CZ}}(\overline{y})\mu _{\mathrm{CZ}}(\overline{x})$$
(2.14)
admitting a free $``$-action if $`xy`$. Moreover, the moduli space
$$\widehat{}(\overline{x},\overline{y};J,H):=(\overline{x},\overline{y};J,H)/$$
(2.15)
is a finite set for relative index one, i.e. if $`\mu _{\mathrm{CZ}}(\overline{y})\mu _{\mathrm{CZ}}(\overline{x})=1`$, and it can be compactified by adding broken solutions for relative index two.
For relative index one we denote the (mod 2) number of elements by
$$n(\overline{x},\overline{y}):=\mathrm{\#}_2\left(\widehat{}(\overline{x},\overline{y};J,H)\right)$$
(2.16)
and define the boundary operator of the Floer complex on generators by
$`:\mathrm{CF}_k(H)`$ $`\mathrm{CF}_{k1}(H)`$ (2.17)
$`\overline{y}`$ $`\overline{y}:={\displaystyle \underset{\mu _{\mathrm{CZ}}(\overline{x})=\mu _{\mathrm{CZ}}(\overline{y})1}{}}n(\overline{x},\overline{y})\overline{x}.`$
We extend it linearly to a $`/2`$โvector space homomorphism. The Floer equation yields
$$n(\overline{x},\overline{y})0๐_H(\overline{x})๐_H(\overline{y}),$$
(2.18)
i.e. our conventions imply that the value of the action functional and the Conley-Zehnder index increase along Floer cylinders. Furthermore, the following compactness result holds
$$\underset{c_1(A)=0,\omega (A)c}{}\mathrm{\#}\widehat{}(\overline{x},\overline{y}\mathrm{\#}A;J,H)<\mathrm{}$$
(2.19)
for all $`\overline{x},\overline{y}\stackrel{~}{๐ซ}(H)`$ with $`\mu _{\mathrm{CZ}}(\overline{y})\mu _{\mathrm{CZ}}(\overline{x})=1`$ and every $`c`$ (cf. ). These two facts imply that $`\overline{y}`$ is a well-defined element in $`\mathrm{CF}_{k1}(H)`$. The identity $`n(\overline{x}\mathrm{\#}A,\overline{y}\mathrm{\#}A)=n(\overline{x},\overline{y})`$ implies that $``$ actually is a $`\mathrm{\Lambda }_0(\mathrm{\Gamma })`$-linear homomorphism. Floerโs fundamental theorem asserts $`=0`$, so that Floer homology groups are defined and form $`\mathrm{\Lambda }_0(\mathrm{\Gamma })`$ vector spaces
$$\mathrm{HF}_{}(H):=\mathrm{H}_{}(\mathrm{CF}_{}(H),).$$
(2.20)
###### Example (monotone case).
We recall that if $`(M,\omega )`$ is monotone, then
$$\mathrm{\Gamma }_0=\{0\}\text{and}\mathrm{\Lambda }_0(\mathrm{\Gamma })=/2.$$
(2.21)
In particular, $`\stackrel{~}{๐ซ}_k(H)`$ is finite, although in general $`\stackrel{~}{๐ซ}(H)`$ will be infinite. Furthermore, each Floer chain group $`\mathrm{CF}_k(H)`$ is a finite-dimensional vector space over the field $`\mathrm{\Lambda }_0(\mathrm{\Gamma })=/2`$. In this case the Novikov conditions are empty. Nonetheless the full Novikov ring $`\mathrm{\Lambda }(\mathrm{\Gamma })\left(/2\right)[[q]][q^1]`$ appears for instance in the following theorem by Hofer-Salamon.
For the moment let us denote the Floer homology over the full Novikov ring $`\mathrm{\Lambda }(\mathrm{\Gamma })`$ by $`\overline{\mathrm{HF}}_{}(H)`$, which is not $``$ but $`/(2N_M)`$-graded, where $`N_M`$ is the minimal Chern number of $`M`$ (cf. definition 2.2). Both homology groups are related by the isomorphism
$$\overline{\mathrm{HF}}_{\overline{k}}(H)\mathrm{HF}_k(H)_{\mathrm{\Lambda }_0(\mathrm{\Gamma })}\mathrm{\Lambda }(\mathrm{\Gamma })\overline{k}/(2N_M).$$
(2.22)
###### Theorem (Hofer-Salamon ).
There exists a $`\mathrm{\Lambda }(\mathrm{\Gamma })`$-module isomorphism
$$\overline{\mathrm{HF}}_{\overline{k}}(H)_{\mathrm{\Lambda }(\mathrm{\Gamma })}\underset{ik\mathrm{mod}\mathrm{\hspace{0.33em}2}N_M}{}\mathrm{H}^{mi}(M;/2)_{/2}\mathrm{\Lambda }_0(\mathrm{\Gamma }).$$
(2.23)
The Piunikhin-Salamon-Schwarz isomorphism. In this section, we recall the construction of the Piunikhin-Salamon-Schwarz isomorphism, which we call $`\mathrm{PSS}`$-isomorphism for brevity. The construction is presented in for semi-positive symplectic manifolds.
Since the $`\mathrm{PSS}`$-isomorphism plays a crucial role in the basic theorems below, we give here a fairly detailed exposition.
The $`\mathrm{PSS}`$-isomorphism is defined via counting solutions of a special boundary value problem whose solutions we call plumberโs helper solutions. The corresponding moduli space $`^{\mathrm{PSS}}(q,\overline{x};J,H,f,g)`$ consists of pairs $`(\gamma ,u)`$ of maps
$$\gamma :(\mathrm{},0]M\text{and}u:\times S^1M\text{with }E(u)=_{\mathrm{}}^{\mathrm{}}_0^1|_su|^2๐t๐s<\mathrm{}$$
(2.24)
$$\text{solving}\dot{\gamma }+^gf\gamma =0\text{and}_su+J(u)\left(_tu\beta (s)X_H(t,u)\right)=0,$$
(2.25)
where $`\beta :[0,1]`$ is a smooth cut-off function satisfying $`\beta (s)=0`$ for $`s\frac{1}{2}`$ and $`\beta (s)=1`$ for $`s1`$. The function $`f:M`$ is a Morse function on $`M`$ and $`^g`$ is the gradient with respect to a Riemannian metric $`g`$. Again $`J`$ is a compatible almost complex structure. The pair $`(\gamma ,u)`$ is required to satisfy boundary conditions
$$\gamma (\mathrm{})=q,\gamma (0)=u(\mathrm{})\text{and}u(+\mathrm{})=x,$$
(2.26)
(see remark 2.4 below) where $`q\mathrm{Crit}(f)`$ is a critical point of $`f`$ and $`\overline{x}=[x,d_x]\stackrel{~}{๐ซ}(H)`$. Furthermore, we require the homotopy condition $`\omega (u\mathrm{\#}\overline{d}_x)=0=c_1(u\mathrm{\#}\overline{d}_x)`$. For brevity we denote this moduli space by $`^{\mathrm{PSS}}(q,\overline{x})`$. A pair $`(\gamma ,u)`$ forms a plumberโs helper.
###### Remark 2.4.
$`^{\mathrm{PSS}}(q,\overline{x})`$ is composed of gradient flow half-lines $`\gamma `$ for $`f`$ and Floer cylinders $`u`$. Because of the cut-off function $`\beta `$ the cylinder $`u`$ actually is holomorphic for $`s\frac{1}{2}`$. Since we impose the finite energy condition $`E(u)<\mathrm{}`$, the punctured holomorphic disk $`u|_{(\mathrm{},1/2)\times S^1}`$ has a removable singularity at the origin (cf. ), i.e. $`u`$ has a continuous extension $`u(\mathrm{})`$. Therefore, the second boundary condition and the homotopy condition is meaningful.
For generic choices of $`H,J,f`$ and $`g`$ the moduli space $`^{\mathrm{PSS}}(q,\overline{x})`$ is a smooth manifold and
$$dim^{\mathrm{PSS}}(q,\overline{x})=\mu _{\mathrm{CZ}}(\overline{x})+\mu _{\mathrm{Morse}}(q)m,$$
(2.27)
according to our normalization that for an autonomous $`C^2`$-small Morse function $`f`$ we have $`\mu _{\mathrm{CZ}}(d_x)=m\mu _{\mathrm{Morse}}(x)`$ for $`\overline{x}=[x,d_xx]\stackrel{~}{๐ซ}(f)`$, i.e. $`x\mathrm{Crit}(f)`$. As always $`dimM=2m`$.
A standard computation for an element $`(\gamma ,u)^{\mathrm{PSS}}(q,\overline{x})`$ shows
$$0E(u)๐_H(\overline{x})+_0^1\underset{M}{sup}H(t,)dt.$$
(2.28)
This implies a universal energy bound on plumberโs helper solutions. In particular, sequences converge up to breaking and bubbling.
Assumption (SP) ensures compactness of the 0-dimensional components of $`^{\mathrm{PSS}}(q,\overline{x})`$. We denote by $`\mathrm{CM}^m(f)`$ the Morse co-chain complex associated to $`f`$ (and $`g`$), and define
$`\mathrm{PSS}:\mathrm{CF}_{}(H)`$ $`\mathrm{CM}^m(f)`$ (2.29)
$`\overline{x}`$ $`{\displaystyle \underset{\mu _{\mathrm{Morse}}(q)=m\mu _{\mathrm{CZ}}(\overline{x})}{}}\mathrm{\#}_2^{\mathrm{PSS}}(q,\overline{x})q.`$
###### Theorem 2.5 ().
The 1-dimensional components of the moduli space $`^{\mathrm{PSS}}(q,\overline{x})`$ can be compactified by adding either triples $`(\gamma ,u,u^{})`$, where $`(u,u^{})`$ are broken Floer cylinders, or $`(\gamma ^{},\gamma ,u)`$, where $`(\gamma ^{},\gamma )`$ are broken Morse trajectories.
This implies that $`\mathrm{PSS}`$ commutes with the Morse respectively the Floer (co)-boundary operator, i.e. $`\mathrm{PSS}`$ descends to homology
$$\mathrm{PSS}:\mathrm{HF}_{}(H)\mathrm{HM}^m(f).$$
(2.30)
For proving that $`\mathrm{PSS}`$ is an isomorphism, an explicit inverse is constructed in the following. The adequate moduli space $`^{\mathrm{PSS},\mathrm{inv}}(\overline{x},q)`$ is composed of pairs $`(u,\gamma )`$ such that
$$u:\times S^1M\text{ with}E(u)<\mathrm{}\text{and}\gamma :[0,\mathrm{})M$$
(2.31)
solve the equations
$$_su+J(u)\left(_tu\beta (s)X_H(t,u)\right)=0\text{and}\dot{\gamma }+^gf\gamma =0,$$
(2.32)
and fulfill boundary conditions $`u(\mathrm{})=x`$, $`u(\mathrm{})=\gamma (0)`$ and $`\gamma (+\mathrm{})=q`$, where $`q\mathrm{Crit}(f)`$ is a critical point of $`f`$ and $`\overline{x}\stackrel{~}{๐ซ}(H)`$. Finally, we impose the homotopy condition $`\omega (d_x\mathrm{\#}u)=0=c_1(d_x\mathrm{\#}u)`$. As before, for generic choices we obtain a smooth manifold $`^{\mathrm{PSS},\mathrm{inv}}(\overline{x},q)`$ of $`dim^{\mathrm{PSS},\mathrm{inv}}(\overline{x},q)=m\mu _{\mathrm{Morse}}(q)\mu _{\mathrm{CZ}}(\overline{x})`$, which is compact in dimension 0. We define
$`\mathrm{PSS}^1:\mathrm{CM}^{}(f)`$ $`\mathrm{CF}_m(H)`$ (2.33)
$`q`$ $`{\displaystyle \underset{\mu _{\mathrm{CZ}}(\overline{x})=m\mu _{\mathrm{Morse}}(q)}{}}\mathrm{\#}_2^{\mathrm{PSS},\mathrm{inv}}(\overline{x},q)\overline{x}.`$
As above a suitable compactification of the 1-dimensional components of $`^{\mathrm{PSS},\mathrm{inv}}(\overline{x},q)`$ entails
$$\mathrm{PSS}^1:\mathrm{HM}^{}(f)\mathrm{HF}_m(H).$$
(2.34)
Up to here the nomenclature $`\mathrm{PSS}^1`$ is purely formal. In what follows a justification is sketched. The following arguments are taken from .
Composing $`\mathrm{PSS}^1\mathrm{PSS}:\mathrm{HF}_{}(H)\mathrm{HF}_{}(H)`$ amounts to
$$\mathrm{PSS}^1\mathrm{PSS}(\overline{x})=\underset{\overline{y}}{}\underset{q}{}\mathrm{\#}_2^{\mathrm{PSS},\mathrm{inv}}(\overline{y},q)\mathrm{\#}_2^{\mathrm{PSS}}(q,\overline{x})\overline{y},$$
(2.35)
i.e. counting Floer cylinders emanating from $`\overline{x}\stackrel{~}{๐ซ}(H)`$ connecting to gradient flow half-trajectories ending at some critical point $`q\mathrm{Crit}(f)`$. From $`q`$ further gradient flow half-trajectories finally connect to Floer cylinders ending at $`\overline{y}\stackrel{~}{๐ซ}(H)`$. The idea is to form a cobordism between these configurations and the identity as follows:
1. We glue the two gradient flow half-trajectories at $`q`$ to obtain a finite length gradient flow trajectory.
2. Shrink that finite length to zero, i.e. we end up with two Floer cylinders, which meet at the same point.
3. Due to the cut-off functions on the respective ends, these two Floer cylinders are holomorphic near that point. We employ a glueing theorem for holomorphic curves to obtain one Floer cylinder passing from $`\overline{x}`$ to $`\overline{y}`$.
4. The Hamiltonian function for this Floer cylinder is not yet $`H`$ due to the cut-off. Therefore, we choose a (compactly supported) homotopy changing this Hamiltonian function to $`H`$.
5. We end up with counting Floer cylinders for the Hamiltonian function $`H`$ connecting periodic solutions $`\overline{x}`$ and $`\overline{y}`$ of the same Conley-Zehnder index. For dimension reasons generically there are no such solutions unless $`\overline{x}=\overline{y}`$, in which case there is only the constant solution. This defines the identity homomorphism on Floer homology.
For the other composition $`\mathrm{PSS}\mathrm{PSS}^1:\mathrm{HM}_{}(f)\mathrm{HM}_{}(f)`$ we obtain
$$\mathrm{PSS}\mathrm{PSS}^1(q)=\underset{p}{}\underset{\overline{x}}{}\mathrm{\#}_2^{\mathrm{PSS}}(p,\overline{x})\mathrm{\#}_2^{\mathrm{PSS},\mathrm{inv}}(\overline{x},q)p.$$
(2.36)
Here two Floer cylinders meet at a periodic solution $`\overline{x}`$. This leads to the following cobordism:
1. We glue the two Floer cylinders at $`\overline{x}`$ to obtain a sphere obeying Floerโs equation.
2. Since the Hamiltonian term vanishes outside a neighborhood of the equator of this sphere we choose a homotopy of the Hamiltonian term to zero and end up with a holomorphic sphere.
3. Assuming for the moment that $`\omega |_{\pi _2(M)}=0`$, this sphere is constant. We reduced the problem to count gradient trajectories for index difference 0. As above this defines the identity homomorphism.
In the paper and in the book \[11, chapter 12.1\] are series of figures picturise the above ideas. We conclude that our notation is meaningful, i.e. $`\mathrm{PSS}`$ is an isomorphism with inverse
$$(\mathrm{PSS})^1=\mathrm{PSS}^1.$$
(2.37)
in the case $`\omega |_{\pi _2(M)}=0`$. In the semi-positive case (SP) step (3) becomes much more delicate since, in general, there will exist holomorphic spheres. Then an elaborate transversality argument is necessary to allow for a time independent almost complex structure along the sphere. Since only two points on the sphere are fixed by the gradient trajectories, an $`S^1`$-symmetry remains. In particular, solutions come in 1-dimensional families as long as the sphere is non-constant. This contradicts the fact that the moduli space that we started with has dimension 0. We end up with the same conclusion as in (3).
###### Definition 2.6.
An element $`\overline{x}\mathrm{CF}_{}(H)`$ is called essential if there are $`p,q\mathrm{Crit}(f)`$ s.t.
$$^{\mathrm{PSS}}(p,\overline{x})\mathrm{}\text{and}^{\mathrm{PSS},\mathrm{inv}}(\overline{x},q)\mathrm{}.$$
(2.38)
###### Lemma 2.7.
For essential elements $`\overline{x}\mathrm{CF}_{}(H)`$ we can estimate (cf. inequality (2.28))
$$_0^1\underset{M}{sup}Hdt๐_H(\overline{x})_0^1\underset{M}{inf}Hdt.$$
(2.39)
###### Proposition 2.8.
All periodic orbits representing non-zero homology classes are essential.
###### Proof.
Pick $`\overline{x}\mathrm{CF}_{}(H)`$ such that $`\overline{x}=0`$ and $`[\overline{x}]0`$. If $`^{\mathrm{PSS}}(p,\overline{x})=\mathrm{}`$ for all $`p\mathrm{Crit}(f)`$, then $`\mathrm{PSS}([\overline{x}])=0`$.
The existence of $`q`$ such that $`^{\mathrm{PSS},\mathrm{inv}}(\overline{x},q)\mathrm{}`$ follows via Poincarรฉ duality, cf.. โ
The action filtration and the definition of spectral capacities. The action filtration and the action selector are used only to state and prove theorem 5.5. Therefore, we are very brief about the definitions and properties. Detailed expositions can be found in .
We define the action filtration on Floer homology by
$$\mathrm{CF}_k^a(H):=\{\underset{\mu _{\mathrm{CZ}}(\overline{x})=k}{}a_{\overline{x}}\overline{x}a_{\overline{x}}=0\text{ for }๐_H(\overline{x})>a\},$$
(2.40)
and check that $`(\mathrm{CF}_k^a(H),)`$ forms a subcomplex of the Floer complex (cf. ) due to the chosen Novikov condition. In particular, we obtain the long exact sequence
$$\mathrm{}\mathrm{HF}_{}^a(H)\stackrel{i_H^a}{}\mathrm{HF}_{}(H)\stackrel{j_H^a}{}\mathrm{HF}_{}^{(a,\mathrm{}]}(H)\mathrm{}.$$
(2.41)
This gives rise to the action selector
$$c(\alpha ,H):=inf\{a\mathrm{PSS}^1(\alpha )\mathrm{im}i_H^a\}$$
(2.42)
for all $`0\alpha \mathrm{QH}^{}(M)`$. We note here, that the quantum cohomology of the symplectic manifold comes into play and refer the reader to \[11, chapter 11.1\].
The action selector is well-defined, i.e. $`c(\alpha ,H)`$ is a finite number, since each representation of the class $`\mathrm{PSS}^1(\alpha )`$ contains an element $`\overline{x}\stackrel{~}{๐ซ}(H)`$ with maximal action value, due to the Novikov condition (but there might not be such an element with minimal action value).
Properties of the action selector. A full list of properties of the action selector $`c(\alpha ,H)`$ can be found in \[11, chapter 12.4\]. We shall need the following: For $`\alpha ,\beta \mathrm{QH}^{}(M)`$, $`A\mathrm{\Gamma }`$ and Hamiltonian functions $`H,K:S^1\times M`$ we have
(Zero) $`c(\alpha ,0)=0`$ (Novikov action) $`c(\alpha e^A,H)=c(\alpha ,H)\omega (A)`$ (Product) $`c(\alpha \beta ,H\mathrm{\#}K)c(\alpha ,H)+c(\beta ,K)`$ (Poincarรฉ duality) $`c(1,H)c([\omega ^m],H^{(1)})`$ (Non-degeneracy) for each $`H0`$ there exists a $`\delta >0`$ such that $`c(1,H)+c(1,H^{(1)})c(1,H)c([\omega ^m],H)\delta .`$
Here, $`H\mathrm{\#}K(t,m)=H(t,m)+K(t,(\psi _H^t)^1(m))`$ and $`\frac{d}{dt}\psi _H^t=X_H(t,\psi _H^t)`$; furthermore, we set $`H^{(1)}(t,m)=H(t,m)`$ and $``$ denotes the quantum-cup-product in $`\mathrm{QH}^{}(M)`$.
With the help of the action selector we define two norms (cf. )
$$\gamma (H):=c(1,H)c([\omega ^m],H)\text{and}\stackrel{~}{\gamma }(H):=c(1,H)+c(1,H^{(1)}).$$
(2.43)
Non-degeneracy implies $`\stackrel{~}{\gamma }(H)\gamma (H)`$. These norms give rise to two symplectic capacities as follows. For an open subset $`UM`$ we set
$$c_\gamma (U)=sup\{\gamma (H)\mathrm{supp}X_HS^1\times U\}\text{and}\stackrel{~}{c}_{\stackrel{~}{\gamma }}(U)=sup\{\stackrel{~}{\gamma }(H)\mathrm{supp}X_HS^1\times U\}.$$
Again non-degeneracy implies $`\stackrel{~}{c}_{\stackrel{~}{\gamma }}(U)c_\gamma (U)`$. We call the capacity $`c_\gamma (U)`$ the spectral capacities of $`U`$. Most interesting to us is the case $`U=M`$. For a Lagrangian submanifold $`LM`$ we define $`c_\gamma (L)`$ as $`c_\gamma (L):=inf\{c_\gamma (U)U\text{ open},LU\}`$ and $`\stackrel{~}{c}_{\stackrel{~}{\gamma }}(L)`$ accordingly.
## 3. Basic theorems
The basic theorems which all other results rely on are the following theorems 3.1 and 3.3.
###### Theorem 3.1.
Let $`(M^{2m},\omega )`$ be a monotone closed symplectic manifold and $`LM`$ a monotone closed Lagrangian submanifold with $`N_L2`$ (cf. definition 2.2). For each closed submanifold $`SL`$ such that $`s:=dimS`$ satisfies
$$s>dimL+1N_L,\text{i.e.}N_L>\mathrm{codim}_LS+1,$$
(3.1)
the image of $`[S]\mathrm{H}_s(M;/2)`$ under the $`\mathrm{PSS}`$-isomorphism is represented by the cycle
$$\mathrm{PSS}^1\left([S]\right)=\left[\underset{\overline{x}\stackrel{~}{๐ซ}_{s2m}(H)}{}\mathrm{\#}_2\left(_{(L,S)}^+(\overline{x};J,H;\mathrm{๐})\right)\overline{x}\right]\mathrm{HF}^{2ms}(H),$$
(3.2)
in the Floer cohomology of the (generically chosen) Hamiltonian $`H:S^1\times M`$, and dually
$$\mathrm{PSS}^1\left(\mathrm{PD}[S]\right)=\left[\underset{\overline{x}\stackrel{~}{๐ซ}_{s2m}(H)}{}\mathrm{\#}_2\left(_{(L,S)}^{}(\overline{x};J,H;\mathrm{๐})\right)\overline{x}\right]\mathrm{HF}_{s2m}(H).$$
(3.3)
The moduli spaces $`_{(L,S)}^\pm (\overline{x};J,H;\mathrm{๐})`$ are defined with the help of a generically chosen almost complex structures $`J`$ as
$$_{(L,S)}^+(\overline{x};J,H;\mathrm{๐}):=\{u:[0,\mathrm{})\times S^1(M,L)|\begin{array}{cccc}& \overline{}_{J,H}u=0,\hfill & & [u\mathrm{\#}\overline{d}_x]=\mathrm{๐}\hfill \\ & u(+\mathrm{})=x,\hfill & & u(0,0)S\hfill \end{array}\}$$
(3.4)
respectively
$$_{(L,S)}^{}(\overline{x};J,H;\mathrm{๐}):=\{u:(\mathrm{},0]\times S^1(M,L)|\begin{array}{cccc}& \overline{}_{J,H}u=0,\hfill & & [d_x\mathrm{\#}u]=\mathrm{๐}\hfill \\ & u(\mathrm{})=x,\hfill & & u(0,0)S\hfill \end{array}\},$$
(3.5)
where $`\overline{}_{J,H}`$ abbreviates Floerโs equation, namely $`_su+J(u)\left(_tuX_H(t,u)\right)=0`$, and $`u:[0,\mathrm{})\times S^1(M,L)`$ means $`u(\{0\}\times S^1)L`$. Recall that $`\overline{x}=[x,d_x]`$. The homotopical condition $`[d_x\mathrm{\#}u]=\mathrm{๐}`$ in the definition of $`_{(L,S)}^{}(\overline{x};J,H;\mathrm{๐})`$ is meant in the Novikov sense, i.e. $`\omega ([d_x\mathrm{\#}u])=0=\mu _{\mathrm{Maslov}}([d_x\mathrm{\#}u])`$ and accordingly for $`_{(L,S)}^+(\overline{x};J,H;\mathrm{๐})`$.
We count Floer half-cylinders $`u`$ which are asymptotic to a periodic orbit and have boundary on the Lagrangian submanifold, i.e. $`u(0,t)L`$ for all $`tS^1`$. An important feature of the moduli spaces is the *marking* of the half-cylinders: we require them to map the marked point $`(0,0)`$ on the half-cylinder to $`S`$. The closed submanifold $`S`$ can be replaced by any singular chain representing a cycle on $`L`$.
###### Remark 3.2.
If $`S=L`$, we have $`dimL>dimL+1N_L`$, since $`N_L2`$. In particular, theorem 3.1 implies that we always can represent the class $`[L]\mathrm{H}_m(M;/2)`$.
If $`S=\mathrm{pt}`$ and $`N_L>dimL+1`$, we obtain a Floer-theoretic representation of the class $`[\mathrm{pt}]\mathrm{H}_0(M;/2)`$ and its Poincarรฉ dual $`[\omega ^m]\mathrm{H}^{2m}(M;/2)`$. The condition $`N_L>dimL+1`$ excludes all *displaceable* (cf. definition 5.1) Lagrangian submanifolds. This is no coincidence in view of theorem 5.2.
###### Theorem 3.3.
If we assume in addition to the assumptions of theorem 3.1 that the Hamiltonian function $`H:S^1\times M`$ satisfies
$$H\mathrm{min}\{A_M,A_L\}$$
(3.6)
(see definition 2.2 for $`A_M,A_L`$), then the representations (3.2) and (3.3) hold for all closed submanifolds $`SL`$ regardless of their dimension.
The basic theorems 3.1 and 3.3 deal with the extrinsic topology of the Lagrangian submanifold: They provide a Floer theoretic representation of the homomorphism $`\iota :\mathrm{H}_{}(L;/2)\mathrm{H}_{}(M;/2)`$ induced by the inclusion $`\iota :LM`$.
## 4. Proof of the basic theorems
The proofs of theorems 3.1 and 3.3 rely both on the compactness theorem 4.3 stated below for the moduli spaces $`_{(L,S)}^\pm (\overline{x};J,H;\mathrm{๐})`$ in dimensions 0 and 1. The respective assumption on the dimension of $`S`$ or the Hofer norm of $`H`$ prohibits bubbling-off, which in turn results in the compactness of the moduli spaces in question. Before proving this compactness result we make sure that the moduli spaces are smooth manifolds for generic choice of $`H`$ and $`J`$.
###### Theorem 4.1.
Let $`L`$ be a closed Lagrangian submanifold of the closed symplectic manifold $`(M,\omega )`$ and $`SL`$ a closed submanifold. For a generic Hamiltonian function $`H:S^1\times M`$ and a generic almost complex structure $`J`$ the moduli spaces
$$_{(L,S)}^{}(\overline{x};J,H;\mathrm{๐})\text{and}_{(L,S)}^+(\overline{x};J,H;\mathrm{๐})$$
(4.1)
defined in (3.4) and (3.5) are smooth manifolds of dimension
$$dim_{(L,S)}^\pm (\overline{x};J,H;\mathrm{๐})=\pm \mu _{\mathrm{CZ}}(\overline{x})dimL+dimS,$$
(4.2)
where $`\mu _{\mathrm{CZ}}(\overline{x})`$ denotes the Conley-Zehnder index of the periodic orbit $`\overline{x}\stackrel{~}{๐ซ}(H)`$.
###### Proof.
This is a standard result in Floer theory (cf. ) with one minor modification, cf. remark 4.2 (i). First of all, for a non-degenerate $`H`$ we can regard $`\overline{}_{J,H}=_s+J(_tX_H)`$ as a Fredholm-section in a suitable Banach-bundle and identify the moduli spaces with the vanishing locus of this Fredholm-section. For generic $`H`$ and $`J`$, this section will be transverse to the zero-section. In particular, by the implicit function theorem, the moduli spaces are smooth finite-dimensional manifolds. Computing the dimension of the moduli spaces is achieved by the Riemann-Roch theorem and additivity of the Fredholm index. All details within the setting of this paper can be found in . โ
###### Remark 4.2.
1. To prove the transversality result we need to assume, in addition to $`H`$ being non-degenerate, that there are no periodic orbits of $`H`$ lying entirely on $`L`$. This excludes $`s`$-independent solutions of the Floer equation, for which we were not able to prove the necessary transversality statement. In contrast to the construction of Floer homology there is no (immediate) automatic transversality result for half-cylinders. Anyway, this additional assumption on $`H`$ is clearly fulfilled generically.
2. The moduli space $`_{(L,S)}^\pm (\overline{x};J,H;\mathrm{๐})`$ is a submanifold of $`_{(L,L)}^\pm (\overline{x};J,H;\mathrm{๐})`$ for generic almost complex structures, because the evaluation map
$$\mathrm{ev}:_{(L,L)}^\pm (\overline{x};J,H;\mathrm{๐})L\mathrm{ev}(u):=u(0,0)$$
(4.3)
is a submersion if regarded on a universal moduli space, cf. \[11, chapter 3.4\], and
$$_{(L,S)}^\pm (\overline{x};J,H;\mathrm{๐})=\mathrm{ev}^1(S).$$
(4.4)
To complete the proofs of theorems 3.1 and 3.3 we need to show that
1. the 0-dimensional moduli spaces are compact,
2. the chains $`_{\overline{x}}\mathrm{\#}_2_{(L,S)}^\pm (\overline{x};J,H;\mathrm{๐})\overline{x}`$ are cycles in the Floer complex,
3. they represent the claimed (co)-homology classes.
The first two points are the content of the following compactness results, theorem 4.3. Since the assertions of theorem 3.3 are homological it suffices in the second case of theorem 4.3 to consider only essential elements since all elements representing non-vanishing homology classes are essential, cf. proposition 2.8.
###### Theorem 4.3.
Let $`(M^{2m},\omega )`$ be a closed monotone symplectic manifold and $`LM`$ a monotone closed Lagrangian submanifold with $`N_L2`$. If we choose a closed submanifold $`SL`$ and a Hamiltonian function $`H:S^1\times M`$ such that one of the following holds:
1. $`dimS>dimL+1N_L`$ ,
2. $`H\mathrm{min}\{A_M,A_L\}`$ and $`\overline{x}\stackrel{~}{๐ซ}(H)`$ is essential (cf. definition 2.6),
then the 0-dimensional moduli spaces $`_{(L,S)}^\pm (\overline{x};J,H;\mathrm{๐})`$ are compact.
Furthermore, the 1-dimensional moduli spaces are compact up to splitting-off one Floer cylinder, i.e. they can be compactified in such a way that the boundary of the compactification, which we denote by the same symbol, decomposes as follows (cf. formula (2.15))
$`_{(L,S)}^{}(\overline{x};J,H;\mathrm{๐})`$ $`={\displaystyle \underset{\mu _{\mathrm{CZ}}(\overline{y})=\mu _{\mathrm{CZ}}(\overline{x})+1}{}}\widehat{}(\overline{x},\overline{y};J,H)\times _{(L,S)}^{}(\overline{y};J,H;\mathrm{๐}),`$ (4.5)
$`_{(L,S)}^+(\overline{x};J,H;\mathrm{๐})`$ $`={\displaystyle \underset{\mu _{\mathrm{CZ}}(\overline{y})=\mu _{\mathrm{CZ}}(\overline{x})1}{}}_{(L,S)}^+(\overline{y};J,H;\mathrm{๐})\times \widehat{}(\overline{y},\overline{x};J,H).`$
A typical element in $`_{(L,S)}^{}(\overline{x};J,H;\mathrm{๐})`$ is depicted in figure 2.
###### Proof.
The following inequality for elements $`u_{(L,S)}^{}(\overline{x};J,H;\mathrm{๐})`$
$$E(u)=_{\mathrm{}}^0_0^1|_su|^2๐t๐s\omega ([d_x\mathrm{\#}u])๐_H(\overline{x})_0^1\underset{L}{inf}H(t,)dt.$$
(4.6)
is easily derived using Floerโs equation. Let us abbreviate from now on
$$\underset{L}{inf}H:=_0^1\underset{L}{inf}H(t,)dt\text{and}\underset{L}{sup}H:=_0^1\underset{L}{sup}H(t,)dt.$$
(4.7)
By the definition of $`_{(L,S)}^{}(\overline{x};J,H;\mathrm{๐})`$ we have $`\omega ([d_x\mathrm{\#}u])=0`$, therefore (4.6) reads
$$0E(u)๐_H(\overline{x})\underset{L}{inf}H$$
(4.8)
and analogously for elements $`u_{(L,S)}^+(\overline{x};J,H;\mathrm{๐})`$
$$0E(u)+๐_H(\overline{x})+\underset{L}{sup}H.$$
(4.9)
Hence, the energy of elements of the moduli spaces is uniformly bounded.
The case $`S=L`$.
Since the energy is uniformly bounded and the marking $`u(0,0)S=L`$ for an element in the moduli space is obsolete we conclude by the standard compactness results in Floer theory that sequences $`(u_n)_{(L,L)}^{}(\overline{x};J,H;\mathrm{๐})`$ converge in the Gromov-Hausdorff topology
$$u_n(u_{\mathrm{}};v_1,\mathrm{},v_\mathrm{\Gamma };s_1,\mathrm{},s_\mathrm{\Sigma };d_1,\mathrm{},d_\mathrm{\Delta })$$
(4.10)
where
* $`u_{\mathrm{}}_{(L,L)}^\pm (\overline{x}_0;J,H;\mathrm{๐})`$,
* $`v_\gamma \widehat{}(\overline{x}_\gamma ,\overline{x}_{\gamma 1};J,H)`$, where $`\overline{x}_\gamma \stackrel{~}{๐ซ}(H)`$ and $`\overline{x}_\mathrm{\Gamma }=\overline{x}`$,
* $`\{s_\sigma \}`$ are holomorphic spheres,
* $`\{d_\delta \}`$ are holomorphic disks,
i.e. sequences converge to a finite family of adjacent Floer cylinders starting at $`\overline{x}`$ and finally connecting to a half-cylinder in $`_{(L,L)}^{}(\overline{x}_0;J,H;\mathrm{๐})`$. Furthermore, there are finitely many holomorphic spheres and disks attached to this configuration.
The Fredholm index behaves additively with respect to this convergence, i.e. the Fredholm index at an element $`u_n`$, which equals the dimension of $`_{(L,L)}^{}(\overline{x};J,H;\mathrm{๐})`$, is
$$\mathrm{ind}_{u_n}=\mathrm{ind}_u_{\mathrm{}}+\underset{\gamma }{}\mathrm{ind}_{v_\gamma }+\underset{\sigma }{}2c_1(s_\sigma )+\underset{\delta }{}\mu _{\mathrm{Maslov}}(d_\delta ).$$
(4.11)
Since the Lagrangian submanifold and therefore the symplectic manifold is monotone, the Chern class and the Maslov index evaluate positively in case that $`s_\sigma `$ and $`d_\delta `$ are non-constant. Moreover, by assumption the Maslov index is at least 2, and if the Floer cylinders $`v_\gamma `$ are $`s`$-dependent, the Fredholm indices are at least 1. In particular, if the moduli space $`_{(L,L)}^{}(\overline{x};J,H;\mathrm{๐})`$ has dimension 0 all maps but $`u_{\mathrm{}}`$ must be trivial. In other words, the moduli space $`_{(L,L)}^{}(\overline{x};J,H;\mathrm{๐})`$ is compact.
If the moduli space $`_{(L,L)}^{}(\overline{x};J,H;\mathrm{๐})`$ is 1-dimensional, then a sequence converges to the limit half-cylinder $`u_{\mathrm{}}`$ and at most one further Floer cylinder $`v_\gamma `$. In particular, at most one Floer cylinder may split-off. This gives rise to the compactification with the decomposition of the boundary as asserted in the theorem.
Since twice the first Chern class evaluated on a sphere is a multiple of the minimal Maslov number $`N_L`$ (cf. remark 2.3) we obtain the following refined statement. If the dimension of the moduli space is less than the minimal Maslov number
$$dim_{(L,L)}^{}(\overline{x};J,H;\mathrm{๐})=\mu _{\mathrm{CZ}}(\overline{x})<N_L$$
(4.12)
the moduli space $`_{(L,L)}^{}(\overline{x};J,H;\mathrm{๐})`$ is compact up to splitting-off a finite number of Floer cylinders. No holomorphic spheres or disks bubble off. A completely analogous statement holds for the moduli space $`_{(L,L)}^+(\overline{x};J,H;\mathrm{๐})`$.
The case $`dimS>dimL+1N_L`$.
We already noted above that the marked moduli space $`_{(L,S)}^\pm (\overline{x};J,H;\mathrm{๐})`$ is a submanifold of the moduli space $`_{(L,L)}^\pm (\overline{x};J,H;\mathrm{๐})`$ for generic almost complex structures.
If $`_{(L,S)}^\pm (\overline{x};J,H;\mathrm{๐})`$ is 0 or 1-dimensional and in addition we require the condition $`dimS>dimL+1N_L`$ to hold, we conclude that $`dim_{(L,L)}^\pm (\overline{x};J,H;\mathrm{๐})<N_L`$. As explained at the end of the first case we obtain the statement of theorem 4.3 in this case.
The case $`H\mathrm{min}\{A_M,A_L\}`$ and $`\overline{x}\stackrel{~}{๐ซ}(H)`$ is essential.
Without the dimension restriction on the submanifold $`S`$ we cannot expect compactness of all 0 and 1-dimensional moduli spaces $`_{(L,S)}^\pm (\overline{x};J,H;\mathrm{๐})`$. In fact, some of our applications are derived with help of this non-compactness. But as long as all elements in $`_{(L,S)}^\pm (\overline{x};J,H;\mathrm{๐})`$ have energy less than the minimal energy of a non-constant holomorphic sphere or disk, a bubble would take away more energy than is available.
We assume that the Hamiltonian function $`H`$ has sufficiently small Hofer norm, namely $`H\mathrm{min}\{A_M,A_L\}`$, and the periodic orbit $`\overline{x}`$ is essential, cf. definition 2.6. For such a periodic orbit $`\overline{x}`$ we recall from lemma 2.7 the inequalities $`sup_MH๐_H(\overline{x})inf_MH`$. Combining this with inequality (4.8) resp. (4.9) we obtain
$$E(u)๐_H(\overline{x})\underset{L}{inf}H\underset{M}{sup}H\underset{M}{inf}H=H$$
(4.13)
for an element $`u_{(L,S)}^{}(\overline{x};J,H;\mathrm{๐})`$ and analogously
$$E(u)+๐_H(\overline{x})+\underset{L}{sup}H\underset{M}{inf}L+\underset{M}{sup}H=H$$
(4.14)
for an element $`u_{(L,S)}^+(\overline{x};J,H;\mathrm{๐})`$.
In particular, bubbling-off cannot occur for a sequence $`(u_n)_{(L,S)}^\pm (\overline{x};J,H;\mathrm{๐})`$. Indeed, the energy $`E(u_{\mathrm{}})`$ of the limit solution $`u_{\mathrm{}}`$ is positive by the assumption on $`H`$ that no periodic orbits lie entirely on $`L`$. Therefore,
$$E(b)<E(u_{\mathrm{}})+E(b)lim\; infE(u_n)H\mathrm{min}\{A_M,A_L\}.$$
(4.15)
We obtain $`E(b)=0`$. Since we excluded bubbling-off of holomorphic spheres and disks the relevant moduli spaces are compact or can be compactfied as asserted. This finishes the proof of theorem 4.3. โ
###### Remark 4.4.
For later purposes we remark that in the proof of theorem 4.3 the assumption $`H\mathrm{min}\{A_M,A_L\}`$ is used only in combination with inequality (4.13) to derive the crucial inequality $`E(u)\mathrm{min}\{A_M,A_L\}`$. Therefore, if we assume this inequality right away, the assertion of theorem 4.3 still holds.
###### Remark 4.5.
A geometric explanation for the hypothesis $`dimS>dimL+1N_L`$ is that a holomorphic disk bubbling-off may take away the marking. If we disregard the marking in a configuration as depicted in figure 3, it lies in the boundary of the higher dimensional moduli space $`_{(L,L)}^\pm (\overline{x};J,H;\mathrm{๐})`$. The assumption $`dimS>dimL+1N_L`$ excludes such configurations in $`_{(L,L)}^\pm (\overline{x};J,H;\mathrm{๐})`$ for index reasons so that no bubble can take away the marking.
Theorem 4.3 immediately implies that the chain defined in theorem 3.1 resp. 3.3 is well-defined and in fact a cycle. To finish the proof of theorems 3.1 and 3.3 we need to show that the $`\mathrm{PSS}`$-isomorphism maps $`[S]\mathrm{H}_s(M;/2)`$ of the homology class of this cycle.
###### End of the proof of theorems 3.1 and 3.3.
We will define a cobordism relating the counting procedures defining $`\mathrm{PSS}^1([S])`$ and the explicitly given cycle. This shows that they are chain homotopic and therefore agree in homology.
We define the moduli space $`_{(L,S)}^{\mathrm{PSS},}(q;g,f;J,H;\mathrm{๐})`$ as the set of triples $`(R,\gamma ,u)`$, where
$$R0,\gamma :(\mathrm{},0]M,u:(\mathrm{},R]\times S^1M$$
(4.16)
subject to the equations
$$\dot{\gamma }+^gf\gamma =0\text{and}_su+J(u)\left(_tu\beta (s)X_H(t,u)\right)=0,$$
(4.17)
where $`\beta :[0,1]`$ is a smooth cut-off function satisfying $`\beta (s)=0`$ for $`s\frac{1}{2}`$ and $`\beta (s)=1`$ for $`s1`$. The function $`f:M`$ is an auxiliary Morse function and $`^g`$ is the gradient with respect to an auxiliary Riemannian metric $`g`$ on M. Furthermore, $`(R,\gamma ,u)_{(L,S)}^{\mathrm{PSS},}(q;g,f;J,H;\mathrm{๐})`$ has to satisfy the following boundary conditions
$$E(u)<\mathrm{},\gamma (\mathrm{})=q,\gamma (0)=u(\mathrm{})\text{and}u(R,t)LtS^1,u(R,0)S,$$
(4.18)
where $`q\mathrm{Crit}(f)`$ is a critical point of $`f`$. As in the definition of the $`\mathrm{PSS}`$-isomorphism the finiteness of the energy of $`u`$ in conjunction with the cut-off function $`\beta `$ allows for a continuous extension $`u(\mathrm{})`$. In particular, topologically we can think of $`u`$ as a disk. Finally, we impose the homotopy condition $`[u]=0`$ for the disk $`u`$, i.e. $`\omega (u)=0=\mu _{\mathrm{Maslov}}(u)`$.
The geometric ideas behind the definition of $`_{(L,S)}^{\mathrm{PSS},}(q;g,f;J,H;\mathrm{๐})`$ are:
1. The cycle $`\mathrm{PSS}\left(_{\overline{x}\stackrel{~}{๐ซ}(H)}\mathrm{\#}_2\left(_{(L,S)}^{}(\overline{x};J,H;\mathrm{๐})\right)\overline{x}\right)`$ is defined by counting gradient half trajectories attached to Floer cylinders, i.e. plumberโs helper solutions, ending at a periodic orbit $`\overline{x}`$, from where Floer half-cylinders start which end on $`L`$.
2. We glue such half-cylinders and Floer cylinders and obtain new half-cylinders which obey Floerโs equation with respect to a Hamiltonian function different from zero only near $`L`$.
3. We regard this Hamiltonian function as a compactly supported deformation of the zero Hamiltonian. Therefore, after a homotopy, we deal with holomorphic cylinders which in fact are holomorphic disks since they have finite energy.
4. The assumption on the homotopy type of this disk immediately implies that it is constant. Thus, we are left with counting gradient half-trajectories starting at a critical point and ending at $`S`$, due to the marking. It is easy to see that this represents the class $`\mathrm{PD}[S]\mathrm{HM}^s(M)`$ in the Morse cohomology of $`M`$.
Let us come back to our object of study: $`_{(L,S)}^{\mathrm{PSS},}(q;g,f;J,H;\mathrm{๐})`$. Combining the methods from the definition of the $`\mathrm{PSS}`$-isomorphism and from this chapter, it follows that these spaces are smooth manifolds for generic choices of $`g`$, $`J`$ and $`H`$ and of dimension
$`dim_{(L,S)}^{\mathrm{PSS},}(q;g,f;J,H;\mathrm{๐})`$ $`=dimSdimL\left(\frac{1}{2}dimM\mu _{\mathrm{Morse}}(q;f)\right)+1`$ (4.19)
$`=dimSdimM+\mu _{\mathrm{Morse}}(q;f)+1.`$
The $`+1`$ accounts for the variable $`R`$. A straight forward computation implies for an element $`(R,\gamma ,u)_{(L,S)}^{\mathrm{PSS},}(q;g,f;J,H;\mathrm{๐})`$ that
$$E(u)\beta (R)\left(\underset{M}{sup}H\underset{L}{inf}H\right)H.$$
(4.20)
We note, that for $`R=0`$ this inequality implies $`E(u)=0`$.
Exactly the same arguments as in theorem 4.3 imply that the moduli spaces are compact in dimension 0 and can be compactified in dimension 1. We denote again the compactification by the same letter. The boundary of the compactification decomposes as
$`_{(L,S)}^{\mathrm{PSS},}(q;g,f;J,H;\mathrm{๐})=`$ $`{\displaystyle \underset{\begin{array}{c}q^{}\mathrm{Crit}(f)\\ \mu _{\mathrm{Morse}}(q^{})=\mu _{\mathrm{Morse}}(q)1\end{array}}{}}\widehat{}^{\mathrm{Morse}}(q,q^{};f)\times _{(L,S)}^{\mathrm{PSS},}(q^{};g,f;J,H;\mathrm{๐})`$
$``$ $`{\displaystyle \underset{\begin{array}{c}\overline{x}\stackrel{~}{๐ซ}(H)\\ \mu _{\mathrm{CZ}}(\overline{x})=m\mu _{\mathrm{Morse}}(q)\end{array}}{}}^{\mathrm{PSS}}(q,\overline{x})\times _{(L,S)}^{}(\overline{x};J,H;\mathrm{๐})`$
$``$ $`\left\{(R,\gamma ,u)_{(L,S)}^{\mathrm{PSS},}(q;g,f;J,H;\mathrm{๐})\right|R=0\}.`$
If we consider a sequence $`(R_n,\gamma _n,u_n)_n`$, the first union collates breaking of the gradient half-trajectory. In this case the sequence $`(R_n)`$ converges. If the sequence $`R_n\mathrm{}`$, the sequence of half-cylinders $`u_n`$ breaks into a pair consisting of a half-cylinder and a plumberโs helper solution due to the chosen cut-off function $`\beta `$. This makes up the second union. The last union appears for obvious reasons.
Since the moduli spaces $`_{(L,S)}^{\mathrm{PSS},}(q;g,f;J,H;\mathrm{๐})`$ are compact in dimension 0, we can define
$$\theta (S):=\underset{\mu _{\mathrm{Morse}}(q)=dimMdimS1}{}\mathrm{\#}_2_{(L,S)}^{\mathrm{PSS},}(q;g,f;J,H;\mathrm{๐})q\mathrm{CM}^{}(f;/2),$$
(4.21)
which is *not* a cycle but merely a chain. Furthermore, we abbreviate the set
$$_{(L,S)}^{\mathrm{PSS},}(q;R=0):=\left\{(R,\gamma ,u)_{(L,S)}^{\mathrm{PSS},}(q;g,f;J,H;\mathrm{๐})\right|R=0\}$$
(4.22)
and define the cycles
$$\rho (S):=\underset{\mu _{\mathrm{Morse}}(q)=dimMdimS}{}\mathrm{\#}_2_{(L,S)}^{\mathrm{PSS},}(q;R=0)q\mathrm{CM}^{}(f;/2)$$
(4.23)
$$\text{and}\mathrm{\Phi }(S):=\underset{\overline{x}\stackrel{~}{๐ซ}_{s2m}(H)}{}\mathrm{\#}_2\left(_{(L,S)}^{}(\overline{x};J,H;\mathrm{๐})\right)\overline{x}\mathrm{CF}_{}(H).$$
(4.24)
From the decomposition of moduli space $`_{(L,S)}^{\mathrm{PSS},}(q;g,f;J,H;\mathrm{๐})`$ we conclude
$$\delta ^{\mathrm{Morse}}\left(\theta (S)\right)=\mathrm{PSS}\left(\mathrm{\Phi }(S)\right)\rho (S),$$
(4.25)
where the minus sign is arbitrary as long as we are working with $`/2`$-coefficients. We conclude that in cohomology $`[\mathrm{\Phi }(S)]=\mathrm{PSS}^1\left([\rho (S)]\right)`$. To finish the proof of theorems 3.1 and 3.3 we need to show $`\rho (S)=\mathrm{PD}[S]`$.
Either from equation (4.20) or directly from the definition we conclude that for an element $`(0,\gamma ,u)_{(L,S)}^{\mathrm{PSS},}(q;R=0)`$, the map $`u`$ is constant. Namely, since $`R=0`$ the Hamiltonian term in the Floer equation is 0, i.e. $`u`$ is a holomorphic disk satisfying $`E(u)=0`$ due to the assumption $`[u]=0`$. We conclude that $`_{(L,S)}^{\mathrm{PSS},}(q;R=0)`$ is in bijection to the space containing gradient half-trajectories $`\gamma :(\mathrm{},0]M`$ such that $`\gamma (\mathrm{})=q`$ and $`\gamma (0)S`$.
Choosing for instance a Morse function on the manifold $`S`$ and adding a quadratic function in the normal direction this is easily seen to define the Poincarรฉ dual of the class $`[S]\mathrm{H}_s(M;/2)`$ of $`S`$, i.e. $`\rho (S)=\mathrm{PD}[S]`$. The proof is immediately adjusted to the Poincarรฉ dual case. This concludes the proof of theorems 3.1 and 3.3. โ
## 5. Applications
### 5.1. Extrinsic topology of Lagrangian submanifolds
###### Definition 5.1.
We call a closed Lagrangian submanifold $`L`$ displaceable in the closed symplectic manifold $`(M,\omega )`$ if there exists a Hamiltonian diffeomorphism $`\mathrm{\Phi }_H\mathrm{Ham}(M,\omega )`$ such that $`L\mathrm{\Phi }_H(L)=\mathrm{}`$. In particular, $`\mathrm{\Phi }_H`$ displaces a sufficiently small neighborhood $`๐ฐ(L)`$ of $`L`$.
The displacement energy of $`L`$ is $`e(L):=inf\{HL\mathrm{\Phi }_H(L)=\mathrm{}\}`$, i.e. $`L`$ is displaceable iff $`e(L)<\mathrm{}`$.
###### Theorem 5.2.
In the situation of theorem 3.1 we denote by $`\iota _k:\mathrm{H}_k(L;/2)\mathrm{H}_k(M;/2)`$ the homomorphism induced by the inclusion $`\iota :LM`$. If the Lagrangian submanifold $`L`$ is displaceable, then the homomorphism $`\iota _k`$ vanishes for degrees $`k>dimL+1N_L`$.
The homomorphism $`\iota _0`$ does not vanish, so that $`N_LdimL+1`$. This is well-known for displaceable Lagrangian submanifolds.
###### Corollary 5.3.
For any displaceable monotone Lagrangian submanifold $`L`$ with $`N_L2`$,
$$[L]=0\mathrm{H}_m(M;/2).$$
(5.1)
###### Proof of theorem 5.2.
This is an application of theorem 3.1, namely we take a non-zero class $`0[S]\mathrm{H}_k(L;/2)`$ and represent its image $`\iota _k([S])\mathrm{H}_k(M;/2)`$ Floer theoretically with the help of theorem 3.1. Using the fact that $`L`$ is displaceable we will prove that this cycle in Floer homology vanishes simply by the fact that all moduli spaces involved are empty for a certain class of Hamiltonian functions. This shows that $`\mathrm{PSS}^1\left(\iota _k([S])\right)=0`$ and proves the assertion of theorem 5.2.
In contrast to theorem 3.1 we denote here by $`[S]`$ a class in $`\mathrm{H}_k(L;/2)`$ and by $`\iota _k([S])`$ its image in $`\mathrm{H}_k(M;/2)`$. We claim that under the assumption that $`L`$ is displaceable there exists a Hamiltonian function $`G:S^1\times M`$ such that $`_{(L,S)}^{}(\overline{x};J,G;\mathrm{๐})=\mathrm{}`$ for all $`\overline{x}\stackrel{~}{๐ซ}(G)`$ of Conley-Zehnder index $`\mu _{\mathrm{CZ}}(\overline{x})<N_L`$ and analogously for $`_{(L,S)}^+(\overline{x};J,G;\mathrm{๐})`$.
By assumption $`L`$ is displaceable by the time-1-map $`\mathrm{\Phi }_K\mathrm{Ham}(M,\omega )`$ associated to a Hamiltonian function $`K:S^1\times M`$, say. In particular, $`\mathrm{\Phi }_K\left(๐ฐ(L)\right)๐ฐ(L)=\mathrm{}`$ for some small neighborhood $`๐ฐ(L)`$ of $`L`$. We choose an autonomous Hamiltonian function $`H:M_0`$ such that the support of $`X_H`$ is contained in $`๐ฐ(L)`$ and $`H|_L>0`$. Since $`\mathrm{\Phi }_K`$ displaces the support of $`X_H`$, we draw the standard conclusion
$$\stackrel{~}{๐ซ}(\rho H\mathrm{\#}K)=\stackrel{~}{๐ซ}(K)$$
(5.2)
for all constants $`\rho 0`$ (cf. ). Furthermore, $`๐_{\rho H\mathrm{\#}K}(\overline{x})=๐_K(\overline{x})`$, since by assumption $`H0`$ outside $`๐ฐ(L)`$ and all elements of $`\stackrel{~}{๐ซ}(K)`$ and $`\stackrel{~}{๐ซ}(\rho H\mathrm{\#}K)`$ lie in $`M๐ฐ(L)`$.
We employ inequality (4.8) to obtain the inequality
$$0E(u)๐_{\rho H\mathrm{\#}K}(\overline{x})\underset{L}{inf}\left(\rho H\mathrm{\#}K\right)๐_K(\overline{x})\rho \underset{L}{inf}H\underset{L}{inf}K$$
(5.3)
for an element $`u_{(L,S)}^{}(\overline{x};J,\rho H\mathrm{\#}K;\mathrm{๐})`$. We conclude that for $`\rho 0`$ the right hand side becomes negative, since $`H|_L>0`$. In particular, all moduli spaces $`_{(L,S)}^{}(\overline{x};J,\rho H\mathrm{\#}K;\mathrm{๐})`$ are empty for sufficiently large $`\rho `$. This concludes the argument. โ
In contrast to theorem 5.5 below it is essential that we assume that $`L`$ is displaced by a Hamiltonian diffeomorphism and not a symplectic diffeomorphism as the examples $`S^1\times \{\mathrm{pt}\}S^1\times S^1`$ shows.
Theorem 5.2 was an application of theorem 3.1. Using the same idea we present a new proof of Chekanovโs result within our set-up as an application of theorem 3.3. We should mention that in this result is proved under the sole assumption that $`M`$ is geometrically bounded.
###### Theorem 5.4 (Chekanov ).
Under the same assumptions as in theorem 3.1 the displacement energy $`e(L)`$ of a monotone closed Lagrangian submanifold $`L`$ with $`N_L2`$ is bounded below by the minimal area of a non-constant holomorphic disk or sphere: $`e(L)\mathrm{min}\{A_M,A_L\}`$.
###### Proof.
This goes along the same lines as the proof of the preceding theorem. There is nothing to prove in case $`e(L)=\mathrm{}`$. Therefore, we assume from now on that $`L`$ is displaceable.
We choose a Hamiltonian function $`K`$ such that such that $`\mathrm{\Phi }_K^1`$ displaces a small neighborhood $`๐ฐ(L)`$ of $`L`$. Let us assume that the assertion of the theorem is false, i.e. $`K<\mathrm{min}\{A_M,A_L\}`$.
Again we choose a Hamiltonian function $`H:M_0`$ such that the support of $`X_H`$ is contained in $`๐ฐ(L)`$ and $`H|_L>0`$. For an element $`u_{(L,\mathrm{pt})}^{}(\overline{x};J,\rho H\mathrm{\#}K;\mathrm{๐})`$ we combine inequality (4.8) and the estimate of lemma 2.7 for $`\rho 0`$ to obtain
$`E(u)`$ $`๐_{\rho H\mathrm{\#}K}(\overline{x})\underset{L}{inf}(\rho H\mathrm{\#}K)`$ (5.4)
$`๐_K(\overline{x})\underset{L}{inf}K\rho \underset{L}{inf}H`$
$`\underset{M}{sup}K\underset{L}{inf}K\rho \underset{L}{inf}H`$
$`\underset{M}{sup}K\underset{M}{inf}K=K`$
$`<\mathrm{min}\{A_M,A_L\}.`$
The second inequality is valid by exactly the same reasoning as in the proof of theorem 5.2.
Now we want to apply theorem 4.3. Although the assumption $`\rho H\mathrm{\#}K<\mathrm{min}\{A_M,A_L\}`$ does not hold for large values of $`\rho `$ we established the crucial inequality $`E(u)<\mathrm{min}\{A_M,A_L\}`$ directly. As explained in remark 4.4 we can still apply theorem 3.3 and thus represent the class $`[\mathrm{pt}]\mathrm{H}_0(M;/2)`$ in the Floer homology $`\mathrm{HF}_{}(\rho H\mathrm{\#}K)`$.
As in the proof of theorem 5.2 we conclude from the displaceability of $`L`$ that the moduli spaces $`_{(L,\mathrm{pt})}^{}(\overline{x};J,\rho H\mathrm{\#}K;\mathrm{๐})`$ are empty if $`\rho `$ is sufficiently large. This contradicts the fact that counting the number of elements of $`_{(L,\mathrm{pt})}^{}(\overline{x};J,\rho H\mathrm{\#}K;\mathrm{๐})`$ defines the class $`[\mathrm{pt}]0`$. Therefore, our assumption $`K<\mathrm{min}\{A_M,A_L\}`$ was false. This proves the theorem. โ
### 5.2. Lagrangian intersections and spectral capacities
If we apply theorem 3.1 in the special case $`S=\mathrm{pt}`$ to spectral capacities, we obtain.
###### Theorem 5.5.
Let $`L_0,L_1`$ be two monotone closed Lagrangian submanifolds of the closed monotone symplectic manifold $`(M^{2m},\omega )`$ of minimal Maslov number $`N_{L_i}>dimL_i+1`$, $`i=0,1`$. If the spectral capacity $`c_\gamma (M)`$ of $`M`$ is finite, then $`L_0`$ and $`L_1`$ intersect.
###### Remark 5.6.
We note that theorem 5.5 is concerned with any pair of Lagrangian submanifolds. For example $`L_1`$ is not assumed to be a Hamiltonian deformation of $`L_0`$.
###### Proof.
As explained in remark 3.2, due to the assumption $`N_{L_i}>dimL_i+1`$, theorem 3.1 provides an explicit representation of $`\mathrm{PSS}^1([\omega ^m])`$ in the Floer chain complex. We can use both, $`L_0`$ and $`L_1`$, for such a representation.
By definition, $`c([\omega ^m],H)`$ is the smallest action value of a representant of $`[\omega ^m]`$ in Floer homology. In particular, from inequality (4.8) we conclude
$$c([\omega ^m],H)\underset{L_i}{inf}H\text{and}c(1,H)\underset{L_i}{sup}H,i=0,1$$
(5.5)
using the property Poincarรฉ duality for the action selector in the second case. We obtain
$$\gamma (H)\mathrm{max}\{\left(\underset{L_1}{sup}H+\underset{L_0}{inf}H\right),\left(\underset{L_0}{sup}H+\underset{L_1}{inf}H\right)\}.$$
(5.6)
If the Lagrangian submanifolds do not intersect, $`L_0L_1=\mathrm{}`$, we can choose a sequence $`\{H_n\}`$ of Hamiltonian functions such that $`H_n|_{L_0}=n`$ and $`H_n|_{L_1}=0`$. In particular, $`\gamma (H_n)n`$ and thus $`c_\gamma (M)=sup\{\gamma (H)\}=\mathrm{}.`$ This proves the theorem. โ
###### Corollary 5.7.
All monotone closed Lagrangian submanifolds $`L`$ in a monotone closed symplectic manifold of finite spectral capacity with minimal Maslov number $`N_L>dimL+1`$ are connected. Furthermore, $`L\phi (L)\mathrm{}`$ $`\phi \mathrm{Symp}(M,\omega )`$.
###### Proof.
Each connected component of $`L`$ has minimal Maslov number at least $`N_L`$. Therefore, two such components of $`L`$ intersect. For the second assertion is obvious from theorem 5.5. โ
###### Remark 5.8.
Theorem 5.5 says that if the monotone symplectic $`(M,\omega )`$ has two disjoint Lagrangian submanifolds of sufficiently large minimal Maslov number, then it has infinite spectral capacity.
The theorem allows for the following refinement: An arbitrarily small neighborhood of each such Lagrangian has infinite spectral capacity, namely choose the sequence $`\{H_n\}`$ from the above proof such that $`H_n`$ is supported in the small neighborhood. In particular, $`c_\gamma (L_i)=\mathrm{}`$.
###### Remark 5.9.
1. The assumption of theorem 5.5 that the symplectic manifold has finite spectral capacity is crucial as the example $`S^1\times \{\mathrm{pt}\}S^1\times S^1`$ shows: $`S^1\times \{\mathrm{pt}\}`$ is monotone and $`N_{S^1\times \{\mathrm{pt}\}}=\mathrm{}`$, but $`c_\gamma (S^1\times S^1)=\mathrm{}.`$
2. The assumption on the minimal Maslov number is necessary, since otherwise the Lagrangian submanifold might be displaceable. Indeed, there exist displaceable Lagrangian spheres $`L`$ in symplectic manifolds of the form $`X^{n+1}\times \mathrm{P}^n`$, cf. . Since $`L`$ is simply connected it is monotone and $`N_L=dimL+1`$. Furthermore, an analogous calculation as in lemma 5.12 shows $`c_\gamma (X\times \mathrm{P}^n)=c_\gamma (\mathrm{P}^n)<\mathrm{}`$ if $`\omega |_{\pi _2(X)}=0`$.
3. We do not know whether the Lagrangian submanifold $`L`$ has to be monotone but we suspect that theorem 5.5 does not generalize to non-monotone $`L`$. A counterexample could consist of an analog of two $`S^1S^2`$ where one of them is not an equator.
###### Theorem 5.10.
Any two monotone Lagrangian submanifolds $`L_0,L_1`$ of $`\mathrm{P}^n\times \mathrm{P}^n`$ with minimal Maslov number $`N_{L_i}>2n+1i=0,1`$ have to intersect.
###### Proof.
In view of theorem 5.5 we have to proof that $`\mathrm{P}^n\times \mathrm{P}^n`$ has finite spectral capacity. This is content of lemma 5.12 and is proved with help of lemma 5.11. โ
###### Lemma 5.11.
The monotone symplectic manifold $`(\mathrm{P}^n,\omega _{\mathrm{FS}})`$ has finite spectral capacity
$$\stackrel{~}{c}_{\stackrel{~}{\gamma }}(\mathrm{P}^n)=c_\gamma (\mathrm{P}^n)\omega _{\mathrm{FS}}(\mathrm{P}^1).$$
(5.7)
###### Proof.
The equality is due to the fact that the minimal Chern number $`N_{\mathrm{P}^n}`$ is sufficiently large, namely $`2N_{\mathrm{P}^n}>dim\mathrm{P}^n`$. We prove only the inequality.
The quantum cohomology ring of $`\mathrm{P}^n`$ over $`\mathrm{\Lambda }_0(\mathrm{\Gamma })=/2`$ is given by
$$\mathrm{QH}^{}(\mathrm{P}^n)\frac{(/2)[[q]][q^1][p]}{<p^{n+1}=q>},$$
(5.8)
i.e. Laurent series in $`q`$ and polynomials in $`p`$ up to degree $`n`$, where the class of the symplectic form corresponds to $`p`$. Furthermore, $`q`$ corresponds to the class $`[\mathrm{P}^1]`$ of the holomorphic sphere $`\mathrm{P}^1`$. In particular, the relation $`p^{n+1}=q`$ reads
$$[\omega _{\mathrm{FS}}][\omega _{\mathrm{FS}}^n]=1_{\mathrm{P}^n}e^{[\mathrm{P}^1]}.$$
(5.9)
See for details on the quantum cohomology ring of $`\mathrm{P}^n`$.
Using the properties Novikov and Product of the action selector we find
$`c([\omega _{\mathrm{FS}}][\omega _{\mathrm{FS}}^n],H)`$ $`=c(1_{\mathrm{P}^n}e^{[\mathrm{P}^1]},H)=c(1_{\mathrm{P}^n},H)\omega _{\mathrm{FS}}(\mathrm{P}^1)\text{and}`$ (5.10)
$`c([\omega _{\mathrm{FS}}][\omega _{\mathrm{FS}}^n],H)`$ $`c([\omega _{\mathrm{FS}}],0)+c([\omega _{\mathrm{FS}}^n],H)=c([\omega _{\mathrm{FS}}^n],H).`$
This implies $`\gamma (H)=c(1_{\mathrm{P}^n},H)c([\omega _{\mathrm{FS}}^n],H)\omega _{\mathrm{FS}}(\mathrm{P}^1)`$, and thus $`c_\gamma (\mathrm{P}^n)\omega _{\mathrm{FS}}(\mathrm{P}^1)`$. โ
###### Lemma 5.12.
The spectral capacity of $`(\mathrm{P}^n\times \mathrm{P}^n,\omega _{\mathrm{FS}}\omega _{\mathrm{FS}})`$ is finite:
$$c_\gamma (\mathrm{P}^n\times \mathrm{P}^n)2\omega _{\mathrm{FS}}(\mathrm{P}^1).$$
(5.11)
###### Proof.
We abbreviate
$$a:=[\omega _{\mathrm{FS}}]0,b:=0[\omega _{\mathrm{FS}}],A:=\mathrm{P}^1\times \{\mathrm{pt}\}\text{and}B:=\{\mathrm{pt}\}\times \mathrm{P}^1.$$
(5.12)
The class $`a^nb^n:=(a\mathrm{}a)(b\mathrm{}b)`$ represents (up to a factor $`\left(\genfrac{}{}{0pt}{}{2n}{n}\right)`$) the class $`[(\omega _{\mathrm{FS}}\omega _{\mathrm{FS}})^{2n}]`$. Since both factors in $`\mathrm{P}^n\times \mathrm{P}^n`$ have the same symplectic form we can use the Kuenneth formula for quantum cohomology \[11, chapter 11.1\]. Therefore, together with lemma 5.11 we compute in the quantum cohomology of $`\mathrm{P}^n\times \mathrm{P}^n`$:
$$(ab)(a^nb^n)=(aa^n)(bb^n)=(1_{\mathrm{P}^n}e^A)(1_{\mathrm{P}^n}e^B)=\mathrm{\hspace{0.33em}1}_{\mathrm{P}^n\times \mathrm{P}^n}e^{(A+B)}.$$
Now the reasoning is as above, namely
$`c((ab)(a^nb^n),H)`$ $`=c(1_{\mathrm{P}^n\times \mathrm{P}^n}e^{(A+B)},H)=c(1_{\mathrm{P}^n\times \mathrm{P}^n},H)2\omega _{\mathrm{FS}}(\mathrm{P}^1)`$ (5.13)
$`c((ab)(a^nb^n),H)`$ $`c(ab,0)+c(a^nb^n,H)=c(a^nb^n,H).`$
Note, that $`c(\left(\genfrac{}{}{0pt}{}{2n}{n}\right)a^nb^n,H)=c(a^nb^n,H)`$. The lemma follows. โ
###### Corollary 5.13.
Any two simply connected Lagrangian submanifolds in $`\mathrm{P}^n\times \mathrm{P}^n`$ intersect each other. Furthermore, they are all connected.
###### Proof.
Since a simply connected Lagrangian submanifold $`L`$ in a monotone symplectic manifold is monotone with minimal Maslov number equal to twice the minimal Chern number, we conclude $`N_L=2N_{\mathrm{P}^n\times \mathrm{P}^n}=2N_{\mathrm{P}^n}=2(n+1)>2n+1`$ and can apply theorem 5.10. โ
###### Example 5.14.
An example of a simply connected Lagrangian submanifold of $`\mathrm{P}^n\times \mathrm{P}^n`$ is the anti-diagonal $`\overline{\mathrm{\Delta }}=\{(\overline{z},z)z\mathrm{P}^n\}`$. Unfortunately, we are not aware of any other examples of simply connected Lagrangian submanifolds of $`\mathrm{P}^n\times \mathrm{P}^n`$ or of non-simply connected monotone Lagrangian submanifolds of sufficiently large minimal Maslov number.
###### Remark 5.15.
Biran informed us that using his techniques he can prove corollary 5.13, too. Using tools adapted to four dimensions, Hind proves that all Lagrangian spheres in $`S^2\times S^2`$ are Hamiltonianly isotopic to the anti-diagonal. Since the minimal Maslov number of such spheres equals 4 they are not displaceable and corollary 5.13 follows in the case $`n=1`$.
Cornea reported on an ongoing project with Lalonde at the โSMS 2004โ in Montreal<sup>1</sup><sup>1</sup>1see http://www.dms.umontreal.ca/sms/index.html. The following theorem is a special case of their results. We give here a new proof since it demonstrates our method quite nicely. Another proof of theorem 5.16, in spirit of Gromovโs theorem asserting that $`\mathrm{H}^1(L;)0`$ for all closed Lagrangian $`L^{2n}`$, was explained to us by Salamon.
###### Theorem 5.16.
Let $`(M,\omega )`$ be a symplectically aspherical manifold and $`LM`$ a monotone closed Lagrangian submanifold with $`N_L2`$. If $`L`$ is displaceable, then through *each* point of $`L`$ passes a non-constant holomorphic disk with boundary lying entirely on $`L`$ and which is of Maslov index $`\mu _{\mathrm{Maslov}}dimL+1`$.
###### Proof.
We argue by contradiction and assume that there is a point $`p_0L`$ through which passes no holomorphic disk of Maslov index $`\mu _{\mathrm{Maslov}}dimL+1`$. Then inspection of the proof of theorem 4.3 shows that the moduli spaces $`_{(L,\{p_0\})}^\pm (\overline{x};J,H;\mathrm{๐})`$ are compact up to splitting of multiple Floer cylinders for all Hamiltonian functions $`H`$. In other words, no bubbling-off of holomorphic disks can occur. Furthermore, no holomorphic spheres exist since $`M`$ is assumed to be symplectically aspherical.
Now we argue as in the proof of theorem 5.2 or 5.4. Namely, since $`L`$ is displaceable, we can show that the moduli spaces $`_{(L,\{p_0\})}^\pm (\overline{x};J,H;\mathrm{๐})`$ are empty for appropriately chosen Hamiltonian function $`H`$. On the other hand, we can represent the non-vanishing class $`[p_0]\mathrm{H}_0(M;/2)`$ by counting the number of solutions. This contradiction finishes the proof. โ
Since the dimension of $`_{(L,\{p_0\})}^\pm (\overline{x};J,H;\mathrm{๐})`$ equals $`\pm \mu _{\mathrm{CZ}}(\overline{x})dimL`$ we can relax the assumption of $`(M,\omega )`$ being symplectically aspherical to $`N_M>dimL`$. In this case there might be holomorphic spheres but the index formula (4.11) shows that the moduli spaces $`_{(L,\{p_0\})}^\pm (\overline{x};J,H;\mathrm{๐})`$ are compact for all relevant $`\overline{x}`$, more precisely for $`\mu _{\mathrm{CZ}}(\overline{x})=\pm dimM`$. For instance the theorem holds for $`\mathrm{P}^n`$.
Combining the proofs of theorem 5.5 and 5.16 we obtain the following assertion for a symplectic manifold $`(M,\omega )`$ of finite spectral capacity.
If there exists a point on the monotone Lagrangian submanifold $`LM`$ through which passes no non-constant holomorphic disk of Maslov index $`\mu _{\mathrm{Maslov}}dimL+1`$, e.g. $`N_L>dimL+1`$, then $`L\phi (L)\mathrm{}`$ for all $`\phi \mathrm{Symp}(M,\omega )`$.
The example $`S^1S^2`$ shows that this statement canโt be reversed. Through each point of an equator passes a holomorphic disk of Maslov index 2 and each image under a symplectomorphism intersect the equator again. More generally, this holds true for the Clifford torus in $`\mathrm{P}^n`$, see .
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# Optimal Quantum Filtering and Quantum Feedback Control
## 1 Introduction
With technological advances now allowing the possibility of continuous monitoring and rapid manipulations of systems at the quantum level , there is an increasing awareness of the applications and importance of quantum feedback control. Such applications include the engineering of quantum states, stability theory, quantum error correction and substantial applications in quantum computation . This current interest marks quantum control theory as a highly rewarding branch of control theory for study and as such there is a growing number of recent publications on the subject . In particular, contains a useful introduction to quantum probability and along with gives a comprehensive discussion on the comparison of classical and quantum control techniques and we refer the unfamiliar reader to these articles and references within.
The main ingredients of quantum control are essentially the same as in the classical case. One controls the system by coupling to an external control field which modifies the system in a desirable manner. The desired objectives of the control can be encoded into a *cost function* along with any other stipulations or restrictions on the controls such that the minimization of this cost indicates optimality of the control process. There are two types of dynamical control - *open loop* (or blind) control where the controls are predetermined at the start of the experiment and *closed loop* (or feedback) control where controls can be chosen throughout the experiment and thus is preferable for stochastic dynamics. Previous work on the theory of optimal quantum open loop control includes variational techniques on closed qubit systems , which was also extended to open (dissipative) quantum systems . However, this approach can only seek locally optimal solutions which can often be improved further with measurement and feedback, since an open quantum system inevitably loses information to its surrounding environment.
Quantum feedback control was formally initiated by Belavkin in a series of papers in the 80s. This work was developed as a quantum analogy to the classical theories of nonlinear (Stratonovich) filtering and Bellman dynamic programming. In fact, the separation lemma of classical control theory was shown also to hold in the quantum domain. That is, the problem of optimal quantum feedback control is separated into quantum filtering which provides optimal estimates of the stochastic quantum variables (operators) and then an optimal control problem based on the output of the quantum filter. The quantum noise which we filter out comes from the disturbance to the system due to the quantum measurement. Unlike classical systems, this is an unavoidable feature of quantum measurement since the quantum system is not directly observable. The quantum filter describes a classical stochastic process, albeit on the space of quantum states, so Belavkin showed how one can progress using tools from classical feedback control theory when applied to sufficient coordinates of the system . However, the lack of urgency for such a theory and the complexity of the mathematical language at the time left this work relatively undiscovered only to be rediscovered recently in the physics and engineering community.
The purpose of this paper is to build on the original work of Belavkin and present an accessible account of the theory of nonlinear optimal quantum feedback control. Firstly we introduce the necessary concepts from modern quantum theory including quantum probability, non-demolition measurement, quantum stochastic calculus and quantum filtering. Next the quantum Bellman equation for optimal feedback control with diffusive non demolition measurement is derived. Often in optimal control problems of this nature, the separation lemma is assumed and the control objectives are defined in terms of posterior sufficient coordinates . In this paper, we show how the general Bellman equation is applied with the same effect by application to the many dimensional quantum LQG problem. Next a physcial example of LQG control is given and we conclude with a discussion on the results with comparison to the corresponding classical control problem.
## 2 Optimal quantum measurement and filtering
This section highlights the differences between quantum and classical systems and introduces the problem of quantum measurement. After the appropriate setting is given, the measurement problem is then restated as a problem of optimal estimation of the output of a noisy quantum channel. Finally, the quantum filtering equation describing the dynamical least squares estimator is given.
### 2.1 Quantum Probability
Quantum physics which deals with the unavoidable random nature of the microworld requires a new, more general, noncommutative probability theory than the classical one based on Kolmogorovโs axioms. It was developed through the 70s and 80s by Accardi, Belavkin, Gardiner, Holevo, Hudson and Parthasarthy amongst others.
The essential difference between classical and quantum probability is that classically, Kolmogorovโs probability axioms allow the occurrence of simultaneous events only. This is because the classical events are described by indicator functions $`1_\mathrm{\Delta }(\omega )`$ of the measurable subsets $`\mathrm{\Delta }\mathrm{\Omega }`$ on the space of point states $`\mathrm{\Omega }`$. They are the building blocks for the classical random variables described by measurable functions $`x:\mathrm{\Omega }`$ as linear combinations (integrals) of the indicator functions $`1_\mathrm{\Delta }`$. Such classical essentially bounded variables represented by operators of multiplication by the corresponding functions, form an abelian (commutative) von Neumann algebra on the Hilbert space $`L^2(\mathrm{\Omega },)`$ of square-integrable random functions with respect to a probability measure $``$.
In quantum probability, there are some events which cannot occur simultaneously, so we must generalize the framework of classical probability to incorporate these features. This is done by considering quantum events as self adjoint orthoprojectors $`P^2=P=P^{}`$ (where $``$ denotes the Hermitian adjoint) acting in some Hilbert space $``$ not only by multiplications on the indicator functions $`1_\mathrm{\Delta }(\omega )`$. Quantum random variables are also built from events as linear (integral) combinations of their projectors $`P`$. The events are incompatible if the corresponding projectors do not commute, i.e. $`[P_i,P_j]:=P_iP_jP_jP_i0`$ and therefore cannot be represented classically by the indicator functions which always commute.
One can form the non commutative von Neumann algebra $`๐`$ of bounded quantum random variables generated by the self adjoint projectors $`\{P^1,..,P^m\}`$. This algebra is equal to its double commutant $`๐:=\{P^1,\mathrm{},P^m\}^{\prime \prime }`$ where the commutant of a set $`S()`$ in the algebra $`()`$ of all bounded operators on $``$ is defined by $`S^{}:=\{X^{}()\text{ s.t. }[X,X^{}]=0`$ $`XS\}`$. The quantum state on $`๐`$, given by a positive operator $`\rho =\rho ^{}0`$ with unit trace $`\text{Tr}[\rho ]=1`$, defines all expectations
$$X=\text{Tr}[\rho X]=\rho ,X$$
(1)
for operators $`X๐`$. So we describe a quantum probability space by the pair $`(๐,\rho )`$. In the case where $`๐`$ is an abelian von Neumann operator algebra, there is a natural isomorphism $`(๐,\rho )L^{\mathrm{}}(\mathrm{\Omega },)`$ with bounded functions on the classical probability space $`(\mathrm{\Omega },)`$ and so we recover the classical statistics.
The incompatibility of quantum events means that after one has observed an event, the state of the system needs to be updated to account for the change to the system or *back-action* affecting the expectations of all other incompatible events. This state change was traditionally described by the normalized projection postulate
$$\rho \rho _i=\frac{P_i\rho P_i}{\text{Tr}[\rho P_i]}$$
(2)
which also ensures instantaneous repeatability of the observed event corresponding to the projection $`P_i`$. However, it has long been known that this phenomological description is inadequate, since it fails to describe continuous measurements and experimentally it is not possible to perform a direct measurement of eigenstates of such a quantum operator. Instead we must consider an indirect measurement of operators in a coupled semi-classical field and describe the state change $`\rho \rho _i`$ by an optimal estimator based on the results of measurements in this field. Let $``$ denote the Hilbert space of the field, which we view as a noisy measurement channel in the initial vacuum state $`\varphi `$. We only observe compatible events in the channel (corresponding to output meter readings for example). So we describe these events by commuting projectors $`\{P_\omega \}_{\omega \mathrm{\Omega }}`$ which can be represented by classical indicator functions and generate the abelian subalgebra $`()`$ where $`\mathrm{\Omega }`$ is now the space of measurement results (eigenvalues) for these commuting operators. So the field operators $`W`$ which are linear combinations of the commuting projectors are in one-to-one correspondence with classical random variables as functions $`w:\mathrm{\Omega }`$ on the data space $`\mathrm{\Omega }`$. In the quantum noise model, we consider input quantum noises as quantum random variables represented by operators in the full field algebra $`()`$ of bounded operators on $``$ which perturb the quantum system in such a way to allow a classical correlated output. This interaction between the open quantum system and the semi-classical field is described on the composite system by a unitary operator $`U`$, which for an initial state $`\varphi `$ of the field gives the state evolution
$$\rho U(\rho \varphi )U^{}$$
called the *prior state*. The reduced conditional evolution can then be described by the nonlinear map
$$\rho \rho _\omega =\frac{\text{Tr}_{}[U(\rho \varphi )U^{}(IP_\omega )]}{\text{Tr}[U(\rho \varphi )U^{}(IP_\omega )]}$$
(3)
called the *posterior state* which is the Bayes law of conditioning for the measurement result $`\omega \mathrm{\Omega }`$, normalized with respect to the output probabilities $`(\omega )=\text{Tr}[U(\rho \varphi )U^{}(IP_\omega )]`$ and $`\text{Tr}_{}`$ denotes the partial trace over $``$.
We denote the posterior state as a classical random variable $`\rho _{}:\mathrm{\Omega }๐_{}`$ taking values $`\rho _\omega `$ in the space $`๐_{}`$ of states on $`๐`$. The posterior state gives the conditional expectation
$$๐ผ[X^{}|Y]=\rho _{},X$$
which is the least squares estimator of the system operator $`X^{}=U^{}(XI)U`$ after interaction, with respect to the output operators $`Y:=U^{}(IW)U`$. We now describe the appropriate model for the dynamical coupling between the open quantum system and the field.
### 2.2 The Quantum Vacuum Noise Model and Markov Approximation
The indirect measurement of the quantum system is via a coupled measurement channel, playing the role of a quantum noise bath. It is modelled by the symmetric Fock space $``$ over the single particle space $`L^2(_+๐ข)`$ of square integrable functions from $`[0,\mathrm{})`$ into a Hilbert space $`๐ข`$ of the bath degrees of freedom. Having in mind the vacuum noise model of the bath, let $`๐ฒ:=()`$ denote the quantum noise algebra of bounded operators on $``$ initially in the vacuum state $`\varphi `$. From the divisibility property of the symmetric Fock space, we can factorize the noise algebra
$$๐ฒ=๐ฒ_0^t๐ฒ_t^{\mathrm{}},=_{[0,t)}_{[t,\mathrm{})}$$
for arbitrary $`t>0`$ where $`๐ฒ_a^b=(_{[a,b)})`$ and $`_{[a,b)}`$ is the symmetric Fock space over $`L^2([a,b)๐ข)`$ for $`0a<b`$. This tensor independence implies compatibility for operators belonging to the disjoint time intervals of the noise algebra. The time evolution of the quantum system and the quantum noise bath (which together form a closed composite quantum system) can be described in the interaction representation by a family $`\{U_t\}_{t_+}`$ of unitary operators $`U_t:_{[0,t)}_{[0,t)}`$. In the weak coupling limit , (short bath memory), they describe the Markovian *flow* $`j_t:๐๐๐ฒ_0^t`$ by $`j_t(X):=U_t^{}(XI)U_t`$ for operators $`X๐`$ (we use the symbol $`X๐`$ to denote that $`X`$ is an element of $`๐`$, or that its spectral projectors belong in $`๐`$ for the case of unbounded $`X`$). We complete the description of the joint system and field evolution by introducing the unitary shift operator $`S_t:_{[0,s)}_{[t,s+t)}`$ which models the free evolution in the field. Thus the combined evolution and interaction on the composite system is given by a family of endomorphisms $`\{\gamma _t\}`$ on $`๐๐ฒ`$ such that $`\gamma _t(XW)=\widehat{U}_t^{}(XW)\widehat{U}_t`$ for unitaries $`\widehat{U}_t:=(IS_t)U_t`$. This gives the cocyle identity $`U_{t+s}=S_sU_tS_sU_s`$ for the interaction unitaries $`\{U_t\}`$. Note that for ease of presentation, we avoid the repetition of tensoring with the identity on $``$ and $`_{[t,\mathrm{})}`$ and assume the domain of the operators is clear from the context.
We now briefly discuss quantum stochastic calculus, a necessary tool when developing a time-continuous theory of quantum stochastic evolution.
### 2.3 Quantum Stochastic Calculus
In this paper we consider feedback control based on a homodyne detection scheme. This is the quantum analogue of measurement of the Wiener process in the field and is described by the field quadrature $`W_t=A_t+A_t^{}`$ where $`A_t๐ฒ_0^t`$ is called the *annihilation* operator on $``$. The properties of $`A_t`$ are such that $`W_t^\theta :=\mathrm{exp}(\text{i}\theta )A_t+\mathrm{exp}(\text{i}\theta )A_t^{}`$ is equivalent to the classical Wiener process for each $`\theta [0,2\pi )`$, however they do not commute for different $`\theta `$, so by considering solely the measurement of $`W_t`$, we restrict ourselves to a chosen classical diffusive measurement process corresponding to $`\theta =0`$.
Hudson and Parthasarathy , developed the theory of quantum stochastic calculus using the annihilation process and its adjoint, the creation process $`A_t^{}`$ as the fundamental diffusive adapted processes and defined the interaction unitaries $`\{U_t\}`$ as the unique solutions to the quantum stochastic differential equation which we chose of the simple form
$$dU_t+KU_tdt=LU_tdA_t^{}L^{}U_tdA_t.$$
(4)
with $`U_0=I`$. Here $`K=\frac{\text{i}}{\mathrm{}}H+\frac{1}{2}L^{}L`$, $`H`$ is the Hamiltonian of the quantum system and $`L`$ is the operator describing the coupling of the system to the measurement channel. The increments $`dt`$, $`dA_t`$, $`dA_t^{}`$ are considered as operators acting in $`_{[t,t+dt)}`$ and define stochastic Itรด calculus using the product rule
$$d(M_tN_t)=d(M_t)N_t+M_td(N_t)+d(M_t)d(N_t)$$
for adapted quantum stochastic processes $`M_t`$, $`N_t`$ where the quantum Itรด correction term (last term) is calculated using the multiplication table
$$\begin{array}{c}(dt)^2=0,dtdA_t=0=dtdA_t^{},\hfill \\ dA_t^{}dA_t=0,dA_tdA_t^{}=dt.\hfill \end{array}$$
(5)
### 2.4 Quantum Langevin Equations and Non-demolition Measurements
From the quantum Itรด formula applied to $`X_t=U_t^{}(XI)U_t`$ and the quantum Itรด multiplication table (5), we obtain the quantum Langevin equation
$$dX_t=_t[X_t]dt+[X_t,L_t]dA_t^{}[X_t,L_t^{}]dA_t.$$
(6)
Here $`_t[X_t]=j_t([X])`$ is the time evolved Lindblad (or Gorini-Kossakovski-Sudarshan) generator
$$[X]=\frac{\text{i}}{\mathrm{}}[H,X]+\frac{1}{2}(L^{}[X,L]+[L^{},X]L)$$
(7)
for the semigroup of completely positive maps describing the dissipative evolution in the Markovian limit. The dual $`^{}`$ of this map describes the unconditional dissipative evolution of states
$$\frac{d}{dt}\rho ^t=\frac{\text{i}}{\mathrm{}}[H,\rho ^t]+\frac{1}{2}(L[\rho ^t,L^{}]+[L,\rho ^t]L^{})$$
(8)
called the *master equation* which is the quantum analogue of the Focker-Plank equation. A time continuous measurement of the field quadrature $`W_t`$ in the output channel represents an indirect measurement of the evolved generalized coordinate $`L_t+L_t^{}๐_t`$ as can be seen from the quantum Itรด formula applied to the output operators $`Y_t=U_t^{}(IW_t)U_t`$:
$$dY_t=(L_t+L_t^{})dt+IdW_t.$$
(9)
Note that the output process $`Y_t`$ is directly observable as it is a commutative family of self-adjoint operators $`\{Y_s\}_{st}`$ unitary equivalent to the family $`\{W_s\}_{st}`$ for each $`t`$. This simply follows from the following lemma which was first observed by Belavkin in ,.
###### Lemma 1.
The input and output operators satisfy the quantum non-demolition (QND) condition
$$[X_t,Y_s]=0[Y_t,Y_s]=00st$$
(10)
###### Proof.
Let $`t=s+r,r>0`$, then from the cocycle identity we get
$$\begin{array}{c}U_{s+r}^{}(IW_s)U_{s+r}=\hfill \\ U_s^{}(S_sU_rS_s)^{}(IW_s)(S_sU_rS_s)U_s=\hfill \\ U_s^{}(IW_s)U_s=Y_s\hfill \end{array}$$
where the last step uses the commutativity of $`S_sU_rS_s๐๐ฒ_s^{s+r}`$ and $`W_s๐ฒ_0^s`$. So $`[X_t,Y_s]=U_t^{}[X,W_s]U_t=0`$ and $`[Y_t,Y_s]=U_t^{}[W_t,W_s]U_t=0`$ follows from the tensor independence of $`X`$, $`W_s`$ and $`W_t`$ for all $`st`$. โ
### 2.5 Quantum Filtering
Classically, filtering equations are used when we need to estimate the value of dynamical variables about which we have incomplete knowledge due to an indirect observation. For example, the Kalman-Bucy filter , gives a continuous least-squares estimator for a Gaussian classical random variable with linear dynamics when we only have access to a correlated, noisy output signal. Since closed quantum systems are fundamentally unobservable unless they are open, e.g. disturbed by quantum noise processes (c.f. (6),(9)), filtering of quantum noise plays an important role in quantum measurement. Belavkin was the first to realize that an optimal estimation without further disturbance is possible in the Markovian limit and is based on an output nondemolition measurement ,,. He constructed the quantum filtering equation which describes the evolution of the optimal estimate given by the density matrix conditioned on a classical output of the noisy quantum channel. This is used to estimate arbitrary input operators $`X_t๐_t`$ which are driven by environmental quantum noises. The previous lemma shows that the expectation of $`X_t`$ is not disturbed when we measure $`Y_s`$ for $`0st`$. This is necessary for the existence of a well defined conditional expectation of $`X_t`$ with respect to past measurement results of $`Y_s`$.
Let $`๐_s^t:=\{Y_s^t\}^{\prime \prime }`$ be the abelian von Neumann algebra generated by the output operators $`Y_s^t:=\{Y_r|srt\}`$ (or their spectral projectors in the case of unbounded $`Y_r`$). Also let $`๐_s^t=\{X_r|srt\}^{\prime \prime }`$ denote the von Neumann algebra generated by the system operators $`X_r๐_r`$. From the QND condition, $`๐_0^t`$ lies in the center of (i.e. it is a subalgebra commuting with the whole of) $`_t^T๐๐ฒ_0^T`$, where $`_t^T:=๐_t^T๐_0^T`$ is the smallest von Neumann algebra containing $`๐_t^T`$ and $`๐_0^T`$ as subalgebras for $`0tT`$. This gives the necessary conditions for the existence of a conditional expectation , defined as a linear, normcontractive projection $`E_0^T:_t^T๐_0^t`$.
The conditional expectation $`๐ผ[X_t|Y_0^t]:=E_0^t[X_t]`$ gives the least squares estimator $`\widehat{X}_t`$ of an operator $`X_t๐_t`$ conditional on the output operators $`Y_0^t`$ and so is equivalent to a classical random variable on the space of measurement trajectories $`\mathrm{\Omega }_0^t:=\{\omega _s|0st`$ s.t. $`\omega _s`$ is an eigenvalue of $`Y_s\}`$. This conditional expectation is most conveniently written in the Schrรถdinger picture $`E_0^t[X_t]=\rho _{}^t,X`$ for the solution $`\rho _{}^t`$ to the classical stochastic nonlinear differential equation
$$d\rho _{}^t=^{}[\rho _{}^t]dt+\sigma (\rho _{}^t)(dY_t\rho _{}^t,L+L^{}dt)$$
(11)
often called the Belavkin quantum filtering equation, where
$$\sigma (\rho _{}^t)=\rho _{}^tL^{}+L\rho _{}^t\rho _{}^t,L^{}+L\rho _{}^t$$
is the nonlinear *fluctuation coefficient*.
We can generalize the filtering equation to the case where we couple the open quantum system to $`d`$ independent measurement channels. If we assume no scattering between the channels, then the family of unitary operators $`\{U_t\}_{t_+}`$ describing the evolution in the interaction picture $`U_t:_{[0,t)}^d_{[0,t)}^d`$ satisfy
$$dU_t+KU_tdt=\underset{i=1}{\overset{d}{}}[L_iU_tdA_{i,t}^{}L_i^{}U_tdA_{i,t}]$$
where $`L_i`$ describes the coupling to the $`i`$th channel and $`K=\frac{\text{i}}{\mathrm{}}H+\frac{1}{2}_{i=1}^dL_i^{}L_i`$. Throughout this paper we reserve the Roman character i to denote the imaginary unit $`\text{i}:=\sqrt{1}`$, whereas italic $`i`$ is freely used as an index. Note that we have tensor independence of the annihilation increments $`dA_{i,t},dA_{j,t}`$ for $`ij`$, so the quantum vacuum noises commute for different channels. The Belavkin filtering equation for a simultaneous diffusive measurement of $`Y_{i,t}=U_t^{}(IW_{i,t})U_t`$ gives
$$d\rho _{}^t=^{}[\rho _{}^t]dt+\underset{i=1}{\overset{d}{}}\sigma _i(\rho _{}^t)(dY_{i,t}\rho _{}^t,L_i+L_i^{}dt)$$
(12)
for $`W_{i,t}=(A_{i,t}+A_{i,t}^{})`$.
## 3 Optimal Quantum Control
We now couple the system to a control field. If we assume no scattering between the measurement and control fields and assume a weak coupling such that information is not lost into the control field, then this effectively replaces the Hamiltonian $`H`$ of the system with a controlled Hamiltonian $`H(u_s)`$ for admissible real valued control functions $`u_s`$ say, at time $`s`$. This Hamiltonian generates the controlled unitaries $`U_t(u_0^t)`$ giving the controlled flow
$$j_t(u_0^t)[X]:=U_t^{}(u_0^t)(XI)U_t(u_0^t)$$
where $`u_0^t:=\{u_s|0s<t\}`$ is the control process over the interval $`[0,t)`$. The controlled posterior density operator $`\rho _{}^t(u_0^t)`$ can then be obtained from (12) with the controlled Hamiltonian $`H(u_t)`$ which appears in the controlled Lindblad term $`(u_t)`$.
In classical control, we can allow complete observability of the controllable system, so that feedback controls are determined by the system variables $`x_tu_t(x_t)`$. However, in quantum systems, we do not have the point states $`x_t`$ due to joint non observability of the system operators $`X_t`$, so the stochastic feedback controls should be given by a function of the stochastic output process $`Y_0^t`$ which is associated with a classical random variable $`u_t()`$ on $`\mathrm{\Omega }_0^t`$. I.e. the measurement trajectory is fed into the control $`\omega _0^tu_t(\omega _0^t)`$. Thus the feedback controlled flow is a map $`j_t(u_0^t(Y_0^t))`$ from $`๐`$ to $`๐_t๐_0^t`$.
The optimality of control is judged by the expected cost associated to the admissible control process $`u_0^T`$ for the finite duration $`T`$ of the experiment. Admissible control strategies are defined as those $`u_0^T`$ for which the operator valued cost integral
$$J(u_0^T)=_0^Tj_s(u_0^s)[C(u_s)]๐s+j_T(u_0^T)[S]$$
(13)
exists in the strong operator topology for self adjoint positive operators $`C(u_s),S๐`$ giving the expected cost by the expectation
$$\rho \varphi ,J(u_0^T).$$
(14)
An optimal feedback control strategy $`u_0^T()`$ for nondemolition measurements of the output operators $`Y_0^T`$ is one which minimizes the expected posterior cost-to-go
$$\rho \varphi ,J(u_0^T())=\underset{u_0^T()U_0^T()}{\mathrm{min}}\rho \varphi ,J(u_0^T())$$
where $`U_0^T()`$ is the space of admissible stochastic control strategies $`u_0^T()`$. This dynamical optimization problem is considerably simplified by the following Lemma first observed by Bellman.
###### Lemma 2 (Principle of Optimality).
If $`u_0^T()`$ is an optimal strategy for the cost function (13) given the initial state $`\rho \varphi `$, then its restriction $`u_t^T()`$ to the interval $`[t,T)`$ is optimal for the *cost-to-go*
$$\begin{array}{c}J_t(u_t^T())=_t^Tj_s(u_t^s())[C(u_s())]๐s\hfill \\ +j_T(u_t^T())[S]\hfill \end{array}$$
(15)
given the state $`\rho _{}^t(u_0^t())`$ at time $`t`$.
We can now reduce the dynamics to the observable output algebra and rewrite the expectation as a conditional one
$$\rho \varphi ,J_t(u_t^T(Y_0^t))=\varphi _0^t,E_0^t[J_t(u_t^T(Y_0^t))]$$
where $`E_0^t:_t^T๐_0^t`$ is the conditional expectation on $`_t^T=๐_t^T๐_0^T`$ which defines the feedback controlled posterior density operator by $`\rho _{}^t(u_0^t()),X=E_0^tj_t(u_0^t())[X]`$ at time $`t`$.
###### Theorem 1.
The posterior cost-to-go from state $`\rho `$ at time $`t`$ satisfies
$$E_0^t[J_t(u_t^T())]=๐ผ_0^t[\text{J}(t,u_t^T(),\rho )]$$
(16)
where
$`\text{J}(t,u_t^T(),\rho )=`$ $`{\displaystyle _t^T}\rho _{}^s(u_t^s()),C(u_s())๐s`$
$`+\rho ^T(u_t^T()),S`$
is a random variable on $`\mathrm{\Omega }_0^T`$ and $`\rho _{}^s(u_t^s())`$ is the solution to the controlled filtering equation for $`st`$ with $`\rho =\rho _{}^t(u_0^t())`$.
###### Proof.
The โquantumโ conditional expectation $`E_0^t`$ acting on future operators gives
$$E_0^tj_s(u_t^s())[X]=๐ผ_0^t[\rho _{}^s(u_t^s()),X]$$
for $`X๐`$, where $`๐ผ_0^t:๐_0^T๐_0^t`$ is the โclassicalโ conditional expectation on $`๐_0^T`$ satisfying the tower property $`๐ผ_0^t๐ผ_0^s=๐ผ_0^t`$ for $`tsT`$. โ
Let us denote the minimum posterior cost-to-go
$$\text{S}(t,\rho ):=\underset{u_t^T()U_t^T()}{\mathrm{min}}๐ผ_0^t[\text{J}(t,u_t^T(),\rho )].$$
(17)
###### Theorem 2.
The minimum posterior cost-to-go satisfies the Bellman equation
$$\frac{}{t}\text{S}(t,\rho )+\frac{1}{2}\underset{i=1}{\overset{d}{}}\sigma _i(\rho )\sigma _i(\rho ),(\delta \delta )\text{S}(t,\rho )$$
$$+\underset{u_t()}{\mathrm{min}}\left\{\rho ,C(u_t())+(u_t())[\delta \text{S}(t,\rho )]\right\}=0$$
(18)
where $`\delta \text{S}(t,\rho )๐`$ denotes the derivation of $`\text{S}(t,\rho )`$ with respect to $`\rho `$ and $`\sigma _i(\rho )`$ is the non-linear fluctuation coefficient in the filtering equation (12).
###### Proof.
From the definition of $`\text{S}(t,\rho )`$ and $`\text{J}(t,u_t^T(),\rho )`$, we have
$$\text{S}(t,\rho ^t)=\underset{u_t^T()}{\mathrm{min}}๐ผ_0^t\left\{\begin{array}{c}_t^{t+ฯต}\rho _{}^s(u_t^s()),C(u_s())๐s\hfill \\ +\text{J}(t+ฯต,u_{t+ฯต}^T(),\rho ^{t+ฯต})\hfill \end{array}\right\}$$
So when $`ฯตdt`$ becomes sufficiently small, we approximate this by
$$\text{S}(t,\rho ^t)=\underset{u_t()}{\mathrm{min}}๐ผ_0^t\left\{\begin{array}{c}\rho ^t,C(u_t())dt\hfill \\ +\text{S}(t+dt,\rho ^{t+dt})\hfill \end{array}\right\}$$
(19)
where we use the tower property of the classical conditional expectation. Assuming that $`\text{S}(t,\rho ^t)`$ is sufficiently differentiable, we use the Taylor expansion
$$\begin{array}{c}\text{S}(t+dt,\rho ^{t+dt})=\hfill \\ \text{S}(t,\rho ^t)+\frac{}{t}\text{S}(t,\rho ^t)dt+d\rho ^t,\delta \text{S}(t,\rho ^t)+\hfill \\ \frac{1}{2}_{i=1}^d\sigma _i(\rho ^t)\sigma _i(\rho ^t),(\delta \delta )\text{S}(t,\rho ^t)dt\hfill \end{array}$$
where $`\delta \text{S}(t,\rho ):=\frac{\delta }{\delta \rho }\text{S}(t,\rho )`$ denotes the derivation of $`\text{S}(t,\rho )`$ with respect to $`\rho `$. Using this expansion in (19) gives the Bellman equation (18) when we observe that $`\text{S}(t,\rho )+\frac{}{t}\text{S}(t,\rho )`$ does not depend on $`u_t`$ and $`๐ผ_0^t[d\stackrel{~}{Y}_{i,t}]=0`$ for the innovation process $`d\stackrel{~}{Y}_{i,t}=dY_{i,t}\rho ^t,L_i+L_i^{}dt`$. โ
## 4 Application of Results to a Linear Quantum Dynamical System
We illustrate the ideas of quantum filtering and control described above by application to the multidimensional quantum LQG control problem. LQG control is well studied in classical control theory and we shall see many similarities between quantum and classical LQG control theory.
### 4.1 Quantum Filtering of Linear, Gaussian Dynamics
Let $`๐ฟ`$ be the phase space vector of self adjoint operators $`X^i`$, $`i=1,\mathrm{},m`$ satisfying the canonical commutation relations (CCRs)
$$[X^i,X^j]=X^iX^jX^jX^i=\text{i}\mathrm{}J^{ij}I$$
for $`i,j=1,\mathrm{}m`$ where $`I`$ is the identity operator on $``$. The CCRs can be written in vector form as
$$[๐ฟ,๐ฟ^{}]:=๐ฟ๐ฟ^{}(๐ฟ๐ฟ^{})^{}=\text{i}\mathrm{}๐I$$
where $`๐ฟ^{}=(X^1,\mathrm{},X^m)`$ is the row vector transpose of $`๐ฟ`$ and $`๐=(J^{ij})`$ is an anti-symmetric real valued matrix which is assumed to be nondegenerate for an even $`m=2d`$ say. We couple the open quantum system to $`d`$ measurement channels via the operator vector $`๐ณ=๐ฒ๐ฟ`$, where $`๐ฒ`$ is a $`d\times m`$ matrix of complex-valued coefficients. Let us place it in a controllable potential which is described by the Hamiltonian
$$H(๐_t)=\frac{1}{2}๐ฟ^{}๐๐ฟ+๐ฟ^{}๐๐_t+๐_t^{}๐^{}๐ฟ$$
(20)
for real vector valued control parameters $`๐_t^d`$, where $`๐`$ is a real symmetric $`m\times m`$ matrix and $`๐`$ is a complex $`m\times d`$ matrix. We shall use $`๐ฒ^{}`$ to denote complex conjugation $`(๐ฒ^{})_{ij}=\mathrm{\Lambda }_{ij}^{}`$ and $`๐ฒ^{}=(๐ฒ^{})^{}`$ the Hermitian conjugate.
These definitions allow us to calculate the components of the controlled Lindblad generator from (7) with the controlled Hamiltonian (20) which we write here in vector form
$$(๐_t)[๐ฟ]=๐(๐+\mathrm{}\mathrm{}(๐ฒ^{}๐ฒ))๐ฟ+๐(๐+๐^{})๐_t$$
omitting the identity $`I`$ for notational convenience where $`2i\mathrm{}(๐ฒ^{}๐ฒ)=๐ฒ^{}๐ฒ๐ฒ^{}๐ฒ^{}`$. So from (6) and (9) we obtain the following quantum linear Langevin vector equation
$$d๐ฟ_t=(๐๐ฟ_t+๐๐_t)dt+d๐ฝ_t$$
(21)
and linear output equation
$$d๐_t=๐๐ฟ_tdt+d๐พ_t$$
(22)
where $`๐:=๐(๐+\mathrm{}\mathrm{}(๐ฒ^{}๐ฒ))`$, $`๐:=๐(๐+๐^{})`$, $`๐:=๐ฒ+๐ฒ^{}`$. The quantum noise increments are given by vectors
$`d๐ฝ_t`$ $`=`$ $`i\mathrm{}๐(๐ฒ^{}d๐จ_t^{}๐ฒ^{}d๐จ_t)`$
$`d๐พ_t`$ $`=`$ $`d๐จ_t+d๐จ_t^{}`$
for $`(๐จ_t)_i=A_{i,t}`$ the annihilation operator on the $`i`$th coupled independent measurement channel.
Let us denote the initial mean $`\overline{๐ฟ}`$ of the phase space operator vector by the component wise expectation $`(\overline{๐ฟ})^i=\overline{X}^i=\rho ,X^i`$ and symmetric covariance
$$\mathrm{\Sigma }^{ij}:=\frac{1}{2}\rho ,X^iX^j+X^jX^i\overline{X}^i\overline{X}^j.$$
which is given by a real positive definite matrix $`๐บ=(\mathrm{\Sigma }^{ij})`$ satisfying the Heisenberg uncertainty principle
$$๐บ\pm \frac{\text{i}\mathrm{}}{2}๐$$
(23)
The filtering equation (12) preserves the Gaussian nature of the posterior state , so the posterior mean $`(\widehat{๐ฟ_t})^i=\widehat{X}_t^i=\rho _{}^t,X^i`$ and symmetric error covariances
$$\mathrm{\Sigma }_t^{ij}:=\frac{1}{2}\rho _{}^t,X^iX^j+X^jX^i\widehat{X}_t^i\widehat{X}_t^j$$
form a set of sufficient coordinates for the quantum LQG system and agree with the initial mean and covariance for $`\rho _{}^0=\rho `$. Using (12), the posterior expectation of $`๐ฟ_t`$ for non demolition measurement of the output operators $`๐_t`$ is given in vector form
$`d\widehat{๐ฟ}_t`$ $`=`$ $`(๐\widehat{๐ฟ}_t+๐๐_t)dt+\stackrel{~}{๐}_td\stackrel{~}{๐}_t`$ (24)
$`\stackrel{~}{๐}_t`$ $`=`$ $`๐บ_t๐^{}+๐`$ (25)
where $`d\stackrel{~}{๐}_t=d๐_t๐\widehat{๐ฟ}_tdt`$ is the innovating martingale which describes the information gain from measurement of the output vector operator $`๐_t`$.
The symmetric error covariance $`๐บ_t`$ satisfies the matrix Ricatti equation
$$\begin{array}{cc}\hfill \frac{d}{dt}๐บ_t=& ๐๐บ_t+๐บ_t๐^{}+๐\hfill \\ & (๐บ_t๐^{}+๐)(๐บ_t๐^{}+๐)^{}\hfill \\ \hfill ๐บ_0=& ๐บ\hfill \end{array}$$
(26)
where
$$๐=\frac{1}{2}\mathrm{}^2๐(๐ฒ^{}๐ฒ+๐ฒ^{}๐ฒ^{})๐^{}$$
is the intensity (symmetric covariance) matrix of the quantum noise increment $`d๐ฝ_t`$ and
$$๐=\frac{i}{2}\mathrm{}๐(๐ฒ^{}๐ฒ^{})$$
is the covariance matrix of the noise increments $`d๐ฝ_t`$ and $`d๐พ_t`$.
### 4.2 Quantum LQG Control
We aim to control the phase space operator whilst constraining the amplitude of the controlling force for energy considerations. Thus, our control objectives and restraints can be described by the operator valued risk (13) with quadratic parameters
$$C(๐_s)=๐ฟ^{}๐
๐ฟ+๐ฟ^{}๐^{}๐_s+๐_s^{}๐๐ฟ+๐_s^{}๐_s$$
and $`S=๐ฟ^{}๐๐ฟ`$ for positive real symmetric $`m\times m`$ matrices $`๐,๐
`$ and a real $`d\times m`$ matrix $`G`$.
Since $`\widehat{๐ฟ}`$ and $`๐บ`$ form a set of sufficient coordinates, they describe the full probability distribution given by $`\rho `$, so we may consider the derivation of $`\text{S}(t,\rho )`$ as partial derivatives of $`\text{S}(t,\widehat{๐ฟ},๐บ)`$. So from (21) and the Gaussian nature of the system, we obtain
$$\begin{array}{c}\rho ,(๐_t)[\delta \text{S}(t,\widehat{๐ฟ},๐บ)]=\hfill \\ \frac{1}{2}(๐\widehat{๐ฟ}+๐๐_t)^{}_{\widehat{๐ฟ}}\text{S}+\frac{1}{2}_{\widehat{๐ฟ}}\text{S}^{}(๐\widehat{๐ฟ}+๐๐_t)\hfill \\ +(๐๐บ+๐บ๐^{}+๐,_๐บ\text{S})\hfill \end{array}$$
$$\begin{array}{c}_{j=1}^d\sigma _j(\rho )\sigma _j(\rho ),(\delta \delta )\text{S}(t,\widehat{๐ฟ},๐บ)=\hfill \\ ((๐บ๐^{}+๐)(๐บ๐^{}+๐)^{},_{\widehat{๐ฟ}}^2\text{S}2_๐บ\text{S})\hfill \end{array}$$
where $`(๐,๐):=\text{Tr}[๐^{}๐]`$ is the Hilbert-Schmidt inner product on the vector space of complex-valued $`m\times m`$ matrices. We denote the partial derivatives by $`(_{\widehat{๐ฟ}}\text{S})_i=\frac{}{\widehat{๐ฟ}^i}\text{S}(t,\widehat{๐ฟ},๐บ)`$, $`(_๐บ\text{S})_{ij}=\frac{}{๐บ_{ij}}S(t,\widehat{๐ฟ},๐บ)`$ and $`(_{\widehat{๐ฟ}}^2\text{S})_{ij}=\frac{}{\widehat{๐ฟ}^i}\frac{}{\widehat{๐ฟ}^j}\text{S}(t,\widehat{๐ฟ},๐บ)`$. Inserting into the Bellman equation (18) and minimizing gives $`๐_t=(\frac{1}{2}๐^{}_{\widehat{๐ฟ}}\text{S}+๐\widehat{๐ฟ})`$ where $`\text{S}(t,\widehat{๐ฟ},๐บ)`$ now satisfies the nonlinear partial differential equation
$$\begin{array}{c}\frac{}{t}\text{S}(t,\widehat{๐ฟ},๐บ)=\hfill \\ \frac{1}{2}(\widehat{๐ฟ}^{}๐^{}_{\widehat{๐ฟ}}\text{S}+_{\widehat{๐ฟ}}\text{S}^{}๐\widehat{๐ฟ})+\widehat{๐ฟ}^{}๐
\widehat{๐ฟ}\hfill \\ +(๐๐บ+๐บ๐^{}+๐,_๐บ\text{S})+(๐บ,๐
)\hfill \\ (\frac{1}{2}๐^{}_{\widehat{๐ฟ}}\text{S}+๐\widehat{๐ฟ})^{}(\frac{1}{2}๐^{}_{\widehat{๐ฟ}}\text{S}+๐\widehat{๐ฟ})\hfill \\ +((๐บ๐^{}+๐)(๐บ๐^{}+๐)^{},\frac{1}{2}_{\widehat{๐ฟ}}^2\text{S}_๐บ\text{S})\hfill \end{array}$$
(27)
which is called the Hamilton-Jacobi-Bellman (HJB) equation for this example.
It is well known from classical control theory that LQG control gives a posterior cost-to-go which is quadratic in the posterior mean. So we use the ansatz
$$\text{S}(t,\widehat{๐ฟ},๐บ)=\widehat{๐ฟ}^{}๐_t\widehat{๐ฟ}+๐_t,๐บ+\alpha _t$$
in the HJB equation (27). This gives the optimal feedback control strategy
$`๐_t`$ $`=`$ $`\stackrel{~}{๐}_t\widehat{๐ฟ}_t`$ (28)
$`\stackrel{~}{๐}_t`$ $`=`$ $`๐^{}๐_t+๐`$ (29)
which is linear in the solution to the filtering equation $`\widehat{๐ฟ}_t`$ at time $`t`$ where $`๐_t`$ satisfies the backwards matrix Ricatti equation
$$\begin{array}{cc}\hfill \frac{d}{dt}๐_t=& ๐_t๐+๐^{}๐_t+๐
\hfill \\ & (๐^{}๐_t+๐)^{}(๐^{}๐_t+๐)\hfill \\ \hfill ๐_T=& ๐\hfill \end{array}$$
(30)
and $`\alpha _t`$ satisfies
$$\begin{array}{cc}\hfill \frac{d}{dt}\alpha _t=& ((๐^{}๐+๐)^{}(๐^{}๐+๐),๐บ_t)\hfill \\ & +(๐_t,๐)\hfill \\ \hfill \alpha _T=& 0.\hfill \end{array}$$
(31)
From this we obtain the total minimal cost
$$\begin{array}{c}\text{S}(0,\overline{๐ฟ},๐บ)=\hfill \\ \overline{๐ฟ}^{}๐_0\overline{๐ฟ}+\text{Tr}[๐_0๐บ]+_0^T\text{Tr}[๐_t๐]๐t\hfill \\ +_0^T\text{Tr}[(๐^{}๐+๐)^{}(๐^{}๐+๐)๐บ_t]๐t\hfill \end{array}$$
(32)
where $`๐_0`$ is the solution to (30) at time $`t=0`$.
### 4.3 Duality
The example of the quantum LQG control problem is important since it is one of the few exactly solvable control problems and emphasizes the similarities between the two components of optimal quantum feedback control, namely quantum filtering and optimal control. The duality between the solutions of filtering (24)-(26) and control (28)-(30) is summarized in the duality table
$$\begin{array}{ccccccc}\text{Filtering}& ๐บ_t& \stackrel{~}{๐}_t& ๐& ๐& ๐& ๐\\ & & & & & & \\ \text{Control}& ๐_{Tt}& \stackrel{~}{๐}_{Tt}^{}& ๐^{}& ๐^{}& ๐
& ๐^{}\end{array}$$
(33)
which allows us to formulate and solve the dual control problem given the filtering parameters. The duality can be understood when we examine the nature of each of the methods used. Both methods involve the minimization of a quadratic function for linear, Gaussian systems, (i.e. the least squares error for filtering and the quadratic cost for control). The time reversal in the dual picture is explained by the forward (backward) induction used in the dynamical minimization problem for the filtering (control) problem.
### 4.4 Optimal feedback control of continuously observed quantum free particle
We give a more physical interpretation of the above results by application to an explicit example of LQG control where the duality between filtering and control is preserved. The example of the complex Gaussian oscillator was given in , however we may now use the multidimensional quantum LQG control solutions derived above for application on higher dimensional systems which do not have such complex representation. The optimal control of a continuously observed quantum free particle with quadratic cost is the simplest such example.
Let $`๐ฟ^{}=(Q,P)`$ be the phase space vector operator consisting of the position $`Q`$ and momentum $`P`$ operators of the free particle having the initial expectations $`\overline{Q}`$ and $`\overline{P}`$ respectively. Let us also denote the initial dispersions by $`\sigma _Q`$ and $`\sigma _P`$ respectively and the initial covariance of $`Q`$ and $`P`$ by $`\sigma _{QP}=\sigma _{PQ}`$. We can perform a continuous observation of the particle by coupling the position operator to the measurement channel $`L=Q`$ in which we measure the classical Wiener process $`W_t=A_t+A_t^{}`$ and the particle is controlled using the linear potential $`V(u_t)=u_tQ`$ for $`u_t`$. The Hamiltonian of this simple system is then given by $`H(u_t)=\frac{1}{2M}P^2u_tQ`$ where $`M`$ is the mass of the particle and the corresponding Langevin equations are
$$\frac{d}{dt}Q_t=\frac{1}{M}P_t\frac{d}{dt}P_t=u_t+\dot{V}_t$$
(34)
where $`\dot{V}_t`$ is the time derivative of the Wiener process $`V_t=\mathrm{}W_t^{\pi /2}`$ and represents the system process noise due to the interaction with the coupled noise bath. The operators $`Y_t`$ satisfy the linear output equation
$$\frac{d}{dt}Y_t=2Q_t+\dot{W}_t$$
(35)
which is perturbed by measurement noises represented by the time derivative of the Wiener process $`W_t`$.
The optimal estimates of the position and momentum based on a non demolition measurement of $`Y_t`$ are then given by the quantum Kalman Bucy filter (24)
$`d\widehat{Q}_t`$ $`=`$ $`{\displaystyle \frac{1}{M}}\widehat{P}_tdt+2\sigma _{Q,t}d\stackrel{~}{Y}_t`$ (36)
$`d\widehat{P}_t`$ $`=`$ $`u_tdt+2\sigma _{QP,t}d\stackrel{~}{Y}_t`$ (37)
where the innovation process $`\stackrel{~}{Y}_t`$ describes the gain of information due to measurement of $`Y_t`$ given by
$$\stackrel{~}{Y}_t=Y_t2\widehat{Q}_t.$$
In practice, for a continuous observation, it is the measurement current $`I_t:=dY_t/dt`$ which we observe and so we write the filtering equations in the form
$`{\displaystyle \frac{d}{dt}}\widehat{Q}_t`$ $`=`$ $`{\displaystyle \frac{1}{M}}\widehat{P}_t+2\sigma _{Q,t}(I_t2\widehat{Q}_t)`$ (38)
$`{\displaystyle \frac{d}{dt}}\widehat{P}_t`$ $`=`$ $`u_t+2\sigma _{QP,t}(I_t2\widehat{Q}_t)`$ (39)
where the error covariances satisfy the Ricatti equations
$$\begin{array}{ccc}\frac{d}{dt}\sigma _{Q,t}\hfill & =\hfill & \frac{2}{M}\sigma _{QP,t}4(\sigma _{Q,t})^2\hfill \\ \frac{d}{dt}\sigma _{QP,t}\hfill & =\hfill & \frac{1}{M}\sigma _{P,t}4\sigma _{Q,t}\sigma _{QP,t}\hfill \\ \frac{d}{dt}\sigma _{P,t}\hfill & =\hfill & \mathrm{}^24(\sigma _{QP,t})^2\hfill \end{array}$$
(40)
with initial conditions
$$\sigma _{Q,0}=\sigma _Q,\sigma _{QP,0}=\sigma _{QP},\sigma _{P,0}=\sigma _P.$$
The Ricatti equations for the error covariance in the filtered free particle dynamics have an exact solution , however we will simply comment on the stationary solutions which are the solutions obtained by setting the LHS of (40) to zero, giving the asymptotic behaviour of the posterior dispersions for $`t\mathrm{}`$
$$\sigma _{Q,t}\frac{1}{2}\sqrt{\frac{\mathrm{}}{M}},\sigma _{P,t}\mathrm{}\sqrt{\mathrm{}M},\sigma _{PQ,t}\frac{\mathrm{}}{2}.$$
(41)
This proper treatment dispels the paradoxical quantum Zeno effect which insists that a quantum state is frozen in time by a continuous observation. Instead we can describe the continuous observation as an optimal estimation with posterior dispersions tending to a finite limit satisfying the Heisenberg uncertainty relation
$$\mathrm{\Delta }_{Q,t}\mathrm{\Delta }_{P,t}=\sqrt{\sigma _{Q,t}\sigma _{P,t}}\mathrm{}/\sqrt{2}\mathrm{}/2.$$
In contrast, for the case without conditioning (where the measurement results are ignored or averaged over) the Ricatti equations for the dispersions become linear which have solutions tending to infinity like $`t^3`$. This is faster than the $`t^2`$ spreading of the wavefunction due to the closed evolution described by Schrรถdingerโs equation as one would expect since the coupled noise bath only serves to increase the dispersion.
The dual optimal control problem can be found by identifying the corresponding dual matrices from the table (33) which give the quadratic control parameters
$`C(u_t)`$ $`=`$ $`\beta Q^2+u_t^2`$
$`S`$ $`=`$ $`\omega _QQ^2+\omega _{QP}(PQ+QP)+\omega _PP^2`$
which for the linear Gaussian system (34) gives the optimal control strategy
$$u_t=2(\omega _{PQ,t}\widehat{P}_t+\omega _{P,t}\widehat{Q}_t)$$
(42)
where the coefficients are the solutions to the Ricatti equations
$$\begin{array}{ccc}\frac{d}{dt}\omega _{P,t}\hfill & =\hfill & \frac{2}{M}\omega _{QP,t}4(\omega _{P,t})^2\hfill \\ \frac{d}{dt}\omega _{QP,t}\hfill & =\hfill & \frac{1}{M}\omega _{Q,t}4\omega _{P,t}\omega _{QP,t}\hfill \\ \frac{d}{dt}\omega _{Q,t}\hfill & =\hfill & \beta 4(\omega _{QP,t})^2\hfill \end{array}$$
(43)
with terminal solutions
$$\omega _{P,0}=\omega _P,\omega _{QP,0}=\omega _{QP},\omega _{Q,0}=\omega _Q.$$
Note that in this example, as well as identifying the dual matrices by transposition and time reversal according to the duality table (33), one must also interchange the coordinates $`PQ`$. This is because the matrix of coefficients $`๐`$ is non-symmetric and nilpotent, so it is dual to its transpose only when we interchange the coordinates in the dual picture. Thus the optimal coefficients $`\{\omega _{P,t},\omega _{QP,t},\omega _{Q,t}\}`$ in the quadratic cost-to-go correspond to the minimal error covariances $`\{\sigma _{Q,Tt},\sigma _{QP,Tt},\sigma _{P,Tt}\}`$ in the dual picture.
The minimal total cost for the experiment can be obtained from (32) by substitution of these solutions
$$\begin{array}{c}\text{S}=\omega _{Q,0}(\overline{Q}^2+\sigma _Q)+2\omega _{QP,0}(\overline{Q}\overline{P}+\sigma _{QP})\hfill \\ +\omega _{P,0}(\overline{P}^2+\sigma _P)+_0^T(\mathrm{}^2\omega _{P,t}+\omega _{PQ,t}^2\sigma _{Q,t})๐t\hfill \\ +_0^T(\omega _{P,t}^2\sigma _{P,t}+2\omega _{QP,t}\omega _{P,t}\sigma _{PQ,t})๐t\hfill \end{array}$$
(44)
## 5 Discussion
We have shown that the optimal quantum feedback control problem reduces to an optimal estimation problem followed by an optimal control problem based on this optimal estimator. The optimal (least-squares) estimator for quantum random variables (operators) given a classical nondemolition output measurement process is the conditional expectation which is given by the result of the filtering equation (12). The resulting optimal control problem is then defined on the output of this filter, which reduces to a classical control problem on the space of quantum states. For cost functions that are linear in the state, the optimal feedback control strategy is given by the solution to the Bellman equation (18).
In the LQG example, the space of quantum states are restricted to the class of Gaussian states so the probability distribution is parameterized by the mean and covariance of the generating operators. However, due to non commutativity of these quantum operators there are many different definitions of the covariance matrices. For direct comparison to classical LQG control theory, we choose the symmetric representation of the covariance matrices, although unlike the classical case, the Heisenberg uncertainty principle (23) places a positive lower bound on the covariances. In particular, this prohibits the common classical assumption of uncorrelated process and measurement noise if the coupling to the noise bath is complex.
## Acknowledgements
This work has been supported by the EPSRC under the programme Mathfit (grant no RA2273). VPB also acknowledges support from the EC under the programme ATESIT (contract no IST-2000-29681).
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# Evaluating deterministic policies in two-player iterated games
## 1 Introduction
Game theory has been formalized by Neumann and Morgenstern in 1944, . Their objective was to introduce into the language of economic theory some mathematical tools for the quantitative analysis of the behaviour of economic agents without a central authority. One of the Neumann and Morgenstern arguments in favour of the usefulness of a theory of games is based on the intrisic limited knowledge about the facts which economists deal with. This argument has also been applied to the description of some physical systems and to evolutionary theories in biology.
In the context of evolutionary biology and in order to analyse the logic of animal conflict, Maynard-Smith introduced a game theory approach to describe some of the evolutionary features of organisms. In the framework of sociology, Axelrod gave several examples where the game theoretical framework is useful. In economics, there is today a vast literature on the applicability of the game theoretical approach to economic decision , and . More recently, the same type of formalism has been applied to quantum mechanics .
In a game, a policy is a rule of decision for each player, and policies can be deterministic, depending of the previous choices of one or of both players, or can be stochastic. In a two-player non-cooperative game with a finite number of choices or pure strategies, both players know the payoffs, make their choices independently of each other, and know the past history of their choices. It is also assumed that each player maximizes its payoff after a finite or an infinite number of choices or moves.
An important problem is game theory is to determine which policies perform better than others. In this context, Axelrod, , proposed the following problem: โUnder what conditions will cooperation emerge in a world of egoists without central authority?โ. To help to answer this question a computer tournament has been settled to decide which policy would perform better in an iterated Prisonerโs Dilemma game, introduced by Dresher, Flood and Tucker . The tournament has been won by the tit-for-tat (TFT) policy, submitted by Rapoport, \[3, pp. 31\]. The TFT policy consists in a simple rule that says that oneโs actual move is equal to what the other player did in the previous move.
In fact, several approaches have been developed in order to decide which policies perform better than others in infinitely iterated games. One of these approaches relies on the concept of mixed strategy. In a game with several possible choices or moves, a player has a mixed strategy if he has a probability profile associated to all the possible moves of the game. Based on the concept of mixed strategy, the replicator dynamics approach, and , postulates an evolution equation for the probability profiles of each playerโs move. This evolution equation implies a precise type of rationality of the players, and the mixed strategy concept has a subjacent infinite memory associated to the choices of the players.
The formal construction of game theory depends on the relation between players and from whom they receive their payoffs. For example, we can formalize a two-player game in such a way that the payoffs won by one player are the losses of the other, . Another approach is to consider that the playerโs payoffs are obtained from external sources. Our construction applies to the second case and applies to games describing the global behaviour of systems from economy, sociology and evolutionary biology.
The aim of this paper is to derive the mean payoffs of the โrepresentative playersโ of a game, and to formulate the problem of deciding which policy performs better than another in iterated non-cooperative games.
This paper is organized as follows. In Section 2, we introduce some of the definitions that will be used along this paper, and we analyse and interpret iterated non-cooperative games from the point of view of dynamical system theory. In order to evaluate games and deterministic strategies with finite memory, we take the point of view of uniform ensembles of statistical mechanics and we introduce the concept of representative ensemble of a game. In this context, the players of the infinite set of games are substituted by the โrepresentative agentsโ of the game.
In Section 3, we consider the case of a uniform ensemble of games, where in each subensemble we have two players playing the same game. The mean value of the payoffs per move taken over the uniform ensemble is calculated, and gives information about the performance of a game. In Section 4, we evaluate the performance of deterministic strategies with finite memory length. In the case where in each subensemble a player has a deterministic strategy and the other makes his choices with equal probabilities, we calculate the ensemble averages of the payoffs per move. The main results of sections 3 and 4 are summarized in Theorems 3.1 and 4.1, and Corollary 4.2. In particular, we show that if one of the players has a generalized tit-for-tat policy, the mean payoff per move of both players is the same. Therefore, generalized tit-for-tat is the best policy against exploitation.
In Section 5, we consider the case where the opponent players have deterministic policies, within the same memory class. In this case, the game dynamics is a deterministic process, and the mean payoffs per move depend on the initial moves of both players and on the policy functions of both players. Comparing all the possible deterministic strategies with memory length $`1`$, we prove that, in dilemmatic games, the generalized tit-for-tat policy together with the condition of not being the first to defect, leads to the highest possible mean payoffs per move for the players.
In Section 6, we apply the formalism developed in this paper to the Prisonerโs Dilemma and to the Hawk-Dove games, and we analyse their state space structure. Finally, in Section 7, we summarize the main conclusions of the paper.
## 2 Formalism and definitions
We take a two-player game with two possible choices or moves โ two pure strategies. At times $`n1`$, each player chooses, independently of the other, one of the two possible pure strategies. These pure strategies are represented by the symbols โ0โ and โ1โ. We denote by $`๐=\{0,1\}`$ the set of pure strategies, and by $`P`$ and $`Q`$ the two players. After a move, each player owns a profit or payoff that is dependent of the opponent move. The payoff matrices of the game are:
$$A=\left(\begin{array}{cc}A_{00}& A_{01}\\ A_{10}& A_{11}\end{array}\right),B=\left(\begin{array}{cc}B_{00}& B_{01}\\ B_{10}& B_{11}\end{array}\right)$$
where the payoff of $`P`$ is $`A_{ij}`$ if player $`P`$ plays $`i`$ and $`Q`$ plays $`j`$. In the same move, $`Q`$ has payoff $`B_{ij}`$. If each player makes its choice independently of the other, we are in the context of non-cooperative games. If $`B=A^T`$, the two-player game is symmetric. In the following, we analyze only the case of symmetric and non-cooperative games.
In a two-player symmetric game, we say that a pure strategy $`i๐`$ is dominant, , if $`A_{ij}A_{kj}`$, for every $`j=0,1`$, and $`A_{ij}>A_{kj}`$ for some $`j`$ and $`ki`$. For example, the symmetric and non-cooperative games with payoff matrices,
$$A_1=\left(\begin{array}{cc}3& 0\\ 5& 1\end{array}\right),A_2=\left(\begin{array}{cc}1& 0\\ 5& 3\end{array}\right)$$
have โ1โ as dominant strategy ($`A_{10}>A_{00}`$ and $`A_{11}>A_{01}`$). In the first game, if the two players choose both the dominant strategy โ1โ, their payoffs is $`1`$. In the second game, the payoff of each players is $`3`$, and the dominant strategy is the right choice for both players. However, as $`A_{00}>A_{11}`$ for the first game, if both players choose the non-dominant strategy, their individual payoffs per move is higher when compared with the choice of the dominant strategy by both players.
These two examples suggest the following definition: A symmetric two-player game is dilemmatic, if either,
$$A_{10}>A_{00}>A_{11}>A_{01}\text{and}2A_{00}>A_{01}+A_{10}$$
$`(2.1)`$
or,
$$A_{01}>A_{11}>A_{00}>A_{10}\text{and}2A_{11}>A_{10}+A_{01}$$
$`(2.2)`$
where the second inequalities in (2.1) and (2.2) have been introduced in order to favour the non-dominant strategy.
In the first case of a dilemmatic game, (2.1), the strategy โ1โ is dominant. In the second case, (2.2), โ0โ is the dominant strategy. If both players choose the dominant strategy in one move, they get smaller payoffs than the ones they could have obtained if both had chosen the non-dominant strategy.
In an iterated game with a fixed payoff matrix, players are always playing the same game, and their payoffs accumulate. Therefore, a two-player iterated game is described by the two sequences of pure strategies of each player,
$$\begin{array}{c}\mu =(\mu _1,\mu _2,\mathrm{},\mu _n,\mathrm{})\hfill \\ \sigma =(\sigma _1,\sigma _2,\mathrm{},\sigma _n,\mathrm{})\hfill \end{array}$$
$`(2.3)`$
where $`\mu _n`$ and $`\sigma _n`$ represent the choices of the players $`P`$ and $`Q`$, respectively, at discrete time $`n1`$, and $`\mu _n,\sigma _n๐`$. The sequences (2.3), completely specify the accumulated payoffs of both players. We call $`\mu `$ and $`\sigma `$ the game record sequences. In an infinitely iterated game, the accumulated payoff of the players can be infinite. The mean payoffs per move are always finite and, for a symmetric game, they are given by,
$$\begin{array}{c}G_p=lim_n\mathrm{}\frac{1}{n}_{i=1}^nA_{\mu _i\sigma _i}\hfill \\ G_q=lim_n\mathrm{}\frac{1}{n}_{i=1}^nA_{\sigma _i\mu _i}\hfill \end{array}$$
$`(2.4)`$
where $`G_p`$ and $`G_q`$ are the mean payoffs of players $`P`$ and $`Q`$, respectively.
An example of a symmetric, non-cooperative, and dilemmatic game is the Prisonerโs Dilemma game. In this game, we have two players with two possible pure strategies, โ0โ and โ1โ, and we have chosen the payoff matrix,
$$A=\left(\begin{array}{cc}3& 9\\ 11& 5\end{array}\right)$$
$`(2.5)`$
As, $`A_{10}>A_{00}>A_{11}>A_{01}`$ and $`2A_{00}>A_{01}+A_{10}`$, the Prisonerโs Dilemma game is dilemmatic with โ1โ as dominant pure strategy. The pure strategy โ0โ corresponds to cooperation and the pure strategy โ1โ to defection. For a discussion about the importance of dilemmatic games and the Prisonerโs Dilemma game, see the discussion in Axelrod, .
Following Neumann and Morgenstern , a strategy or policy is a set of rules that tells each participant how to behave in every situation which may arise. The only sources of information available to the players is the set of all possible moves, their possible payoffs, and the history of the previous moves of both players. To describe a rule of decision, policy, or strategy we can adopt the Neumann-Morgenstern view where a rule of decision is specified through the knowledge of a function of the $`m`$ previous moves.
Deterministic strategy: In an iterated two-player game, with game records $`\mu `$ and $`\sigma `$ for players $`P`$ and $`Q`$, respectively, a rule of decision or a deterministic strategy with memory length $`m1`$ for player $`P`$ is a function $`f:๐^m๐`$ such that,
$$\mu _i=f(\sigma _{im},\mathrm{},\sigma _{i1})$$
$`(2.6)`$
for every $`i>m`$. Analogously, player $`Q`$ has a deterministic strategy with memory length $`n1`$, if there exists a function $`g:๐^n๐`$ such that,
$$\sigma _i=g(\mu _{in},\mathrm{},\mu _{i1})$$
for every $`i>n`$.
In the following, deterministic policy and deterministic strategy have the same meaning. In some game theory texts, the word โstrategyโ refers to โpure strategyโ, an element of the set $`๐=\{0,1\}`$, and in other contexts if refers to policies, as in โtit-for-tat strategyโ.
By definition, the outcome of a playerโs choice or move at time $`im+1`$ is determined by a finite number of previous moves of the other. In general, we can take the functions $`f:๐^{m+n}๐`$ and $`g:๐^{r+s}๐`$, and set,
$$\begin{array}{c}\mu _i=f(\sigma _{im},\mathrm{},\sigma _{i1},\mu _{in},\mathrm{},\mu _{i1})\hfill \\ \sigma _i=g(\mu _{ir},\mathrm{},\mu _{i1},\sigma _{is},\mathrm{},\sigma _{i1})\hfill \end{array}$$
In the following we will only analyze the case where each playerโs choice depends on a finite number of previous moves of the other, and $`m=n`$.
For example, adopting the definition of the tit-for-tat (TFT) strategy given in the introduction, a TFT strategy with memory length $`m=1`$ is described by the boolean identity function $`f:๐๐`$, defined by,
$$f(0)=0\text{and}f(1)=1$$
Generalized tit-for-tat strategy (GTFT): We say that $`f:๐^m๐`$ is a generalized tit-for-tat strategy with memory length $`m`$, if the number of โ0โ and โ1โ in $`f(\sigma ^{(m)})`$, when $`\sigma ^{(m)}`$ runs over the set $`๐^m`$, are equal. More formally, $`f:๐^m๐`$ is a generalized tit-for-tat strategy with memory length $`m`$, if,
$$\mathrm{\#}\{\sigma ^{(m)}๐^m:f(\sigma ^{(m)})=0\}=\mathrm{\#}\{\sigma ^{(m)}๐^m:f(\sigma ^{(m)})=1\}$$
where $`\sigma ^{(m)}=(\sigma _1,\mathrm{},\sigma _m)`$, and $`\sigma _i๐`$.
To solve a game it is meant to find a procedure to determine for each playerโs choice which is the most favourable result (, and ). In this context, the concepts of mixed strategy and equilibrium state of a game are fundamental tools in game theory.
A mixed strategy is a collection of probabilities associated to each player and its pure strategies. The players $`P`$ and $`Q`$ have mixed strategies $`s_p=(s_{0p},s_{1p})`$ and $`s_q=(s_{0q},s_{1q})`$, if each player plays pure strategy $`i`$ with probability $`s_i`$. Obviously, $`s_0+s_1=1`$.
In a symmetric game with two pure strategies and mixed strategies $`s_p`$ and $`s_q`$ for players $`P`$ and $`Q`$, respectively, the mean payoffs per move of players $`P`$ and $`Q`$ are,
$$\begin{array}{cc}P:E(s_p|s_q)\hfill & =s_{0p}(s_{0q}A_{00}+s_{1q}A_{01})+s_{1p}(s_{0q}A_{10}+s_{1q}A_{11})\hfill \\ & =s_{0p}s_{0q}A_{00}+s_{0p}s_{1q}A_{01}+s_{1p}s_{0q}A_{10}+s_{1p}s_{1q}A_{11}\hfill \\ Q:E(s_q|s_p)\hfill & =s_{0q}(s_{0p}A_{00}+s_{1p}A_{01})+s_{1q}(s_{0p}A_{10}+s_{1p}A_{11})\hfill \\ & =s_{0q}s_{0p}A_{00}+s_{0q}s_{1p}A_{01}+s_{1q}s_{0p}A_{10}+s_{1q}s_{1p}A_{11}\hfill \end{array}$$
$`(2.7)`$
The time evolution of a game with mixed strategies $`s_p`$ and $`s_q`$ can be seen as a stochastic processes with two independent random variables $`X`$ and $`Y`$. The random variables $`X`$ and $`Y`$, associated to players $`P`$ and $`Q`$, respectively, have mean values given by (2.7). More precisely, $`X`$ can assume the values $`A_{00}`$, $`A_{01}`$, $`A_{10}`$, $`A_{11}`$ with probabilities $`s_{0p}s_{0q}`$, $`s_{0p}s_{1q}`$, $`s_{1p}s_{0q}`$ and $`s_{1p}s_{1q}`$, respectively. Analogously, $`Y`$ takes values in the same set, with probabilities: $`s_{0q}s_{0p}`$, $`s_{0q}s_{1p}`$, $`s_{1q}s_{0p}`$ and $`s_{1q}s_{1p}`$. Therefore, the deviations from the mean payoffs per move of the players, or the fluctuations from the mean values, are characterised by the variances,
$$\begin{array}{cc}\sigma _p^2(s_p|s_q)\hfill & =s_{0p}s_{0q}(A_{00}E(s_p|s_q))^2+s_{0p}s_{1q}(A_{01}E(s_p|s_q))^2\hfill \\ & +s_{1p}s_{0q}(A_{10}E(s_p|s_q))^2+s_{1p}s_{1q}(A_{11}E(s_p|s_q))^2\hfill \\ \sigma _q^2(s_q|s_p)\hfill & =s_{0q}s_{0p}(A_{00}E(s_q|s_p))^2+s_{0q}s_{1p}(A_{01}E(s_q|s_p))^2\hfill \\ & +s_{1q}s_{0p}(A_{10}E(s_q|s_p))^2+s_{1q}s_{1p}(A_{11}E(s_q|s_p))^2\hfill \end{array}$$
$`(2.8)`$
In general, $`E(s_p|s_q)E(s_q|s_p)`$, but, by a straightforward calculation, $`\sigma _p^2(s_p|s_q)=\sigma _q^2(s_q|s_p)`$.
Imposing the condition $`E(s_p|s_q)=E(s_q|s_p)`$ in (2.7), a game or a mixed strategy is equalitarian, if either, $`A_{01}=A_{10}`$, or $`s_p=s_q`$.
To characterise the dynamics of an iterated symmetric game, we introduce the concept of phase or state space of a game. The state space of a two-player game is the convex closure of the points $`(A_{00},A_{00})`$, $`(A_{11},A_{11})`$, $`(A_{01},A_{10})`$ and $`(A_{10},A_{01})`$, in the two-dimensional space of the payoffs of players $`P`$ and $`Q`$. Let us denote by $`๐ฆ`$ the state space of a game. As $`s_{0p},s_{0q}[0,1]`$, then $`(E(s_p|s_q),E(s_q|s_p))๐ฆ`$. In Fig. 1, we show the state space $`๐ฆ`$ for the Prisonerโs Dilemma game with payoff matrix (2.5).
In an iterated game, the mean initial (at time $`n=1`$) payoffs per move of both players is $`(x_1,y_1)=(A_{\mu _1\sigma _1},A_{\sigma _1\mu _1})๐ฆ`$. By (2.4) and after $`n+1`$ moves, the mean payoffs per move of both players is,
$$\begin{array}{cc}(x_{n+1},y_{n+1})\hfill & =(\frac{1}{n+1}_{i=1}^{n+1}A_{\mu _i\sigma _i},\frac{1}{n+1}_{i=1}^{n+1}A_{\sigma _i\mu _i})\hfill \\ & =(\frac{n}{n+1}x_n+\frac{1}{n+1}A_{\mu _{n+1}\sigma _{n+1}},\frac{n}{n+1}y_n+\frac{1}{n+1}A_{\sigma _{n+1}\mu _{n+1}})๐ฆ\hfill \end{array}$$
$`(2.9)`$
and the iterated two-player game is dynamically described by a (non-deterministic) one-to-many map, $`\beta :๐ฆ๐ฆ`$, . The equilibrium point or equilibrium solution of a game is the point $`lim_n\mathrm{}(x_n,y_n)`$.
For a given mixed strategy profile $`s_p`$ and $`s_q`$ of players $`P`$ and $`Q`$, the iterated two-player game has the equilibrium point, $`(E(s_p|s_q),E(s_q|s_p))๐ฆ`$. As $`s_{0p},s_{0q}[0,1]`$, the set of equilibrium states of the map $`\beta :๐ฆ๐ฆ`$ span all the state space of a game.
For example, in a two-player game with the mixed strategy profiles $`s_p=(1/2,1/2)`$ and $`s_q=(1/2,1/2)`$, by (2.7), the equilibrium point of the game is,
$$(\frac{1}{4}(A_{00}+A_{01}+A_{10}+A_{11}),\frac{1}{4}(A_{00}+A_{01}+A_{10}+A_{11}))$$
$`(2.10)`$
By (2.8), the fluctuations around the equilibrium are,
$$\begin{array}{cc}& \sigma _p^2((1/2,1/2)|(1/2,1/2))=\sigma _q^2((1/2,1/2)|(1/2,1/2))\hfill \\ & =\frac{1}{4^2}(3A_{00}^2+3A_{01}^2+3A_{10}^2+3A_{11}^2\hfill \\ & 2A_{00}(A_{01}+A_{10}+A_{11})2A_{11}(A_{01}+A_{10})2A_{01}A_{10})\hfill \end{array}$$
$`(2.11)`$
In Fig. 1, we represent several iterates of the map $`\beta `$ for the Prisonerโs Dilemma game with payoff matrix (2.5), and mixed strategies $`s_p=s_p=(1/2,1/2)`$. In the limit $`n\mathrm{}`$, $`(x_n,y_n)(0,0)`$. The fluctuations from equilibrium have standard deviations $`\sigma _p=\sigma _q=\sqrt{59}=7.68`$.
A mixed strategy is a strict Nash equilibrium solution of a game if $`P`$ and $`Q`$ maximizes their payoffs per move independently of each other. A mixed strategy is a Nash bargain equilibrium solution of a game if $`E(s_p|s_p)`$ is a maximum.
In the case of the of the Prisonerโs Dilemma game with payoff matrix (2.5), by (2.7), we have,
$$\begin{array}{cc}P\hfill & :E(s_p|s_q)=54s_{0p}+16s_{0q}4s_{0p}s_{0q}\hfill \\ Q\hfill & :E(s_q|s_p)=54s_{0q}+16s_{0p}4s_{0p}s_{0q}\hfill \end{array}$$
$`(2.12)`$
Maximizing $`E(s_p|s_q)`$ in order to $`s_{0p}`$ and $`E(s_q|s_p)`$ in order to $`s_{0q}`$, the strict Nash equilibrium of the game is obtained when both player choose the mixed strategy $`(s_{0p},s_{1p})=(s_{0p},s_{1p})=(0,1)`$. In this case, the Nash equilibrium state of the map $`\beta :๐ฆ๐ฆ`$ is the point $`(5,5)๐ฆ`$. The Nash bargain solution of the game is obtained from (2.12) with $`s_p=s_q`$, and is the point $`(3,3)๐ฆ`$, Fig. 1. As both Nash solutions correspond to the choices of pure strategies with probability 1, by (2.8), the fluctuations of the iterated game have zero standard deviations. In general, a n-person non-cooperative game has always a strict Nash equilibrium, .
The choice of a mixed strategy profile for a game has the advantage that the iterates of the map $`\beta :๐ฆ๐ฆ`$ converges to the equilibrium solution $`E(s_p|s_q)`$. However, the choice of a mixed strategy profile implies that both players have infinite memory, which, in real situations, is difficult or even impossible to fulfil.
On the other hand, in some game theory approaches describing the global behaviour of economic, social and evolutionary systems, there are a large number of agents or players in mutual interaction. These individual agents interact with the same rules and can also change partners along time. These situations are difficult to interpret under the infinite memory hypothesis, implicitly associated to the concept of mixed strategies.
Following this point of view, to evaluate a non-cooperative and symmetric game and their possible deterministic strategies (short memory), we adopt the point of view of the statistical ensembles of statistical mechanics. We suppose first that we have an infinite system composed by independent subensembles, where in each subensemble we have two players playing the same game with payoff matrix $`A`$. We call this ensemble of independent games the uniform ensemble (\[15, pp. 56\]) of the game. This uniform ensemble is characterized by the payoff matrix $`A`$, and the players $`P`$ and $`Q`$ are the representative agents of the ensemble of the game.
In each subensemble, a game with payoff matrix $`A`$ is played, and subensembles are characterized by the mean payoffs per move of both players. The global properties of the game will be described by the mean payoffs per move averaged over all the subensembles. We say that the representative players $`P`$ and $`Q`$ of the game have mean payoffs per move $`\overline{G}_p`$ and $`\overline{G}_q`$, respectively, where the average is taken over all the subensembles.
To evaluate a game, we first consider that each player chooses its pure strategies with equal probabilities, and each subensemble is characterized by the two sequences of pure strategies $`\mu `$ and $`\sigma `$. The properties of the game are determined by $`\overline{G}_p`$ and $`\overline{G}_q`$.
To evaluate a deterministic strategy or policy, we consider that in each subensemble game, $`P`$ plays with the deterministic policy $`f`$, and $`Q`$ has a game record $`\sigma ๐^{\text{}}`$. Defining an ensemble probability density function $`\rho _q`$ for the occurrence of game record $`\sigma `$ for player $`Q`$, the ensemble of games will by characterized by the mean payoff per move of both players averaged over the set of all allowed sequences $`\sigma ๐^{\text{}}`$ with probability measure $`\rho _q`$. These averages depend on $`f`$ and $`\rho _q`$, and we can compare the performance of a policy with the case where the players have no policies. Within the same memory class, we use these ensemble averages to compare the mean payoffs for different policies.
When both players have a deterministic policy, the mean payoffs per move of the players depend on the finite number of initial conditions of the game.
## 3 The uniform ensemble of a game
We consider an ensemble of subsystems, where in each subsystem there are two players playing the same game. We denote the game record of players $`P`$ and $`Q`$ by $`\mu `$ and $`\sigma `$, respectively. As $`\mu `$ and $`\sigma `$ are infinite sequences of โ0โ and โ1โ, we can identify $`\mu `$ and $`\sigma `$ as real numbers in the interval $`[0,1]`$ through,
$$x=\underset{i=1}{\overset{\mathrm{}}{}}\frac{\mu _i}{2^i},y=\underset{i=1}{\overset{\mathrm{}}{}}\frac{\sigma _i}{2^i}$$
$`(3.1)`$
Relations (3.1) define a map $`\varphi :๐^{\text{}}[0,1]`$. The map $`\varphi `$ is an isomorphism, except when $`(\sigma _1,\sigma _2,\mathrm{})`$ or $`(\mu _1,\mu _2,\mathrm{})`$ represents dyadic rational numbers, . As the set of dyadic rationals has zero Lebesgue measure, the infinite sequence of moves of both players can be represented, almost everywhere, by two real numbers $`x,y[0,1]`$. Therefore, the interval $`[0,1]`$ is naturally the space of game records.
Making this identification between game records and real numbers, in each subensemble game, by (2.4), the mean payoffs per move of the players are,
$$\begin{array}{cc}G_p(\mu ,\sigma )\hfill & =lim_n\mathrm{}\frac{1}{n}_{i=1}^nA_{\mu _i\sigma _i}:=G_p(x,y)\hfill \\ G_q(\mu ,\sigma )\hfill & =lim_n\mathrm{}\frac{1}{n}_{i=1}^nA_{\sigma _i\mu _i}:=G_q(x,y)\hfill \end{array}$$
$`(3.2)`$
where $`x,y[0,1]`$.
As each subensemble game is independent of the other, and each playerโs move is independent of the history of the game, we can assign ensemble probability density functions to the game records. Let $`\rho _p(x)`$ and $`\rho _q(y)`$ be the ensemble probability density functions of game records of the representative players $`P`$ and $`Q`$, respectively. For example, $`\rho _p(x)dx`$ is the probability of finding a subensemble with player $`P`$ with a game record in an interval of length $`dx`$ centred around $`x`$.
Assuming further that all the game records are equally probable, $`\rho _p(x)=1`$ and $`\rho _q(y)=1`$, the ensemble averages of the mean payoffs per move are,
$$\begin{array}{cc}\overline{G}_p\hfill & =_0^1_0^1G_p(x,y)\rho _p(x)\rho _q(y)๐x๐y=_0^1_0^1G_p(x,y)๐x๐y\hfill \\ \overline{G}_q\hfill & =_0^1_0^1G_q(x,y)\rho _p(x)\rho _q(y)๐x๐y=_0^1_0^1G_q(x,y)๐x๐y\hfill \end{array}$$
$`(3.3)`$
To characterize the statistical ensemble of a non-cooperative and symmetric game with payoff matrix $`A`$, we now calculate the integrals in (3.3). We consider the sequences of functions,
$$\begin{array}{cc}G_p^n(\mu ,\sigma )\hfill & =\frac{1}{n}_{i=1}^nA_{\mu _i\sigma _i}\hfill \\ G_q^n(\mu ,\sigma )\hfill & =\frac{1}{n}_{i=1}^nA_{\sigma _i\mu _i}\hfill \end{array}$$
$`(3.4)`$
As $`n\mathrm{}`$, $`G_p^n(\mu ,\sigma )G_p(\mu ,\sigma )`$, and $`G_q^n(\mu ,\sigma )G_q(\mu ,\sigma )`$. In the sense of Lebesgue integration, the integrals in (3.3) can be calculated as the limits of the integrals of the functions $`G_p^n(\mu ,\sigma )`$ and $`G_q^n(\mu ,\sigma )`$.
Let us first take $`n=1`$. By (3.2), (3.3) and (3.4), we have,
$$\begin{array}{cc}\overline{G}_p^1\hfill & =_0^1_0^1A_{\mu _1\sigma _1}๐x๐y\hfill \\ \overline{G}_q^1\hfill & =_0^1_0^1A_{\sigma _1\mu _1}๐x๐y\hfill \end{array}$$
$`(3.5)`$
where $`\mu _1`$ and $`\sigma _1`$ are the first digits in the binary developments of $`x`$ and $`y`$, both in the interval $`[0,1]`$. Therefore, the functions $`A_{\mu _1\sigma _1}A_{\mu _1\sigma _1}(x,y)`$ and $`A_{\sigma _1\mu _1}A_{\sigma _1\mu _1}(x,y)`$ are piecewise constant in the unit square, and the integrals in (3.5) are straightforwardly evaluated to,
$$\begin{array}{cc}\overline{G}_p^1\hfill & =_0^1_0^1A_{\mu _1\sigma _1}๐x๐y=\frac{1}{2^2}(A_{00}+A_{01}+A_{10}+A_{11})\hfill \\ \overline{G}_q^1\hfill & =_0^1_0^1A_{\sigma _1\mu _1}๐x๐y=\frac{1}{2^2}(A_{00}+A_{01}+A_{10}+A_{11})\hfill \end{array}$$
$`(3.6)`$
Note that, the functions $`A_{\mu _1\sigma _1}(x,y)`$ and $`A_{\sigma _1\mu _1}(x,y)`$ are piecewise constant functions from $`[0,1]\times [0,1]`$ to the set $`\{A_{00},A_{01},A_{10},A_{11}\}`$.
In general, by (3.4) and (3.3),
$$\begin{array}{cc}\overline{G}_p^{n+1}\hfill & =\frac{n}{n+1}\overline{G}_p^n+\frac{1}{n+1}_0^1_0^1A_{\mu _{n+1}\sigma _{n+1}}๐x๐y\hfill \\ \overline{G}_q^{n+1}\hfill & =\frac{n}{n+1}\overline{G}_q^n+\frac{1}{n+1}_0^1_0^1A_{\sigma _{n+1}\mu _{n+1}}๐x๐y\hfill \end{array}$$
$`(3.7)`$
The functions $`A_{\mu _n\sigma _n}A_{\mu _n\sigma _n}(x,y)`$ and $`A_{\sigma _n\mu _n}A_{\sigma _n\mu _n}(x,y)`$ are piecewise constant and assume the constant values $`A_{00}`$, $`A_{01}`$, $`A_{10}`$ and $`A_{11}`$ in squares of side $`1/2^n`$. As, for each pair of indices $`(\sigma _n,\mu _n)`$, the domain where $`A_{\sigma _n\mu _n}(x,y)`$ is constant is composed by $`2^{2(n1)}`$ disjoint squares, we have,
$$\begin{array}{cc}_0^1_0^1A_{\mu _n\sigma _n}๐x๐y\hfill & =_0^1_0^1A_{\sigma _n\mu _n}๐x๐y\hfill \\ & =\frac{2^{2(n1)}}{2^{2n}}(A_{00}+A_{01}+A_{10}+A_{11})\hfill \\ & =\frac{1}{2^2}(A_{00}+A_{01}+A_{10}+A_{11})\hfill \end{array}$$
$`(3.8)`$
Introducing (3.8) into (3.7), by induction, and taking the limit $`n\mathrm{}`$, we obtain the values of the ensemble average of the mean payoffs per move of each player:
###### Theorem 3.1.
We consider an ensemble of non-cooperative and symmetric two-player game, where in each subensemble we have two players making their choices with equal probabilities. Assume that each playerโs move is independent of the history of the game and that the ensemble probability density functions of each representative player are uniform in the interval $`[0,1]`$ of the game records. Then, the mean payoffs per move of the representative players of the game are equal and are given by,
$$\overline{G}_p=\overline{G}_q=\frac{1}{2^2}(A_{00}+A_{01}+A_{10}+A_{11})$$
where the $`A_{ij}`$ are the entries of the payoff matrix.
In the uniform statistical ensemble of a non-cooperative and symmetric game with all the players choosing their pure strategies with equal probabilities, the average payoff per move is equal to the average value of the entries of the payoff matrix $`A`$.
These elementary results can be straightforwardly generalized to non-cooperative and non-symmetric $`n`$-player games.
## 4 Evaluating deterministic strategies
To evaluate the performance of a deterministic strategy in an iterated game, we first enumerate the class of all the boolean functions $`f:๐^m๐`$, where $`๐=\{0,1\}`$. These boolean functions describe all the possible deterministic strategies.
For each class of functions with memory length $`m`$, there are exactly $`2^{2^m}`$ different functions. To enumerate a deterministic policy function within a memory class $`m`$, $`f(\sigma ^{(m)})`$, where $`\sigma ^{(m)}=(\sigma _1,\mathrm{},\sigma _m)\{0,1\}^m`$, we introduce an additional index $`n`$. Within each memory class $`m`$, each possible policy function will be denoted by $`f_{m,n}`$, where $`n=_{i=0}^{2^m1}f_{m,n}(\sigma _i^{(m)})2^i`$ is the policy number, $`\sigma _{i+1}^{(m)}=\sigma _i^{(m)}+(0,0,\mathrm{},1)`$, $`\sigma _0^{(m)}=(0,0,\mathrm{},0)`$, and the โplusโ symbols must be understood in the sense of binary arithmetic. For example, in Table 1, we show all the possible deterministic policy functions with memory length $`m=1`$.
In this case, the deterministic TFT policy corresponds to the boolean function $`f_{1,2}`$. The functions $`f_{1,2}`$ and $`f_{1,1}`$ are GTFT policies with memory length $`m=1`$.
Suppose now that the representative player $`P`$ of a game has a policy $`f_{m,n}`$ and the opponent player $`Q`$ can have any game sequence $`\sigma =(\sigma _1,\sigma _2,\mathrm{})`$. Then, by (2.4), for the infinitely iterated game, the mean payoff per move for each player is,
$$\begin{array}{cc}G_p(\sigma )\hfill & =lim_M\mathrm{}\frac{1}{M}_{i=1}^MA_{\mu _i\sigma _i}\hfill \\ & =lim_M\mathrm{}\frac{1}{Mm}_{i=m+1}^MA_{f_{m,n}(\sigma _{im},\mathrm{},\sigma _{i1})\sigma _i}\hfill \\ G_q(\sigma )\hfill & =lim_M\mathrm{}\frac{1}{Mm}_{i=m+1}^MA_{\sigma _if_{m,n}(\sigma _{im},\mathrm{},\sigma _{i1})}\hfill \end{array}$$
$`(4.1)`$
and $`G_p`$ and $`G_q`$ are functions of $`\sigma =(\sigma _1,\sigma _2,\mathrm{})`$ and $`f_{m,n}`$. In the first $`m`$ iterations of the game, the accumulated mean payoffs depend on the initial choices of the players. However, in the limit $`M\mathrm{}`$, the dependence on the initial choices vanishes.
Let us take the infinite sequence $`(\sigma _1,\sigma _2,\mathrm{})๐^{\text{}}`$ characterizing one of the possible outcomes of the choices of the player $`Q`$, and define the real number,
$$y=\underset{i=1}{\overset{\mathrm{}}{}}\frac{\sigma _i}{2^i}$$
$`(4.2)`$
With this identification between infinite sequences of zeros and ones with real numbers in the interval $`[0,1]`$, we write the mean payoffs as,
$$\begin{array}{c}G_p(\sigma )G_p(y;f_{m,n}):=P_{m,n}(y)\hfill \\ G_q(\sigma )G_q(y;f_{m,n}):=Q_{m,n}(y)\hfill \end{array}$$
$`(4.3)`$
Let us suppose now that we are in framework of statistical mechanics and we have an ensemble or collectivity of players $`P`$ and $`Q`$. In each subensemble of the collectivity, the player $`P`$ plays according to strategy $`f_{m,n}`$ and $`Q`$ has some game record $`y[0,1]`$. Suppose additionally that all the subensembles of the collectivity are independent.
As each member of the collectivity is independent of the others, we can assign an ensemble density function $`\rho _q(y)`$ to the collectivity. The function $`\rho _q(y)`$ is the probability density of the game record $`y`$ of player $`Q`$. If $`\rho _q(y)=1`$, all the game records of $`Q`$ are equally probable. The uniform ensemble of the game can then be characterized by the ensemble mean payoffs per move and per player,
$$\begin{array}{cc}\overline{P}_{m,n}\hfill & =_0^1P_{m,n}(y)\rho _q(y)๐y=_0^1P_{m,n}(y)๐y\hfill \\ \overline{Q}_{m,n}\hfill & =_0^1Q_{m,n}(y)\rho _q(y)๐y=_0^1Q_{m,n}(y)๐y\hfill \end{array}$$
$`(4.4)`$
In Fig. 2, we show the mean payoff functions $`P_{m,n}(y)`$ and $`Q_{m,n}(y)`$ for the TFT policy $`f_{1,2}`$ and payoff matrix (2.5) of the Prisonerโs Dilemma game. These functions have been calculated numerically from (4.1), (4.2) and (4.3).
To calculate the mean payoffs $`\overline{P}_{m,n}`$ and $`\overline{Q}_{m,n}`$ given by (4.4), we first approximate the functions $`P_{m,n}(y)`$ and $`Q_{m,n}(y)`$ by sequences of piecewise constant functions. By (4.1)-(4.3), we define the sequences of functions, $`\{P_{m,n}^M(y)\}_{Mm+1}`$ and $`\{Q_{m,n}^M(y)\}_{Mm+1}`$, as,
$$\begin{array}{cc}P_{m,n}^M(y)\hfill & =\frac{1}{Mm}_{i=m+1}^MA_{f_{m,n}(\sigma _{im},\mathrm{},\sigma _{i1})\sigma _i}\hfill \\ Q_{m,n}^M(y)\hfill & =\frac{1}{Mm}_{i=m+1}^MA_{\sigma _if_{m,n}(\sigma _{im},\mathrm{},\sigma _{i1})}\hfill \end{array}$$
$`(4.5)`$
where $`(\sigma _1,\sigma _2,\mathrm{})`$ is the binary development of $`y`$. For $`Mm+1`$, $`P_{m,n}^M(y)`$ and $`Q_{m,n}^M(y)`$ are piecewise constant functions in the interval $`[0,1]`$, and, in the limit $`M\mathrm{}`$, they converge almost everywhere to $`P_{m,n}(y)`$ an $`Q_{m,n}(y)`$, Fig. 2. In the sense of Lebesgue integrations, this implies that,
$$\begin{array}{c}lim_M\mathrm{}_0^1P_{m,n}^M(y)๐y=\overline{P}_{m,n}\hfill \\ lim_M\mathrm{}_0^1Q_{m,n}^M(y)๐y=\overline{Q}_{m,n}\hfill \end{array}$$
$`(4.6)`$
Let us now calculate the integrals in (4.6). For $`M=m+1`$, by (4.5), we have,
$$\begin{array}{cc}P_{m,n}^{m+1}(y)\hfill & =A_{f_{m,n}(\sigma _1,\mathrm{},\sigma _m)\sigma _{m+1}}\hfill \\ Q_{m,n}^{m+1}(y)\hfill & =A_{\sigma _{m+1}f_{m,n}(\sigma _1,\mathrm{},\sigma _m)}\hfill \end{array}$$
$`(4.7)`$
As $`(\sigma _1,\mathrm{},\sigma _{m+1})`$ represents the first $`m+1`$ terms of the binary development of $`y`$, the functions in (4.7) assume constant values in subintervals of $`[0,1]`$ of length $`1/2^{m+1}`$. In each of these subintervals, $`P_{m,n}^{m+1}(y)`$ and $`Q_{m,n}^{m+1}(y)`$ assume one of the four values: $`A_{00}`$, $`A_{01}`$, $`A_{10}`$ and $`A_{11}`$.
Associated to each deterministic policy function $`f_{m,n}`$, we define the numbers,
$$\begin{array}{cc}n_0^{m,n}\hfill & =\mathrm{\#}\{\sigma ^{(m)}\{0,1\}^m:f_{m,n}(\sigma ^{(m)})=0\}\hfill \\ n_1^{m,n}\hfill & =\mathrm{\#}\{\sigma ^{(m)}\{0,1\}^m:f_{m,n}(\sigma ^{(m)})=1\}\hfill \end{array}$$
$`(4.8)`$
and, $`n_0^{m,n}+n_1^{m,n}=2^m`$.
Under these conditions, by (4.7) and (4.8), we have,
$$\begin{array}{cc}_0^1P_{m,n}^{m+1}(y)๐y\hfill & =\frac{1}{2^m}\left(n_0^{m,n}\frac{(A_{00}+A_{01})}{2}+n_1^{m,n}\frac{(A_{10}+A_{11})}{2}\right)\hfill \\ & =\frac{1}{2^{m+1}}\left(n_0^{m,n}(A_{00}+A_{01})+n_1^{m,n}(A_{10}+A_{11})\right)\hfill \\ _0^1Q_{m,n}^{m+1}(y)๐y\hfill & =\frac{1}{2^{m+1}}\left(n_0^{m,n}(A_{00}+A_{10})+n_1^{m,n}(A_{01}+A_{11})\right)\hfill \end{array}$$
$`(4.9)`$
For $`M>m+1`$, by (4.5), we have,
$$\begin{array}{cc}P_{m,n}^{M+1}(y)\hfill & =\frac{Mm}{Mm+1}P_{m,n}^M(y)+\frac{1}{M+1m}A_{f_{m,n}(\sigma _{M+1m},\mathrm{},\sigma _M)\sigma _{M+1}}\hfill \\ Q_{m,n}^{M+1}(y)\hfill & =\frac{Mm}{Mm+1}Q_{m,n}^M(y)+\frac{1}{M+1m}A_{\sigma _{M+1}f_{m,n}(\sigma _{M+1m},\mathrm{},\sigma _M)}\hfill \end{array}$$
$`(4.10)`$
and, as in (4.9),
$$\begin{array}{cc}_0^1P_{m,n}^{M+1}(y)๐y=\hfill & \frac{Mm}{Mm+1}_0^1P_{m,n}^M(y)๐y\hfill \\ & +\frac{1}{2^{m+1}(M+1m)}\left(n_0^{m,n}(A_{00}+A_{01})+n_1^{m,n}(A_{10}+A_{11})\right)\hfill \\ _0^1Q_{m,n}^{M+1}(y)๐y=\hfill & \frac{Mm}{Mm+1}_0^1Q_{m,n}^M(y)๐y\hfill \\ & +\frac{1}{2^{m+1}(M+1m)}\left(n_0^{m,n}(A_{00}+A_{10})+n_1^{m,n}(A_{01}+A_{11})\right)\hfill \end{array}$$
$`(4.11)`$
By (4.9) and by induction from (4.11), we obtain,
$$\begin{array}{cc}_0^1P_{m,n}^{M+1}(y)๐y=\hfill & \frac{1}{2^{m+1}}\left(n_0^{m,n}(A_{00}+A_{01})+n_1^{m,n}(A_{10}+A_{11})\right)\hfill \\ _0^1Q_{m,n}^{M+1}(y)๐y=\hfill & \frac{1}{2^{m+1}}\left(n_0^{m,n}(A_{00}+A_{10})+n_1^{m,n}(A_{01}+A_{11})\right)\hfill \end{array}$$
$`(4.12)`$
As the integrals in (4.12) are independent of $`M`$, by (4.6), we have proved:
###### Theorem 4.1.
We consider an ensemble of non-cooperative and symmetric two-player games, where in each subensemble we have a player $`P`$ playing with deterministic policy $`f_{m,n}`$, and a player $`Q`$ making the choices of pure strategies with equal probabilities. Then, the mean payoffs per move of the representative players of the game depend on the payoff matrix $`A`$ and on the strategy $`f_{m,n}`$, and the mean payoffs per move are,
$$\begin{array}{cc}\overline{P}_{m,n}=\hfill & \frac{1}{2^{m+1}}\left(n_0^{m,n}(A_{00}+A_{01})+n_1^{m,n}(A_{10}+A_{11})\right)\hfill \\ \overline{Q}_{m,n}=\hfill & \frac{1}{2^{m+1}}\left(n_0^{m,n}(A_{00}+A_{10})+n_1^{m,n}(A_{01}+A_{11})\right)\hfill \end{array}$$
where the $`A_{ij}`$ are the entries of the payoff matrix, $`n_0^{m,n}=\mathrm{\#}\{\sigma ^{(m)}\{0,1\}^m:f_{m,n}(\sigma ^{(m)})=0\}`$, and $`n_1^{m,n}=\mathrm{\#}\{\sigma ^{(m)}\{0,1\}^m:f_{m,n}(\sigma ^{(m)})=1\}`$.
This theorem has a direct consequence. With the definitions of Section 2, a policy or strategy $`f_{m,n}`$ is equalitarian if the mean payoffs of the representative players are equal. Imposing the equality between $`\overline{P}_{m,n}`$ and $`\overline{Q}_{m,n}`$ in Theorem 4.1, we obtain,
$$n_0^{m,n}(A_{01}A_{10})+n_1^{m,n}(A_{10}A_{01})=0$$
$`(4.13)`$
From (4.13) it follows that a policy is equalitarian if either $`n_0^{m,n}=n_1^{m,n}`$ or, $`A_{01}=A_{10}`$. In the first case, we have the class of all GTFT policies, independently of the values of the entries of the payoff matrix $`A`$.
If $`A_{01}=A_{10}`$, it follows from Theorem 4.1 and (4.13), that,
$$\overline{P}_{m,n}=\overline{Q}_{m,n}=\frac{n_0^{m,n}}{2^{m+1}}(A_{00}A_{11})+\frac{1}{2}A_{11}+\frac{1}{2}A_{01}$$
$`(4.14)`$
where we have introduced the relation $`n_1^{m,n}=2^mn_0^{m,n}`$. Therefore, we have:
###### Corollary 4.2.
We consider an ensemble of non-cooperative and symmetric two-player games with payoff matrix $`A`$, where in each subensemble we have a player $`P`$ playing strategy $`f_{m,n}`$, and a player $`Q`$ making the choices of pure strategies with equal probabilities. Then the policy $`f_{m,n}`$ is equalitarian if either, it is GTFT or, $`A_{01}=A_{10}`$. Moreover, the payoffs per move of GTFT policies are given by,
$$\overline{P}_{m,n}=\overline{Q}_{m,n}=\frac{1}{2^2}\left(A_{00}+A_{01}+A_{10}+A_{11}\right)$$
For example, in games with memory length $`m=1`$, independently of the payoff matrix $`A`$, the equalitarian strategies are $`f_{1,1}`$ and $`f_{1,2}`$, both GTFT. From the point of view of the ensemble mean payoffs per move, all the GTFT strategies are equivalent to ensemble games where all player play randomly with equal probability.
We determine now the best policy for a player $`P`$ with an opponent $`Q`$ choosing pure strategies with equal probabilities. By Theorem 4.1, and with $`n_1^{m,n}=2^mn_0^{m,n}`$, we obtain,
$$\overline{P}_{m,n}=\frac{n_0^{m,n}}{2^{m+1}}(A_{00}+A_{01}A_{11}A_{10})+\frac{1}{2}(A_{11}+A_{10})$$
$`(4.15)`$
Therefore, in the sense of ensemble average and for a given memory length $`m`$, the best policies for the player $`P`$ are the ones that maximise (4.15), for all the choices of the integers $`n_0^{m,n}=0,\mathrm{},2^m`$.
## 5 Both players have deterministic strategies
When the two representative players $`P`$ and $`Q`$ have deterministic strategies within the same memory class, their game records become dependent of the first $`m`$ moves of the players. As we have two players and $`2^m`$ different initial conditions for each player, for each choice of a pair of deterministic strategies, there are at most $`2^{m+1}`$ different payoffs per move for both players. As there are $`2^{2^m}`$ different boolean functions of memory length $`m`$, the maximum number of equilibrium states is, $`2^{2^{m+1}}\times 2^{m+1}`$, which, for $`m=1`$, is $`64`$.
Let us analyse now in detail the case of memory length $`m=1`$. If $`\mu _1`$ and $`\sigma _1`$ represent the choices for the first move of players $`P`$ and $`Q`$, and $`P`$ and $`Q`$ have policies $`ff_{1,r}`$ and $`gf_{1,s}`$, respectively, their game records are,
$$\begin{array}{cc}P:\hfill & (\mu _1,f(\sigma _1),fg(\mu _1),fgf(\sigma _1),\mathrm{})\hfill \\ Q:\hfill & (\sigma _1,g(\mu _1),gf(\sigma _1),gfg(\mu _1),\mathrm{})\hfill \end{array}$$
where $`fg(\mu _1)=f(g(\mu _1))`$. After a few moves, the game records become periodic. Therefore, the mean payoff per move of each player can be calculated by the periodic sequences which depend on the initial moves and on the policies. For example, with $`ff_{1,2}`$ and $`gf_{1,2}`$, and initial moves $`\mu _1=0`$ and $`\sigma _1=1`$, we obtain the game records,
$$\begin{array}{cc}P:\hfill & (0,1,0,1,\mathrm{})\hfill \\ Q:\hfill & (1,0,1,0,\mathrm{})\hfill \end{array}$$
and the mean payoff per move of both players is $`\overline{P}=\overline{Q}=(A_{01}+A_{10})/2`$. But for the initial moves $`\mu _1=0`$ and $`\sigma _1=0`$, we have, $`\overline{P}=\overline{Q}=A_{00}`$.
In Table 2, we show the mean payoffs per move and per player, for all the deterministic policies with memory length $`m=1`$ and all the possible four different initial moves of the players. Counting the different values in the entries in table, we conclude that, for $`m=1`$, the number of equilibrium states is $`7`$. For a given game, the best strategy and initial conditions is obtained by analyzing the entries of Table 2. Clearly, the best strategy depends on the entries of the payoff matrix of the game.
In general, let $`(\mu _1,\mathrm{},\mu _m)`$ and $`(\sigma _1,\mathrm{},\sigma _m)`$ be the first $`m`$ moves of players $`P`$ and $`Q`$, respectively. Suppose further that player $`P`$ and $`Q`$ choose the deterministic strategies $`f_{m,n}`$ and $`g_{m,n}`$, respectively. Iterating the game, after some transient iteration the game record sequences become periodic, and the mean payoffs per move and per player are easily calculated. If $`(\mu _{i+1},\mathrm{},\mu _{i+p})`$ and $`(\sigma _{i+1},\mathrm{},\sigma _{i+p})`$, for some $`i1`$, are the periodic patterns of period $`p`$ of the game record sequences, the mean payoffs per move of the players are,
$$\begin{array}{cc}P:\hfill & \frac{1}{p}(A_{\mu _{i+1}\sigma _{i+1}}+\mathrm{}+A_{\mu _{i+p}\sigma _{i+p}})\hfill \\ Q:\hfill & \frac{1}{p}(A_{\sigma _{i+1}\mu _{i+1}}+\mathrm{}+A_{\sigma _{i+p}\mu _{i+p}})\hfill \end{array}$$
and these mean payoffs are the equilibrium states of the game. For example, for $`m=2`$, we have at most $`2^{2^{2+1}}\times 2^{2+1}=2048`$ equilibrium states.
## 6 Examples and policy analysis
The formalism introduced in the previous sections leads to the evaluation of policies for an iterated game with a given payoff matrix $`A`$. In this context, we can forget the role of players $`P`$ and $`Q`$ and speak about the performance of the game, the performance of a deterministic strategy and the relative performance of two deterministic strategies.
We analyze now two examples, the Prisonerโs Dilemma game and the Hawk-Dove game.
### 6.1 The Prisonerโs Dilemma
In the Prisonerโs Dilemma game with payoff matrix (2.5), if all players make their choices with equal probabilities, by Theorem 3.1, the mean payoff per move and per player is zero. By Corollary 4.2, a player with a GTFT policy against a player choosing its pure strategies randomly has also zero payoffs. This includes the simplest case of the tit-for-tat policy.
From the point of view of the non-deterministic map $`\beta :๐ฆ๐ฆ`$, the situation of Theorem 3.1 corresponds to the equilibrium solution $`(0,0)๐ฆ`$, Fig. 3a).
In the case of Theorem 4.1 and for deterministic strategies with memory length $`m=1`$, there are three equilibrium solutions for the Prisonerโs Dilemma game. These equilibrium solutions are: $`(3,7)๐ฆ`$, $`(0,0)๐ฆ`$ and $`(3,7)๐ฆ`$, Fig. 1b). So, in a uniform collectivity, the players that choose the dominant strategy have a better payoff, provided their partners choose their strategies with equal probabilities.
Suppose now that both representative players $`P`$ and $`Q`$ adopt a policy with memory length $`m=1`$. Analysing the results of Table 2, the best payoff per move for both players is obtained when player $`P`$ and $`Q`$ play tit-for-tat and both choose the initial strategy โ0โ. The best payoff per move is also obtained when one at least of the contenders chooses โ0โ and the other plays according the tit-for-tat policy. In the case of policies of memory length $`m=1`$, the tit-for-tit policy forces cooperation. If one of the player plays tit-for-tat and the other player chooses another strategy, tit-for-tit ensures that the payoffs per move of both players are equal and the second player is not able to increase its payoff per move. If both players choose a tit-for-tat policy, depending on the initial condition, we can have four different equilibrium states of the game, Table 2 and Fig. 1c). In this case, two of them are the strict and the bargain Nash solutions. If one of the players chooses always the strategy โ1โ, it corresponds to the deterministic policy $`f_{1,3}`$, and the outcome of the game against a tit-for-tat corresponds to the Nash strict equilibrium of the game. The seven equilibrium solutions of the Prisonerโs game are plotted in Fig. 3c).
If $`P`$ has to choose a policy against a player $`Q`$ that plays its strategies with equal probabilities, by (4.15) and (2.5), the best policy for $`P`$ is the one that maximizes,
$$\overline{P}_{m,n}=\left(312\frac{n_0^{m,n}}{2^{m+1}}\right)$$
Therefore, in the Prisonerโs Dilemma game, the best policy corresponds to $`n_0^{m,n}=0`$, which corresponds to the policy function $`f_{m,2^m1}`$. In this case, we have, $`\overline{P}_{m,2^m1}=3`$ and $`\overline{Q}_{m,2^m1}=7`$, Fig. 3b).
A more detailed analysis of Fig. 3, shows that Nash bargain solutions and Nash strict solutions only exist when both players have deterministic policies. In the sense of ensemble averages, Nash solutions are not equilibrium solutions of a game.
### 6.2 The Hawk-Dove game
The Hawk-Dove game has been introduced by Maynard-Smith and Price as a game theoretical basic model to describe animal conflicts. They have assumed two pure strategies: Hawk (โ0โ) and Dove (โ1โ). A player chooses Hawk or โ0โ if he acts fiercely, and chooses Dove or โ1โ if he looks fierce and then retires. In the context of evolutionary biology, this game aims to explain the struggle for a territory whose payoff is related with the number of offsprings. The payoff matrix of the Hawk-Dove game is,
$$A=\left(\begin{array}{cc}\frac{1}{2}(rc)& r\\ 0& \frac{r}{2}+\epsilon \end{array}\right)$$
where $`r`$ represents the reproductive value and $`c`$ is the cost of injury. In this game the Hawk strategy is dominant, provided $`c<r`$ and $`\epsilon <r/2`$. If $`\epsilon >0`$ and $`\epsilon <r/2`$, the Hawk-Dove game is also dilemmatic. Globally, the species has advantage if everybody acts Dove, which is the non-dominant strategy.
If all players choose their strategies with equal probability, by Theorem 3.1, the mean payoff per player and move is,
$$\overline{P}=\overline{Q}=\frac{1}{2}(r\frac{1}{4}c)+\epsilon $$
If $`c<4r+8\epsilon `$, $`\overline{P}=\overline{Q}>0`$, the Hawk-Dove game shows advantage for the species. If $`c4r+8\epsilon `$, the cost of injury is too high and globally the mean payoff per player and move is non-positive.
If the representative players of the game choose a generalized tit-for-tat policy with memory length $`m=1`$, and $`c<r`$, both players have a positive mean payoff per move. If the players choose not being the first to play Hawk, they both obtain the highest mean payoffs per move.
In the Hawk-Dove non-cooperative and symmetric game, the tit-for-tat policy or imitation of the adversary move implies a positive payoff for both players, provided the cost of injury is not too high ($`c<4r+8\epsilon `$).
## 7 Conclusions
The dynamics of an iterated game is described by a one-to-many map defined on a state space, . Within this framework, the concept of mixed strategy leads to the definition of the equilibrium solution of a game. This equilibrium solution is obtained as the limit of the iterates of a one-to-many map. In general, for a specific game, the equilibrium solutions associated to the set of all mixed strategies span the state space of the game. The concepts of strict Nash equilibrium and bargain solutions of a game are discussed within this framework.
In applications of game theory to economics, ethics, sociology, biology, physics, etc., it is sometimes easy to identify rules of behaviour and interactions between agents and to make guesses about payoffs. However, it is difficult to argue about the (infinite) memory of all the past choices of the players, and to insure that opponent players remain the same during all the iterated game. Therefore, the way of evaluating a game, or a policy depends on the context in which the game is considered.
In order to evaluate a game, we have introduced the concept of representative ensemble of a game. This technique has been applied to the global evaluation of a game, without any specific considerations about policies. In this evaluation, all the players make their choices of pure strategies with equal probabilities. In this case, we have shown that the mean payoffs per move of the players are the mean value of the entries of the payoff matrices of the game.
To evaluate a deterministic policy with a finite memory length, we have calculated the mean payoffs per move of the players, for the case where one of the players has a deterministic policy and the other player chooses its pure strategies with equal probabilities. In this case, there exists a class of deterministic policies that forces equality of the mean payoffs per move of the players. This class of policies is the class of generalized tit-for-tat policies. When a representative player has a generalized tit-for-tat policy, in the limit of the iterated game, the payoffs of both representative players are equal. If a player tries to increase its payoff by changing its strategy and the other player plays tit-for-tat, the change in the strategy can increase or decrease the payoffs, but the payoffs per move of both players remain equal. Generalized tit-for-tat or imitation strategies force equalitarian payoffs per move. In dilemmatic games, the generalized tit-for-tat policy together with the condition of not being the first to defect, leads to the highest possible mean payoffs per move for the players.
Acknowledgments: This work has been partially supported by the POCTI Project P/FIS/13161/1998, and by Fundaรงรฃo para a Ciรชncia e a Tecnologia, under a plurianual funding grant.
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# A Brighter Past: Galaxy Luminosity Function At High Redshifts
## 1 Introduction
The luminosity function (LF) of field galaxies at the present day can be described using a simple empirical model that involves the dark matter halo mass function and the relation between central galaxy luminosity and the mass of the halo occupied (Cooray & Milosavljevic 2005b; Cooray 2005). The model utilizes an approach based on the conditional luminosity function (CLF; Yang et al. 2003b, 2005), or the luminosity distribution of galaxies as a function of the halo mass, $`\mathrm{\Phi }(L|M)`$, to construct the galaxy LF, $`\mathrm{\Phi }(L)`$. The CLFs are basically an extension of the halo approach to galaxy statistics (Seljak 2000; Scoccimarro et al. 2001; Cooray et al. 2000; see, Cooray & Sheth 2002 for a review) where instead of the halo occupation number as a function of the halo mass, $`N_g(M)`$ โ which describes the average number of galaxies in a dark matter halo โ we consider the number of galaxies in a given dark matter halo conditioned in terms of the galaxy luminosity such that $`\mathrm{\Phi }(L|M)=dN_g(M)/dL`$. Other conditions of the halo occupation number has been suggested (such as in terms of the stellar mass in Zheng et al. 2004), but for comparison with observational data, a halo occupation number conditioned in terms of galaxy luminosity and type is most useful.
As the halo occupation statistics are best described based on a division to central and satellite galaxies (Kravtsov et al. 2003), we also divide the CLF to central galaxies and satellites; central galaxies are assigned a log-normal distribution in the luminosity, centered around the mean central galaxy luminosity given the halo mass (Cooray & Milosavljevic 2005a), while satellite galaxies are assigned a power-law distribution in luminosity. This empirical approach has the main advantage that it can elucidate important aspects associated with the galaxy distribution as measured by surveys of the large scale structure. Statistical measurements from observations mostly include galaxy LFs, both field and cluster galaxies, and statistics related to galaxy clustering such as the large-scale bias factor or the correlation length. Thus, we use CLFs to model same statistics here.
Previous studies using the halo model to describe galaxy statistics concentrated primarily on average clustering properties, such as the galaxy power spectrum (e.g., Seljak 2000; Scoccimarro et al. 2001; Cooray 2002; Berlind et al. 2003); for these statistics the conditional occupation number, either in luminosity or other galaxy property, is not needed as the statistic depends simply on the total number of galaxies in a given halo. With large data sets, on the other hand, measurements can be made with the galaxy sample divided to various physical properties such as the luminosity, color, environment, etc (e.g., Zehavi et al. 2004 who considered clustering of Sloan galaxies as a function of the luminosity). In this scenario, to compare observations with analytical or numerical models, the average halo occupation number must be conditioned in terms of the physical property. Similarly, wide-field surveys at high redshifts have now allowed detailed measurements related to the redshift evolution of the LF. Thus, one must also account for redshift variations in the CLF. Here, we consider the latter application and improve prior analytical models of the LF by discussing redshift dependence of the CLF.
In Cooray & Milosavljeviฤ (2005b), we used the CLF to explain why the LF can be described with the Schechter (1976) form of $`\mathrm{\Phi }(L)(L/L_{})^\alpha \mathrm{exp}(L/L_{})`$. In this approach the main ingredient, in addition to the halo mass function, is the relation between central galaxy luminosity and the halo mass, hereafter called the $`L_c(M)`$ relation. This relation, as appropriate for galaxies at low-redshifts, was established in Cooray & Milosavljevic (2005a) from a combination of weak lensing (e.g., Yang et al. 2003a) and direct measurements of galaxy luminosity and mass in groups and clusters (e.g., Lin et al. 2004). The same relation has been established with a statistical analysis of the 2dFGRS $`b_J`$-band LF (e.g., Norberg et al. 2002) by Vale & Ostriker (2004) and, independently, by Yang et al. (2005) based on the 2dF Galaxy Redshift Survey (2dFGRS; Colles et al. 2001) galaxy group catalog. The shape of the $`L_c(M)`$ relation, where luminosities grow rapidly with increasing mass but flattens at a mass scale around $``$ 10<sup>13</sup> M$`_{}`$ is best explained through dissipationless merging history of central galaxies (Cooray & Milosavljevic 2005a).
In Cooray & Milosavljeviฤ (2005b), we related the two parameters of the Schechter (1976) LF involving the slope at low luminosities, $`\alpha `$, and the exponential cut-off luminosity, $`L_{}`$, to a combination of the mass function slope and the slope of the $`L_c(M)`$ relation. In this analytical model, we expect $`\alpha <1.25`$, consistent with observations that indicate $`\alpha 1.3`$ (Blanton et al. 2004; Huang et al. 2003). The characteristic scale $`L_{}`$ comes about when the luminosity scatter in the $`L_c(M)`$ relation dominates over the increase in the luminosity with mass or when $`d\mathrm{ln}L_c/d\mathrm{ln}M\mathrm{ln}(10)\mathrm{\Sigma }`$ where $`\mathrm{\Sigma }`$ is the dispersion in the $`L_c(M)`$ relation. Given the observed dispersion, we find $`M_{}2\times 10^{13}`$ M$`_{}`$ and $`L_{}=L_c(M_{})`$ to be consistent with observed $`L_{}`$ values from Schechter (1976) function fits to the LF (Cooray & Milosavljevic 2005a in the case of k-band observations and Cooray 2005 in the case of 2dFGRS $`b_J`$ band).
The empirical modeling approach can easily be extended to consider statistics of galaxy types or color as well. For example, in Cooray (2005), we studied the environmental dependence of galaxy colors, broadly divided to blue and red types given the bimodal nature of the color distribution (e.g., Baldry et al. 2004; Balogh et al. 2004). There, we described the conditional type-dependent LFs from 2dFGRS (Croton et al. 2004) as a function of the galaxy overdensity based on an empirical description of blue-to-red galaxy fraction in dark matter halos as a function of the halo mass. With an increasing fraction of early-type, or red, galaxies with increasing halo mass, the simple analytical model considered in Cooray (2005) explain why the LF of galaxies in dense environments are dominated by these galaxies.
In addition to galaxy statistics today from wide-field redshift surveys such as 2dFGRS or Sloan Digital Sky Survey (SDSS; York et al. 2000), various techniques, such as the Lyman drop-out method (e.g., Steidel et al. 1999), have now allowed the study of galaxy LF and related statistics on the galaxy distribution at high redshifts. Spectroscopic redshift surveys, such as DEEP2 (Davis et al. 2003), and photometric data based redshift selections, such COMBO-17 survey (Wolf et al. 2001, 2003), have now allowed detailed studies of galaxy properties at redshifts around unity (e.g., Willmer et al. 2005; Faber et a. 2005 with DEEP2 and Bell et al. 2004 with COMBO-17). Similar photometric redshift based studies extend the LF statistics to higher redshifts using ground-based data (e.g., Gabasch et al. 2004 using ESO VLTโs FORS Deep Field), space-based data (e.g., Bouwens et al. 2004 using HST NICMOS Ultra Deep Field), or a combination of space-based imaging and ground-based spectroscopic followup data (e.g., Giallongo et al. 2005). While the LF of galaxies today can be described through the observationally established $`L_c(M)`$ relation, it is also useful to understand the extent to which the same empirical approach can be applied at high redshifts. In return, using the measured galaxy LFs out to a redshift of 6, we can attempt to extract information on how galaxy luminosities evolve as a function of redshift. The model may then allow one to address if the galaxy formation was efficient in the past and how galaxy properties are different when compared to properties today. When comparing to LF measurements, we ignore any potential systematics in these data such as due to selection effects and biases that may have affected the LF measurements. We make the assumption that, if any biases or selection effects exist, these effects have been considered and that the LFs are properly corrected to account for them. Thus, our model comparisons may only be accurate to the same extent that measurements of high redshift LFs can be considered reliable.
Here, we compare our predictions to the measured rest frame B-band LF of galaxies, including red and blue galaxies, out to a redshift of 1.2 from DEEP2 (Willmer et al. 2005) and COMBO-17 (Bell et al. 2004), out to a redshift of 5 from Gabasch et al. (2004) and out to a redshift of 3.5 from Giallongo et al. (2005). We use the latter data set from Giallongo et al. (2005) to build out models but then perform a detailed model fit to other data sets by varying some of the parameters related to redshift evolution of the $`L_c(M)`$ relation. We also make comparisons to rest-UV LFs of galaxies at redshifts 3 to 6 from Steidel et al. (1999) and Bouwens et al. (2004a), since there surveys provide an additional data sets to compare with models at the highest redshift ranges surveyed so far. While previous studies have measured the redshift dependence of the LF, say in the K-band (e.g., Drory et al. 2003), given that the LF corresponds to different rest wavelengths as a function of redshift, any evolutionary aspects associated with galaxy properties, at a given wavelength, must be distinguished from evolutionary effects resulting from color differences. The B-band LFs considered for modeling here have the advantage that one can directly address how galaxy properties change with redshift at the same band, regardless of color differences.
The paper is organized as follows: In the next section, we will outline the basic ingredients in the empirical model for CLFs and how it is modified to model the LF at high redshifts. We refer the reader to Cooray (2005) and Cooray & Milosavljeviฤ (2005b) for initial discussions related to this empirical modeling approach. In Section 3, we will describe the z-dependent LF and compare with measurements by discussed above. We also compare our models to rest-UV LFs of galaxies between redshifts of 3 to 6, and galaxy clustering bias around the same redshift ranges. We conclude with a summary of our main results and implications related to the galaxy distribution at redshifts $``$ 3 to 6 in ยง 4. Throughout the paper we assume cosmological parameters consistent with observational analyses of LF measurements modeled here and take $`\mathrm{\Omega }_m=0.3`$, $`\mathrm{\Omega }_\mathrm{\Lambda }=0.7`$, and a scaled Hubble constant of $`h=0.7`$ in units of 100 km s<sup>-1</sup> Mpc<sup>-1</sup>. The matter power spectrum is normalized to a $`\sigma _8`$, rms fluctuations at 8 $`h^1`$ Mpc scales, of 0.84 consistent with WMAP data (Spergel et al. 2003).
## 2 Conditional Luminosity Function: Empirical Model
In order to construct the redshift evolution of the luminosity function (LF), we follow Cooray & Milosavljeviฤ (2005b) and Cooray (2005). The redshift-dependent conditional luminosity function (CLF; Yang et al. 2003b, 2005), denoted by $`\mathrm{\Phi }(L|M,z)`$, is the average number of galaxies with luminosities between $`L`$ and $`L+dL`$ that reside in halos of mass $`M`$ at a redshift of $`z`$. As in our application to 2dFGRS at $`z=0`$ (Cooray 2005), the CLF is separated into terms associated with central and satellite galaxies, such that
$`\mathrm{\Phi }(L|M,z)`$ $`=`$ $`\mathrm{\Phi }_\mathrm{c}(L|M,z)+\mathrm{\Phi }_\mathrm{s}(L|M,z)`$
$`\mathrm{\Phi }_\mathrm{c}(L|M,z)`$ $`=`$ $`{\displaystyle \frac{f_\mathrm{c}(M,z)}{\sqrt{2\pi }\mathrm{ln}(10)\mathrm{\Sigma }(z)L}}\times `$
$`\mathrm{exp}\left\{{\displaystyle \frac{\mathrm{log}_{10}[L/L_\mathrm{c}(M,z)]^2}{2\mathrm{\Sigma }^2(z)}}\right\}`$
$`\mathrm{\Phi }_\mathrm{s}(L|M)`$ $`=`$ $`A(M,z)L^{\gamma (M,z)}.`$ (1)
Here $`L_c(M,z)`$ is the relation between central galaxy luminosity of a given dark matter halo and itโs halo mass, taken to be a function of redshift, while $`\mathrm{ln}(10)\mathrm{\Sigma }(z)`$ is the dispersion in this relation, again a function of redshift. In equation (1), $`f_c(M,z)`$ is an additional selection function introduce to account for the efficiency for galaxy formation as a function of the halo mass, given the fact that at low mass halos galaxy formation is inefficient and not all dark matter halos may host a galaxy. Given limited statistics to extract information on $`f_c(M,z)`$, and given that our primary goal is to study the evolution associated with $`L_c(M,z)`$ relation, we set $`f_c(M,z)=1`$ throughout here. This assumption also leads to reasonably well model fits to LFs either at low redshifts as studied in Cooray & Milosavljevic (2005b) or at high redshifts as in the case here.
The central galaxy CLF takes a log-normal form, while the satellite galaxy CLF takes a power-law form in luminosity. Such a separation describes the LF best, with an overall better fit to the data in the K-band as explored by Cooray & Milosavljeviฤ (2005b) and 2dFGRS $`b_J`$-band in Cooray (2005). Our motivation for log-normal distribution also comes from measured conditional LFs, such as galaxy cluster LFs that include bright central galaxies, where a log-normal component, in addition to the Schechter (1976) form, is required to fit the data (e.g., Trentham & Tully 2002). Similarly, the stellar mass function, as a function of halos mass in semi-analytical models, is best described with a log-normal component for central galaxies (Zheng et al. 2004). To simplify the modeling approach, we assume $`\mathrm{\Sigma }(z)=\mathrm{\Sigma }(0)=0.17`$, and $`\gamma (M,z)=0.5`$; The value for $`\mathrm{\Sigma }`$ comes from a comparison to low-redshift LF in 2dFGRS (Norberg et al. 2002), while the latter ignores the mass dependence of the satellite luminosity distribution suggested in Cooray (2005).
To describe galaxies as a function of color in this analytical description, we must further divide central and satellite galaxies as a function of their color given the luminosity. Here, motivated by the bimodality of color (e.g., Baldry et al. 2004) that extends out to high redshifts (e.g., Giallongo et al. 2005), we consider models in terms of galaxy types. The description in terms of galaxy types is also useful since measurements at high redshifts, so far, involve the division of galaxy samples to two broad categories involving early-type, or red, and late-type, or blue, galaxies. Thus, in the case of early type galaxies, we write CLFs as
$`\mathrm{\Phi }_{\mathrm{early}\mathrm{cen}}(L|M,z)`$ $`=`$ $`\mathrm{\Phi }_\mathrm{c}(L|M,z)f_{\mathrm{early}\mathrm{cen}}(M,z)`$
$`\mathrm{\Phi }_{\mathrm{early}\mathrm{sat}}(L|M)`$ $`=`$ $`\mathrm{\Phi }_\mathrm{s}(L|M,z)f_{\mathrm{early}\mathrm{sat}}(M,L,z),`$ (2)
where the two functions that divide between early- and late-types are taken to be functions of mass, in the case of central galaxies, and both mass and luminosity in the case of satellites. As fractions are defined with respect to the total galaxy number of a halo, late-type fractions are simply $`[1f_{\mathrm{early}\mathrm{cen}}(M,z)]`$ and $`[1f_{\mathrm{early}\mathrm{sat}}(M,L,z)]`$ for central and satellite galaxies, respectively. These functions, as appropriate for 2dFGRS $`b_J`$-band galaxy type LFs, are described in Cooray (2005).
Here, we will use these functions to describe the type dependence of the high redshift galaxy LF and ignore all redshift dependences and assume that regardless of the redshift, the fraction of of early-to-late type galaxies, in a halo of fixed mass, is same as the fraction of that halo mass today. While this assumption may seem contradictory with observations at the first instance, given that observations indicate that ellipticals are older than spirals or the fact that mass correlates with age, this is not necessarily the case given the strong evolution associated with the halo mass distribution. While our assumption is that the ratio of red-to-blue galaxies, of a given halo mass, at a redshift $`z`$ is same as the fraction today, the mass function evolves such that at high redshifts, the universe is dominated by small mass halos while the density of high mass halos increases to low redshifts. Given that in the model of Cooray (2005) the fraction of early-type galaxies increases with increasing mass, the relative fraction of early-to-late type galaxies, when averaged over the mass distribution of halos at a given redshift, is not the same, but increases to lower redshifts. Thus, while we have not assumed or specified important astrophysical processes such as galaxy mergers that may be responsible for galaxy types, given that the halo mass function evolves through merging, galaxies must undergo merging as well. While the specific process may remain hidden, what we can extract with this empirical modeling approach is the exact fraction of early-to-late type galaxies as a function of halo mass; It remains a task for numerical simulators and other approaches to understand galaxy distribution, such as semi-analytical models of the galaxy formation, to explain the results extracted from the empirical model here when compared to observations.
After various simplifications, which can be ignored as observational statistics at high redshifts improve, the only redshift dependence in our empirical model comes from redshift variations associated with the $`L_c(M,z)`$ relation, and when describing satellites, $`L_{\mathrm{tot}}(M,z)`$, the total luminosity of galaxies as a function of the halo mass, relation; the latter, however, is not an important ingredient since the LF is primarily determined by statistics of central galaxies (Cooray & Milosavljevic 2005b; Cooray 2005).
### 2.1 Central Galaxy Luminosity-Halo Mass Relation
For $`L_\mathrm{c}(M,z=0)`$ relation, here we make use of the relation derived in Vale & Ostriker (2004). These authors established this relation by inverting the 2dFGRS luminosity function given an analytical description for the sub-halo mass function of the Universe (e.g., De Lucia et al. 2004; Oguri & Lee 2004). We used this relation to describe the 2dFGRS $`b_J`$-band LF in Cooray (2005), and we will use the same relation, at $`z=0`$, as an approximation to describe the B-band LF. The relation is described with a general fitting formula given by
$$L(M,z=0)=L_0\frac{(M/M_1)^a}{[b+(M/M_1)^{cd}]^{1/d}}.$$
(3)
For central galaxy luminosities, the parameters are $`L_0=5.7\times 10^9L_{}`$, $`M_1=10^{11}M_{}`$, $`a=4.0`$, $`b=0.57`$, $`c=3.72`$, and $`d=0.23`$ (Vale & Ostriker 2004). For the total galaxy luminosity, as a function of the halo mass, we also use the fitting formula in equation (3), but with $`c=3.57`$. As discussed in Cooray (2005), the overall shape of the LF is strongly sensitive to the shape of the $`L_\mathrm{c}`$$`M`$ relation, and itโs scatter, and less on details related to the $`L_{\mathrm{tot}}`$$`M`$ relation.
To describe the redshift evolution, first, we consider two possibilities, but using a large sample of datasets will combine them to consider a general model fit. First, we describe the high-z LFs with $`L(M,z)=L(M,z=0)(1+z)^\alpha `$. This is a scenario in where all luminosities either increase or decrease depending on the value and sign of $`\alpha `$, which we take to be mass independent. Such an evolution provides an acceptable description of the LFs, say of Giallongo et al. (2005), though the increase in luminosity of galaxies in less massive dark matter halos overestimates the LF at the faint-end at high redshifts. At $`z6`$, this overestimate becomes significant and even the LF at the bright-end is overestimated relative to the UV LF measured by Bouwens et al. (2004a). A preferred description may be a case where low luminosity end of the $`L_c(M,z)`$ relation remains independent of the redshift, while the bright-end increases with increasing redshift. To describe this behavior, we take parameters $`c`$ and $`d`$ in equation (3) to be dependent on the redshift with $`c(z)=c(z=0)(1+z)^\beta `$ and $`d(z)=d(z=0)(1+z)^\eta `$, where $`\beta `$ and $`\eta `$ are taken to be free parameters. When combined, as we discuss later, we find that the redshift evolution related to $`d(z)`$ is not strongly constrained by observational data while $`\alpha `$ and $`\beta `$ are. Using the LF measurements from DEEP2, COMBO-17, and those extending to $`z5`$ by Gabasch et al. (2004), we will provide general constraints on $`\alpha `$ and $`\beta `$.
In Figure 1, we show the $`L_c(M,z)`$ relation as a function of the halo mass and for redshifts from 0 to 6. For comparison, we also show measurements from the 2dFGRS galaxy group catalog from Yang et al. (2005). In the left panel, we show the pure-luminosity evolution case with $`\alpha =0.75`$ and in the right-panel, we show the evolution of the bright-end with $`\beta =0.07`$ and $`\eta =0.05`$; These numerical values were selected based on a comparison to the high-redshift LFs of Giallongo et al. (2005).
Note that we are assuming here that the central galaxy luminosity of a given halo increases with redshift. This assumption does not violate the fact that the halo occupation number is not changing with redshift (e.g., Yan et al. 2003; Coil et al. 2004), since the integral of $`\mathrm{\Phi }(L|M)`$ over luminosities remain the same; all we have done is to shift the mean of the log-normal distribution that describes central galaxy luminosity distribution to a higher luminosity when compared to the value today. The fact that the halo occupation number is the same at high redshifts, when compared to today, should not be considered as a statement that galaxy properties do not change with redshift.
## 3 High-Redshift Luminosity Functions
Given our model for CLFs, we can now construct the LF by averaging CLFs over the halo mass distribution given by the mass function. Here, we use the Sheth & Tormen (1999; ST) mass function $`dn/dM`$ for dark matter halos. This mass function is in better agreement with numerical simulations (Jenkins et al. 2001), when compared to the more familiar Press-Schechter (PS; Press & Schechter 1974) mass function. While there are differences in the ST mass function and the numerically simulated mass functions at the high mass end, these differences do not affect these results as statistics of the LF at the present day are dominated by galaxies in halos around $`10^{13}`$ M$`_{}`$. As one moves to a higher redshift, statistics become dominated by lower mass halos than today; at $`z3`$, statistics are dominated by halos with mass in the range between few times $`10^{11}`$ M$`_{}`$ and few times $`10^{12}`$ M$`_{}`$. Here, we make the assumption that the ST mass function is the correct description for halo masses out to $`z`$ 6 and above; any differences in the evolution of the mass function, relative to ST description, could affect our conclusions regarding the luminosity evolution. The same can also said of our assumption related to underlying cosmological model; a different cosmology than the one considered here impacts the redshift evolution of the mass function differently than the one we assume here. These differences, however, are at the few percent level, at most, and do not strongly affect our results given the few tens percent uncertainties in the LF at high redshifts.
Given the mass function, the galaxy LF as a function of $`z`$ is
$$\mathrm{\Phi }_i(L,z)=_0^{\mathrm{}}\mathrm{\Phi }_i(L|M,z)\frac{dn}{dM}(z)๐M,$$
(4)
where $`i`$ is an index for early and late type galaxies. The conditional luminosity function for each type involves the sum of central and satellites. To compare with Giallongo et al. (2005), Willmer et al. (2005; DEEP2) and Bell et al. (2004; COMBO-17) measurements, we will plot these two divisions, as well as the sum separately for late-type (blue) and early-type (red) galaxies. In the case of Gabasch et al. (2004), since the measurements did not consider the division to galaxy types, we will only consider the total LF.
In Figure 2, we show the LF of galaxies at redshifts out to 3.5 from Giallongo et al. (2005). In Figure 2(a), we present a comparison to a model of the $`z=0`$ LF of 2dFGRS data from Cooray (2005). In Figure 2(b), we assume no redshift evolution in the $`L_c(M,z)`$ relation. The resulting LFs are then affected only by the redshift evolution of the dark matter halo mass function, $`dn/dM(z)`$. At low redshifts, the dark matter halo mass function evolves such that the number density of massive halos is rapidly decreasing while the density of halos at masses around and below $`10^{12}`$ M$`_{}`$ is slightly increasing relative to the mass function at $`z=0`$ (see, Reed et al. 2003). This evolution in the halo mass function is directly reflected on the high redshift galaxy luminosity function. Given that the $`L_c(M,z=0)`$ relation flattens at mass scales $``$ 10<sup>13</sup> M$`_{}`$ (Cooray & Milosavljeviฤ 2005a), the rapid decline in the number density of massive halos, corresponding to groups and clusters, does not lead to the same fractional decline in the bright-end of the galaxy LF relative to values today. Models based on the mass function evolution alone, however, suggest a decline in the density of bright galaxies at high-redshifts when compared to densities measured today.
On the other hand, the observed galaxy LF at high redshifts indicates that the bright-end density is, in fact, increasing as one moves to $`z3`$ from today. To compensate for the decline in the number density of dark matter halos that host galaxies, associated with the redshift evolution of the mass function, the only possibility to increase the density of luminous galaxies is to consider positive luminosity evolution; The galaxies must brighten at high redshifts relative to luminosity values today. This brightening can be accomplished in several ways. First, galaxies, regardless of the host halo mass, can brighten by a constant factor; This description can be considered as a scenario involving pure luminosity evolution. We consider this possibility by scaling the $`L_c(Mz,z=0)`$ relation by $`(1+z)^\alpha `$ with $`\beta `$ and $`\eta `$ both set to zero. With $`\alpha =0.75`$, our model descriptions are plotted in Figure 3. For comparison, we also plot the measured LFs by Giallongo et al. (2005) where we show the total sample as well as the division to galaxy types. The $`L_c(M,z)`$ relations related to this pure luminosity evolution scenario is shown in Figure 1(a). The model can provide an adequate description, though one underestimates the density of most luminous galaxies shown with the bright-end data points in Figures 3(b), (c) and (d), while overestimating the faint-end density, for example, in Figure 3(c). To obtain a better fit to the bright-end density, one can increase $`\alpha `$. This, however, comes at the expense of overestimating the number density of galaxies at the faint-end further.
A better description of the Giallongo et al. (2005) data may be that the luminosity evolution is mass dependent. In Figure 1(b), we plot $`L_c(M,z)`$ relations under this alternative description, where we allow the luminosity of halos above the flattening mass scale to grow rapidly with redshift. This is accomplished with non-zero values for two parameters $`\beta `$ and $`\eta `$, with $`\alpha =0`$, though, alternative model descriptions that increase luminosities of central galaxies in massive dark matter halos, while keeping the luminosity at the low-end of the mass distribution essentially the same as today, can also be considered. As shown in Figure 4, where we also plot measurements of Giallongo et al. (2005), with the mass-dependent luminosity evolution description, the density of bright galaxies is increased while keeping the faint-end density similar to values measured out to $`z3`$. It is likely that this model provides a more accurate description of the data, though given various uncertainties in LFs out to a redshift of 3.5, we cannot distinguish reliably between the mass-dependent luminosity evolution and the pure luminosity evolution description from this data set alone.
### 3.1 $`L_c(M,z)`$ evolution through model fits to data
While Giallongo et al. (2005) data alone may not allow us to differentiate between the mass-dependent luminosity evolution and the pure luminosity evolution the DEEP2 (Willmer et al. 2005), COMBO-17 (Bell et al. 2004) and Gabasch et al. (2004) LFs may provide adequate statistics to consider detailed model fits in combination. For this purpose, we describe the $`L_c(M,z)`$ relation as
$$L_c(M,z)=L_0(1+z)^\alpha \frac{(M/M_1)^a}{[b+(M/M_1)^{cd(1+z)^\beta }]^{1/d}},$$
(5)
where we have combined the two model descriptions to one by allowing variations in $`\alpha `$, pure luminosity evolution, and $`\beta `$, mass-dependent evolution, but setting $`\eta =1`$ as $`L_c(M,z)`$ relation is not strongly sensitive to $`\eta `$ given the appearance of $`d`$ and $`1/d`$ in the numerator of equation (5). We vary $`\alpha `$ and $`\beta `$ and consider likelihood model fits to LFs from DEEP2, COMBO-17, and Gabash et al. (2004); given that DEEP2 and COMBO-17 probe the same redshift range, the models are distinguished mostly based on a comparison of those datasets with Gabasch et al. (2004) B-band LFs out to $`z5`$.
In Figure 5, we summarize our results in the $`(\alpha ,\beta )`$ plane where we show the $`1\sigma `$, $`2\sigma `$, and $`3\sigma `$ allowed values of these parameters based on a comparison to DEEP2/COMBO-17 LFs and Gabasch et al. (2004) LFs separately as well the combined constraints. When model fitting to the data, in the case of DEEP2 and COMBO-17 LFs, we only consider the total LF and do not use the LFs of galaxy types; In an upcoming paper, we will consider an analysis of galaxy statistics related to types including an analysis of galaxy type-dependent clustering at $`z0`$ and the redshift dependence based on galaxy type LFs from DEEP2 and COMBO-17. While DEEP2/COMBO-17 allow $`\beta >0`$, in combination with high redshift LFs, we can establish that at the 3 $`\sigma `$ level, $`\beta <0`$ with a preferred value around -0.1; As shown in Figure 1(b), $`\beta <0`$ leads to an increase in the luminosities of central galaxies at large halo masses while the luminosities at low masses remain the same. While $`\alpha =0`$ is consistent with the data, the constraint in the $`(\alpha ,\beta )`$ plane lies in a degeneracy line that allow for large negative values for $`\alpha `$, leading to a decrease in the luminosity of galaxies under the $`L_c(M,z)`$ relation, compensated by the increase in luminosity in the $`L_c(M,z)`$ relation associated with negative values for $`\beta `$. In either case, our models suggest one strong conclusion: while pure-luminosity evolution is consistent, the high redshifts LFs require halo mass-dependent evolution in the $`L_c(M,z)`$ relation.
In Figs. 6 and 7 we show the DEEP2 and COMBO-17 LFs and a comparison to model fits based on these constraints, respectively. Here, for comparison, we also show the LF of galaxy types. In addition to the best fit model parameters for $`\alpha `$ and $`\beta `$ based on the total sample, we also show the total LF related to model fits using these data alone. The difference here is minor since out to $`z1`$, the difference between overall best-fit parameters $`\alpha `$ and $`\beta `$ and the same parameters that best describe these datasets best individually is minor. In Figure 8, we consider the LF out to $`z5`$ from Gabasch et al. (2004). Here, again, we show the overall best-fit model as well the model that describe Gabasch et al. (2004) LFs best. In this case, the difference is significant since small variations in $`\alpha `$ and $`\beta `$ can lead to large differences at high redshifts. For example, the overall best fit model, does not describe adequately the Gabasch et al. (2004) LF at the highest redshift (results in an underestimate of the number density of galaxies at the bright-end), though, this difference is accounted by the best fit model preferred by these data alone.
While our models generally suggest that the central galaxy luminosity increases with redshift, we do not find strong evidence for a pure-luminosity increase in the B-band, but rather a mass dependent increase. It will be useful to see if the same increase repeats in other bands, such as the in K-band. Observations suggest that in such redder bands, galaxy luminosities increase first with increasing redshift (passive evolution), but then begin to decrease beyond a particular redshift that correspond to the formation epoch of galaxies where most stellar mass was assembled (e.g., Drory et al. 2003). In our models, the formation of galaxies are strictly linked to the assembly of dark matter halos. In the mass-dependent luminosity evolution description, galaxies in dark matter halos below a certain characteristic halo mass scale, M$`{}_{c}{}^{}10^{11}`$ M$`_{}`$ in Figure 1(b), do not increase in luminosity with an increase in the redshift. Thus, a a decrease in galaxy luminosity beyond a certain redshift could be associated with the transition where average halo mass at that epoch, given the halo mass function corresponding to that redshift, becomes less than $`M_c`$. While we have not seen clear indications for such a transition in the B-band high-z LFs, we plan to investigate this possibility in detail in near future using high-z LFs in redder bands (e.g., Drory et al. 2003), where the $`L_c(M,z)`$ relation may become more sensitive to the mass-dependent luminosity evolution needed to explain the observations.
In addition to the redshift evolution of the total galaxy LF, the model description built using Cooray (2005) model, but with an accounting of the redshift evolution in the $`L_c(M,z)`$ relation is in good agreement with high redshift LFs of galaxy types. Note that we have not allowed for a redshift variation in the fraction of red and blue galaxies, as a function of the halo mass. As shown in Figures 3 and 4 in the case of Giallongo et al. (2005) data and Figures 6 and 7 with DEEP2 and COMBO-17 data, these models generally overestimate the LF of red galaxies at the faint-end. At $`z>1`$, the fraction of blue galaxies increases; in Cooray (2005), the late-type (blue) fraction of central galaxies increases at low-mass halos. Thus, at high redshifts, with the rapid disappearance of halos with masses above few times $`10^{13}`$ M$`_{}`$, the blue fraction begins to dominate. This, however, does not mean that all galaxies at redshifts out to 3 or so is late-types. We certainly expect $``$ 50% of bright galaxies at redshifts $``$ 3 to be early type, red galaxies. The exact observed fraction of late- vs. early-type galaxies, as a function of redshift, can eventually be used to update this model and, especially, to understand whether there is a redshift evolution in the galaxy type fraction relative to values seen today.
### 3.2 $`z6`$ LF
While there is no measured LF of galaxies at $`z6`$ in the rest B-band, ignoring complications resulting from differences in the rest wavelength, we can also compare our model predictions with the UV LF at $`z6`$ from Bouwens et al. (2004a). In Figure 9, we summarize our results. Note that the pure luminosity evolution scenario over predicts the LF at all luminosities of interest, while the mass-dependent luminosity evolution provides a reasonable description of the data; note that any agreement or disagreement between these models and measurements must be considered with the difference in color in mind. This is due to the fact that, while the $`L_c(M,z=0)`$ relation is constructed as appropriate for the B-band, the measurements at $`z6`$ is in the rest UV band.
Given differences in model predictions compared to measurements in Figure 9, however, we believe that the two luminosity evolution descriptions, pure luminosity evolution at all mass scales or mass-dependent luminosity evolution only at the high mass end, may be distinguished with $`z6`$ LFs, though these two models give equally acceptable description of LFs out to a redshift of 3.5. Thus, based on the UV LF at $`z6`$, we suggest that a favorable description of the high-redshift galaxies may be the evolution of luminosities based on the host halo mass. As shown in Figure 9, the $`z6`$ LF is clearly dominated by late-type blue galaxies; However this does not mean that there are no early-type galaxies at these high redshifts. The early-type red galaxies will only appear, in the statistical sense, at the bright-end with a density close to that of blue galaxies. At the faint-end, for each red-type galaxy, one should statistically expect a factor of 5 to 10 more blue galaxies, depending on the luminosity. The existence of bright and red galaxies at redshifts $``$ 6, as found with Spitzer (Yan et al. 2005; see, review in Stark & Ellis 2005) does not contradict these models, as we do expect early-type galaxies to be present even at these redshifts. While currently the observed sample is a few galaxies, a precise determination of their number density, or the luminosity function, could strongly constrain the redshift evolutionary properties of the early-to-late type fraction, which we have assumed to be a constant in our models, so far.
### 3.3 $`z3`$ to 6 halo masses from LFs
Given that our model for the LF is constructed from CLFs โ the number of galaxies as a function of the halo mass โ we can directly address an important question as to what mass dark matter halos host galaxies seen at $`z3`$ to 6. We show the mass dependence of the LF in Figure 10. At $`z3`$, galaxies that are brighter than $`M_{\mathrm{AB}}`$ -22 are hosted in dark matter halos with masses in the range of $`10^{12}`$ to $`10^{13}`$ $`h_{70}^1`$ M$`_{}`$. To statistically detect dark matter halos at $``$ $`10^{11}`$ M$`_{}`$, one must study the LF down to an absolute magnitude of -18. In the case of $`z6`$ galaxy LF, all galaxies in the luminosity range corresponding to absolute magnitudes between -22 and -19 are hosted in dark matter halos between $`10^{11}`$ to $`10^{12}`$ $`h_{70}^1`$ M$`_{}`$; the bright-end of the $`z6`$ LF corresponds to the upper-end of this mass range.
While these are approximate mass ranges, using CLFs, we can quantify the mass distribution of $`z3`$ to $`6`$ galaxies exactly. Here, we calculate the conditional probability distribution $`P(M|L,z)`$ that a galaxy of a given luminosity $`L`$ at redshift $`z`$ is in a halo of mass $`M`$ (Yang et al. 2003b, Cooray 2005):
$$P(M|L,z)dM=\frac{\mathrm{\Phi }(L|M,z)}{\mathrm{\Phi }(L,z)}\frac{dn(z)}{dM}dM.$$
(6)
In Figure 11, we summarize our results, where we plot probabilities at luminosities that correspond to absolute magnitudes of -18 to -24. At the bright end of $`M_{\mathrm{AB}}=24`$, at $`z3`$, galaxies are primarily in dark matter halos of mass $``$ 10<sup>13</sup> M$`_{}`$. In comparison, such galaxies are central galaxies in groups and clusters today with masses above 10<sup>14</sup> M$`_{}`$. At $`z6`$, $`M_{\mathrm{AB}}=24`$ galaxies are primarily in dark matter halos with mass $`5\times 10^{12}`$ M$`_{}`$. Similarly, $`M_{\mathrm{AB}}=18`$ galaxies at $`z3`$ are found in dark matter halos with mass 10<sup>11</sup> M$`_{}`$, though one finds a few percent probability that some of these galaxies are satellites of dark matter halos with masses between 10<sup>12</sup> M$`_{}`$ and 10<sup>13</sup> M$`_{}`$. The four panels, when combined, show the mass-dependent redshift evolution of the galaxy luminosity. Luminous galaxies at high redshifts are found at lower mass halos than dark matter halo masses that corresponds to the same galaxy luminosity today. At the faint-end, $`M_{\mathrm{AB}}>20`$, regardless of the redshift, faint galaxies are essentially found in dark matter halos with a similar range in mass, though at low redshifts, a 30% or more fraction of low-luminous galaxies could be satellites in more massive halos.
While we have simply used the LF to establish the mass scale of $`z3`$ to 6 galaxies, previous attempts have also been made to establish the dark matter halo masses associated with these galaxies. These estimates on halo masses were primarily based on observed galaxy clustering at these redshifts (e.g., Steidel et al. 1998; Bullock et al. 2002; Moustakas & Somerville 2002). To compare with these halo mass estimates, which essentially lead to an estimate of the average halo mass, we calculate the probability distribution of halo mass associated with high-redshift galaxies:
$$P(M|z)dM=\frac{_{L_{\mathrm{min}}}๐L\mathrm{\Phi }(L|M,z)}{_{L_{\mathrm{min}}}\mathrm{\Phi }(L,z)}\frac{dn(z)}{dM}dM,$$
(7)
with the low-end of luminosity integral set at $`L_{\mathrm{min}}`$. For example, to compare with Moustakas & Somerville (2002) estimate on the average $`z3`$ Lyman Break Galaxy (LBG) halo mass, we set $`L_{\mathrm{min}}`$ to be that corresponding to $`M_{\mathrm{AB}}20.8`$; at this magnitude level, the number density of $`z3`$ galaxies is $`5\times 10^3`$ h<sup>3</sup> Mpc<sup>-3</sup>, comparable to the number density of LBG galaxies used in Moustakas & Somerville (2002) together with clustering statistics (correlation length and bias) of galaxies down to this density.
Figure 12 shows the probability distribution of mass related to this number density of galaxies at $`z3`$. The probability peaks around $``$ 7 $`\times 10^{11}`$ h<sup>-1</sup> M$`_{}`$ with the 1$`\sigma `$ range of $`(4`$$`21)\times 10^{11}`$ h<sup>-1</sup> M$`_{}`$, which is consistent with the average halo mass of $`5.5\times 10^{11}`$ $`h^1`$ M$`_{}`$ suggested in Moustakas & Somerville (2002) based on detailed halo model descriptions of the bias factor; A straight forward description of the bias, assuming a single LBG per each dark matter halo, leads to a halo mass of $`8\times 10^{11}`$ h<sup>-1</sup> M$`_{}`$ (Adelberger et al. 1998), with mass values similar to this by others (Baugh et al. 1998; Giavalisco & Dickinson 2001). These values, based on the two-point correlation function, are in good agreement with the mass estimate here based on model fits to the $`z3`$ LF.
This agreement between mass estimates from two-point clustering statistic and the one-point LF of the LBG distribution suggests that the LBG distribution follows the standard hierarchical clustering model. Motivated by a slight discrepancy in the mass estimates from the clustering argument and an estimate based on the velocity dispersion of spectroscopic lines (Pettini et all. 2001), Scannapieco & Thacker (2003) suggested a modified model for the clustering of LBGs that involve a time-dependent correction associated with the merging history (also, Furlanetto & Kamionkowski 2005). The agreement between mass estimates suggested here argues against an additional correction to the bias factor. Recent estimates of the LBG halo masses based on H$`\alpha `$ spectroscopic observations suggest that LBG halo masses at $`z3`$ is $`M>3\times 10^{11}`$ M$`_{}`$ (Erb et al. 2003), in good agreement with previous clustering based estimates as well as the estimates based here from the LF. While the approach based on galaxy clustering allows the average halo mass scale to be established, while velocity dispersion measurements only lead reliably to a lower limit on the mass scale, making use of CLFs, here, we have quantified the dark matter halo mass of $`z3`$ and 6 galaxies as a probability distribution function. This probability distribution function captures not only the average mass function, but also the dispersion related to the average mass. In addition to the average halo mass, a proper measurement of this dispersion must be considered before discrepancies are highlighted.
### 3.4 $`z>3`$ UV LFs
In Figure 13, we compare our predictions with several measurements in the literature on the UV LF at high redshifts. At $`z3`$, Lyman-break galaxy (LBG) LF is measured by Steidel et al. (1999). We find reasonable agreement, though we emphasize that our models are constructed for the rest B-band instead of rest UV-band related to these observations. For reference, we also show the LF of galaxies at redshifts 8 and 10; while no measurements currently exist at these high redshifts, the agreement, at least out to $`z6`$ suggests that our predictions may be directly testable in the near future using deep IR images in near-IR wavelengths. Relative to $`z6`$ LF, the $`z10`$ LF predicts a factor of $``$ 10 lower number density of galaxies at $`M_{\mathrm{AB}}20`$. Based on $`i`$-band dropouts, the surface density of $`z6`$ galaxies, is roughly $`0.5\pm 0.2`$ galaxies per square arcmin. To detect a $`z`$ 10 galaxy, it may be that one must search over an area of $``$ 15 to 30 square arcmins, on average. The strong clustering of high-z galaxies, discussed below, may affect search for these galaxies.
In Figure 14, for comparison with existing measurements, we plot the cosmic luminosity density as a function of redshifts. These luminosity densities are calculated via $`\rho _L(z)_{L_{\mathrm{min}}}L\mathrm{\Phi }(L,z)๐L`$, where we consider two values for the minimum luminosity. In Figure 14(a), we consider a fainter cut off, at $`M_{\mathrm{AB}}16`$, and compare with measurements of the luminosity density at rest B-band from the literature. Our predictions generally agree at low redshifts, though, over predicts the density measured by Dahlen et al. (2005) using LFs constructed from GOODS data. In Figure 14(b), we set the low luminosity end of the integral to be roughly $`0.3L_{}`$ ($`M_{\mathrm{AB}}20`$) at $`z3`$ to be consistent with most measurements by Bouwens et al. (2004a, 2004b, 2005). In this panel, we compare with measurements of the luminosity density at high redshifts in the rest UV-band. Note that our underlying model here is designed for rest B-band and we do not attempt to include any corrections due to color differences when comparing with observations at rest UV wavelengths. Our models suggest that the luminosity density at $`z`$ 10 should be an order of magnitude below what is suggested in Bouwens et al. (2005), under the assumption that they detect 3 $`z10`$ dropouts in deep HST NICMOS fields in a search area over $`z15`$ square arcmins. Based on our LFs, we expect at most a single $`z10`$ galaxy in such a small survey area. If a significant density of $`z10`$ galaxies were to exist, one would require a sharp increase in the evolution of galaxy luminosities at halo mass scales around $`10^{11}`$ M$`_{}`$ than the luminosity evolution we have suggested so far to explain $`z3`$ to 6 galaxy LFs.
### 3.5 Galaxy Bias
While our models can generally describe the LF of galaxies at redshifts out to 6, another useful quantity to compare with observed data is the galaxy bias, as a function of the luminosity. Using the conditional LFs, we calculate the luminosity and redshift dependent galaxy bias as
$$b(L,z)=b_{\mathrm{halo}}(M,z)\frac{\mathrm{\Phi }_i(L|M,z)}{\mathrm{\Phi }_i(L,z)}\frac{dn}{dM}(z)๐M,$$
(8)
where $`b_{\mathrm{halo}}(M,z)`$ is the halo bias with respect to the linear density field (Sheth, Mo & Tormen 2001; also, Efstathiou et al. 1988; Cole & Kaiser 1989; Mo et al. 1997) and $`i`$ denotes the galaxy type.
In Figure 15, we show the galaxy bias as a function of the luminosity. We also divide the sample to galaxy types and redshifts at $`z=0`$, 3 and 6. At redshift of 3, we plot several estimates of the LBG bias; we convert the biasโnumber density relations, from e.g., Bullock et al. 2002, to plot bias as a function of luminosity based on the expected number density of $`z3`$ galaxies, down to the given luminosity, given our model description for the LF. The $`b(L)`$ relation provides a more direct approach to compare how galaxy bias evolves with redshift, than using the bias factor as a function of the number density, though the latter is what is measured from the data; With adequate statistics, in fact, it should be possible to measure $`b(L)`$ relation at high redshifts directly from the data as has been at $`z=0`$ with SDSS (Zehavi et al. 2004) and with 2dFGRS (Norberg et al. 2002b). As shown in Figure 11, the $`z6`$ galaxies are biased by factors of $``$ 3 or higher relative to the linear density field. While there are no published measurements of galaxy clustering at $`z6`$, the existence of clustered large-scale structures has been noted through searches for Lyman-$`\alpha `$ emitters at these redshifts (e.g., Ouchi et al. 2005; Wang et al. 2005). While our models were developed to discuss the LF and clustering of galaxies, one can easily modify the current prescriptions to describe statistics of sources such as Lyman-$`\alpha `$ emitters; we plan to model Ly-$`\alpha `$ galaxy statistics at redshifts 4 to 7 in an upcoming paper.
## 4 Summary and Conclusions
To summarize our discussion involving high-redshift galaxy LFs, our main results are:
(1) Galaxy luminosities must evolve with redshift; while to describe rest B-band LFs out to $`z3`$, one can consider either a mass-independent evolution, with luminosities increasing as $`(1+z)^\alpha `$, or a mass-dependent evolution scenario (see, Figure 1b); galaxy LFs from DEEP2 (Willmer et al.2005), COMBO-17 (Bell et al. 2004) and Gabasch et al. (2004) are more compatible with a mass-dependent evolution model with the model parameter $`\beta <0`$ at the 3$`\sigma `$ confidence level. In this scenario, galaxies that are present in halos above $`10^{12}`$ M$`_{}`$ brighten by factor of 4 to 6 between now and $`z6`$. This suggests that the star formation rate, per given dark matter halo mass, was increasing to high redshifts and that the star formation was more efficient in the past relative to low redshifts. A conclusion similar to what we generally suggest here was also reached by Dahlen et al. (2005) based on the rest B-band LF constructed from GOODS data.
(2) While the bright-end density of galaxies increases with redshift, the faint-end density remains essentially the same out to $`z3`$. The suggested evolution in then $`L_c(M)`$ relation with redshift compensates for the decrease in dark matter halo density. Another important reason why the number density of faint galaxies does not decrease rapidly out to a $`z2`$ is that $`L_c(M)`$ relation flattens at a halo mass scale around $`10^{13}`$ M$`_{}`$; at low redshifts, the redshift evolution of the dark matter halo mass function mostly results in a decrease in the number density of massive halos, while the number density around $`10^{13}`$ M$`_{}`$ is not affected. At $`z>2`$, the exponential cut-off associated with the halo mass function moves to mass scales around and below $`10^{13}`$ M$`_{}`$. This leads to a decrease in the number density of halos hosting central galaxies over the range in luminosity of interest. Thus, the turn over in the LF at the bright-end between redshifts 2 to 3 is a reflection on the lack of an adequate density of dark matter halos with masses around $`10^{13}`$ M$`_{}`$ to host bright galaxies. One does not need to invoke different or multiple scenarios to explain why the density of bright galaxies first increases and then decreases at redshifts greater than 2 to 3. Even in the presence of a decrease in the number density of bright galaxies as a function of redshift, the suggested luminosity evolution continues to be present.
(3) The mass scale of galaxies at redshifts $``$ 3 is distributed between $`4\times `$ 10<sup>11</sup> M$`_{}`$ to $`2\times `$ 10<sup>12</sup> M$`_{}`$ (Figure 12). The brightest galaxies at $`z3`$ (with $`M24`$) are found in dark matter halos with masses slightly below 10<sup>13</sup> M$`_{}`$. Today, these galaxies are found in dark matter halos with masses around 10<sup>15</sup> M$`_{}`$, or as central galaxies in massive clusters. At a redshift of $`6`$, these galaxies were in dark matter halos of mass few times 10<sup>12</sup> M$`_{}`$. The $`z3`$ halo masses based on the one-point LF, as determined with model fits here, agree with previous estimates in the literature based on the two-point correlation function and the bias factor of these galaxies (e.g., Adelberger et al. 1998; Moustakas & Somerville 2002) as well as recent estimates of the halo masses for $`z3`$ LBGs based on spectroscopic observations (Erb et al. 2003) and that a disagreement alluded in the literature (e.g., Scannapieco & Thacker 2003; Furlanetto & Kamionkowski 2005) is not present. These agreements, especially between the LF and galaxy bias, also suggest that no major correction to the bias factor of LBGs is necessary from additional effects such as merging (though such corrections may be necessary to interpret other galaxy samples).
(4) The number density of $`z10`$ galaxies are roughly a factor of 10 lower than the density of galaxies at $`z6`$; if a higher density is found, the $`L_c(M,z)`$ relation must evolve rapidly at redshifts above 6 than suggested here to explain $`z3`$ to 6 LFs. If our predictions our correct, one must search roughly an area of $``$ 30 sqr. arcmins to find a galaxy at a redshift of 10. The high redshift galaxies are extremely biased with respect to the linear dark matter density field; the bias factor of galaxies at $`z6`$ varies from less than 3 to 10 for brightest galaxies. Our predictions for clustering bias factors, as a function of luminosity, are in general agreement with estimated values based on the correlation length (Figure 11), at least out to $`z3`$.
To conclude, we have presented a description of the LF of galaxies as a function of redshift, in the rest B-band, using a simple empirical model. The approach has the main advantage that one can extract underlying reasons associated with the observed evolution of the LF. Here, we have characterized this evolution in terms of the $`L_c(M,z)`$ relation โ the luminosity of central galaxies of dark matter halos as a function of halo mass and redshift. Given the rapid decline in the number density of dark matter halos, in order to explain the presence of bright galaxies, we have suggested that the $`L_c(M,z)`$ relation must evolve such that galaxies have a higher luminosity, when compared to today, at the high mass end of the halo distribution.
The $`z=0`$, $`L_c(M)`$ relation was explained in Cooray & Milosavljeviฤ (2005a) in terms of dissipationless merging of galaxies in dark matter halos centers. The flattening of the $`L_c(M)`$ relation at halo mass scales $`10^{13}`$ M$`_{}`$ was suggested as due to a decrease in the efficiency at which dynamical friction decays satellite orbits to merge with the central galaxy; the dynamical friction time scale for the merging of satellites in halos with masses above few times $`10^{13}`$ M$`_{}`$ is more than Hubble time. As one moves to a high redshift, the flattening mass scale is expected to move to a lower mass scale given that the age of the Universe, or the age of the halo since it formed, is lower than the age today. Such a behavior is captured in our empirical model for the $`L(M,z)`$ relation based on mass-dependent luminosity evolution. While our approach has allowed us to capture the evolution associated with the $`L_c(M,z)`$ relation, what physical processes govern this behavior is yet to be understood. Any reasonable explanation for the brightening of galaxies at high redshifts must consider the merging history in addition to the underlying astrophysical reason why galaxies were brighter in the past when compared to today. The $`L_c(M,z)`$ relation is certainly a key aspect in galaxy formation and evolution, and we expect approaches involving numerical simulations (e.g., Kay et al. 2002) and semi-analytic models (e.g., Benson et al. 2003) will be used to understand if what we have suggested on the evolution of this relation is consistent or not with astrophysical processes associated with galaxy formation and evolution.
Acknowledgments: The author thanks Emanuele Giallongo, Chris Willmer, and Armin Gabasch for electronic tables of the high-redshift LFs described in this paper and members of Cosmology and Theoretical Astrophysics groups at Caltech and UC Irvine for useful discussions. This paper was completed while the author was at the Aspen Center for Physics in Summer of 2005.
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# Fixing Two-Nucleon Weak-Axial Coupling ๐ฟ_{1,๐ด} From ๐โปโข๐ Capture
## Abstract
We calculate the muon capture rate on the deuteron to next-to-next-to-leading order in the pionless effective field theory. The result can be used to constrain the two-nucleon isovector axial coupling $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{L}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}}`$ to $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\pm }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}`$ fm<sup>3</sup> if the muon capture rate is measured to 2% level. From this, one can determine the neutrino-deuteron break up reactions and the $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{p}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{p}$}`$-fusion cross section in the sun to a same level of accuracy.
preprint: DOE/ER/40762-322preprint: MU-PP#05-006
The strong evidence of neutrino oscillations observed at the Sudbury Neutrino Observatory (SNO) SNO is based on detecting the <sup>8</sup>B solar neutrino flux through the following three reactions:
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$$\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\nu }$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{x}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{+}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{e}$}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\nu }$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{x}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{+}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{e}$}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{ES}}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{.}$}$$
(1)
The charged current (CC) reaction involves only the electron neutrinos, while the neutral current (NC) reaction and elastic scattering (ES) involve all the active neutrinos $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\left(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{x}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{e}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mu }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\tau }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\right)}$}`$. The $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\nu }$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{e}$}`$ and $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\nu }$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{x}$}`$ fluxes are found to be significantly differentSNO . Further detailed measurements of the fluxes could sharpen the constraints to neutrino oscillation parameters and provide precision tests to the standard solar model SSM . However, while the ES cross section is known to high accuracy, the CC and NC cross sections have hadronic uncertainties. As shown by Butler, Chen and Kong BCK , the dominant uncertainties in low energy CC and NC cross sections comes from the coupling of a two-body isovector axial current, $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{L}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}}`$, in pionless effective field theory ($`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{EFT}}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\pi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{/}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}`$). The potential model results of Refs. potentialmodel1 and potentialmodel2 can be reproduced by different choices of $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{L}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}}`$, indicating that the $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{5}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\%}$}`$ difference between the models comes from the different assumptions about the short-distance nuclear physics. There are other interesting weak reactions involving the same two-body current, for example, the $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{p}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{p}$}`$ and $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{p}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{e}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{p}$}`$ fusion processes ($`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{p}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{p}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{d}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{e}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{+}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\nu }$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{e}$}`$, $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{p}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{p}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{e}$}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{d}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\nu }$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{e}$}`$) which power the sun fusion1 ; fusion2 . It is one of the great current interests to measure these neutrino fluxes to further test the standard solar model.
Recently much effort has been going into determining the effective two-body axial current interaction fixing1 . Butler, Chen, and Vogel attempted to fix $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{L}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}}`$ from reactor antineutrino-deuteron breakup reactions, and they found $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{L}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{3.6}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\pm }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{5.5}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{fm}}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{3}$}`$ fixing1 . Chen, Heeger, and Robertson obtained $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{L}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{4.0}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\pm }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{6.3}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{fm}}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{3}$}`$ by using SNOโs CC and NC data, calibrated by the ES events of SNO and Super-Kamiokande(SK) fixing2 . Schiavilla et al.โs idea fixing3 of using the tritium $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\beta }$}`$ decay rate to control the strength of the two-body current was adopted by Park et al. in their hybrid EFT calculation fix dr , and the pp fusion rate was predicted with a small error. When compared with the $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{EFT}}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\pi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{/}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}`$ calculation fusion2 , their result yields $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{L}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{4.2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\pm }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2.5}$}`$fm$`^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{3}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{.}$}`$
In this paper we aim to make a high-precision determination of $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{L}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}}`$ from the $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mu }$}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{d}$}`$ capture process
$$\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mu }$}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{+}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{d}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\nu }$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mu }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{+}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{n}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{+}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{n}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}$$
(2)
by calculating the rate to next-to-next-to leading order (N<sup>2</sup>LO) in $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{EFT}}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\pi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{/}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}`$. The $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mu }$}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{d}$}`$ capture rate has been measured previously by different groups with rather different results $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{\Gamma }}$}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{exp}}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{470}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\pm }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{29}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{s}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}`$ mu-d-exp1 and $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{\Gamma }}$}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{exp}}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{409}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\pm }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{40}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{s}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}`$ mu-d-exp2 . A measurement of this rate with 1% precision is under investigation at PSI Kammel . An earlier potential model calculation tkk gave $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{\Gamma }}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{397}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{400}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{s}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}`$. More recently, the hybrid approach mentioned above gave $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{\Gamma }}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{386}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\pm }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{5}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{s}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}`$ Ando .
A concern in applying $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{EFT}}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\pi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{/}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}`$ to the $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mu }$}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{d}$}`$ capture is that the energy transfer into the hadronic system might be too large to apply $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{EFT}}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\pi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{/}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}`$. However, as shown in Ref.Ando and also in this calculation, the contribution to the total rate from high-energy neutrons is small, and it is possible to impose a neutron energy cut to isolate the low-energy ($`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{10}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{20}$}`$ MeV) neutron events without significantly increasing the statistical errors fixing1 ; Kammel .
Effective field theory is useful when low and high energy scales in the problem are widely separated. For low-energy processes, short-distance physics can be taken into account by local operators in an effective lagrangian involving only low-energy degrees of freedom. For $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\nu }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\overline{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\nu }$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{d}$}`$ scattering with neutrino energy below 20 MeV and $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mu }$}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{d}$}`$ capture with small final-state neutron energy, the pion and other meson exchanges are not dynamical degrees of freedom, and their physics can be captured by contact interactions involving nucleons and the external currents. To make predictions with controlled precision, calculations are done with the perturbative expansion parameter $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{Q}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{/}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{a}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{n}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{n}$}}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{S}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{0}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\gamma }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{p}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{/}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{\Lambda }}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}`$ which is the ratio of light to heavy scales. The light scales include the inverse S-wave neutron-neutron scattering length $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{/}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{a}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{n}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{n}$}}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{S}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{0}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{10.6}$}`$ MeV in the $`{}_{}{}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{S}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{0}$}}^{}`$ channel, the deuteron binding momentum $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\gamma }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{45.7}$}`$ MeV in the $`{}_{}{}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{3}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{S}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}}^{}`$ channel, and typical nucleon momentum $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{p}$}`$ in the system. The heavy scale $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{\Lambda }}$}`$ is set by the pion mass $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{m}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\pi }$}`$. This $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{EFT}}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\pi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{/}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}`$ (see e.g. pionless ) and its dibaryon version dibaryon1 ; dibaryon2 ; dibaryon3 have been applied successfully to many processes involving the deuteron, including electro-magnetic processes such as Compton scattering $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\gamma }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{d}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\gamma }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{d}$}`$ Compton1 ; Compton2 , $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{n}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{p}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{d}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\gamma }$}`$ relevant to the big-bang nucleosynthesis Nsynthesis1 ; Nsynthesis2 , weak processes such as $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\nu }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{d}$}`$ reactions for SNO physicsBCK , the solar $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{p}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{p}$}`$ fusion process fusion1 ; fusion2 , and parity violating observables pv . For reviews on three-body systems, see 3body .
The effective Lagrangian for the CC weak interaction is given by $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{C}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{C}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\sqrt{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\omega }$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{i}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{G}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{F}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{V}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{u}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{d}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{l}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{+}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mu }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{J}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mu }$}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{/}$}\sqrt{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{+}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{h}}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{.}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{c}}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{.}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}`$ where $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{G}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{F}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1.166}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\times }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{10}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{5}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{GeV}}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}`$ is the Fermi coupling constant and $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\omega }$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{i}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1.024}$}`$ takes into account the inner electroweak radiative correction radiative . $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{l}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{+}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mu }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\overline{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\nu }$}}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mu }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\gamma }$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mu }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\gamma }$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{5}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mu }$}`$ is the leptonic current. The quark current $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{J}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mu }$}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{V}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mu }$}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mu }$}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{V}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mu }$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mu }$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{i}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{V}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mu }$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mu }$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}`$ contains both vector and axial-vector interactions, where the superscripts 1 and 2 are the isospin indices. At the scale relevant to nuclear physics, the quark current need be matched to a hadronic current which in general contains one-nucleon, two-nucleon, etc., operators.
Up to the order of our interest, the one-nucleon isovector vector and axial vector currents are
$`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{V}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{0}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{a}$}}`$ $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}`$ $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{N}$}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}}{\displaystyle \frac{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\tau }$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{a}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\left(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{+}$}{\displaystyle \frac{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{6}$}}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{r}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{c}$}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{I}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}}\overline{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\right)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{N}$}`$
$`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{V}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{k}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{a}$}}`$ $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}`$ $`{\displaystyle \frac{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{i}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{M}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{N}$}}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{N}$}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}}\stackrel{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{k}$}{\displaystyle \frac{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\tau }$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{a}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{N}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}{\displaystyle \frac{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mu }$}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{M}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{N}$}}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{ฯต}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{k}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{i}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{j}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{N}$}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\sigma }$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{i}$}\overline{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{j}$}{\displaystyle \frac{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\tau }$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{a}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{N}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}`$
$`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{0}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{a}$}}`$ $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}`$ $`{\displaystyle \frac{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{i}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{g}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{M}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{N}$}}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{N}$}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}}\stackrel{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\sigma }$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\stackrel{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}}{\displaystyle \frac{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\tau }$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{a}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{N}$}`$
$`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}{\displaystyle \frac{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{4}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{M}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{N}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{N}$}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{G}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{p}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\left(}$}\overline{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\right)}$}\overline{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{0}$}}\stackrel{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\sigma }$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\overline{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}}{\displaystyle \frac{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\tau }$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{a}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{N}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}`$
$`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{k}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{a}$}}`$ $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}`$ $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{g}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{N}$}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\sigma }$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{k}$}{\displaystyle \frac{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\tau }$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{a}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\left(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{+}$}{\displaystyle \frac{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{6}$}}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{r}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{I}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}}\overline{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\right)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{N}$}`$
$`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{+}$}{\displaystyle \frac{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{4}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{M}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{N}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{N}$}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{G}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{p}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\left(}$}\overline{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\right)}$}\overline{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{k}$}\stackrel{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\sigma }$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\overline{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}}{\displaystyle \frac{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\tau }$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{a}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{N}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}`$
where $`\overline{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\stackrel{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{+}$}\stackrel{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}}`$, $`\overline{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{0}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\stackrel{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{0}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{+}$}\stackrel{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{0}$}}`$, and $`\stackrel{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\stackrel{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\stackrel{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{.}$}`$ The superscript $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{a}$}`$ is the isospin index, $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mu }$}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mu }$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{p}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mu }$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{n}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{/}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2.353}$}`$, is the isovector magnetic moment. Isovector Dirac charge radius $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{r}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{c}$}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{I}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{r}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{c}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{p}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{r}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{c}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{n}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{0.873}$}`$fm<sup>2</sup>, and the isovector axial-charge radius $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{r}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{I}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{0.45}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{fm}}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}`$. We have neglected terms of order $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{p}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{/}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{M}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{N}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}`$ or even $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mu }$}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{p}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{/}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{M}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{N}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}`$. The pseudoscalar form factor is, to a good approximation, dominated by the pion-pole $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{G}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{p}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{q}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\frac{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{4}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{M}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{N}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{g}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{M}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\pi }$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{q}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}}`$ whose $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{q}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}`$ dependence will not be expanded because the momentum transfer $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{|}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{๐ช}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{|}$}`$ is of order muon mass $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{M}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mu }$}`$ with low energy final state neutrons. The $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{G}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{p}$}`$ contribution to the axial current is counted of order $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{Q}$}`$.
The lowest dimensional two-nucleon isovector currents, in the dibaryon version of $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{EFT}}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\pi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{/}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}`$, relevant to the $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mu }$}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{d}$}`$ capture process are
$`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{V}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{k}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{a}$}}`$ $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}`$ $`{\displaystyle \frac{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{L}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{d}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{b}$}}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{M}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{N}$}\sqrt{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{r}$}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{3}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{S}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{r}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{n}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{n}$}}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{S}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{0}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}}}}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{ฯต}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{k}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{i}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{j}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{t}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{i}$}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}}\overline{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{j}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{s}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{a}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{+}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{h}}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{.}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{c}}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{.}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}`$
$`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{k}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{a}$}}`$ $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}`$ $`{\displaystyle \frac{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{L}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{d}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{b}$}}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{M}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{N}$}\sqrt{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{r}$}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{3}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{S}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{r}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{n}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{n}$}}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{S}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{0}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}}}}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{t}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{k}$}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{s}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{a}$}`$
$`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{+}$}{\displaystyle \frac{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{G}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{p}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{4}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{M}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{N}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{g}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{t}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{j}$}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}}\overline{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{j}$}\overline{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{k}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{s}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{a}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{+}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{h}}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{.}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{c}}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{.}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}`$
where $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{t}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{i}$}`$ and $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{s}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{a}$}`$ are dibaryon fields for the two-nucleon $`{}_{}{}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{3}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{S}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}}^{}`$ and $`{}_{}{}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{S}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{0}$}}^{}`$ states, respectively. The second term in $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{k}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{a}$}}`$ is induced by the $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{G}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{p}$}`$ term in the one-nucleon current. $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{r}$}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{3}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{S}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1.764}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{fm}}$}`$ and $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{r}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{n}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{n}$}}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{S}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{0}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2.8}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{fm}}$}`$ are the effective ranges in triplet and two-neutron singlet channels, respectively. $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{L}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{d}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{b}$}}`$ and $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{L}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{d}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{b}$}}`$ are coupling constants in dibaryon formalism. The vector current is N<sup>2</sup>LO and its coupling $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{L}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{d}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{b}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{4.08}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{fm}}$}`$ has been determined by the rate of $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{n}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{+}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{p}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{d}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{+}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\gamma }$}`$ near threshold. The axial current is NLO, and its coupling $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{L}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{d}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{b}$}}`$ is proportional to the renormalization-scale-$`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mu }$}`$-independent $`\stackrel{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{~}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{L}$}}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}}`$ in Ref. BCK as $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{L}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{d}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{b}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\frac{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{M}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{N}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\pi }$}}\stackrel{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{~}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{L}$}}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}}`$, through which $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{L}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{d}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{b}$}}`$ is related to the $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mu }$}`$-dependent $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{L}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mu }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}`$ in Ref. BCK . The numerical relation between $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{L}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{d}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{b}$}}`$ and $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{L}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mu }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}`$ is $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{L}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{d}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{b}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{13.8}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{+}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{0.28}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{L}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{M}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\pi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}`$, where $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{L}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{M}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\pi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}`$ is in units of $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{fm}}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{3}$}`$ and has a natural size $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{6}$}`$ fm<sup>3</sup>.
The $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mu }$}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{d}$}`$ atom has a ground state with a hyperfine structure, corresponding to the total angular momentum $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{F}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{3}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{/}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}`$ and $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{F}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{/}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}`$. The $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mu }$}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{d}$}`$ capture process is known to take place almost uniquely from the doublet $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{F}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{/}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}`$ state. The differential capture rate for muon and deuteron in their specific polarization states can be written in terms of leptonic tensor $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{l}$}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mu }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\nu }$}}`$ and hadronic tensor $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{W}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mu }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\nu }$}}`$ as
$$\frac{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{d}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{\Gamma }}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{S}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\widehat{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\xi }$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{d}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{d}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{\Omega }}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\frac{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\omega }$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{i}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{G}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{F}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{|}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{V}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{u}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{d}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{|}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{|}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\psi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{0}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{|}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{32}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\pi }$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{M}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mu }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{+}$}\frac{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{M}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mu }$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{M}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{d}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{l}$}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mu }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\nu }$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{S}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{W}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mu }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\nu }$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\widehat{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\xi }$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}$$
(5)
where the $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{|}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\psi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{0}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{|}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\frac{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\pi }$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\left(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\alpha }$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{em}}$}}\frac{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{M}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mu }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{M}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{d}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{M}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mu }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{+}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{M}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{d}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\right)}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{3}$}`$ is the $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{S}$}`$-state wave-function-at-origin-squared, $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}`$ is the outgoing neutrino energy, and $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{M}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{d}$}`$ is the mass of the deuteron.
The capture rate depends on the polarization vector of the muon $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{S}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mu }$}`$ and the deuteron polarization vector $`\widehat{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\xi }$}}`$. The leptonic tensor is given by $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{l}$}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mu }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\nu }$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{S}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{4}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{k}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mu }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{k}$}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}}{}_{}{}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\nu }$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{+}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{k}$}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}}{}_{}{}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mu }$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{k}$}_{}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\nu }$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{k}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{k}$}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{g}$}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mu }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\nu }$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{+}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{i}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{ฯต}$}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mu }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\nu }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\rho }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\sigma }$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{k}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\rho }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{k}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\sigma }$}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{4}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{M}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mu }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{S}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mu }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{k}$}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}}{}_{}{}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\nu }$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{+}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{k}$}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}}{}_{}{}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mu }$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{S}$}_{}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\nu }$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{k}$}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{S}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{g}$}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mu }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\nu }$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{+}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{i}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{ฯต}$}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mu }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\nu }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\rho }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\sigma }$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{S}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\rho }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{k}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\sigma }$}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}`$, where $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{k}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{M}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mu }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\stackrel{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{0}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}`$ and $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{k}$}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\stackrel{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{k}$}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}`$, with $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{|}$}\stackrel{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{k}$}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{|}$}`$, are the four-momenta of initial muon and final $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\nu }$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mu }$}`$, respectively. The hadronic tensor is
$$\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{W}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mu }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\nu }$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\widehat{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\xi }$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\frac{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\pi }$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{Im}}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\left[}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{d}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{4}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{x}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{e}$}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{i}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{q}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{x}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{T}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{d}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{|}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{J}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mu }$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{+}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{x}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{J}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\nu }$}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{0}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{|}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{d}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\right]}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}$$
(6)
where $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{q}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mu }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{k}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mu }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{k}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mu }$}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}}`$, and $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{|}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{d}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{|}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{d}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{P}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\widehat{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\xi }$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}`$ is the deuteron state with momentum $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{P}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{M}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{d}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\stackrel{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{0}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}`$ and polarization $`\widehat{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\xi }$}}`$.
The diagrams contributing to $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{W}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mu }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\nu }$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\widehat{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\xi }$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}`$ up to N<sup>2</sup>LO are shown in Fig. 1. A straightforward calculation finds
$`{\displaystyle \frac{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{d}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{\Gamma }}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{d}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}{\displaystyle \frac{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{G}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{F}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{|}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{V}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{u}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{d}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{|}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{|}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\psi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{0}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{|}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\pi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{+}$}\frac{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{M}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mu }$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{M}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{d}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}}}{\displaystyle \frac{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\epsilon }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\gamma }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{r}$}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{3}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{S}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}}}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{[}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{F}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}{\displaystyle \frac{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{F}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}`$
$`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{3}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{G}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{V}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{G}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{V}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{G}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\epsilon }$}{\displaystyle \frac{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{4}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{g}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{g}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{M}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mu }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{3}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{M}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\pi }$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{q}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}}}`$
$`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{+}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\epsilon }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\kappa }$}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}}{\displaystyle \frac{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{8}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{g}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{3}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{M}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{N}$}}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{+}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{F}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}{\displaystyle \frac{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{F}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{6}$}}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{+}$}{\displaystyle \frac{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{3}$}}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{F}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{4}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{a}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\left)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\right(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{9}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{G}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}`$
$`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\epsilon }$}{\displaystyle \frac{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{6}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{g}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{M}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mu }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{M}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\pi }$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{q}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{+}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\epsilon }$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}{\displaystyle \frac{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{g}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{4}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\left(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{M}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\pi }$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{q}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\right)}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}`$
$`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{+}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\left(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{10}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{F}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{F}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{+}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{8}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{F}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{4}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{b}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\right)}$}{\displaystyle \frac{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\epsilon }$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\kappa }$}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{3}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{M}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{N}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}}}`$
$`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{+}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\left(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{F}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}{\displaystyle \frac{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{F}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{6}$}}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{+}$}{\displaystyle \frac{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{3}$}}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{F}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{4}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{c}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\right)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{12}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\epsilon }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{g}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\kappa }$}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}}{\displaystyle \frac{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{M}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{N}$}}}`$
$`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\left(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{F}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{+}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{F}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{4}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{c}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\right)}$}{\displaystyle \frac{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{8}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\epsilon }$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{g}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\kappa }$}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{3}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{3}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{M}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{N}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{M}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\pi }$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{q}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}}}`$
$`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{+}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{F}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{5}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\epsilon }$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}{\displaystyle \frac{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{14}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{3}$}}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{g}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{+}$}{\displaystyle \frac{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{16}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{3}$}}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{g}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{+}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{]}$}`$ (7)
where $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{G}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{V}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{+}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\epsilon }$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\frac{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{6}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{r}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{c}$}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{I}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{q}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}`$ is the Dirac form factor, and $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{G}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{g}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{+}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\epsilon }$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\frac{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{6}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{r}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{c}$}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{I}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{q}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}`$ is the axial form factor. The $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\epsilon }$}`$ is formally introduced to keep track of the $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{Q}$}`$ expansion. After expanded in $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\epsilon }$}`$ and truncated at $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\epsilon }$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}`$, the N<sup>2</sup>LO result is obtained by setting $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\epsilon }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}`$. The functions $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{F}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{a}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{b}$}}`$, $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{F}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}`$, and $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{F}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{3}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{a}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{3}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{b}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{3}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{c}$}}`$ are from diagrams $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{a}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}`$, $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{b}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}`$, and $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{c}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}`$, respectively, in Fig. 1
$`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{F}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{a}$}}`$ $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}`$ $`{\displaystyle \frac{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{M}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{N}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{p}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\gamma }$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\pi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{M}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{N}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{\Delta }}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{p}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}}}`$
$`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{F}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{b}$}}`$ $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}`$ $`{\displaystyle \frac{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\gamma }$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\pi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{3}$}}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\left[}$}{\displaystyle \frac{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{M}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{N}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{p}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{\Delta }}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{M}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{N}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{\Delta }}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{p}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{+}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{ln}}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\left(}$}{\displaystyle \frac{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{M}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{N}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{\Delta }}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{p}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{M}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{N}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{\Delta }}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{+}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{p}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\right)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\right]}$}`$
$`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{F}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}`$ $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}`$ $`{\displaystyle \frac{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{4}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\gamma }$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\pi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{\Delta }}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{tanh}}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\left(}$}{\displaystyle \frac{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{p}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{M}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{N}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{\Delta }}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\right)}$}`$
$`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{F}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{3}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{a}$}}`$ $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}`$ $`{\displaystyle \frac{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\gamma }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{M}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{N}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\pi }$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{Im}}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\left\{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{G}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{S}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{0}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{p}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\right\}}$}`$
$`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{F}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{3}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{b}$}}`$ $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}`$ $`{\displaystyle \frac{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\gamma }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{M}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{N}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\pi }$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{Im}}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\left\{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{G}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{S}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{0}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{p}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\right\}}$}`$
$`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{F}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{3}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{c}$}}`$ $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}`$ $`{\displaystyle \frac{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\gamma }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{M}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{N}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\pi }$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{Im}}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\left\{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{G}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{G}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{S}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{0}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{p}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\right\}}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}`$ (8)
with the re-scattering amplitude in the singlet channel as $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{S}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{0}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{p}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\frac{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{4}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\pi }$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{M}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{N}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\frac{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{a}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{n}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{n}$}}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{S}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{0}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\epsilon }$}\frac{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{r}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{n}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{n}$}}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{S}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{0}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{p}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{+}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{i}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{p}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}`$, and functions $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{G}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}}`$ are defined as $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{G}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{tan}}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\frac{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\gamma }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{i}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{p}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{+}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\epsilon }$}\frac{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{4}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{M}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{N}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{g}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{L}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{d}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{b}$}}`$ and $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{G}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{tan}}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\frac{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\gamma }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{i}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{p}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{+}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\epsilon }$}\frac{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{4}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mu }$}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{L}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{d}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{b}$}}`$. The energy injection into the two-neutron system is $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{\Delta }}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{M}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mu }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{M}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{n}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{M}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{p}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}`$. The relative momentum between the two final-state neutrons is $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{p}$}`$, with $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{p}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\sqrt{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{M}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{N}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{\Delta }}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\gamma }$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{/}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{4}$}}`$.
Power counting for the present calculation is rather tricky, and it would be misleading, for example, to just use the $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{p}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{/}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{m}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\pi }$}`$ to estimate the accuracy of the expansion. For example, it is well-know that in the nucleon-nucleon scattering, the effective theory without pion works rather well at the nucleon momentum on the order of 100 MeV, a value close to the pion mass. This can also be seen in the present calculation because the NLO result at large neutron momentum does not completely modify the leading order result, whereas the naive power counting would indicate otherwise. Fortunately, for muon capture, the most of the events occur when the neutron momentum is about half of the muon mass, that is, about 50 MeV.
In Fig. 2 we show the differential rate $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{d}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{\Gamma }}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{/}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{d}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{n}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{n}$}}`$ in terms of the relative motion energy $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{n}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{n}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\sqrt{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{M}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{N}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{+}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{p}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{M}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{N}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}`$ of two final-state neutrons in the region where the $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{EFT}}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\pi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{/}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}`$ calculation is most reliable. It is clear from the figure that the differential rate in the energy region $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{n}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{n}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{10}$}`$ MeV is very small, and is negligible for $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{n}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{n}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{15}$}`$ MeV. By comparing the results of LO, NLO, and N<sup>2</sup>LO, we find good convergence of the expansion.
In the case that a neutron energy cut can be imposed on experimental data Kammel , it is possible to define and measure the integrated capture rate up to a threshold energy $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{n}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{n}$}}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{th}}$}}`$
$$\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{\Gamma }}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{n}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{n}$}}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{th}}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{0}$}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{n}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{n}$}}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{th}}$}}}\frac{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{d}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{\Gamma }}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{d}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{n}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{n}$}}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{๐}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{n}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{n}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{.}$}$$
(9)
The result up to 30 MeV is shown in Fig. 3. In the whole energy region, the NLO contribution is less than $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{10}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\%}$}`$ of the LO contribution, while the N<sup>2</sup>LO contribution is less than $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\%}$}`$. This small size of N<sup>2</sup>LO contribution is accidental and does not happen for unpolarized and $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{F}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{3}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{/}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}`$ rates. For example, the unpolarized rate has the expansion (for $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{L}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{6}$}`$ fm<sup>3</sup>)
$$\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{\Gamma }}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{u}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{n}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{p}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{o}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{l}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{.}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{\Gamma }}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{u}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{n}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{p}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{o}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{l}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{.}$}}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{L}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{O}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{+}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{5.3}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\%}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{+}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{4.9}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\%}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}$$
(10)
where the NLO correction is abnormally small, and NNLO is of the normal size. A similar expansion is obtained for $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{L}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{6}$}`$ fm<sup>3</sup>:
$$\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{\Gamma }}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{u}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{n}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{p}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{o}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{l}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{.}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{\Gamma }}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{u}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{n}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{p}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{o}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{l}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{.}$}}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{L}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{O}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{+}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{10.9}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\%}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{+}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{5.2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\%}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}$$
(11)
which shows a nice convergence pattern. Based on this trend, we assign a 2-3% correction at NNNLO, corresponding to an error in $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{L}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{fm}}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{3}$}`$. This is consistent with the naive estimation of 3% if the small expansion parameter is $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{/}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{3}$}`$. A calculation shows the N<sup>3</sup>LO final state P-wave re-scattering contributes only $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\%}$}`$. Furthermore, the result is insensitive to the uncertainty in $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{a}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{n}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{n}$}}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{S}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{0}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}}`$. Choosing $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{L}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{5.6}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{fm}}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{3}$}`$, the energy dependence of our result matches the previous hybrid calculation very well Ando .
To extract $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{L}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}}`$ from experimental data, it is useful to provide the dependence of the rate on $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{L}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}}`$
$$\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{\Gamma }}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{n}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{n}$}}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{th}}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{a}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{n}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{n}$}}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{th}}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{+}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{b}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{n}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{n}$}}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{th}}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{L}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}$$
(12)
where $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{L}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}}`$ is in unit of $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{fm}}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{3}$}`$. The energy-cut dependent functions $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{a}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{n}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{n}$}}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{th}}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}`$ and $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{b}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{n}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{n}$}}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{th}}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}`$ for a set of $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{n}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{n}$}}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{th}}$}}`$s are listed in the Table 1, from which we observe that, for the whole range of $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{n}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{n}$}}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{th}}$}}`$, the size of $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{b}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{n}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{n}$}}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{th}}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}`$ is about $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1.3}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1.5}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\%}$}`$ of the size of $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{a}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{n}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{n}$}}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{th}}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}`$. This shows how an error in capture rate is translated into an uncertainty of $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{L}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}}`$.
Table 1:
Coefficients functions $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{a}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{n}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{n}$}}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{th}}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}`$ and $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{b}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{n}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{n}$}}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{th}}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}`$ for specific values of two-neutron relative energy $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{n}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{n}$}}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{th}}$}}`$ from the $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{EFT}}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\pi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{/}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}`$ calculation.
| $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{n}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{n}$}}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{th}}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{MeV}}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}`$ | 5.0 | 10.0 | 15.0 | 20.0 |
| --- | --- | --- | --- | --- |
| $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{a}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{n}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{n}$}}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{th}}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{s}}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}`$ | 239.2 | 308.0 | 332.0 | 342.3 |
| $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{b}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{n}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{n}$}}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{th}}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{s}}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{fm}}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{3}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}`$ | 3.3 | 4.2 | 4.7 | 4.9 |
In summary, we calculated the $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mu }$}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{d}$}`$ capture rate using $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{EFT}}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\pi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{/}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}`$ to N<sup>2</sup>LO. The major goal is to fix the two-nucleon isovector axial coupling constant $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{L}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}}`$ from future precision experimental data. An experimental result on the integrated rate up to some neutron energy $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{n}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{n}$}}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{th}}$}}`$ with a $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\%}$}`$ error should be able to, through comparison with our calculation with theoretical error 2-3$`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\%}$}`$, fix the $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{L}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}}`$ with error $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\pm }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2.0}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{fm}}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{3}$}`$. This in turn allows us to determine the neutrino deuteron breakup cross section and the pp fusion rate in the sun to 2-3%.
This work was supported by the U. S. Department of Energy via grant DE-FG02-93ER-40762 and the NSC of Taiwan. We thank S. Ando, D.W. Hertzog, P. Kammel, K. Kubodera, and T.-S. Park for helpful discussions.
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# Accurate polarization within a unified Wannier function formalism
## I Introduction
The representation of the one-particle electronic structure of molecules and solids in terms of localized Wannier wannier orbitals is nowadays โenjoying a revivalโ rmmartin as a useful tool for many applications mv ; hydrogen ; sharma ; cangiani ; fernandez ; posternak ; dovesi ; pavarini ; fabris ; feliciano ; danish . The main impetus for this renewal of interest was given by the establishment, by King-Smith and Vanderbilt (KSV), of a formally exact relationship between the sum of the Wannier function (WF) centers and a gauge-invariant Berry phase, in the context of the modern theory of polarization kingsmith . However, the intrinsic nonuniqueness in the Wannier function definition, and the difficulty in defining their centers within a periodic cell calculation, limited their practical use, until a particularly elegant method due to Marzari and Vanderbilt mv (MV) became available some years ago. Their scheme allows one to obtain, in a given isolated or extended system, a unique set of maximally localized Wannier functions that minimizes a well-defined spread functional. Moreover, the MV formalism provides as an important byproduct the positions of the WF centers, whose sum gives direct access to the macroscopic polarization of the physical system. The MV scheme, which became instantly popular, presents nevertheless an inconvenience, in that crystalline solids are treated on a different footing with respect to the case of, e.g., large disordered systems simulated at the $`\mathrm{\Gamma }`$ point only. The two prescriptions are indeed equivalent in the thermodynamic limit, but they formally differ when discrete Brillouin zone (BZ) samplings (or finite supercells for isolated objects) are used, which is necessarily the case in any practical calculation. In the first part of this work we show that there is nothing fundamental in this discrimination, i.e. that a given choice for a spread functional in $`\mathrm{\Gamma }`$-sampled cells dictates unambiguously the mathematical expression in *discrete* $`k`$-point space, and that invariance under โBZ foldingโ is the guideline which estabilishes the link. The resulting formalism is completely general and, while being similar in spirit to the original (MV) one, presents a much simpler algebra.
A more relevant issue affects more directly the modern theory of polarization, and concerns the asymptotic convergence with respect to BZ sampling. It was already shown formally and numerically mv ; puma that both methods (KSV and MV) for calculating the polarization of a molecule or a crystal in periodic boundary conditions (PBC) are plagued by a slow $`๐ช(L^2)`$ convergence, where $`L`$ is the linear dimension of the supercell containing the isolated molecule, or alternatively the resolution of the $`k`$-point mesh. The problem is usually addressed, in the context of the Berry-phase KSV approach, by refining the $`k`$-point grid along โstripesโ in the Brillouin zone within a separate, non-selfconsistent calculation kingsmith . For finite systems, an extrapolation technique was recently proposed sagui , in which the $`๐ช(L^2)`$ error is removed from the Wannier multipoles by performing a series of calculations in cubic supercells of increasing size <sup>1</sup><sup>1</sup>1A similar extrapolation technique was also proposed in the context of the discrete Berry-phase method davidextr . Both solutions are somewhat unsatisfactory, in that they require many calculations to be performed on the same system, with a cost that is higher than what is normally needed to converge total energies and densities. In the second part of this work we propose a simple โrefiningโ procedure, which is able to provide an extremely accurate value for the center and spread once a well-localized set of maximally localized Wannier functions is available. We show first formally and then by numerical examples that this technique, while requiring a very minor computational effort, is able to outperform in terms of accuracy both the standard KSV Berry-phase approach and the alternative formula based on โunrefinedโ WF centers <sup>2</sup><sup>2</sup>2 Another group, working in a different WF formalism dovesiwannier , already noticed dovesiknbo3 incidentally that the WF-based expression for polarization can potentially provide better numerical convergence than the Berry-phase approach.. Finally, our derivation also provides a novel, intuitive interpretation of the position/localization operator in periodic boundary conditions and of its relationship to the corresponding well-known free-space operators.
## II Method
The theoretical basis for the MV approach rests on a continuum formulation, in which the space is infinitely extended in all directions; this translates to an infinitely dense Brillouin zone sampling in the case of crystalline solids. For practical calculations a finite sampling (or finite simulation supercell) is necessary, and MV give detailed prescriptions for the โdiscretizationโ of the relevant mathematical objects (gradients and Laplacians in $`k`$-space). We start here our alternative derivation from a slightly different viewpoint, i.e. we โdiscretizeโ the problem from the very beginning by choosing an appropriate spread functional in the $`\mathrm{\Gamma }`$-point case, and then work out the formulas in $`k`$-space without making any further approximation. This approach leads automatically to a general and size-consistent formalism, that is invariant under BZ unfolding.
We assume a Born-von Kรกrmรกn (BvK) supercell of volume $`V_{BvK}`$, which is a multiple of the primitive (P) unit cell of the crystal (of volume $`V_P`$) under study. In this system with periodic boundary conditions (PBC) there are $`N`$ allowed Bloch vectors, where $`N`$ is given by the ratio between the volumes:
$$N=\frac{V_{BvK}}{V_P}.$$
The generalized Bloch orbitals (which are not necessarily eigenstates of the Hamiltonian) are orthonormal on the primitive cell:
$$_P\psi _{m๐ค}^{}(๐ซ)\psi _{n๐ค}(๐ซ)๐๐ซ=\delta _{mn},$$
and can be written as usual:
$$\psi _{n๐ค}(๐ซ)=e^{i๐ค.๐ซ}u_{n๐ค}(๐ซ),$$
where the $`u_{n๐ค}`$ are periodic functions, and can be represented on the reciprocal lattice of the P cell:
$$u_{n๐ค}(๐ซ)=\frac{1}{\sqrt{V_P}}\underset{|๐+๐ค|^2<E_{cut}}{}e^{i๐.๐ซ}\stackrel{~}{u}_{n๐ค}(๐).$$
$`E_{cut}`$ represents the plane-wave cutoff, while $`\stackrel{~}{u}_{n๐ค}(๐)`$ is the Fourier coefficient of the lattice-periodic part of the Bloch function:
$$\stackrel{~}{u}_{n๐ค}(๐)=\frac{1}{\sqrt{V_P}}_Pe^{i๐.๐ซ}u_{n๐ค}(๐ซ)๐๐ซ.$$
We will use the BvK supercell for representing our Wannier functions:
$$w_n(๐ซ)=\frac{1}{N}\underset{๐ค}{}\psi _{n๐ค}(๐ซ),$$
where the normalization constant is chosen so that these $`w_n`$ are orthonormal on the BvK supercell. A particularly simple relationship holds in reciprocal space:
$$\stackrel{~}{w}_n(๐+๐ค)=\frac{1}{\sqrt{N}}\stackrel{~}{u}_{n๐ค}(๐).$$
(1)
We remind the reader that the reciprocal lattice of the BvK cell is spanned by all vectors of type $`๐=๐+๐ค`$, which we will call $`๐`$ in the following, to distinguish them from the $`๐`$ vectors of the primitive reciprocal lattice.
Eq. 1 does not define a unique set of Wannier functions, because of the gauge arbitrariness in the choice of the unitary representation of the Bloch vectors. This indeterminacy can be solved by defining a spread functional $`\mathrm{\Omega }`$ which depends explicitly on the gauge, so that the minimization of $`\mathrm{\Omega }`$ leads to a well defined set of localized orbitals with the desired properties. Berghold and coworkers berghold proposed a particularly simple and appealing expression for $`\mathrm{\Omega }`$ and the related Wannier centers $`\overline{๐ซ}_n`$, which is valid for $`\mathrm{\Gamma }`$-only BZ sampling in a lattice of general symmetry. Since our BvK supercell is sampled at $`\mathrm{\Gamma }`$ by construction, we can use the same expressions as Berghold, that in our notations read:
$`\overline{๐ซ}_๐ง`$ $`=`$ $`{\displaystyle \underset{i}{}}\overline{\mathrm{w}}_i๐_i\mathrm{Im}\mathrm{ln}z_n^{(i)}`$ (2a)
$`\mathrm{\Omega }`$ $`=`$ $`{\displaystyle \underset{n}{}}{\displaystyle \underset{i}{}}\overline{\mathrm{w}}_i\mathrm{\hspace{0.17em}2}\left(1|z_n^{(i)}|\right).`$ (2b)
Here $`z_n^{(i)}`$ are dimensionless complex numbers given by:
$$z_n^{(i)}=w_n|e^{i๐_i.๐ซ}|w_n=|z_n^{(i)}|e^{i\varphi _n^{(i)}},$$
and {$`๐_i,\overline{\mathrm{w}}_i`$} represents a small set of reciprocal lattice vectors $`๐_i`$ with weights $`\overline{\mathrm{w}}_i`$. In the case of a cubic BvK supercell of edge $`L`$ these quantities reduce to the $`i=1,\mathrm{}3`$ primitive reciprocal-space vectors of the BvK cell and the weights are all equal:
$$๐_i=\frac{2\pi }{L}\widehat{๐ข},\overline{\mathrm{w}}_i=\left(\frac{L}{2\pi }\right)^2,$$
while in the most general case of a triclinic cell a maximum number of six independent vectors and weights must be used, according to the prescriptions given in Ref. mv, and silvestrelli, .
With these notations and conventions in hand, we are now ready to write down a $`k`$-space expression for $`\overline{๐ซ}_n`$ and $`\mathrm{\Omega }`$. Both quantities depend directly on $`z_n^{(i)}`$, and the key of the derivation is then the โBrillouin-zone unfoldingโ of this latter quantity. Using the same notation as MV:
$$M_{mn}^{(๐ค,๐_๐ข)}=u_{m๐ค}|u_{n๐ค+๐_i},$$
it is straightforward to derive a very simple expression for $`z_n^{(i)}`$:
$$z_n^{(i)}=\frac{1}{N}\underset{๐ค}{}M_{nn}^{(๐ค,๐_i)}.$$
With this formula we can write our operational definitions of position and quadratic spread in $`k`$-space:
$`\overline{๐ซ}_๐ง`$ $`=`$ $`{\displaystyle \underset{i}{}}\overline{\mathrm{w}}_i๐_i\mathrm{Im}\mathrm{ln}\left[{\displaystyle \frac{1}{N}}{\displaystyle \underset{๐ค}{}}M_{nn}^{(๐ค,๐_i)}\right],`$ (3a)
$`\mathrm{\Omega }`$ $`=`$ $`{\displaystyle \underset{n}{}}{\displaystyle \underset{i}{}}\overline{\mathrm{w}}_i\mathrm{\hspace{0.17em}2}\left(1\left|{\displaystyle \frac{1}{N}}{\displaystyle \underset{๐ค}{}}M_{nn}^{(๐ค,๐_i)}\right|\right)`$ (3b)
It is interesting to notice the strikingly close similarity between our expression for the centers (Eq. 3a) and Eq. 31 of MV, the only difference being the order in which the complex logarithm and the average over $`k`$-points is taken <sup>3</sup><sup>3</sup>3 Our expression for the spread (Eq. 3b) is also very close to Eq. 23 of MV (which is further discussed in posternak, ), and coincides with it when $`\overline{๐ซ}_๐ง=\mathrm{๐}`$. . We argue that the one proposed here is a more natural choice, since it retains the correct translational properties of their formula, while strictly enforcing size consistency. Size consistency means that the formalism gives mathematically identical answers for the $`k`$-point representation and for the equivalent BvK real-space $`\mathrm{\Gamma }`$-point representation. Our formula is correct by construction, and extends exactly to the case of isolated systems with $`\mathrm{\Gamma }`$-point sampling without any further algebra.
Another advantage of our scheme is its simplicity, which becomes evident when taking the gradient of the spread functional with respect to an infinitesimal unitary rotation in a given $`๐ค`$ subspace. Thus, we consider the transformation:
$$u_{n๐ค}^{}(๐ซ)=\underset{m}{}u_{m๐ค}(๐ซ)U_{mn}^{(๐ค)},$$
where the rotation matrices $`U^{(๐ค)}`$ are obtained by adding an infinitesimal antiHermitian matrix $`dW`$ to the identity:
$$U^{(๐ค)}1+dW^{(๐ค)}.$$
The variation of the total spread with respect to this transformation is readily obtained in terms of the $`M^{(๐ค,๐_i)}`$ matrices and the *phases* of the $`z_n^{(i)}`$ complex numbers:
$$\left(\frac{d\mathrm{\Omega }}{dW^{(๐ค)}}\right)_{mn}=\frac{1}{N}\underset{i}{}\overline{\mathrm{w}}_i(M_{mn}^{(๐ค,๐_i)}C_n^{(i)}+$$
$$+M_{nm}^{(๐ค๐_i,๐_i)}C_n^{(i)})H.c.,$$
(4)
where $`H.c.`$ stays for Hermitian-conjugate, and $`C_n^{(i)}`$ is the phase:
$$C_n^{(i)}=e^{i\varphi _n^{(i)}}=\frac{z_n^{(i)}}{|z_n^{(i)}|}.$$
This expression for the gradient can easily be obtained by observing that, in Eq. 2b, one can write $`|z|=ze^{i\varphi }`$.
We note that the for the spread functional (Eq. 2b) several possibilities exist, which are all equivalent in the thermodynamic limit berghold . In the Appendix we briefly consider these alternatives, and we provide a formal derivation of the gauge-invariant part of the spread mv , which further evidences the close relationship of our formulation to the original MV scheme. Because of the exact mapping between the BvK supercell and the primitive one, we find it particularly natural to choose our Wannier functions to be real. Even if there is no formal proof that at the global minimum of $`\mathrm{\Omega }`$ the Wannier functions are real, this is nevertheless a very reasonable assumption mv , and allows one to fully take advantage of the time-reversal symmetry, with significant gain in computational efficiency.
For the minimization of $`\mathrm{\Omega }`$ with respect to the $`U_{mn}^{(๐ค)}`$ degrees of freedom many efficient schemes are available berghold . We decided in this work to implement a damped dynamics algorithm, which allows for good control over the process, at the expense of requiring more human input for the optimal tuning of the two independent parameters (time step and friction). In antiferromagnetic MnO, a case that is known posternak to be difficult to converge, we were able to obtain this way a very accurate and symmetrical minimum (to machine precision) in a couple of thousand time steps, which required only a few minutes on a modern workstation. An even more appealing feature of the dynamical scheme is the availability of a mathematically conserved constant of motion, which provides a very stringent test on the accuracy of the implementation.
## III Convergence properties
Since the Wannier functions in an insulator are known to be *exponentially* localized in space vanderbiltprl , similar convergence properties can be expected for any physical quantity that is extracted from this particular representation of the electronic structure. Instead, as we pointed out at the beginning, both the sum of Wannier centers and the Berry phase (which is formally related to the Wannier centers by the derivation in KSV) converge only as $`๐ช(L^2)`$, and need special treatment whenever accurate values are needed.
We will show in this section that this slow convergence is indeed not a *intrinsic* feature of the ground state electronic structure of an extended system, and can be dramatically improved by a simple, inexpensive and very general procedure. Before explaining our correction in detail, we will first provide an intuitive picture of the position operator in PBC, which, as Resta showed restaprl , is the โkernelโ of both Berry phase and maximally localized Wannier function calculations.
Letโs consider a one-dimensional system of one single electronic state $`|\psi `$, which we will assume to be *well localized* within a periodic cell of length $`L`$. The expression for the fundamental, dimensionless complex number $`z`$ is very similar to the 3D expression:
$$z=\psi |e^{i\frac{2\pi }{L}x}|\psi =|z|e^{i\varphi }.$$
(5)
The average value of the position operator (Eq. 2a) becomes:
$$\overline{x}=\frac{L}{2\pi }\mathrm{Im}\mathrm{ln}z=\frac{L}{2\pi }\varphi ,$$
(6)
while the quadratic spread (Eq. 2b) reduces to:
$$\mathrm{\Omega }=\left(\frac{L}{2\pi }\right)^22(1|z|).$$
By defining the charge density $`\rho (x)=|\psi (x)|^2`$, it is easy to see that the following is true <sup>4</sup><sup>4</sup>4 The generalization of Eq. 7a and Eq. 7b to a general 3D lattice is straightforward by using the set of $`๐_i`$ vectors and weights defined in the text.:
$$_0^L\rho (x)\mathrm{sin}[\frac{2\pi }{L}(x\overline{x})]๐x=0$$
(7a)
$$\mathrm{\Omega }=\left(\frac{L}{2\pi }\right)^2_0^L\rho (x)\left\{22\mathrm{cos}[\frac{2\pi }{L}(x\overline{x})]\right\}๐x$$
(7b)
These equations can be directly compared to the elementary textbook definitions of the position and quadratic spread for a square-integrable electronic state in one dimension (i.e. without PBC, the superscript $`F`$ stays for โfree-spaceโ):
$$_{\mathrm{}}^{\mathrm{}}\rho (x)(x\overline{x}^F)๐x=0$$
(8a)
$$\mathrm{\Omega }^F=_{\mathrm{}}^{\mathrm{}}\rho (x)(x\overline{x}^F)^2๐x$$
(8b)
The resemblance is indeed striking, the only difference being the replacement of the $`x`$ and $`x^2`$ operators with trigonometric functions that are periodic on the cell. This relationship between the (polynomial) free-space operators and the (trigonometric) PBC ones is made evident in Fig. 1, where they are plotted together in order to show their close matching in a region surrounding the localized state. Indeed, by a Taylor expansion one obtains:
$$\left(\frac{L}{2\pi }\right)\mathrm{sin}(\frac{2\pi }{L}x)x\left(\frac{2\pi }{L}\right)^2\frac{x^3}{3!}+\mathrm{}$$
$$\left(\frac{L}{2\pi }\right)^2\left[22\mathrm{cos}(\frac{2\pi }{L}x)\right]x^22\left(\frac{2\pi }{L}\right)^2\frac{x^4}{4!}+\mathrm{}$$
Thus, we arrived from a different starting point at the same $`๐ช(L^2)`$ convergence of the position and spread, which has already been discussed in the literature mv ; puma . It is particularly clear from this derivation that the intrinsic property of the periodic lattice and the localized state are by no means responsible for the slow convergence, which is instead determined exclusively by the mathematical form of the PBC position operator.
To end this section, we note that Eq. 7a alone is not sufficient to *define* the center $`\overline{x}`$, since also $`\overline{x}+\frac{L}{2}`$ satisfies the same requirement. In the context of Equations 7a and 7b, a correct definition of $`\overline{x}`$ can be given as the points in the lattice which minimize Eq. 7b (it is easy to show that this definition is identical to the standard one in Eq. 6). Interestingly, from this point of view the position $`\overline{x}`$ can be thought of as an *internal* parameter of the formalism, which is implicit in the definition of the spread.
## IV Correction scheme
Since we are working with Wannier functions which are expected to be well localized in space (as the 1D state depicted in Fig. 1), there is actually no need to insist on using the PBC formulas for calculating Wannier centers. One could argue here that our โWannier functionsโ are formally still periodic (although represented on a large BvK supercell), and since their Hilbert space is defined within PBC, only the action of PBC-allowed operators is justified on them. Actually, another point of view can be used. We recall that *true* Wannier functions are continuous functions in the full 3D (reciprocal) $`q`$-space. The mean value of a local operator $`V(๐ซ)`$ in real space can be written
$$w_n|V|w_n=_{Allspace}|w_n(๐ซ)|^2V(๐ซ)๐๐ซ$$
(9)
or, equivalently in $`q`$-space ($`\stackrel{~}{\rho }_n(๐ช)`$ is the continuous Fourier transform of the Wannier density $`\rho _n(๐ซ)=|w_n(๐ซ)|^2`$):
$$w_n|V|w_n=_{|๐ช|<E_{cut}}\stackrel{~}{V}(๐ช)^{}\stackrel{~}{\rho }_n(๐ช)๐๐ช$$
(10)
If the Wannier function is localized (i.e. zero beyond a given distance from its center), the integral in Eq. 9 can be limited to a finite region of space, for example a cubic box centered around the region where the Wannier density is nonzero. The $`q`$-space integral in Eq. 10 can then be recast to a sum over a *discrete* set of reciprocal space vectors, which is also *finite* because of the plane-wave cut off, and the result is still *exact*.
If the integration box is chosen to be smaller than the region where the Wannier density is nonzero, then the reciprocal-space sum carries an error which is due to the overlap between the tails of the Wannier functions and their (artificially) repeated images. This overlap depends on the decay properties of the localized state, and in particular it goes *exponentially* to zero for increasing integration box size whenever $`|w_n`$ is exponentially localized.
In a standard DFT simulation of a periodic crystal, the discrete set of reciprocal-space Wannier function coefficients are defined by Eq. 1, and converge to their thermodynamic limit as soon as the total charge density is converged. Then, the only effect of a further refinement of the $`k`$-points mesh is an increase in the BvK cell volume, which leads to the progressive reduction of the overlap term discussed above. Therefore, assuming exponential decay for the Wannier functions, our technique allows for an *exponential* convergence of the calculated expectation value of any free-space operator. The natural โbounding boxโ for the integration domain in real space is, for a general lattice, a Wigner-Seitz BvK cell aligned on the Wannier center. With this choice, the discrete Fourier representation of a given local free-space operator (we use here again the standard conventions for normalizations and Fourier transforms) is:
$$\stackrel{~}{V}(๐)=\frac{1}{V_{BvK}}_{WignerSeitz}e^{i๐.๐ซ}V(๐ซ)๐๐ซ,$$
and the expectation value is simply given as
$$w_n|V|w_n=V_{BvK}\underset{๐}{}\stackrel{~}{V}^{}(๐)\stackrel{~}{\rho }_n(๐)$$
Starting from a well-localized set of Wannier functions we can now define a โrefinedโ spread operator $`\mathrm{\Omega }^{}=_n\mathrm{\Omega }_n^{}`$, where the contribution from the individual WF is:
$$\mathrm{\Omega }_n^{}=_{WignerSeitz}|๐ซ\overline{๐ซ}_n^{}|^2\rho _n(๐ซ)๐๐ซ.$$
The $`b`$-space expression for this formula can be derived starting from the Fourier series of a parabola in a one dimensional box:
$$\frac{1}{L}_{\frac{L}{2}}^{\frac{L}{2}}\mathrm{cos}(\frac{2\pi k}{L}x)x^2๐x=\left(\frac{L}{2\pi }\right)^2\frac{2(1)^k}{k^2}(k>0)$$
$$\frac{1}{L}_{\frac{L}{2}}^{\frac{L}{2}}x^2๐x=\frac{L^2}{4},$$
and is readily generalized to three dimensions using the same set of vectors and weights {$`๐_i,\overline{\mathrm{w}}_i`$} introduced in Section 2:
$$\mathrm{\Omega }_n^{}=2V_{BvK}\underset{i,k>0}{}\overline{\mathrm{w}}_i\mathrm{Re}\left[\frac{2(1)^k}{k^2}e^{ik๐_i.\overline{๐ซ}_n^{}}\stackrel{~}{\rho }_n(k๐_i)\right]+$$
$$+\frac{\pi ^2}{3}\underset{i}{}\overline{\mathrm{w}}_i$$
(11)
The โrefinedโ position $`\overline{๐ซ}_n^{}`$ which appears in Eq. 11 is then again an *internal parameter*, which is defined by the minimum of $`\mathrm{\Omega }_n^{}`$ for a given $`\stackrel{~}{\rho }_n`$ (see the discussion at the end of Section 3). By taking the gradient of $`\mathrm{\Omega }_n^{}`$ with respect to $`\overline{๐ซ}_n^{}`$ one obtains that, at the minimum, the integral:
$$\mathrm{\Delta }๐ซ_n=_{WignerSeitz}(๐ซ\overline{๐ซ}_n^{})\rho _n(๐ซ)๐๐ซ$$
vanishes. Consistently with the definition of the spread, this condition has to be enforced in reciprocal space, where this integral becomes:
$$\mathrm{\Delta }๐ซ_n=2V_{BvK}\underset{i,k>0}{}\overline{\mathrm{w}}_i๐_i\mathrm{Re}\left[i\frac{(1)^k}{k}e^{ik๐_i.\overline{๐ซ}_n^{}}\stackrel{~}{\rho }_n(k๐_i)\right],$$
(12)
The stationary point can be obtained iteratively starting from a set of maximally localized Wannier functions and unrefined centers, by updating at every iteration $`\overline{๐ซ}_n^{}`$ through the addition of $`\mathrm{\Delta }๐ซ_n`$ as calculated in Eq. 12 until convergence is reached. If the Wannier function is exactly zero in a region surrounding the boundary of the Wigner-Seitz cell, one iteration is sufficient to provide the *exact* value of the center, while for less converged cases up to ten iterations may be necessary to achieve machine precision. These iterations have anyway negligible cost, since the Fourier transform of the Wannier function on the BvK cell has to be evaluated only at the beginning (twice for each Wannier function to get the density in reciprocal space).
Both expressions $`\mathrm{\Omega }`$ and $`\mathrm{\Omega }^{}`$ are in fact particular cases of a class of localization criteria which rely on individual Wannier densities only, through some generalized spread functional $`S`$:
$$\mathrm{\Omega }=\underset{n}{}S[\rho _n]$$
A similar generalised, density-dependent spread can be used in practice to explore alternative localization criteria, like e.g. the maximal Coulomb self-repulsion of Edmiston and Ruedenberg edmiston , or the orbital self-interaction as defined by Perdew and Zunger perdew . An article comparing such alternatives is under preparation inprep .
Since the present free-space-like expressions for position and spread are more accurate than those derived in the first part of this work, one could wonder why we did not use them from the beginning. The reason is exclusively related to computational efficiency. In Eqns. 3 the localization algorithm involves operations on small $`J\times J`$ matrices only, where $`J`$ is the number of bands in the primitive cell (the computationally intensive calculation of the $`M^{(๐ค,๐_i)}`$ matrices has to be performed only once at the beginning of the iterative minimization). If, instead, the refined spread (or one of the alternative localization criteria discussed above) is used directly for localizing the Wannier functions, several Fourier transforms on the full Wannier (BvK) grid are required for each iteration, at a substantially higher cost. This expensive procedure is anyway not necessary for the scope of the present work, since the actual set of Wannier wavefunctions obtained from one localization method or the other coincide in practice to a high degree of accuracy hydrogen (in particular, the decay properties are expected to be very similar). Therefore, we find it most convenient to use this refinement step in a one-shot fashion once a set of maximally localized Wannier functions is obtained within the more efficient localization functional (Eqns. 3).
## V Numerical tests
To demonstrate the effectiveness of our method we have chosen two examples which have been extensively studied in the literature: (i) the dynamical Born effective charge of oxygen atoms in rocksalt MgO, and (ii) the spontaneous polarization of the ferroelectric, tetragonally distorted phase of KNbO<sub>3</sub>. Our calculations were performed within the local density approximation perdew , by using norm-conserving Troullier and Martins tm pseudopotentials in the Kleinman and Bylander kb form. A non-linear core correction louie was adopted for the Mg pseudopotential, while the K pseudopotential was generated in the $`4s^0`$ ionized configuration with the semicore $`3s,3p`$ orbitals included in the valence. We used the experimental lattice constants and atomic positions ($`a_0`$=7.96 a.u. for MgO mgoexp , and the structural data for KNbO<sub>3</sub> from Ref. baldoknbo3, ). We expanded the electronic ground state on a plane-wave basis up to a cutoff of 70 Ry. The BZ sampling was performed with $`\mathrm{\Gamma }`$-centered simple cubic (orthorombic) grids in reciprocal space for MgO (KNbO<sub>3</sub>), by taking into account the time-reversal symmetry only.
We will compare the results as a function of $`k`$-mesh resolution for three different methods for calculating the polarization: (i) the sum of Wannier centers as obtained by Eq. 3a; (ii) the sum of *refined* Wannier centers as described in the previous section; (iii) the Berry-phase approach. We note that the Berry-phase result can be readily obtained from the quantities that are already available in the localization formalism:
$$๐ซ_{Berry}=\frac{1}{N}\underset{i}{}\overline{\mathrm{w}}_i๐_i\underset{๐ค}{}\mathrm{Im}\mathrm{ln}detM^{(๐ค,๐_i)}.$$
(13)
### V.1 MgO
The dynamical Born effective charge ($`Z^{}`$) of oxygen was calculated by the finite-difference method, i.e. by considering the difference in total polarization between the ideal centrosymmetric ground state and an atomic configuration where the oxygen sublattice was displaced by 1% of the cubic lattice constant along the $`x`$ direction. The atomic coordinates were prepared in such a way that, in the ideal lattice, the O atom sits at the origin, and in this case the electronic contribution to the polarization is exactly zero modulo a polarization quantum (all four Wannier functions are symmetrical about the O in this case). We compare in Fig. 2 the resulting value for $`Z^{}`$ as calculated by the three different methods (i-iii). The results clearly show that the sum of the *unrefined* Wannier centers can be very inaccurate, and even for the finest mesh the error is still large. MgO is probably a very unfortunate case in that each $`sp^3`$-like Wannier function has a strongly asymmetric shape, and the errors in the individual centers do not cancel out efficiently in the strained configuration, so that the total polarization carries an important deviation from the exact value. The Berry-phase calculation is a much better estimate, but in the inset it can be seen that the convergence is still relatively slow. As we explained in the preceding sections, it was already shown that the Berry-phase result converges only as $`๐ช(L^2)`$, i.e. it shares the same asymptotic behaviour as the sum of the unrefined Wannier centers (*albeit* with a quite different prefactor in this particular case). The sum of the *refined* Wannier centers instead shows an extremely fast convergence, and gives a very accurate result already for a $`2\times 2\times 2`$ mesh. The value of $`Z^{}`$ we obtain is -1.95, which is in excellent agreement with previous experimental and theoretical investigations pumaprl .
To complement our methodological test, we calculated also the refined value for the total quadratic spread as a function of $`k`$-mesh resolution, and the results are reported in Fig. 3. It is clear that this quantity shares the same, excellent, convergence properties as the position operator (upper panel). In the lower panel of the same figure we report for comparison the results of an analogous calculation of the total spread in bulk silicon. The convergence is slower than in MgO, as can be expected from the very different character of this covalent compound as compared to the highly ionic magnesium oxide, but the benefit that can be obtained through the use of the more accurate free-space definition of the spread is still very clear. The โunrefinedโ value of the spread is also compared to the alternative, very similar prescriptions discussed in the Appendix.
### V.2 KNbO<sub>3</sub>
We present in Fig. 4 our results for the spontaneous polarization of KNbO<sub>3</sub>. The sum of the unrefined Wannier centers is less inaccurate in this case, and is fairly close to the values obtained within the Berry phase formalism. The sum of the refined centers has, again, much better convergence properties than the two traditional methods. By increasing the mesh from $`6\times 6\times 6`$ to $`7\times 7\times 7`$ the value of the spontaneous polarization increases by 0.02 %, while within the same $`7\times 7\times 7`$ mesh the traditional techniques carry an error which is two orders of magnitude higher. Extrapolating the $`๐ช(L^2)`$ trend one can guess that $`70`$ $`k`$-points along the reciprocal space stripes kingsmith would be needed to achieve similar accuracy within the Berry-phase formalism. The final value we obtain, 0.38 C/m<sup>2</sup>, compares very well with experimental data and previous theoretical investigations baldoknbo3 ; dovesiknbo3 ; restaknbo3 ; vanknbo3 .
## VI Conclusions
In conclusion, we have derived a simple and general formalism for the computation of maximally localized Wannier functions. We provide an intuitive picture of the convergence properties of this scheme and similar ones, relating them to the Taylor expansion of elementary trigonometric functions. We show that the convergence can be dramatically improved by a simple strategy based on the exponential localization of the Wannier functions in insulating materials. We expect our scheme to open the way to both accurate and efficient calculations of polarization properties in a wide range of physical systems, making the expensive linear-response approach or the relatively cumbersome non-self consistent calculation of โstripesโ in reciprocal space unnecessary.
As a final remark, we note that the Wannier function-based theory of polarization is becoming increasingly important especially in disordered systems, where not only the global polarization but also the *local* bonding properties and dipole moments may be interesting to follow during, e.g. a molecular dynamics simulation sharma . In these applications the improved accuracy provided by our method could be an extremely valuable tool.
We wish to thank David Vanderbilt for insightful comments on the manuscript. This work was supported by the National Science Foundationโs Division of Materials Research through the Information Technology Research program, grant number DMR-0312407, and made use of MRL Central Facilities supported by the National Science Foundation under award No. DMR-0080034
*
## Appendix A Decomposition into invariant, off-diagonal and diagonal parts
The form 2b for the spread functional was chosen mainly because of its simplicity, and because it allows for a direct interpretation as the integral of the Wannier density multiplied by a real function on the BvK cell (see the discussion in Sec. 3). Unfortunately this expression does not lead to an elegant separation into invariant and non-invariant parts. However, this issue is readily solved by choosing an alternative definition of the spread:
$$\mathrm{\Omega }_{MV}=\left(\frac{L}{2\pi }\right)^2(1|z|^2),$$
(14)
which coincides with the $`\mathrm{\Gamma }`$-point prescription of MV and which *does* allow for an exact separation of the invariant part. This choice still allows for the simple interpretation based on cosine-like functions. If we define a function of $`x_0`$:
$$f(x_0)=_0^L\rho (x)\mathrm{cos}[\frac{2\pi }{L}(xx_0)]๐x$$
(15)
it is clear that when $`x_0`$ maximizes $`f`$, $`x_0`$ is automatically the Wannier center of Eq. 2a. Both expressions for the spread ($`\mathrm{\Omega }`$ as in Eq. 2b and $`\mathrm{\Omega }_{MV}`$ discussed here) are consistent with the same value of $`x_0`$ at the minimum:
$$\mathrm{\Omega }=\left(\frac{L}{2\pi }\right)^2\underset{x_0}{\mathrm{min}}2[1f(x_0)]$$
$$\mathrm{\Omega }_{MV}=\left(\frac{L}{2\pi }\right)^2\underset{x_0}{\mathrm{min}}[1f^2(x_0)]$$
Moving on to 3D, the operational definition of the spread becomes:
$$\mathrm{\Omega }_{MV}=\underset{n}{}\underset{i}{}\overline{\mathrm{w}}_i\left(1|z_n^{(i)}|^2\right),$$
where it is easy to see that $`z_n^{(i)}`$ are nothing other than the matrix elements indicated as $`X_{nn}`$, $`Y_{nn}`$, $`Z_{nn}`$ in MV.
Now, โfoldingโ this expression in $`k`$-space leads to a formula which is similar to Eq. 3b:
$$\mathrm{\Omega }=\underset{n}{}\underset{i}{}\overline{\mathrm{w}}_i\left(1\left|\frac{1}{N}\underset{๐ค}{}M_{nn}^{(๐ค,๐_i)}\right|^2\right)$$
Thinking in terms of the big BvK cell, this can be written equivalently as:
$$\mathrm{\Omega }=\frac{1}{N}\underset{i}{}\overline{\mathrm{w}}_i(NJ\underset{๐,n}{}|๐n|e^{i๐_i.๐ซ}|๐n|^2),$$
where the leading factor $`1/N`$ gives the spread per *primitive* cell, and the same notations as MV for the $`n`$-th Wannier function at the $`๐`$ site, $`|๐n`$ are used. From this expression it is clear how to construct an obvious invariant quantity, $`\mathrm{\Omega }_I`$ ($`J`$ is the number of bands in the primitive cell):
$$\mathrm{\Omega }_I=\frac{1}{N}\underset{i}{}\overline{\mathrm{w}}_i(NJ\underset{\mathrm{๐๐}^{},nm}{}|๐n|e^{i๐_i.๐ซ}|๐^{}m|^2),$$
and what remains to do is to โunfoldโ this formula in $`k`$-space. A first simplification is trivial:
$$\mathrm{\Omega }_I=\underset{i}{}\overline{\mathrm{w}}_i(J\underset{๐,nm}{}|๐n|e^{i๐_i.๐ซ}|\mathrm{๐}m|^2).$$
A second simplification is obtained by reversing the formula between Eq. 5 and 6 of MV, leading to:
$$\mathrm{\Omega }_I=\underset{i}{}\overline{\mathrm{w}}_i(J\underset{๐,nm}{}\left|\frac{1}{N}\underset{๐ค}{}e^{i๐ค.๐}u_{n๐ค}|e^{i๐_i.๐ซ}|u_{m๐ค+๐}\right|^2).$$
By writing explicitly $`|z|^2=z^{}z`$ and noticing that
$$\underset{๐}{}e^{i(๐ค๐ค^{}).๐}=N\delta _{๐ค,๐ค^{}},$$
we obtain the final expression in $`k`$-space:
$$\mathrm{\Omega }_I=\underset{i}{}\overline{\mathrm{w}}_i(J\underset{mn}{}\frac{1}{N}\underset{๐ค}{}\left|M_{nm}^{(๐ค,๐_i)}\right|^2),$$
(16)
which is exactly Eq. 34 of the MV paper.
It is interesting to work out the remaining terms, $`\mathrm{\Omega }_D`$ and $`\mathrm{\Omega }_{OD}`$, which are indicated as โdiagonalโ and โoff-diagonalโ parts in MV (we recall that $`\mathrm{\Omega }_D`$ vanishes for a centrosymmetric system mv ). The easiest way is to first solve the expression for
$$\mathrm{\Omega }_{MV}\mathrm{\Omega }_D=\underset{i}{}\overline{\mathrm{w}}_i(J\underset{๐,n}{}|๐n|e^{i๐_i.๐ซ}|\mathrm{๐}n|^2).$$
By using an analogous algebra we readily arrive at the formula in $`k`$-space:
$$\mathrm{\Omega }_{MV}\mathrm{\Omega }_D=\underset{i}{}\overline{\mathrm{w}}_i(J\underset{n}{}\frac{1}{N}\underset{๐ค}{}\left|M_{nn}^{(๐ค,๐_i)}\right|^2),$$
from which it is very easy to evaluate $`\mathrm{\Omega }_{OD}`$:
$$\mathrm{\Omega }_{OD}=\underset{i}{}\overline{\mathrm{w}}_i\underset{mn}{}\frac{1}{N}\underset{๐ค}{}\left|M_{mn}^{(๐ค,๐_i)}\right|^2.$$
This means that $`\mathrm{\Omega }_D`$ is given by the following difference:
$$\mathrm{\Omega }_D=\underset{i,n}{}\overline{\mathrm{w}}_i\left(\left|\frac{1}{N}\underset{๐ค}{}M_{nn}^{(๐ค,๐_i)}\right|^2\frac{1}{N}\underset{๐ค}{}\left|M_{nn}^{(๐ค,๐_i)}\right|^2\right).$$
Comparing this formalism with the MV one, it is clear that $`\mathrm{\Omega }_I`$ and $`\mathrm{\Omega }_{OD}`$ are identical, while the terms $`\mathrm{\Omega }`$ and $`\mathrm{\Omega }_D`$ differ. This derivation provide in a certain sense a โunificationโ of the two, formerly distinct, MV prescriptions for the $`\mathrm{\Gamma }`$-point case and in $`k`$-space. The gradients with respect to unitary rotations of the Bloch orbitals are simply given by setting $`C=z`$ (instead of $`C=z/|z|`$) in Eq. 4.
To complete our discussion, we note that a third form of the one-particle quadratic spread was proposed by Resta and Sorella rs , which leads to yet another operational definition for the localization criterion:
$$\mathrm{\Omega }_{RS}=\underset{n}{}\underset{i}{}\overline{\mathrm{w}}_i\mathrm{ln}|z_n^{(i)}|^2.$$
The $`k`$-space folding of this formula is straightforward, while the gradient is again given by Eq. 4, with $`C=z/|z|^2`$. All functionals $`\mathrm{\Omega }`$, $`\mathrm{\Omega }_{MV}`$ and $`\mathrm{\Omega }_{RS}`$ are identical in the thermodynamic limit. For finite BvK cells they all share the same definition of the Wannier center. The resulting maximally localized Wannier functions themselves are identical in cases where $`|z_n|`$ are equal for all $`n=1,\mathrm{},J`$ bands (e.g. bulk Si, centrosymmetric MgO crystal). The numerical value of the spread can differ slightly, because the higher orders in the Taylor expansion are different. Some examples concerning this discrepancy are reported in the main text (see, e.g., Fig. 3).
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# The Molecular Accretion Flow in G10.6-0.4
## 1 Introduction
G10.6-0.4 is a bright, 2.5 Jy at 23 GHz (Keto et al., 1987a), ultra-compact (UC) Hii region (Wood & Churchwell, 1989) at a distance of 6.0 kpc (Downes et al., 1980), in an area of active star formation. The associated IRAS point source, IRAS 18075-1956, has a luminosity of $`9.2\times 10^5`$ L (Casoli et al., 1986) and has colors that meet the criteria of Wood & Churchwell (1989) for an UCHii region. G10.6 is known to be embedded in a hot molecular core (HMC) (Braz & Epchtein, 1983; Ho & Haschick, 1986; Plume et al., 1992). The core is thought to contain 1200 M of gas within a radius of 0.2 pc, based on an analysis of a variety of dust continuum measurements (Mueller et al., 2002), and 3300 M within 1.1 pc based on C<sup>18</sup>O and C<sup>17</sup>O measurements (Hofner et al., 2000). Previous studies of the inversion lines of NH<sub>3</sub> have determined that the molecular core is rotating and collapsing inward toward the UCHii region (Ho & Haschick, 1986; Keto et al., 1987b, 1988; Keto, 1990). In these studies, using the NH<sub>3</sub>(1,1) and NH<sub>3</sub>(3,3) lines, rotation is seen at size scales from 1 pc down to 0.08 pc, and infall is detected in the form of red-shifted absorption seen against the continuum source. CH<sub>3</sub>OH and H<sub>2</sub>O masers are seen distributed linearly in the plane of the rotation (Walsh et al., 1998; Hofner & Churchwell, 1996), while OH masers seem to lie along the axis of rotation (Argon et al., 2000). In C<sup>18</sup>O $`(J=10)`$, Ho et al. (1994) see $`10^3`$ M of dense (n$`10^6`$cm<sup>-3</sup>), rotating gas in a flattened (0.3$`\times `$0.1 pc) disk-like structure. At the highest resolution achieved in earlier work, infall and rotation in the molecular gas were seen simultaneously in absorption, showing that the molecular gas was spiraling inward on size scales comparable to the size of the UCHii region.
Recent observations of G10.6 hinted that it might represent a previously unobserved mode of high mass star formation. Observations of H66$`\alpha `$ from the ionized gas within the UCHii region indicate that the ionized gas is also spiraling inward toward the stars at the center of the UCHii region (Keto, 2002a). Subsequent theoretical work showed that in small Hii regions, the gravitational effect of the central star(s) can overcome the thermal pressure of the ionized gas causing the molecular accretion flow to pass through the Hii region boundary and continue inward as an ionized accretion flow (Keto, 2002b). In this model, the Hii region boundary exists as a standing R-type ionization front within a continuous accretion flow. These results differ from classical treatments of the pressure driven expansion of Hii regions, which predict outward motion of the ionized gas as soon as the Hii region is formed (Strรถmgren, 1939; Spitzer, 1978). In the classical model for pressure driven expansion, the Hii region boundary, after a very short phase as a moving R-type front, will develop a characteristic double front structure composed of an isothermal shock followed by a moving D-type ionization front. As the Hii region expands, most of the displaced molecular material remains between the shock and the ionization front as a dense outward moving shell, which snow-plows ahead of the Hii region. If, however, the accretion flow passes through a standing R-type ionization front at the Hii region boundary and continues toward the star(s) as an ionized flow, as suggested by Keto (2002b), there will be no dense molecular layer at the boundary, and all the molecular gas will be moving inward.
Sollins et al. (2005) did preliminary analysis of the data presented here, and showed that G10.6 is accreting through its UCHii region. The velocities of the molecular gas showed clear evidence of both infall and rotation. Based on geometrical arguments, Sollins et al. (2005) concluded that the infalling layer proceeded directly up to the ionization front. They also showed a non-detection of any expanding molecular gas. The non-detection placed such a stringent upper limit on the mass of any expanding molecular shell that might be present, that it was concluded that no such expansion was taking place, and that the accretion in the ionized and molecular gas were all part of a single accretion flow which continues across a stalled ionization front.
In this paper we present a more detailed analysis of the NH<sub>3</sub> data. We conclude that, while the accretion flow is spherical at large radii, the rotation does cause it to flatten somewhat on the size scale of the UCHii region, so that the highest column densities appear in a thin strip. We also conclude that the axis of rotation is inclined away from the observer in the northeast. We find that G10.6 is quite different than other young high mass stars in which disk-like molecular structures have been observed (Zhang et al., 1998, 2002; Shepherd & Kurtz, 1999; Chini et al., 2004).
## 2 Observations
We observed the UCHii region G10.6 with the NRAO Very Large Array (VLA)<sup>1</sup><sup>1</sup>1The National Radio Astronomy Observatory is a facility of the National Science Foundation operated under cooperative agreement by Associated Universities, Inc. on February 1, 2002, with the phase center at $`\alpha (2000)=18^\mathrm{h}10^\mathrm{m}28^\mathrm{s}.683,\delta (2000)=19^\mathrm{o}55^{}49^{\prime \prime }.07`$. We observed the (3,3) inversion line of NH<sub>3</sub> at 23.870130 GHz with 63 spectral channels of width 48.828 kHz (0.61 km s<sup>-1</sup>) for a total bandwidth of 3.125 MHz (38.7 km s<sup>-1</sup>) centered on $`\mathrm{v}_{\mathrm{lsr}}`$$`=10`$ km s<sup>-1</sup>, and 1.3 cm continuum with a bandwidth of 15.6 MHz. The array was in the A-configuration, yielding a uniform-weighted synthesized beam of width $`0.^{\prime \prime }12\times 0.^{\prime \prime }072`$ for a physical resolution of 0.0034$`\times `$0.0021 pc or 700$`\times `$430 AU.
We observed the quasars 3C286, 3C273 and 1733-130 for flux, bandpass and phase calibration respectively. Self-calibration of the source amplitudes and phases resulted in a noise level of 0.18 (0.14) mJy beam<sup>-1</sup> in the uniform (natural) weighted continuum map, and 1.9 (1.5) mJy beam<sup>-1</sup> in each uniform (natural) weighted channel map, about 3 times the thermal noise limit. The images were deconvolved by CLEANing in the usual way with the AIPS task IMAGR. No special steps were taken to deal with the fact that the emission area is much larger than the synthesized beam. The total flux in the resulting natural weighted continuum map is 2.5 Jy, which is consistent with the total flux detected in earlier lower resolution maps of 23 GHz continuum (Keto et al., 1987b). For this reason, we believe that the continuum map is missing very little flux due to the lack of short baselines.
Expressed as a temperature, our sensitivity in a natural weighted continuum map is about 25 K and in a natural weighted channel map is about 280 K. The physical temperature of the molecular gas around G10.6 is estimated to be only 110 K at the ionization front (Keto, 1990). Thermal line emission always has a brightness temperature less than the temperature of the gas. Thus, the brightest possible thermal emission from the molecular gas would be undetectable, less than $`0.4\sigma `$. Absorption, however, should be detectable at a wide range of optical depths. The continuum has a peak brightness temperature of 6900 K, and since the noise level is a channel map is 280 K, the molecular line absorption should be detectable at up to $`25\sigma `$. The quality of the self-calibration solutions and improvements in the K-band receiver system at the VLA have resulted in 25 times better sensitivity in our NH<sub>3</sub>(3,3) channel maps than in the previous best existing NH<sub>3</sub>(3,3) data for this source (Keto et al., 1988), with 3 times better spatial resolution and 2 times better velocity resolution. A sample spectrum is shown in Figure 1.
It should be noted that we have achieved the highest possible angular resolution in studying the problem of high mass star formation. With the VLA in its most extended array configuration, at a frequency of 23 GHz, our spatial resolution is $`0^{\prime \prime }.1`$. Imaging thermal (i.e., non-masing) molecular gas at that resolution is only possible toward sources with strong continuum emission, and only at wavelengths which include spectral lines useful for studying the dense gas surrounding UCHii regions. Using radiation with wavelengths around a centimeter is ideal because it is near the peak of the continuum emission from many interesting UCHii regions. At lower frequency, the spatial resolution decreases. At higher frequency the brightness temperature of the continuum emission declines rapidly. So for optimal backlighting from the UCHii region, centimeter wavelengths are ideal. There are few thermal lines associated with high density gas in the centimeter wavelength regime apart from the inversion lines of ammonia which have a critical density of roughly $`10^4`$cm<sup>-3</sup>(Ho & Townes, 1983). In addition to its fortuitous wavelength, the hyperfine structure (one main line, four satellite lines) of the ammonia inversion transitions is extremely useful. Because the different hyperfine components have well-known intrinsic line-strengths, the ratio of the main line to a satellite line can be used to directly calculate the optical depth (what we call hyperfine optical depth) and column density of the ammonia in the rotational state in question, in this case (J,K) = (3,3). While the optical depth of any absorption line can be calculated directly by comparing the depth of the absorption to the strength of the background continuum, this apparent optical depth has limitations. If the filling factor of the absorbing gas is less than one, the apparent optical depth decreases. Also, for deeply embedded objects like G10.6, the main hyperfine component easily saturates. For this reason, the satellite lines, which are much more optically thin and do not saturate easily, are invaluable in investigating the highest optical depths and column densities. Since the hyperfine optical depth is calculated from the ratio of the main line absorption to that of the satellite line, the hyperfine optical depth is accurate even when the main line saturates, and also does not include the effects of the filling factor. The high spatial resolution and sensitivity to high column density gas achieved here have not yet been possible for this sort of object in the millimeter, IR, optical, or X-ray regimes.
## 3 Results
We report 6 key observational results. First, we find that the absorbing molecular gas is less spatially extensive than the continuum emission, and, on average, northeast of the average position of the continuum emission. Second, we find that the highest column density gas is localized in clumps, while lower column density gas is seen over the entire face of the UCHii region. Third, we find that the characteristic size scale of the infall-and-rotation kinematic pattern of the molecular gas is larger than the characteristic size scale of the structure in either the optical depth maps, or continuum map. Fourth, we find that the highest optical depth gas is missing on the southwest side of the UCHii region, reinforcing the idea that the absorption is preferentially located in the northeast. Fifth, on size scales smaller than the synthesized beam, i.e. less than $`0^{\prime \prime }.1`$ or 500 AU, in more than 75% of the pixels where absorption is detected, the filling factor of the absorbing gas is greater than 0.7. Sixth, none of the sharp edges clearly seen in the continuum emission are seen in maps of the optical depth.
We find that the absorbing molecular gas, as located by the actual line absorption, the apparent optical depth and the hyperfine optical depth, is on average northeast of the continuum emission, and is spatially narrower than the continuum emission in the direction of the minor axis, with the hyperfine optical depth being the most skewed to the northeast, and the most narrow. We have taken slices through the maps shown in Figure 2. The slices run northeast-southwest, and are separated by a synthesized beam-width ($`0^{\prime \prime }.14`$). For each slice we have calculated the first moment to determine the mean position of the flux in the maps, and the second moment to determine the width of the flux in the maps. Figures 3 and 4 show the mean positions and widths as a function of position. The origin was chosen so that the mean positions in the continuum map have an average of zero. In the other three maps, the mean positions are generally northeast of the continuum positions. For the hyperfine optical depth map, the mean positions are on average $`0^{\prime \prime }.07`$ northeast of the continuum mean positions. Only in slices which pass through the large line-width clump, noted in Figure 5, are the mean positions of the line absorption or apparent or hyperfine optical depth southwest of the continuum positions. The continuum is the widest of the three maps, while the hyperfine optical depth is the narrowest.
We find that the highest column density gas is localized in clumps, while lower column density gas obscures the entire UCHii region. Comparing the two optical depth maps in Figure 2, we note that the hyperfine optical depth map looks โclumpierโ. The flux in the map is mostly collected into peaks about $`0^{\prime \prime }.25`$ to $`0^{\prime \prime }.5`$ (1500 to 3000 AU) in size, while the apparent optical depth map shows more extended absorption over the whole face of the continuum source. The large line-width clump at the western edge is visible in both maps, as is the clump at the eastern edge. Since the hyperfine optical depth is sensitive to much larger optical depths than the apparent optical depth, we interpret the difference in the maps to imply that, while there is an extended high density envelope, the highest column densities are achieved only in highly localized areas.
We find that the characteristic size scale of the kinematic pattern is much larger than the characteristic size scales of fluctuations in the optical depth or continuum maps. Figure 5 shows the first and second moments of the line absorption in the main hyperfine component. Sollins et al. (2005) has interpreted the bulls-eye pattern in the first moment map as showing simultaneous infall and rotation in a rotating, quasi-spherical, molecular accretion flow. The velocity field appears smooth across the absorption region, with the velocity gradient varying slowly. The size scale of this velocity pattern is visibly much larger than either the size of the continuum structures, or the size of the clumpiness of the optical depth. The clumpiness in the optical depth seems to have no effect on the velocity pattern. The velocity pattern is established for the core as a whole, while the optical depth appears to be picking out over-densities which do not depart from the general flow. The map of the second moment shows that the line width is fairly constant across most of the face of the UCHii region, at around 1.8 km s<sup>-1</sup> (FWHM = 4.2 km s<sup>-1</sup>). One spot on the western edge of the UCHii region shows a much larger width, $`>3`$km s<sup>-1</sup>(FWHM$`>7`$km s<sup>-1</sup>). Interestingly, the bulls-eye pattern in the first moment map shows no real effect of the anomalous large line width clump. That location does not stand out at all in the first moment map, which means that the broadening at this point must be symmetric around a central velocity which fits with the overall velocity field.
The position-velocity cuts highlight both the kinematic pattern seen in the first moment, and also the lack of high optical depth gas in the southwest of the UCHii region. Figure 6 shows two position velocity cuts through the cube of apparent optical depths. The upper panel shows a cut from northwest (negative position) to southeast (positive position). The lower panel shows a cut from southwest (negative position) to northeast (positive position). The largely spherical infall can be seen in the lower panel as a backwards โCโ shape. Only the front side of the infall can be detected since we are seeing the line in absorption. In the upper panel, the backwards โCโ shape is seen again, but this time tilted, showing the effect of rotation. The lower panel is a cut along the axis of rotation, so only infall is seen. The upper panel is a cut in the plane rotation, so both infall and rotation are seen. In addition, it should be noted that the satellite line can be seen tracing all of the absorption in the upper panel. The satellite appears everywhere in the plane of rotation. But in the lower panel, the satellite fades out as the cut approaches the southwest side of the UCHii region. The satellite line has an intrinsic line strength of roughly 3% of the main line, so it traces only the highest optical depth gas. Thus the highest optical depth gas is missing on the southwest side of the nebula.
The filling factor of the absorbing gas is large on size scales smaller than a synthesized beam, despite the apparent clumpiness in the optical depth maps. We have noted above that in many pixels, the main line saturates, i.e., in the central channels of the absorption line, there is no detectable flux. When absorbing gas has a non-zero filling factor, the ratio of the depth of the absorption to the continuum strength is related to the optical depth and the filling factor by
$$\frac{T_{line}}{T_{cont}}=\mathrm{\Phi }(1e^\tau )<\mathrm{\Phi }$$
(1)
where $`T_{line}`$ is the depth of the absorption, $`T_{cont}`$ is the strength of the continuum emission, $`\mathrm{\Phi }`$ is the filling factor, and $`\tau `$ is the optical depth. It is impossible to tell the true optical depth when the line has saturated, only a lower limit can be determined. However, the above inequality shows that, no matter what the optical depth is, saturation is only possible where the filling factor is close to one. Furthermore, every measurement of $`\frac{T_{line}}{T_{cont}}`$ puts a lower limit on $`\mathrm{\Phi }`$. The lower limit on the filling factor is 0.7 in more than 75% of all the points in which main line absorption is detected, and more than 0.9 in 55% of those points.
The sharp-edges of the emission seen in the continuum map are absent in the optical depth maps. The โVโ shaped cavity on the northeast side of the UCHii region has very sharp edges, as does the spur to the south. These were interpreted as being the sides of an outflow cavity. The arcs of continuum emission to the east also have sharp edges on the sides facing the UCHii region. These arcs were interpreted by Sollins et al. (2005) as ionized edges of clumps of molecular material. Photons leaking out of the central UCHii region could ionize these clumps externally, naturally creating the arcs, all of which have sharp edges pointing back toward the central source. By contrast, none of the structures in the molecular material have such sharp edges.
## 4 Discussion
We draw two conclusions from the observational results. First, while the kinematics of the accretion flow are quasi-spherical with slow rotation, the density structure appears flattened and disk-like. Second, the plane of the flattening is inclined to the line of sight.
### 4.1 The Molecular โDiskโ
The densest part of the accretion flow is clearly flattened. The flattening is clear when one compares the apparent optical depth map to the map of the hyperfine optical depth. Figure 2 shows both maps. The hyperfine optical depth, which is sensitive to much higher column densities than the apparent optical depth, is distributed quite narrowly along a line perpendicular to the axis of rotation. Only the large line-width clump in the west deviates from the mid-plane. We have calculated the width of the hyperfine optical depth and the apparent optical depth along 15 slices parallel to the axis of rotation (as described above). Figure 4 shows the widths perpendicular to the disk plane for each of the four maps, continuum, velocity integrated absorption, apparent optical depth, and hyperfine optical depth. Again, the slices which include the large line-width clump stand out. Otherwise, the velocity integrated apparent optical depth, which is sensitive only to lower column density gas, is distributed more broadly, while the velocity integrated hyperfine optical depth, which is sensitive to much higher column densities is narrower. The average 2nd moment of the slices of the continuum map is $`0^{\prime \prime }.36`$, for the apparent optical depth map is $`0^{\prime \prime }.23`$, and for the hyperfine optical depth map it is $`0^{\prime \prime }.19`$. We conclude that the highest density gas is collected in a flattened structure in the mid-plane.
We make three specific predictions of what would be observed if there were a geometrically thin, optically thick accretion disk around G10.6, like those accretion disks seen around low-mass stars. Imagine such a disk-UCHii region system schematically like the planet Saturn and its rings, with the equator inclined relative to the line of sight so that the south pole is visible. (We will discuss the inclination of the rotation axis in G10.6 below) The planet is the UCHii region, the disk is the rings. The northern hemisphere is obscured by the rings, and the southern hemisphere is not. There is a sharp edge to the obscuration, not a gradual edge, where the planet emerges from behind the rings. In the case of a thin-thick molecular disk around a UCHii region, we expect that the molecular absorption will be strong on one side and much weaker on the other. Unlike the case of Saturnโs rings, we do not expect the obscuration to be zero in the southern hemisphere where the disk is behind the UCHii region, because the whole object is embedded in a molecular cloud. But the difference in absorption above and below the disk should be great. Also, if the disk is really geometrically thin compared to the size of the UCHii region, we expect there to be a sharp dividing line between the obscured side and the unobscured side, just like in the Saturnโs rings analogy. Departing from the planet-ring analogy, we can also predict that in a disk-UCHii region system, the absorbing gas in the molecular disk would be well homogenized. Such a disk would only form if the gas were rotationally supported. So the rotation time-scale would be much smaller than the infall time-scale, and any inhomogeneities entering the disk would be quickly smoothed out by differential rotation. The predictions for a thin-thick molecular disk surrounding an UCHii region are a large difference in optical depths from one side to the other, a sharp dividing line between the two sides, and structurally smooth absorbing material.
G10.6 does not show evidence for a geometrically thin, optically thick disk, and in fact it fits none of our three observational predictions for a thin-thick disk. First, Figure 2 shows that there is very high optical depth gas over most of the face of the UCHii region. While the southwestern edge has less high optical depth gas than the rest, Figure 2 shows no directional preference at all for the apparent optical depth. The only continuum emission without detectable absorption is the southern spur (to which we will return below). Second, there is clearly no sharp dividing line, just a general trend of the higher optical depth gas to be thinner. Third, the absorbing material is inhomogeneous. At our full resolution, the optical depth varies greatly on size scales ($`0^{\prime \prime }.25`$) much smaller than the size of the UCHii region ($`12^{\prime \prime }`$), and larger than synthesized beam ($`0^{\prime \prime }.1`$). The existence of the arcuate structures to the east has been interpreted as evidence that the surrounding molecular medium is clumpy, and that the arcs are caused by ionizing photons leaking out from the central UCHii region (Sollins et al., 2005). We can confirm the clumpiness of the accretion flow with the optical depth maps, which show variations in integrated optical depth of as much as a factor of eight in a projected distance of less than 1000 AU. While not all of the variations need to be attributed to variations in column density, the other factors which contribute to the optical depth, excitation temperature and ammonia abundance, might not be expected to vary greatly in the molecular gas. The gas distribution in G10.6 is not a geometrically thin, optically thick accretion disk. Instead, the gas is in a flattened, slowly rotating, molecular accretion flow.
Compared to the central mass, the mass of the accretion flow is appreciable. Based on the hyperfine optical depth, and assuming a temperature and ammonia abundance, we can calculate the total molecular mass seen in absorption in the accretion flow. Using 150 K for the excitation temperature of the gas (Sollins et al., 2005), the peak column density of ammonia is greater than $`1.2\times 10^{17}`$ cm<sup>-2</sup>. This a lower limit because there are points at which the absorption in the satellite line saturates. Our integral over the entire map therefore gives a lower limit. Assuming the ammonia abundance is $`10^7`$ relative to $`\mathrm{H}_2`$ (van Dishoeck & Blake, 1998), and adding a factor of 2 since we only detect the front half of the accretion flow in absorption, the total molecular gas mass in the molecular accretion flow is greater than 72 M. The assumed abundance is the largest source of error here and could be wrong by as much as an order of magnitude in either direction. Using the radius of the UCHii region to set the size scale, the implied mass accretion rate is 0.02 Myr<sup>-1</sup>. Sollins et al. (2005) calculate that the central mass responsible for the infall is roughly 150 M. It is entirely possible that the mass of the accretion flow is comparable to the central stellar mass.
The total continuum flux is 2.44 Jy, so assuming constant density, and electron temperature of 10,000 K and a physical size of 8500 AU, the mass of the ionized gas is 0.22 M. Just by looking at the continuum map it is clear that the density is not uniform, and since we know there is ongoing accretion, the density profile should be proportional to $`r^{3/2}`$. The mass, however, depends strongly on the total size of the region in question. So the more spread out ionized gas will dominate the mass. For example, even if we associate all the emission from the marginally resolved peak of the continuum emission with a single density enhancement, the most ionized gas mass we can possibly associate with the peak is 0.0035 M. Keto (2003) pointed out that, when estimating the Lyman continuum flux necessary to achieve ionization balance, the density gradient can be very important. This is because the recombination rate is proportional to density squared. Because mass is proportional to density only to the first power, the total mass is less sensitive to small high density pockets, and will be dominated by the larger scale structures. Only for a density profile steeper than $`r^2`$ will the mass be dominated by smaller radii rather than larger radii.
### 4.2 Inclination of the โDiskโ
Sollins et al. (2005) determined that the rotation axis of the molecular accretion flow points northeast when projected into the plane of the sky. Based on our data, the axis of rotation appears to be tipped away from the observer in the northeast. Even though the molecular gas is not in a rotationally supported, geometrically thin disk, the gas distribution is flattened, with denser gas collected in the plane of the equator of the system. Since that plane is tipped, we expect the densest gas to be preferentially in the northeast. We have seen a hint of this already in Figure 6, where the absorption from the main hyperfine component extends right down to the southwest edge of the continuum, while the absorption from satellite component does not. To test for this inclination quantitatively, we have analyzed the maps in Figure 2. For each of the four maps we have made fifteen slices parallel to the axis of rotation, southwest-northeast. Then along each slice we calculate a flux-weighted average position, i.e., the first moment of the slice. In the continuum map, the average positions follow the line closely, and do not systematically deviate in one direction or the other. By contrast, in the hyperfine optical depth map, the average positions are all to the northeast, except at the two points where the anomalous large line-width clump has dragged the average to the southwest. The highest density gas is, on average, $`0^{\prime \prime }.08`$ northeast of the projected mid-plane of the continuum, and farther to the northeast if the large line-width clump is excluded. Because the area over which the optical depth can be calculated is defined by the extent of the continuum emission, the mean positions cannot be wildly different. This emphasizes the significance of the offsets in position of the hyperfine optical depth from the continuum. These offsets are direct evidence that the โdiskโ is tilted, and the axis of rotation is inclined. Using the radius of the UCHii region, $`1^{\prime \prime }.1`$, as a lower limit for the radius of the disk, a $`0^{\prime \prime }.08`$ offset is consistent with a tilt of the disk of $`4^{}`$. Other direct evidence for this inclination has been found by Keto & Wood (2005), who have detected red-shifted H66$`\alpha `$ emission to the northeast of the UCHii region, within the notch in the continuum emission on that side. They have interpreted that gas as an outflow.
Another clue as to the inclination of the rotation axis is the strength of the absorption on the narrow spurs of continuum. On the northeast side of the continuum source there is the โVโ-shaped notch mentioned above. Sollins et al. (2005) interpret this as a possible outflow cavity, and Keto & Wood (2005) have confirmed this. On the southwest side, Keto & Wood (2005) detect no corresponding blue-shifted outflow, but the continuum does show a sharp edged spur, reminiscent of the notch in the northeast. While the spur-like structures in the northeast show strong absorption in both apparent and hyperfine optical depth, the spur in the southwest shows only weak absorption in the apparent optical depth map, and no detectable hyperfine optical depth. This is further evidence that the high density gas in front of the continuum source is in the northeast because of a flattened density profile and the rotation axis being tipped away from the observer in the northeast. However we should note a possible alternate explanation. We cannot rule out the possibility that the southern spur is the limb-brightened edge of a photo-ionized molecular clump, just like the arcs to the east. In that case, the clump might be closer along the line of sight, not physically associated with the main UCHii region, and therefore not obscured by the densest molecular gas.
### 4.3 A New Phase of Massive Star Formation
The data on G10.6 are unique in the study of accretion onto massive stars because of their spatial resolution, and because the interpretation is not model dependent. Disk-like structures have been detected around a number of very early B type stars. Chini et al. (2004) detected a 10000 AU, morphologically disk-like structure at 550 AU resolution. Kinematic observations at 13000 AU resolution show that the disk is rotating on that larger size scale. In IRAS 20126+4104, IRAS 18089-1732, AFGL5142, flattening and rotation in dense molecular gas has been detected at roughly 5000 AU resolution. In all three of those cases the sources have infrared luminosities corresponding to early B stars (Zhang et al., 1998, 2002; Beuther et al., 2004). In G192.16-3.82, Shepherd & Kurtz (1999) find a velocity pattern consistent with rotation in water maser spots at a 1000 AU spatial scale, also around an early B star. In all these cases, molecular gas is found to be in rotation, in some cases apparently Keplerian rotation, around early B type or even late O type stars. All are consistent with the existence of rotating, molecular accretion disks which, in many respects, are larger versions of the disks observed around low-mass protostars. In contrast to previous work, in G10.6 we have 500 AU physical resolution in the thermally emitting molecular gas. The spatial resolution is enough to completely resolve the motions involved. The kinematics and optical depths are fairly unambiguous. The densest gas is flattened, and the velocities clearly show infall and slow rotation.
G10.6 itself contrasts with objects from previous massive-star-disk studies in that it cannot be interpreted as a scaled up version of low-mass star formation. All the cases cited above (IRAS 20126, IRAS 18089, G192.16, M17, and AFGL5142) are consistent with the central object being a single stellar system of up to 20 M. In most of these objects the existence of a bipolar outflow indicates ongoing accretion. The analogy to the formation mechanisms of low-mass stars is fairly straightforward, scaled up in size and mass, although a key difference is ratio of disk mass to stellar mass, small for low mass stars, but apparently large for high mass stars. In G10.6, the central source is at least 150 M, close to $`10^6`$ L, and is almost certainly not a single star or binary. At the edge of the UCHii region, at a radius of 5000 AU, the molecular gas is clearly moving inward, and Keto (2002a) detect inward motions in the ionized gas down to radii of less than 1000 AU. While we cannot say how or if the inward moving gas actually accretes onto one or more of the central stars, we can say with great certainty that inward motion continues in the molecular gas from the 0.5 pc scale (Ho & Haschick, 1986; Keto et al., 1988) down to thousands of AU. This is a single continuous accretion flow traceable over two orders of magnitude in size, toward a group of young massive stars. This suggests a completely different phase or mode of star formation than that seen in low mass stars, or in the preceding examples of individual massive young stars.
## 5 Summary
We have utilized the strong 23 GHz radio continuum emission from the UCHii region G10.6-0.4 to serve as a backlight for examining the foreground molecular material seen in absorption. Using the VLA, we have achieved very high angular ($`0^{\prime \prime }.1`$) and spatial (500 AU) resolution. In the past, such resolutions have not been possible for studying the circumstellar environment of massive young stars. Making use of the hyperfine structure of the NH<sub>3</sub>(3,3) inversion line, we are sensitive to optical depths of up to 80. This allowed us to investigate the structure of the densest circumstellar material. We find that in the densest material, the structure is flattened, with an aspect ratio of 5. The structure is displaced with respect to the mean continuum emission, consistent with a tilt of the disk along the line of sight at $`4^{}`$, away from the observer in the northeast. The flattened structure has a mass of 72 M, much larger than the ionized gas in the HII region of 0.2 M. The velocity pattern within the circumstellar material, as well as its clumpiness, suggest a dynamically collapsing structure which is not centrifugally supported. The implied infall rate is very high, on the order of 0.02 Myr<sup>-1</sup>. The kinematics of the circumstellar material, which agree with the kinematics of the ionized gas within the Hii region, suggest that this infalling material continues across the ionization front. Because we do not know how much mass is actually being accreted by the stars, or how much mass is leaving the system in the outflow, it is impossible to know for sure whether G10.6 is in quasi-static equilibrium, or if it is evolving dynamically. However, the very high mass and luminosity involved mean that this is a different type of object than the individual high-mass protostars which have been investigated in the past.
The authors would like to acknowledge Eric Keto and Qizhou Zhang for their helpful comments and suggestions in the preparation of this paper.
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# Geometric transitions and mixed Hodge structures
## 1 Introduction
Large $`N`$ duality is a relation between open and closed string theories on two different Calabi-Yau manifolds connected by an extremal transition. This relation was originally formulated in the context of topological A-model for a local conifold transition and it was extended to more general noncompact toric Calabi-Yau manifolds in .
This paper is part of a long term project aimed at understanding large $`N`$ duality for extremal transitions between compact Calabi-Yau manifolds. A first step in this direction has been made in for topological A-models. Here we will be concerned with large $`N`$ duality in the topological B-model. Open-closed duality for topological B-strings was first developed in for a special class of noncompact toric Calabi-Yau manifolds employing a remarkable relation between holomorphic Chern-Simons theory and random matrix models.
In contrast with the A-model, the topological B-model on compact Calabi-Yau spaces has not been given so far a rigorous mathematical description. However it is well known that the genus zero topological closed string amplitudes can be expressed in Hodge theoretic terms using the formalism of special geometry. On the other hand disc level topological open string amplitudes associated to D-branes wrapping curves in Calabi-Yau threefolds can also be given a geometric interpretation in terms of Abel-Jacobi maps . Higher genus amplitudes do not have a pure geometric interpretation. In principle one would have to quantize Kodaira-Spencer theory coupled to holomorphic Chern-Simons theory on a compact Calabi-Yau space, which is a very hard task at best.
In this paper we will formulate and prove a first order B-model duality statement for general conifold transitions between compact Calabi-Yau spaces. By first order duality we mean a correspondence between topological disc amplitudes on the open string side and first order terms in a suitable expansion of the holomorphic prepotential on the closed string side. The expansion is taken around an appropriate stratum parameterizing nodal Calabi-Yau spaces that admit a projective crepant resolution.
Using special geometry, in section two we show that the first order terms in this expansion admit a intrinsic geometric interpretation in terms of degenerations of Hodge structures. In section three we will show that the first order duality statement follows from a Hodge theoretic result relating two different mixed Hodge structures. The main element in the proof is the Clemens-Schmid exact sequence.
A connection between mixed Hodge structures and B-model topological disc amplitudes on toric Calabi-Yau manifolds has been previously developed in . This machinery has been applied to first order large $`N`$ duality for toric Calabi-Yau manifolds in . Our approach is different and can be used to extend the $`๐`$-model large $`N`$ duality beyond disc level. Some progress along these lines for an interesting class of noncompact transitions is reported in the companion paper .
Acknowledgments. We are very grateful to Bogdan Florea and Antonella Grassi for collaboration at an early stage of the project and many useful discussions. The work of D.-E. D. has been partially supported by an Alfred P. Sloan fellowship. R.D. is partially supported by NSF grant DMS 0104354 and FRG grant 0139799 for โThe Geometry of Superstringsโ. T.P. is partially supported by NSF grants FRG 0139799 and DMS 0403884.
## 2 B-Model transitions and periods
In this section we discuss the first order behavior of B-model geometric transitions associated to conifold singularities of Calabi-Yau spaces. For us a Calabi-Yau space will be a (possibly singular) complex Gorenstein quasi-projective variety $`X`$ which has a trivial canonical class. In addition, in the singular case we will require that $`X`$ has a Kรคhler crepant resolution.
This convention is somewhat broader than the usual notion of a Calabi-Yau space used in physics, where one requires that $`X`$ is a complex analytic space equipped with a Ricci flat Kรคhler metric. We will be primarily interested in moduli spaces parameterizing Calabi-Yau structures. These moduli spaces can have different components corresponding to different sets of values of the topological invariants of $`X`$. A geometric (or extremal) transition is a process connecting two connected components of the moduli space through a degeneration. Schematically an extremal transition is captured in a diagram
where $`X_l`$ is a smooth Calabi-Yau, is a degeneration of $`X_l`$ to a Calabi-Yau variety $`X_m`$ having only ordinary double points, and $`\stackrel{~}{X}_{\stackrel{~}{m}}X_m`$ is a crepant quasi-projective resolution of $`X_m`$.
The standard example of such a situation is the local conifold transition:
###### Example 2.1
Take $`\{X_\mu \}_\mu `$ to be the one parameter family of $`3`$-dimensional affine quadrics $`X_\mu =\{(x,y,z,w)^4|xyzw=\mu \}`$. When $`\mu 0`$ we get a degeneration of $`X_\mu `$ to the $`3`$-dimensional quadratic cone $`xy=zw`$. To complete this degeneration to a transition
we take $`\stackrel{~}{X}_0`$ to be one of the small resolutions of $`X_0`$, i.e. the blow-up of $`X_0`$ along the Weil divisor in $`X_0`$ corresponding to one of the two rulings in the base $`^1\times ^1`$ of the cone. For future reference we note also that as an abstract variety $`\stackrel{~}{X}_0`$ is isomorphic to the total space of $`๐ช(1)๐ช(1)^1`$, and the map $`\stackrel{~}{X}_0X_0`$ is the natural map contracting the zero section of $`๐ช(1)๐ช(1)`$.
The geometric transition from $`X_l`$ to $`\stackrel{~}{X}_{\stackrel{~}{m}}`$ changes the topology. Indeed, the process of degenerating $`X_l`$ to $`X_m`$ collapses some $`3`$-spheres in $`X_l`$ to the singular points of $`X_m`$. These spheres are the vanishing cycles of the degeneration . On the other hand, the small resolution $`\stackrel{~}{X}_{\stackrel{~}{m}}X_m`$ replaces each singular point of $`X_m`$ by a copy of $`^1S^2`$. So in the passage from $`X_l`$ to $`\stackrel{~}{X}_m`$ we deleted some $`3`$ spheres from $`X_l`$ and glued $`2`$-spheres in their place. It is important to note that this process not only changes the topology but also alters the type of geometry of the Calabi-Yau spaces in question. More precisely the transition interchanges holomorphic and symplectic data: the exceptional $`^1`$โs in $`\stackrel{~}{X}_{\stackrel{~}{m}}`$ are holomorphic curves, and conjecturally the vanishing $`3`$-spheres in $`X_l`$ can be chosen to be special Lagrangian submanifolds for the Kรคhler form on $`X_l`$. In particular, in an A-model transition one expects a relationship between open Gromov-Witten invariants on $`X_l`$ (with boundaries on the vanishing $`3`$-spheres) and closed Gromov-Witten invariants on $`\stackrel{~}{X}_m`$. The precise form of such a relationship is the content of the A-model large $`N`$ duality which has been extensively analyzed in the literature, see and references therein. Here we propose a geometric description for the B-model counterpart of the large $`N`$ duality and investigate some of the mathematical and physical consequences of our proposal.
Suppose $`๐ณ`$ is a fixed component of the moduli space of Calabi-Yau threefolds with at most ODP singularities. Given a point $`l๐ณ`$ we write $`X_l`$ for the corresponding Calabi-Yau. We will always assume that for a general $`l๐ณ`$, the variety $`X_l`$ is a smooth (compact or non-compact) Calabi-Yau threefold. In examples we will often take $`X_l`$ to be complete intersections in some toric variety since we want to keep track of the family of $`๐`$-models mirror to the topological $`๐`$-models specified by the $`X_l`$โs.
More precisely, we will look at the subvariety of $`๐ณ`$ parameterizing singular threefolds with ordinary double points which admit a crepant projective resolution. Let $`๐ด`$ be a component of this subvariety and let $`v`$ denote the number of ODPs of $`X_m`$ for a general $`m๐ด`$. In particular on a nearby smooth $`X_l`$ we have a collection of $`v`$ embedded Lagrangian 3-spheres $`L_1,\mathrm{},L_v`$ whose homology classes $`[L_1],\mathrm{},[L_v]H_3(X_m,)`$ vanish under a deformation . Recall that \[18, Theorem 2.9\] in order for $`X_m`$ to admit a projective small resolution, we must have at least one good relation among $`[L_1],\mathrm{},[L{}_{v}{}^{}]`$. That is, in $`H_3(X_l,)`$ we must have a relation of the form
$$\underset{i=1}{\overset{v}{}}c_i[L_i]=0H_3(X_l,),\text{ with }c_i0\text{ for all }i=1,\mathrm{},v.$$
Assuming that this is the case, let $`r1`$ denote the number of relations on the vanishing cycles. Then for a fixed point $`mM`$, $`X_m`$ may have finitely many different projective small resolutions related by flops. This means that the moduli space $`\stackrel{~}{๐ด}`$ of the resolution is a finite to one cover of the component $`๐ด`$. In the following we will denote by $`\rho :\stackrel{~}{๐ด}๐ด`$ the covering map and by $`\stackrel{~}{m}`$ a point of $`\stackrel{~}{๐ด}`$ which projects to $`m๐ด`$. Since above we have restricted our considerations to the moduli space of Calabi-Yau threefolds with at most isolated ODP singularities, this cover is unramified. The branching points of the cover would correspond to singular threefolds with more complicated singularities which have been excluded by our definition of the moduli space $`๐ณ`$.
The exceptional locus of the resolution $`\stackrel{~}{X}_{\stackrel{~}{m}}X_m`$ consists of $`v`$ smooth $`(1,1)`$ rational curves $`C_1,\mathrm{},C_v`$ satisfying $`vr`$ relations in $`H_2(\stackrel{~}{X}_{\stackrel{~}{m}},)`$. Moreover, we have (see e.g. ) the following relations among Betti numbers of $`X_m`$ and $`\stackrel{~}{X}_{\stackrel{~}{m}}`$:
$$\begin{array}{cc}\hfill b_2(\stackrel{~}{X}_{\stackrel{~}{m}})& =b_2(X_m)+r=b_2(X_l)+r,\hfill \\ \hfill b_3(\stackrel{~}{X}_{\stackrel{~}{m}})& =b_3(X_m)(vr)=b_3(X_l)2(vr),\hfill \\ \hfill b_4(\stackrel{~}{X}_{\stackrel{~}{m}})& =b_4(X_m)=b_4(X_l)+r.\hfill \end{array}$$
Large $`N`$ duality conjectures a correspondence between topological string theories defined on the Calabi-Yau manifolds $`X_l`$ and $`\stackrel{~}{X}_{\stackrel{~}{m}}`$ related by a geometric transition. In this paper we will study B-model transitions, in which case the conjecture predicts an equivalence between closed topological strings on $`X_l`$ and open-closed topological strings on $`\stackrel{~}{X}_{\stackrel{~}{m}}`$. In physical terms, the open closed topological string theory on the small resolution $`\stackrel{~}{X}_{\stackrel{~}{m}}`$ is defined by wrapping $`N_i`$ B-branes on the exceptional curves $`C_i`$. It is by now well established that in a rigorous framework B-branes should be described by derived objects. However, for the purposes of the present paper, it suffices to think of a B-brane as an algebraic cycle on $`\stackrel{~}{X}_{\stackrel{~}{m}}`$ of the form $`_{i=1}^vN_iC_i`$. Furthermore we will restrict our considerations to homologically trivial D-brane configurations i.e.
$$\underset{i=1}{\overset{v}{}}N_i[C_i]=0$$
(1)
where $`[C_i]H_2(\stackrel{~}{X}_{\stackrel{~}{m}},)`$ denotes the homology class of $`C_i`$. In principle, the open-closed topological $`๐`$ model should be well defined from a physical point of view for any values of the multiplicities $`N_i`$. However it will become clear later that the above homology condition is required by large $`N`$ duality. While there is no a priori explanation for this condition in topological string theory, in physical superstring theory, this is a direct consequence of the Gauss law for Ramond-Ramond flux. It is quite interesting that the topological version of large $`N`$ duality still requires us to impose the physical Gauss law.
### 2.1 Closed Strings
The central object of study of any topological string theory is the partition function, which is a generating functional for topological string amplitudes. The partition function of the closed topological $`๐`$-model on $`X_l`$ can be written as
$$_{X_l}^{cl}=\underset{g=0}{\overset{\mathrm{}}{}}g_s^{2g2}_{X_l}^g.$$
The genus $`g`$ free energy $`_{X_l}^g`$ is heuristically defined in terms of functional integrals over moduli spaces of maps from compact genus $`g`$ Riemann surfaces to $`X_l`$. In the B-model the functional integral receives contributions only from degenerate maps, which sit on the boundary of the moduli space. For genus zero, the degenerate maps in question are constant maps, and the functional integral reduces to an ordinary integral on $`X_l`$ . Moreover the genus zero free energy depends holomorphically on the complex structure parameters of $`X_l`$. For higher genus, the degenerate maps have a more complicated structure and there is no rigorous mathematical formulation of topological amplitudes. In the following we will restrict ourselves to genus zero topological strings.
Next we will explain the construction of the genus zero free energy $`_{X_l}^0`$ and its relation to the special geometry of the moduli space $`๐ณ`$. Since in the B-model all physical correlators depend on the choice of a global holomorphic three-form, we have to introduce the enlarged moduli space $`๐ณ^{}`$ parameterizing pairs $`(X_l,\mathrm{\Omega }_l)`$ where $`\mathrm{\Omega }_l`$ is a nonzero global holomorphic three-form on $`X_l`$. Note that there is a complex holomorphic line bundle $`๐`$ so that the fiber $`_l`$ is the space of global holomorphic three-forms on $`X_l`$ for any point $`l๐`$. The enlarged moduli space $`๐ณ^{}`$ is isomorphic to the complement of the zero section in the total space of $``$, hence it has the structure of a holomorphic principal $`^\times `$-bundle $`\pi :๐ณ^{}๐ณ`$. Let us denote by $`๐ณ_0`$ the open subspace of $`๐ณ`$ parameterizing smooth varieties $`X_l`$, and by $`๐ณ_0^{}`$ its inverse image in $`๐ณ^{}`$. We also write $`๐ด^{}`$ for the inverse image of $`๐ด`$ in $`๐ณ^{}`$, and $`\stackrel{~}{๐ด}^{}`$ for the enlarged moduli space of the resolution. Note that there is a finite to one unramified cover $`\rho ^{}:\stackrel{~}{๐ด}^{}๐ด^{}`$.
Caution: The previous discussion is somewhat loose. For instance, the moduli $`๐ณ^{}`$ of pairs $`(X_l,\mathrm{\Omega }_l)`$ is the total space of a line bundle $`๐ณ`$ only if we view $`๐ณ`$ as a stack. More importantly we need to make sure that $`๐ณ`$ is a line bundle on $`๐ณ๐ณ_0`$ as well. If we have a universal family $`f:๐ณ๐ณ`$, then $``$ is the pushforward $`f_{}\omega _{๐ณ/๐ณ}`$ of the relative dualizing sheaf, which is locally free by cohomology and base change. Indeed for this we only need to note that for a nodal Calabi-Yau $`X_m`$ we have $`h^0(X_m,K_{X_m})=h^0(X_m,๐ช)=1`$ and $`h^1(X_m,K_{X_m})=h^1(X_m,๐ช)=0`$.
Ideally one would like to define $`_{X_l}^0`$ as an intrinsic global geometric object on the moduli space $`๐ณ^{}`$ which can be locally described as a holomorphic function (for example a section in a certain line bundle.) Unfortunately, there is no such intrinsic construction for $`_{X_l}^0`$. One can construct the three point function, or Yukawa coupling, as a global cubic form on $`๐ณ_0^{}`$ . The genus zero free energy can only be defined locally as a primitive of the Yukawa coupling. This description is of course ambiguous since the Yukawa coupling specifies only the third derivatives of the free energy. Therefore in order to obtain a well defined local function we have to make some choices. Using special geometry (see e.g. ) one can show that a local primitive for the Yukawa coupling is determined by a choice of splitting of the third homology $`H_3(X_l,)`$ into a direct sum $`AB`$ of maximal Lagrangian sublattices.
Recall that we have chosen $`๐ณ`$ to be a component of a moduli space of Calabi-Yau threefolds with at most isolated ODP singularities. $`๐ด`$ is a subvariety of the discriminant parameterizing threefolds with a fixed number $`v`$ of ODPs which admit a crepant projective resolution. In generic situations, the $`(vr)`$ codimensional subvariety $`๐ด`$ of the discriminant can be locally represented as the intersection locus of $`v`$ (local) branches of the discriminant. Therefore we can choose an open subset $`๐ฐ๐ณ`$ so that $`๐ฐ๐ด`$ is the intersection locus of a collection of Weil divisors $`๐ซ_1,\mathrm{},๐ซ_v`$ in $`๐ฐ`$ so that $`vr`$ of them intersect transversely along $`๐ด`$. Moreover $`X_l`$ with $`l๐ฐ`$ is singular if and only if $`l`$ is a point on
$$๐ซ=๐ซ_1+\mathrm{}+๐ซ_v.$$
Let $`๐ฐ^{}`$ denote the inverse image of $`๐ฐ`$ in $`๐ณ^{}`$. In order to write down a local expression for the genus zero free energy we have to introduce special coordinates $`z^\alpha `$ on $`๐ฐ^{}`$ by choosing a symplectic basis of three-cycles $`\{\gamma _{\alpha ,l},\gamma _l^\alpha \}`$, $`\alpha =1,\mathrm{}h^{1,2}(X_l)+1`$ on each threefold $`X_l`$, with $`l๐ฐ๐ซ`$. The symplectic basis of cycles determines a splitting $`H_3(X_l,)=A_lB_l`$ where $`A_l,B_l`$ are complementary maximal Lagrangian sublattices spanned by the cycles $`\{\gamma _{\alpha ,l}\}`$ and $`\{\gamma _l^\alpha \}`$ respectively. Note that we have a monodromy transformation
$$T_i:H_3(X_l,)H_3(X_l,),T_i(\mathrm{\Gamma })=\mathrm{\Gamma }+\mathrm{\Gamma },\xi _i\xi _i$$
(2)
associated to each component $`๐ซ_i`$, $`i=1,\mathrm{},v`$ of the discriminant in $`๐ฐ`$, where $`\xi _iH_3(X_l,)`$ is the corresponding vanishing cycle. We have denoted by $`,`$ the intersection pairing in $`H_3(X_l,)`$. In order to obtain single valued coordinates on $`๐ฐ^{}`$, we have to choose the lattices of $`A`$-cycles so that $`A_l`$ is contained in the fixed locus of $`T_i`$ for each $`i=1,\mathrm{},v`$. Then the special projective coordinates are given by
$$z^\alpha =_{\gamma _{\alpha ,l}}\mathrm{\Omega }_{X_l}$$
where $`\mathrm{\Omega }_{X_l}`$ is a global holomorphic three-form on $`X_l`$. Since the cycles $`\gamma _{\alpha ,l}`$ are fixed by the monodromy transformations (2), the coordinates $`z^\alpha `$ extend as single valued holomorphic functions over the entire open set $`๐ฐ^{}`$.
The genus zero free energy โ or, adopting special geometry terminology, the holomorphic prepotential โ is a multivalued holomorphic function on $`๐ฐ^{}๐ณ_0^{}`$ given by
$$^0=\frac{1}{2}\underset{\alpha =1}{\overset{h^{1,2}(X_l)+1}{}}z^\alpha \mathrm{\Pi }_\alpha $$
(3)
where
$$\mathrm{\Pi }_\alpha =_{\gamma _l^\alpha }\mathrm{\Omega }_{X_l}$$
are the periods of the holomorphic three-form on the $`B`$-cycles $`\gamma _l^\alpha `$. We also have the special geometry relations
$$\mathrm{\Pi }_\alpha =\frac{^0}{z^\alpha }.$$
(4)
which will be useful later in the paper.
### 2.2 Open Strings
Let us now discuss the topological open-closed B-model on the small resolution $`\stackrel{~}{X}_{\stackrel{~}{m}}`$. In open-closed topological string theory, one would like to integrate over maps from genus $`g`$ Riemann surfaces with $`h`$ boundary components to $`\stackrel{~}{X}_{\stackrel{~}{m}}`$ mapping the boundary components to the exceptional curves $`C_i`$. The partition function of the theory is a generating functional of the form
$$_{\stackrel{~}{X}_{\stackrel{~}{m}},\{C_i\},\{N_i\}}=\underset{g=0}{\overset{\mathrm{}}{}}\underset{h_1,\mathrm{},h_v=0}{\overset{\mathrm{}}{}}g_s^{2g2+h}_{\stackrel{~}{X}_{\stackrel{~}{m}},\{C_i\}}^{g,\{h_i\}}\underset{i=1}{\overset{v}{}}N_i^{h_i}$$
(5)
where $`h=_{i=1}^vh_i`$. The coefficients $`_{\stackrel{~}{X}_{\stackrel{~}{m}},C_i}^{g,\{h_i\}}`$ are heuristically defined in terms of functional integrals over maps from genus $`g`$ Riemann surfaces with $`h`$ boundary components to $`\stackrel{~}{X}_m`$ mapping $`h_i`$ boundary components to the curve $`C_i`$. In the B-model they are expected to depend on complex moduli for the data $`(\stackrel{~}{X}_{\stackrel{~}{m}},C_i)`$. Since the exceptional curves are rigid in the threefold, it follows that they only depend on the complex structure parameters of $`\stackrel{~}{X}_{\stackrel{~}{m}}`$.
According to , the open-string path integral localizes on degenerate maps, just as in the closed string situation. However, degenerate open string maps collapse a Riemann surface with boundary considered as a ribbon graph, to the corresponding graph embedded in $`\stackrel{~}{X}_{\stackrel{~}{m}}`$. In particular, the degenerate maps can have a nontrivial structure even at genus zero. For this reason there is no rigorous mathematical formulation of open-closed amplitudes except for $`g=0`$ and $`h=0,1`$ when degenerate maps are constant maps. In this paper we will focus only on these two cases.
The term corresponding to $`g=h=0`$ in (5) is the genus zero closed string free energy $`_{\stackrel{~}{X}_{\stackrel{~}{m}}}^0`$ which was introduced in section 2.1. The term corresponding to $`g=0`$, $`h=1`$ represents the disc open string free energy which is determined by $`3`$ dimensional chain integrals of the holomorphic three-form $`\mathrm{\Omega }_{\stackrel{~}{X}_{\stackrel{~}{m}}}`$. To explain this construction recall that we are considering D-brane configurations satisfying the homology constraint (1). For a fixed $`\stackrel{~}{m}`$, the solutions to this equation are in one-to-one correspondence with points in the lattice
$$\mathrm{\Lambda }_{\stackrel{~}{m}}=\text{ker}\left(H_2(C)H_2(\stackrel{~}{X}_{\stackrel{~}{m}})\right)$$
where $`C=_{i=1}^vC_i`$. Recall that in our set-up, the points in $`๐ด`$ parameterize threefolds with isolated ODP singularities, so that the cover $`\stackrel{~}{๐ด}๐ด`$ is unramified. Then the lattices $`\mathrm{\Lambda }_{\stackrel{~}{m}}`$ span a locally constant sheaf when $`\stackrel{~}{m}`$ varies in $`\stackrel{~}{๐ด}`$. A flat section of this sheaf parameterizes a pair $`(\stackrel{~}{X}_{\stackrel{~}{m}},C_{N,\stackrel{~}{m}})`$ where $`C_{N,\stackrel{~}{m}}`$ is a homologically trivial algebraic cycle on $`\stackrel{~}{X}_{\stackrel{~}{m}}`$ of the form $`C_{N,\stackrel{~}{m}}=_{i=1}^vN_iC_i`$.
Regarding $`C_{N,\stackrel{~}{m}}`$ as a topological brane on $`\stackrel{~}{X}_{\stackrel{~}{m}}`$, the disc partition function is a period of the holomorphic three-form $`\mathrm{\Omega }_{\stackrel{~}{X}_{\stackrel{~}{m}}}`$ over a 3-chain, defined as follows. Consider the relative homology sequence for a pair $`(\stackrel{~}{X}_{\stackrel{~}{m}},C)`$, with $`\stackrel{~}{X}_{\stackrel{~}{m}}`$ smooth
$$0H_3(\stackrel{~}{X}_{\stackrel{~}{m}},)H_3(\stackrel{~}{X}_{\stackrel{~}{m}},C,)\stackrel{\delta }{}\mathrm{\Lambda }_{\stackrel{~}{m}}0.$$
Given a homologically trivial algebraic cycle $`C_{N,\stackrel{~}{m}}`$ on $`\stackrel{~}{X}_{\stackrel{~}{m}}`$, we can find a relative cycle $`\stackrel{~}{\mathrm{\Gamma }}_{N,\stackrel{~}{m}}H_3(\stackrel{~}{X}_{\stackrel{~}{m}},C)`$ so that
$$\delta \stackrel{~}{\mathrm{\Gamma }}_{N,\stackrel{~}{m}}=C_{N,\stackrel{~}{m}}$$
According to , for fixed $`N_i`$ satisfying the condition (1) the disc partition function is the holomorphic function (= the Griffiths normal function associated to $`C_{N,\stackrel{~}{m}}`$) on the enlarged moduli space $`\stackrel{~}{๐ด}^{}`$ given by
$$\stackrel{~}{}_{N,\stackrel{~}{m}}=\frac{1}{g_s}\underset{i=1}{\overset{v}{}}_{\stackrel{~}{X}_{\stackrel{~}{m}},\{C_i\}}^{0,e_i}N_i=_{\stackrel{~}{\mathrm{\Gamma }}_{N,\stackrel{~}{m}}}\mathrm{\Omega }_{\stackrel{~}{X}_{\stackrel{~}{m}}}.$$
(6)
Here as usual we write $`e_i`$ for the $`i`$-th vector in the standard basis of the lattice $`^v`$.
Note that the relative cycle $`\stackrel{~}{\mathrm{\Gamma }}_N`$ is only determined modulo elements in $`H_3(\stackrel{~}{X}_{\stackrel{~}{m}},)`$, therefore the superpotential (6) is defined modulo periods. In general there is no preferred choice for a three-chain $`\stackrel{~}{\mathrm{\Gamma }}_N`$, hence we will regard this as a discrete ambiguity in the disc partition function.
### 2.3 Duality
Large $`N`$ duality predicts a correspondence between topological open and closed strings on a pair of Calabi-Yau threefolds related by an extremal transition. A precise mathematical statement has not been formulated for $`๐`$-model transitions between compact Calabi-Yau threefolds. On general grounds, large $`N`$ duality is expected to relate genus zero open string amplitudes with $`h`$ boundary components on $`\stackrel{~}{๐ด}^{}`$ to terms of order $`h`$ in the Taylor expansion of the closed string partition function on $`๐ณ^{}`$ near $`๐ด^{}`$. In this subsection our goal is to formulate a precise mathematical statement for genus zero terms with $`h=0,1`$.
We will show that for a certain choice of special coordinates $`z^\alpha `$, $`\alpha =1,\mathrm{},h^{1,2}(X_l)+1`$ in a neighborhood of $`๐ด^{}`$ in $`๐ณ^{}`$, the closed string prepotential $`^0`$ and its derivatives $`\left(\frac{^0}{z^\alpha }\right)`$ have well defined limits $`_๐ด^{}^0`$, $`\left(\frac{^0}{z^\alpha }\right)_๐ด^{}`$ along $`๐ด^{}`$. Moreover, the following relations hold
$$\stackrel{~}{^0}=\rho _{}^{}{}_{}{}^{}(_๐ด^{}^0),\stackrel{~}{}_i=\rho _{}^{}{}_{}{}^{}\left(\frac{^0}{z^i}\right)_๐ด^{}$$
(7)
where $`\stackrel{~}{^0}`$ denotes the closed string prepotential on the moduli space $`\stackrel{~}{๐ด^{}}`$, $`\stackrel{~}{}^i`$, $`i=1,\mathrm{},vr`$ are open string superpotentials of the form (6), and $`i=1,\mathrm{},vr`$ labels the normal directions to $`๐ด^{}`$ in $`๐ณ^{}`$. The main point we would like to make is that these relations are a corollary of the following intrinsic geometric results which will be proved in the next section.
For any $`\stackrel{~}{m}\stackrel{~}{๐ด}`$, and $`m๐ด`$ so that $`m=\rho (\stackrel{~}{m})`$, the contraction map $`(\stackrel{~}{X}_{\stackrel{~}{m}},C)(X_m,\mathrm{Sing}(X_m))`$ induces an isomorphism
$$H_3(\stackrel{~}{X}_{\stackrel{~}{m}},C,)H_3(X_m,).$$
Choose a retraction map $`\mathrm{๐ซ๐๐ญ}:๐ณ๐ณ_{|๐ด}`$ from the universal family on $`๐ณ`$ onto the universal family on $`๐ด`$ (This can always be done in a neighborhood of a point $`m๐ด`$.) Suppose $`\mathrm{\Gamma }`$ is a multivalued section of the sheaf of third homology groups $`H_3(X_l,)`$ on $`๐ณ`$, then the image of $`\mathrm{\Gamma }`$ under $`\mathrm{๐ซ๐๐ญ}`$ induces a multivalued section $`\stackrel{~}{\mathrm{\Gamma }}`$ of the sheaf of relative homology groups $`H_3(\stackrel{~}{X}_{\stackrel{~}{m}},C;)`$ on $`\stackrel{~}{๐ด}`$.
Let
$$\mathrm{\Pi }=_{\mathrm{\Gamma }_l}\mathrm{\Omega }_{X_l},\stackrel{~}{\mathrm{\Pi }}=_{\stackrel{~}{\mathrm{\Gamma }}_{\stackrel{~}{m}}}\mathrm{\Omega }_{\stackrel{~}{X}_{\stackrel{~}{m}}}$$
be the periods of the respective holomorphic 3-forms, viewed as multivalued holomorphic functions on $`๐ณ_0^{}`$ and $`\stackrel{~}{๐ด}^{}`$. Then $`\mathrm{\Pi }`$ induces a multivalued holomorphic function $`\mathrm{\Pi }_๐ด^{}`$ on $`๐ด^{}`$ and
$$\stackrel{~}{\mathrm{\Pi }}=\rho _{}^{}{}_{}{}^{}(\mathrm{\Pi }_๐ด^{}).$$
(8)
For any points $`l๐ณ_0`$, $`\stackrel{~}{m}\stackrel{~}{๐ด}`$, such that $`\mathrm{๐ซ๐๐ญ}(l)=\rho (\stackrel{~}{m})`$, there is a commutative diagram of the form
(9)
where the rows are exact and the vertical arrows are isomorphisms. $`V_l`$ is the lattice of vanishing cycles in $`H_3(X_l)`$ and $`V_l^{}`$ denotes the orthogonal lattice with respect to the intersection pairing. The map $`p:H_3(X_l,)/V_lV_l^{}`$ is given by $`p(\overline{\mathrm{\Gamma }})=\overline{\mathrm{\Gamma }},`$ which is well defined because $`V_l`$ is isotropic with respect to the intersection pairing.
The claim (i) follows from the definition of relative homology groups. The proof of assertions (ii) and (iii) will be given in section 3. In the remaining part of this section we will explain how these statements lead to large $`N`$ duality up to first order. Recall from Section 2.1 that for $`l๐ณ_0`$ the genus zero closed string partition function can be written locally in terms of splitting of the third homology groups $`H_3(X_l,)`$ into a direct sum of two Lagrangian sublattices. Restricting to an open subset $`๐ฐ^{}`$ of $`๐ณ^{}`$, suppose we can choose the Lagrangian sublattices $`A_l,B_l`$ for a smooth $`X_l`$ so that
$$A_l=V_l\stackrel{~}{A}_l,B_l=U_l\stackrel{~}{B}_l$$
(10)
subject to the following conditions:
* $`A_l`$ is contained in the fixed locus of the monodromy transformations (2).
* $`V_l`$ is orthogonal to $`\stackrel{~}{B}_l`$ and $`U_l`$ is orthogonal to $`\stackrel{~}{A}_l`$ with respect to the intersection pairing.
More concretely, this means that we choose the symplectic basis of cycles so that $`\{\gamma _{1,l},\mathrm{},\gamma _{vr,l}\}`$ generate the lattice $`V_l`$ of vanishing cycles. Then $`U_l`$ is generated by the dual cycles $`\{\gamma _l^1,\mathrm{},\gamma _l^{vr}\}`$ and $`V_l^{}`$ is generated by $`\{\gamma _{\beta ,l},\gamma _l^\alpha \}`$ for $`\alpha =vr+1,\mathrm{},h^{1,2}(X_l)+1`$ and $`\beta =1,\mathrm{},h^{1,2}(X_l)+1`$. Although there is no such canonical choice, note that we can obtain a basis of $`V_l`$ by choosing any collection $`๐ซ_{i_1},\mathrm{},๐ซ_{i_{vr}}`$ of local components of the discriminant which intersect transversely along $`๐ด`$. Then we can complete this basis to a symplectic basis of cycles of $`H_3(X_l,)`$.
Given such a basis of cycles on each smooth $`X_l`$, the commutative diagram (9) shows that the cycles $`\{\gamma _{\alpha ,l},\gamma _l^\alpha \}`$ for $`\alpha =vr+1,\mathrm{},h^{1,2}(X_l)+1`$ induce a symplectic basis of cycles on $`\stackrel{~}{X}_{\stackrel{~}{m}}`$ for $`\stackrel{~}{m}`$ in the open set $`\stackrel{~}{๐ฐ}=\rho ^1(๐ฐ๐ด)`$. Moreover, to each basis element $`\gamma _l^i`$, $`i=1,\mathrm{},vr`$ we can associate an exceptional curve $`C_{i,\stackrel{~}{m}}`$ on $`\stackrel{~}{X}_{\stackrel{~}{m}}`$ so that the $`C_{i,\stackrel{~}{m}}`$ generate $`\mathrm{\Lambda }_{\stackrel{~}{m}}`$ for any $`\stackrel{~}{m}\stackrel{~}{๐ฐ}`$. The images of the cycles $`\{\gamma _l^1,\mathrm{},\gamma _l^{vr}\}`$ in $`H_3(V_l)/V_l`$ induce relative three-cycles $`\stackrel{~}{\mathrm{\Gamma }}_{1,\stackrel{~}{m}},\mathrm{},\stackrel{~}{\mathrm{\Gamma }}_{vr,\stackrel{~}{m}}`$ in $`H_3(\stackrel{~}{X}_{\stackrel{~}{m}},C)`$ so that
$$\delta (\stackrel{~}{\mathrm{\Gamma }}_{i,\stackrel{~}{m}})=C_{i,\stackrel{~}{m}}$$
for $`i=1,\mathrm{},vr`$. Then we have a well defined closed string prepotential $`\stackrel{~}{}:\stackrel{~}{๐ฐ}^{}`$ and open string superpotentials $`\stackrel{~}{}_i:\stackrel{~}{๐ฐ}^{}`$
$$\stackrel{~}{}_i=_{\stackrel{~}{\mathrm{\Gamma }}_{i,\stackrel{~}{m}}}\mathrm{\Omega }_{\stackrel{~}{X}_{\stackrel{~}{m}}}$$
for each $`C_{i,\stackrel{~}{m}}\mathrm{\Lambda }_{\stackrel{~}{m}}`$. The special geometry relations (3)-(4) together with statement (ii) above imply that the closed string prepotential $``$ and its derivatives $`\frac{}{z^\alpha }`$ have well defined limits $`_๐ด^{}`$, $`\left(\frac{}{z^\alpha }\right)_๐ด^{}`$ at the points of $`๐ฐ^{}๐ด^{}`$. The relations (7) follow from equation (8). Note that these relations are valid for any choice of a symplectic basis in $`H_3(X_l)`$ satisfying conditions $`(a)`$ and $`(b)`$ above.
## 3 Geometric Transitions and Mixed Hodge Structures
In this section we would like to prove statements (ii) and (iii). We will consider the following abstract situation. Let $`\pi :๐ณ\mathrm{\Delta }`$ be a one parameter family of projective Calabi-Yau threefolds over the unit disc $`\mathrm{\Delta }`$. The generic fiber $`X_\mu `$, $`\mu \mathrm{\Delta }\{0\}`$ is assumed to be smooth and the central fiber $`X_0`$ is a nodal threefold with ordinary double points $`p_1,\mathrm{},p_v`$. We will assume that over the complement of the central fiber, the family $`\pi :๐ณ\mathrm{\Delta }`$ is equipped with a globally defined relative holomorphic 3-form, which restricts to a holomorphic volume form on each fiber. Moreover we will assume that there is a smooth crepant projective resolution $`\stackrel{~}{X}X_0`$. Such a family can be obtained for example by taking $`\mathrm{\Delta }`$ to be a holomorphic disc in the moduli space $`๐ณ`$ intersecting $`๐ด`$ transversely at the origin.
The restriction of the family $`๐ณ\mathrm{\Delta }`$ to the punctured disc $`\mathrm{\Delta }^\times =\mathrm{\Delta }\{0\}`$ is a family of smooth Calabi-Yau threefolds which determines a geometric variation of Hodge structures. To this data we can associate a period map $`\varphi :\mathrm{\Delta }^\times D/๐`$ . where $`D`$ is the classifying space of Hodge structures and $`๐`$ is the monodromy group of the family.
The proof of claim (ii) is Hodge theoretic and is based on the nilpotent orbit theorem and the Clemens-Schmid exact sequence associated to the degeneration $`๐ณ\mathrm{\Delta }`$. In the process of proving claim (ii), we will also show that the Clemens-Schmid exact sequence implies claim (iii).
First observe that the nilpotent orbit theorem implies the existence of a well defined limit of the periods of the holomorphic three-form at the origin $`0\mathrm{\Delta }`$. Suppose $`\mathrm{\Gamma }`$ is a multivalued section of the sheaf of third homology groups $`H_3(X_s)`$ over $`\mathrm{\Delta }^\times `$. Let
$$\mathrm{\Pi }(s)=_{\mathrm{\Gamma }_s}\mathrm{\Omega }_{X_s}$$
be the period of the holomorphic three-form on $`\mathrm{\Gamma }_s`$, for $`s0`$. $`\mathrm{\Pi }`$ is a multivalued holomorphic function on the punctured disc $`\mathrm{\Delta }^\times =\mathrm{\Delta }\{0\}`$. Let
$$\psi :\mathrm{\Delta }^{}D^{},\psi (s)=\text{exp}\left(\frac{\mathrm{log}s}{2\pi i}N\right)\varphi (s),$$
be the modified the period map $`\varphi `$ with values in the compact dual $`D^{}`$ of $`D`$ . Here $`N=\mathrm{log}(T)=T\mathrm{id}`$ is the logarithm of the monodromy transformation about the origin. The identity $`\mathrm{log}(T)=T\mathrm{id}`$ is equivalent to $`(T\mathrm{id})^2=0`$ which holds for any ODP degeneration of threefolds. Indeed, this follows from the observation that if we write $`T\mathrm{id}`$ as a block matrix with respect to the decompositions (10), then the only non-zero block of $`T\mathrm{id}`$ is the one sending $`U_l`$ to $`V_l`$.
According to the nilpotent orbit theorem , the map $`\psi `$ can be extended to a single valued holomorphic map
$$\overline{\psi }:\mathrm{\Delta }D^{}$$
In particular, the multivalued map $`\mathrm{\Pi }:\mathrm{\Delta }^\times `$, which is a matrix coefficient of $`\varphi `$, can be written as
$$\mathrm{\Pi }(s)=\overline{\mathrm{\Pi }}(s)+\frac{1}{2\pi i}(\text{log s})\eta (s)$$
for some $`\overline{\mathrm{\Pi }}(s)`$, $`\eta (s)`$ single valued holomorphic functions of $`s`$ with $`\eta (0)=0`$. Later in the proof of relation (8) we will see that the value $`\overline{\mathrm{\Pi }}(0)`$ admits an intrinsic description of in terms of the mixed Hodge structure of the semistable model of the degeneration $`๐ณ\mathrm{\Delta }`$.
The point $`\overline{\psi }(0)D^{}`$ corresponds to the limiting mixed Hodge structure on the third cohomology group of a smooth fiber $`H^3(X_s)`$, $`s0`$. On the other hand the relative cohomology group $`H^3(\stackrel{~}{X},C)`$ also carries a canonical mixed Hodge structure . The relation (8) that we wish to prove asserts that a period in the limiting mixed Hodge structure on $`H^3(X_s)`$ is equal to the corresponding period in the mixed Hodge structure on $`H^3(\stackrel{~}{X},C)`$. In particular (8) will follow if we can show that the natural map $`H^3(\stackrel{~}{X},C)H^3(X_s)`$ is an inclusion of mixed Hodge structures.
We claim that this follows from the Clemens-Schmid exact sequence. In order to formulate a more precise statement, let us first describe the mixed Hodge structure on $`H^3(\stackrel{~}{X},C)`$ and construct the Clemens-Schmid exact sequence for our degeneration.
Recall that a mixed Hodge structure on a vector space $`H=H_{}`$ is defined by
* a descending Hodge filtration $`\{F^k\}`$, and
* an ascending weight filtration $`\{W_m\}`$, defined over $``$, so that for every $`m`$ the successive quotient $`W_m/W_{m1}`$ has a pure Hodge structure of weight $`m`$ given by the induced filtration
$$F^k(W_m/W_{m1})=(W_mF^k)/(W_{m1}F^k).$$
The limiting mixed Hodge structure on the cohomology of a smooth fiber is defined by the Hodge filtration given by the nilpotent orbit and by the monodromy weight filtration
$$0W_0W_1\mathrm{}W_5W_6=H^3(X_s).$$
which is determined by the monodromy action on cohomology. In our case, $`N^2=0`$, and the monodromy weight filtration is very simple
$$W_3=\text{ker}(N),W_2=\text{im}(N).$$
Additional information about the limiting mixed Hodge structure can be obtained from the Clemens-Schmid exact sequence as we will see below.
To describe the mixed Hodge structure on $`H^3(\stackrel{~}{X},C)`$, consider the long exact sequence of the pair $`(\stackrel{~}{X},C)`$
$$\mathrm{}H^2(\stackrel{~}{X})H^2(C)H^3(\stackrel{~}{X},C)H^3(\stackrel{~}{X})0$$
where all cohomology groups have complex coefficients. The ascending weight filtration on $`H^3(\stackrel{~}{X},C)`$ is defined by
$`W_3=H^3(\stackrel{~}{X},C)`$
$`W_2=\text{Im}\left(H^2(C)H^3(\stackrel{~}{X},C)\right)H^2(C)/\text{im}\left(H^2(\stackrel{~}{X})H^2(C)\right)`$
$`W_1=0.`$
The Hodge filtration is the standard Hodge filtration on the cohomology of the quasi-projective variety $`XC`$. The relation between the mixed Hodge structures on $`H^3(X_s)`$ and $`H^3(\stackrel{~}{X},C)`$ will be extracted from the Clemens-Schmid exact sequence, which is constructed in terms of a semistable model of the degeneration $`๐ณ\mathrm{\Delta }`$.
A semistable model for the family $`๐ณ\mathrm{\Delta }`$ is a new family $`\overline{๐ณ}\mathrm{\Delta }`$ which fits in a commutative diagram of the form
(11)
Here $`\overline{๐ณ}`$ is assumed smooth, $`f:\overline{\mathrm{\Delta }}\mathrm{\Delta }`$ is a finite cover of $`\mathrm{\Delta }`$ branched at the origin and the dashed arrow represents a birational map which is an isomorphism over the punctured disc $`\overline{\mathrm{\Delta }}^\times `$. In addition, the central fiber $`\overline{X}\overline{X}_0`$ is required to be a normal crossing divisor in $`\overline{๐ณ}`$ with smooth reduced irreducible components.
The cohomology of the central fiber $`\overline{X}_0`$ of the semistable degeneration can be equipped with a mixed Hodge structure, which can be described in terms of the Mayer-Vietoris spectral sequence. The Clemens-Schmid exact sequence is an exact sequence of mixed Hodge structures of the form
$$\mathrm{}H_1(\overline{X}_0)H^3(\overline{X}_0)H^3(\overline{X}_{\overline{s}})\stackrel{N}{}H^3(\overline{X}_{\overline{s}})H_3(\overline{X}_0)\mathrm{}$$
The Clemens-Schmid theorem guarantees that all maps in this sequence are morphisms of mixed Hodge structures. The mixed Hodge structure on homology is induced by the mixed Hodge structure on cohomology using the universal coefficient formula.
Now, claim (ii) follows from the following two facts which will be proven in the rest of the section.
* $`H^3(\overline{X}_0)\text{ker}(N:H^3(\overline{X}_{\overline{s}})H^3(\overline{X}_{\overline{s}}))`$ as mixed Hodge structures. This follows from $`H_1(\overline{X}_0)=0`$.
* $`H^3(\stackrel{~}{X},C)H^3(\overline{X}_0)`$ as mixed Hodge structures.
To begin with, let us construct the semistable degeneration for our family $`๐ณ\mathrm{\Delta }`$. Let $`f:\overline{\mathrm{\Delta }}\mathrm{\Delta }`$ be a double cover of the disc defined in local coordinates by $`s=\overline{s}^2`$, and let $`f^{}๐ณ`$ denote the pull-back of the family $`๐ณ\mathrm{\Delta }`$ to $`\overline{\mathrm{\Delta }}`$. The total space of the family $`f^{}๐ณ\overline{\mathrm{\Delta }}`$ has double point singularities at the points $`p_i`$ on the central fiber $`(f^{}๐ณ)_0`$, which is canonically isomorphic to $`X_0`$. We construct a new family $`\overline{๐ณ}\overline{\mathrm{\Delta }}`$ by blowing-up the singular points on the total space of $`f^{}๐ณ`$. Since all the singular points lie on the central fiber, it is an easy check that the fiber $`\overline{X}_{\overline{s}}`$ is isomorphic to $`X_s`$, $`s=\overline{s}^2`$, for $`\overline{s}0`$. The central fiber $`\overline{X}_0`$ is a normal crossing variety consisting of $`v+1`$ smooth reduced irreducible components
$$\overline{X}_0=\widehat{X}Q_1\mathrm{}Q_v$$
where $`\widehat{X}X_0`$ is the blow-up of $`X_0`$ at the $`v`$ singular points and $`Q_1,\mathrm{},Q_v`$ are quadric threefolds. The blow-up $`\widehat{X}X_0`$ replaces each ordinary double point $`p_i`$ with an exceptional divisor $`E_i`$ isomorphic to the Hirzebruch surface $`๐ฝ_0`$. $`\widehat{X}`$ intersects each quadric threefold $`Q_i`$ along $`E_i`$, which is a hyperplane section of $`Q_i`$. Note that the exceptional divisors $`E_i`$ are pairwise disjoint, hence the components $`Q_i`$ have no common points. Therefore $`\overline{X}_0`$ is indeed a normal crossing divisor with smooth reduced irreducible components.
Note that a small resolution of $`X_0`$ can be obtained by contracting $`\widehat{X}`$ along a collection of rulings of the exceptional divisor $`E_i\widehat{X}`$. Since each divisor $`E_i`$ admits two distinct rulings, we can obtain in principle different small resolutions $`\stackrel{~}{X}`$ related by flops. Note however that not all possible contractions result in projective small resolutions. The considerations of this section are valid for any projective contraction $`\widehat{X}\stackrel{~}{X}`$.
Next, we will compute the rational homology of the semi-stable model $`\overline{X}_0`$ using Mayer-Vietoris exact sequences. First we compute the homology of the singular fiber $`X_0`$ in the initial family. From a topological point of view $`X_0`$ can be represented as the cone of a map from a disjoint collection of $`v`$ three-spheres to a smooth fiber $`X_s`$
$$f:_{i=1}^vS_i^3X_s$$
so that $`f`$ maps each $`S_i^3`$ homeomorphically onto a vanishing cycle associated to the $`i`$-th node of $`X_0`$. Therefore $`X_0`$ is homotopy equivalent to the union of $`X_s`$ and $`v`$ closed four-discs $`D_i^4`$, $`i=1,\mathrm{},4`$ so that $`X_sD_i^4=f(S_i^3)`$. The associated homology Mayer-Vietoris sequence reads
$$\mathrm{}_{i=1}^vH_k(S_i^3)H_k(X_s)_{i=1}^vH_k(D^4)H_k(X_0)_{i=1}^vH_{k1}(S_i^3)\mathrm{}$$
(12)
In particular we have
$$0H_4(X_s)H_4(X_0)_{i=1}^vH_3(S_i^3)H_3(X_s)H_3(X_0)0$$
The image of the map $`_{i=1}^vH_3(S_i^3)H_3(X_s)`$ is the $`vr`$ dimensional lattice $`V`$ of vanishing cycles on $`X_s`$, and the kernel is the $`r`$-dimensional lattice $`R`$ of relations among vanishing cycles. Therefore we find two short exact sequences
$`0H_4(X_s)H_4(X_0)R0`$
$`0VH_3(X_s)H_3(X_0)0`$
which yield the relations $`b_3(X_0)=b_3(X_s)(vr)`$, $`b_4(X_0)=b_4(X_s)+r`$. The remaining Betti numbers can be easily determined by writing down the other terms in (12). We record the complete results below, using the notation $`b_kb_k(X_s)`$
| $`i`$ | $`6`$ | $`5`$ | $`4`$ | $`3`$ | $`2`$ | $`1`$ | $`0`$ |
| --- | --- | --- | --- | --- | --- | --- | --- |
| $`b_i(X_s)`$ | $`1`$ | $`0`$ | $`b_2`$ | $`b_3`$ | $`b_2`$ | $`0`$ | $`1`$ |
| $`b_i(X_0)`$ | $`1`$ | $`0`$ | $`b_2+r`$ | $`b_3(vr)`$ | $`b_2`$ | $`0`$ | $`1`$ |
(13)
In order to determine the homology of the blow-up $`\widehat{X}`$, consider the long exact homology sequence for the pair $`(\widehat{X},E)`$
$$\mathrm{}H_k(E)H_k(\widehat{X})H_k(\widehat{X},E)H_{k1}(E)\mathrm{},$$
(14)
where $`E=_{i=1}^vE_i`$ denotes the exceptional divisor. Using the fact that $`E_i^1\times ^1`$ for each $`i=1,\mathrm{},v`$ and that $`X_0`$ is $`\widehat{X}`$ with each $`E_i`$ collapsed to a point, we see that (14) reduces to the following straightforward isomorphisms
$`H_k(\widehat{X})H_k(X_0),k=0,1,5,6`$
and two exact sequences of the form
$`0H_4(E)H_4(\widehat{X})H_4(\widehat{X},E)0`$ (15)
$`0H_3(\widehat{X})H_3(\widehat{X},E)\stackrel{\delta }{}H_2(E)\stackrel{\iota _{}}{}H_2(\widehat{X})H_2(\widehat{X},E)0.`$
The first sequence yields
$$b_4(\widehat{X})=b_4(X_0)+v=b_2(X_s)+v+r.$$
By Poincarรฉ duality we also have
$$b_2(\widehat{X})=b_4(\widehat{X})=b_2(X_s)+v+r.$$
Also taking into account (13) and the second sequence in (15) we find
$$b_3(\widehat{X})=b_3(X_s)2(vr).$$
Moreover, it follows that
$$\text{dim}(\text{ker}(\iota _{}))=vr,\text{dim}(\text{im}(\iota _{}))=v+r$$
where $`\iota _{}:H_2(E)H_2(\widehat{X})`$ is the map on homology determined by the inclusion $`\iota :E\widehat{X}`$. Note that $`\text{ker}(\iota _{})=\text{im}(\delta )`$ is spanned by relative homology three-cycles for the pair $`(\widehat{X},E)`$ modulo cycles in $`\widehat{X}`$. These cycles generate the lattice of relations among curve classes on $`E`$ regarded as homology cycles on $`\widehat{X}`$.
There is a similar exact sequence for the pair $`(\stackrel{~}{X},C)`$, where $`\stackrel{~}{X}`$ is a small projective resolution of $`X_0`$ and $`C`$ is the collection of exceptional curves
$$\mathrm{}H_k(C)H_k(\stackrel{~}{X})H_k(\stackrel{~}{X},C)H_{k1}(C)\mathrm{}$$
(16)
Again, by construction we have $`C_i^1`$ and $`X_0`$ is $`\stackrel{~}{X}`$ with each $`C_i`$ collapsed to a point. The long exact sequence (16) reduces to isomorphisms
$$H_k(\stackrel{~}{X})=H_k(X_0),k=0,1,4,5,6,$$
(17)
and a shorter exact sequence
$`0H_3(\stackrel{~}{X})H_3(\stackrel{~}{X},C)\stackrel{\stackrel{~}{\delta }}{}H_2(C)\stackrel{\stackrel{~}{\iota }_{}}{}H_2(\stackrel{~}{X})H_2(\stackrel{~}{X},C)0.`$ (18)
Then, using (17), (18) and the fact that $`b_2(\stackrel{~}{X})=b_4(\stackrel{~}{X})`$ one gets
$$b_3(\stackrel{~}{X})=b_3(X_s)2(vr).$$
and
$$\text{dim}(\text{ker}(\stackrel{~}{\iota }_{}))=vr,\text{dim}(\text{im}(\stackrel{~}{\iota }_{}))=r$$
where $`\stackrel{~}{\iota }_{}:H_2(C)H_2(\stackrel{~}{X})`$ is the map on homology determined by the inclusion $`\stackrel{~}{\iota }:C\stackrel{~}{X}`$. Note that $`\text{ker}(\stackrel{~}{\iota }_{})\text{im}(\stackrel{~}{\delta })`$ is spanned by relative homology three-cycles on $`(\stackrel{~}{X},C)`$ modulo cycles in $`\stackrel{~}{X}`$. Just as in the blow-up case these cycles generate the lattice of relations among exceptional curve classes on $`\stackrel{~}{X}`$.
For future reference, we analyze the direct relation between the rational homology of $`\widehat{X}`$ and the homology of $`\stackrel{~}{X}`$. Choose a basis of homology two-cycles $`(a_i,b_i)`$ in $`H_2(E_i)`$, $`i=1,\mathrm{},v`$ corresponding to the two rulings so that the contraction map $`q:\widehat{X}\stackrel{~}{X}`$ contracts all rational curves in the classes $`b_i`$, $`i=1,\mathrm{},v`$. A simple intersection computation yields the following relations in the homology ring of $`\widehat{X}`$
$$a_iE_j=b_iE_j=\delta _{ij},i,j=1,\mathrm{},v.$$
(19)
Now let us consider the following commutative diagram of spaces
where $`p:EC`$ is a projection map contracting the rulings $`b_i`$, $`i=1,\mathrm{},v`$. Then we have a natural relation between maps on homology
$$q_{}\iota _{}=\stackrel{~}{\iota }_{}p_{}.$$
This shows that that $`p_{}:H_2(E)H_2(C)`$ induces a map
$$\overline{p}_{}:\text{ker}(\iota _{})\text{ker}(\stackrel{~}{\iota }_{}).$$
Moreover, one can show that $`\overline{p}_{}`$ is injective, since any element in the kernel of $`p_{}`$ is necessarily a linear combination of $`b_iH_2(E)`$. However, any nontrivial linear combination of $`b_i`$ cannot lie in $`\text{ker}(\iota _{})`$ because the intersection pairing on $`(b_i,E_i)`$, $`i=1,\mathrm{},v`$ is nondegenerate according to (19). Therefore $`\text{ker}(\overline{p}_{})=0`$. Since $`\text{ker}(\iota _{}),\text{ker}(\stackrel{~}{\iota }_{})`$ are $``$-vector spaces of equal dimension $`vr`$, it follows that $`\overline{p}_{}`$ is an isomorphism. Furthermore, $`\text{ker}(\iota _{})`$ is contained in the linear span of the homology classes $`a_ib_i`$, $`i=1,\mathrm{},v`$.
Note also that the map $`q_{}:H_3(\widehat{X},E)H_3(\stackrel{~}{X},C)`$ is an isomorphism. This follows from the commutative diagram of pairs
which induces isomorphisms $`H_3(\widehat{X},E)H_3(X_0,X_0^{sing})`$ and $`H_3(\stackrel{~}{X},C)H_3(X_0,X_0^{sing})`$. Therefore we obtain the following commutative diagram of homology groups
(20)
where the rows are exact and the middle and rightmost vertical arrows are isomorphisms. We conclude that the leftmost vertical arrow must be an isomorphism as well, and there is also an isomorphism
$$H_3(\widehat{X},E)/H_3(\widehat{X})H_3(\stackrel{~}{X},C)/H_3(\stackrel{~}{X})$$
To conclude the homology computations, let us determine the homology of the central fiber $`\overline{X}_0`$ of the semistable degeneration. We will employ again a Mayer-Vietoris exact sequence with respect to the closed cover
$$\overline{X}_0=\widehat{X}Q_1,\mathrm{},Q_v$$
where $`Q_i`$ are smooth quadric threefolds intersecting $`\widehat{X}`$ transversely along the exceptional divisors $`E_i`$, $`i=1,\mathrm{},v`$. Let $`Q`$ denote the disjoint union $`Q=Q_1\mathrm{}Q_v`$. Then, we have
$`0`$ $`H_6(\widehat{X})H_6(Q)H_6(\overline{X}_0)H_5(E)H_5(\widehat{X})H_5(Q)H_5(\overline{X})`$ (21)
$`H_4(E)H_4(\widehat{X})H_4(Q)H_4(\overline{X}_0)H_3(E)H_3(\widehat{X})H_3(Q)`$
$`H_3(\overline{X})H_2(E)H_2(\widehat{X})H_2(Q)H_2(\overline{X}_0)H_1(E)`$
$`H_1(\widehat{X})H_1(Q)H_1(\overline{X})H_0(E)H_0(\widehat{X})H_0(Q)H_0(\overline{X}_0)0`$
The homology of $`Q`$ can be computed easily from the Lefschetz hyperplane theorem and an Euler characteristic computation via Gauss-Bonet formula. We have
$$H_k(Q)=,k=0,2,4,6,H_k(Q)=0,k=1,3,5.$$
Moreover, since $`E_i`$ is a hyperplane section of $`Q_i`$, $`i=1,\mathrm{},v`$, it follows that the map
$$H_4(E)H_4(\widehat{X})H_4(Q)$$
is injective. Thus the long exact sequence (21) yields the following straightforward isomorphisms
$`H_6(\overline{X}_0)H_6(\widehat{X})H_6(Q)=^{(1+v)}`$ (22)
$`H_5(\overline{X}_0)H_5(\widehat{X})H_5(Q)=0`$
$`H_0(\overline{X}_0)H_0(\widehat{X})=`$
$`H_1(\overline{X}_0)H_1(\widehat{X})=0`$
The remaining part of (21) splits into two exact sequences
$`0H_4(E)H_4(\widehat{X})H_4(Q)H_4(\overline{X}_0)0`$ (23)
$`0H_3(\widehat{X})H_3(\overline{X}_0)H_2(E)\stackrel{(\iota _{},j_{})}{}H_2(\widehat{X})H_2(Q)H_2(\overline{X}_0)0`$
where the $`j_{}:H_2(E)H_2(Q)`$ is the map on homology determined by the inclusion $`j:EQ`$. The first sequence gives
$$b_4(\overline{X})=b_4(\widehat{X})+vv=b_4(\widehat{X})=b_2(X_s)+v+r.$$
Since $`E`$ is a hyperplane section of $`Q`$, the kernel of the map $`j_{}:H_2(E)H_2(Q)`$ is spanned by the homology classes $`a_ib_i`$, $`i=1,\mathrm{},v`$. On the other hand $`\text{ker}(\iota _{})`$ is a $`vr`$ dimensional subspace of the linear span of $`a_ib_i`$, $`i=1,\mathrm{},v`$. Therefore we find that
$$\text{ker}(\iota _{},j_{})\text{ker}(\iota _{})$$
is also $`vr`$ dimensional. This fixes the remaining Betti numbers
$`b_3(\overline{X})`$ $`=b_3(\widehat{X})+(vr)=b_3(X_s)(vr)`$
$`b_2(\overline{X})`$ $`=b_2(\widehat{X})+v(2vv+r)=b_2(X_s)+v.`$
In summary we can extend the table (13) to include all the spaces appearing in the geometric transition and the semi-stable degeneration:
| $`i`$ | $`6`$ | $`5`$ | $`4`$ | $`3`$ | $`2`$ | $`1`$ | $`0`$ |
| --- | --- | --- | --- | --- | --- | --- | --- |
| $`b_i(X_s)`$ | $`1`$ | $`0`$ | $`b_2`$ | $`b_3`$ | $`b_2`$ | $`0`$ | $`1`$ |
| $`b_i(X_0)`$ | $`1`$ | $`0`$ | $`b_2+r`$ | $`b_3(vr)`$ | $`b_2`$ | $`0`$ | $`1`$ |
| $`b_i(\widehat{X})`$ | $`1`$ | $`0`$ | $`b_2+v+r`$ | $`b_32(vr)`$ | $`b_2+v+r`$ | $`0`$ | $`1`$ |
| $`b_i(\stackrel{~}{X})`$ | $`1`$ | $`0`$ | $`b_2+r`$ | $`b_32(vr)`$ | $`b_2+r`$ | $`0`$ | $`1`$ |
| $`b_i(\overline{X}_0)`$ | $`1+v`$ | $`0`$ | $`b_2+v+r`$ | $`b_3(vr)`$ | $`b_2+v`$ | $`0`$ | $`1`$ |
| $`b_i(\stackrel{~}{X},C)`$ | $`1`$ | $`0`$ | $`b_2+r`$ | $`b_3(vr)`$ | $`b_2`$ | $`v1`$ | $`1`$ |
(24)
With all this information in place, we are ready to finish the proof of (ii).
Recall that
$$\text{ker}(\iota _{})H_3(\widehat{X},E)/H_3(\widehat{X})$$
according to the top exact sequence in (20). Then the second exact sequence in (23) yields a short exact sequence of the form
$$0H_3(\widehat{X})H_3(\overline{X}_0)H_3(\widehat{X},E)/H_3(\widehat{X})0.$$
(25)
Comparing (25) with the top row of (20) we get an isomorphism $`H_3(\overline{X}_0)H_3(\widehat{X},E)/H_3(\widehat{X})`$ of mixed Hodge structures. Combined with the fact that the vertical arrows in (20) are isomorphisms we get an isomorphism $`H_3(\overline{X}_0)H_3(\stackrel{~}{X},C)`$ as mixed Hodge structures. Dualizing and applying the universal coefficients theorem we get the desired isomorphism of mixed Hodge structures $`H^3(\overline{X}_0)H^3(\stackrel{~}{X},C)`$. Since by (22) we have $`H_1(\overline{X}_0)=0`$ this completes the proof of (ii).
At this point we still have to tie a few loose ends and prove claim (iii). Note that the universal coefficients theorem and (24) give that $`H^5(\overline{X}_0)=0`$. Therefore the Clemens-Schmid exact sequence reduces to
$$0H^3(\overline{X}_{\overline{s}})/\text{im}(N)\stackrel{\overline{\beta }}{}H_3(\overline{X}_0)0.$$
The middle map $`\overline{\beta }`$ in this exact sequence is defined by the sequence of maps
$$H^3(\overline{X}_{\overline{s}})\stackrel{PD}{}H_3(\overline{X}_{\overline{s}})H_3(\overline{๐ณ})H_3(\overline{X}_0)$$
where the first map is Poicarรฉ duality on a smooth fiber and the next map is induced by inclusion in the total space $`\overline{๐ณ}`$ of the semistable family. Since the image of $`N=\text{log}(T)=T\text{id}`$ is the space $`V`$ of vanishing cycles, it follows that we have an isomorphism of mixed Hodge structures
$$0H_3(\overline{X}_{\overline{s}})/VH_3(\overline{X}_0)0.$$
(26)
The weight filtration on $`H_3(\overline{X}_{\overline{s}})/V`$ is a one step filtration induced by the monodromy weight filtration of $`H_3(\overline{X}_{\overline{s}})`$
$`W_3`$ $`=H_3(\overline{X}_{\overline{s}})/V`$
$`W_2`$ $`=(\mathrm{ker}N)/V=V^{}/V`$
$`W_1`$ $`=0`$
Therefore we have a three-term exact sequence
$$0V^{}/VH_3(\overline{X}_{\overline{s}})/V\stackrel{p}{}V^{}0$$
where the map $`p`$ is induced by $`N=\text{log}(T)=T\mathrm{id}`$
$$p(\overline{\mathrm{\Gamma }})=\overline{\mathrm{\Gamma }},.$$
The monodromy weight filtration on $`H_3(\overline{X}_0)`$ is also a one step filtration of the form
$`W_3`$ $`=H_3(\overline{X}_0)H_3(\stackrel{~}{X},C)`$
$`W_2`$ $`=H_3(\widehat{X})H_3(\stackrel{~}{X})`$
$`W_1`$ $`=0.`$
Since (26) is an isomorphism of mixed Hodge structures, it is in particular compatible with the weight filtrations. This implies the commutative diagram (9) of claim (iii).
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